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warning/0001/hep-ph0001015.html | ar5iv | text | # hep-ph/0001015FTUV/00-03IFIC/00-03 Status of the MSW Solutions to the Solar Neutrino Problem
In this talk we present the results of an updated global analysis of two-flavor MSW solutions to the solar neutrino problem in terms of conversions of $`\nu _e`$ into active or sterile neutrinos including the the full data set corresponding to the 825-day Super–Kamiokande data sample as well as to Chlorine, GALLEX and SAGE experiments.
It is already three decades since the first detection of solar neutrinos. It was realized from the very beginning that the observed rate at the Homestake experiment was far lower than the theoretical expectation based on the standard solar model with the implicit assumption that neutrinos created in the solar interior reach the earth unchanged, i.e. they are massless and have only standard properties and interactions. From the experimental point of view much progress has been done in recent years. We now have available the results of five experiments, the original Chlorine experiment at Homestake , the radio chemical Gallium experiments on pp neutrinos, GALLEX and SAGE , and the water Cherenkov detectors Kamiokande and Super–Kamiokande which we summarize in Table 1.
Super–Kamiokande has been able not only to confirm the original detection of solar neutrinos at lower rates than predicted by standard solar models, but also to demonstrate directly that the neutrinos come from the sun by showing that recoil electrons are scattered in the direction along the sun-earth axis. We now have good information on the time dependence of the event rates during the day and night, as well as a measurement of the recoil electron energy spectrum. After 825 days of operation, Super–Kamiokande has also presented preliminary results on the seasonal variation of the neutrino event rates an issue which will become important in discriminating the MSW scenario from the possibility of neutrino oscillations in vacuum .
In our study we use the following observables:
* the three measured rates shown in Table 1
* Super–Kamiokande results on the zenith angular dependence of the event rates during 1 day and 5 night periods.
* Recoil e spectrum including the 2 points from obtained with the super low energy threshold below 6.5 GeV and 18 points with 6.5¡E¡15 MeV.
* Seasonal variation of the event rates measured in 8 periods of 1.5 months each.
We obtain the allowed value of the parameters and the corresponding CL for the different scenarios by a $`\chi ^2`$ analysis, details of which can be found in Ref. .
In our calculations of the expected values for these observables we use as SSM the fluxes from Ref. but we also consider departures of the SSM by allowing arbitrary $`{}_{}{}^{8}B`$ and $`hep`$ fluxes. For the Chlorine and Gallium experiments we use improved cross sections $`\sigma _{\alpha ,i}(E)`$ $`(\alpha =e,x)`$ from Ref. . For the Super–Kamiokande experiment we calculate the expected signal with the differential cross section $`d\sigma _\alpha (E_\nu ,T^{})/dT^{}`$, that we take from taking into account the finite energy resolution of the experiment which implies that the measured kinetic energy $`T`$ of the scattered electron is distributed around the true kinetic energy $`T^{}`$ according to a resolution function which we take from .
Using the predicted fluxes from the BP98 model the $`\chi ^2`$ for the fit to the three experimental rates is $`\chi _{SSM}^2=62.4/3`$dof what implies that the probability of the observations to be an statistical fluctuation of the SSM is lower than $`10^8`$!!. One may wonder about the possible dependence of the quality of description on the specific solar model used. In order to address this issue we try to fit the data by allowing a free normalization of the dominant neutrino fluxes $`pp`$ $`{}_{}{}^{7}Be`$ and $`{}_{}{}^{8}B`$ only imposing the constraint that the luminosity of the sun is supplied by nuclear reactions among the light elements what implies a linear relation among the three normalizations. The best fit point corresponds to an unphysical situation with negative $`{}_{}{}^{7}Be`$ neutrino flux. After constraining the fluxes to be positive we obtain that the best fit point occurs at $`{}_{}{}^{7}Be/^7Be_{SSM}=0`$, $`pp/pp_{SSM}=1.08`$ and $`{}_{}{}^{8}B/^8B_{SSM}=0.53`$ with $`\chi _{min}^2=21.4/1`$dof which implies that there is no acceptable fit with a CL better than $`5\times 10^4`$%.
Next we test the possibility of describing the data in terms of an energy independent neutrino conversion probability as expected, for instance in models explaining all evidences for neutrino oscillations (from solar, atmospheric and LSND data) in terms of three massive neutrinos. We find the values listed in Table 2
As seen in the table all scenarios with constant survival probability are ruled out with a CL larger than 99 %.
We now consider the description in terms of matter enhanced neutrino oscillations via the MSW mechanism . We use in this case the neutrino survival probabilities in the presence of matter given in Ref. . We show in Fig .1 the allowed regions for oscillations into active and sterile neutrinos for different combinations of the observables. In Table 3 we give the location of the best fit point for the different regions as well as the corresponding CL.
Several comments are in order. First we see that for active neutrinos there are three allowed regions, the small mixing angle region (SMA), the large mixing angle region (LMA) and the low mass region (LOW). For sterile neutrinos only the SMA solution is allowed. This arises from the fact that unlike active neutrinos which lead to events in the Super–Kamiokande detector by interacting via neutral current with the electrons, sterile neutrinos do not contribute to the Super–Kamiokande event rates. Therefore a larger survival probability for $`{}_{}{}^{8}B`$ neutrinos is needed to accommodate the measured rate. As a consequence a larger contribution from $`{}_{}{}^{8}B`$ neutrinos to the Chlorine and Gallium experiments is expected, so that the small measured rate in Chlorine can only be accommodated if no $`{}_{}{}^{7}Be`$ neutrinos are present in the flux. This is only possible in the SMA solution region, since in the LMA and LOW regions the suppression of $`{}_{}{}^{7}Be`$ neutrinos is not enough.
As for the quality of the different solutions for the oscillations into active neutrinos we find that the different observables favour different solutions. The total rates favours the SMA solution while the inclusion of the zenith angular dependence favours the LMA solution as although small, some effect is observed in the zenith angle dependence which points towards a larger event rate during the night than during the day, and that this difference is constant during the night as expected for the LMA solution. In the SMA solution, however, the enhancement is expected to occur mainly in the fifth night. The spectral information is such that the oscillation hypothesis does not improve considerably the fit to the energy spectrum as compared to the no-oscillation hypothesis as the data is basically consistent with a flat distribution what is also in better agreement with the LMA and LOW solutions while the SMA predicts a continuously raising spectrum. Moreover, the observation of a possible seasonal variation of the higher energy event rates in Super–Kamiokande could also be accommodated in terms of the LMA solution to the solar neutrino problem . So once all the observables are combined we find that all solutions are allowed at the 90 % CL. LMA and SMA give similar descriptions, the LMA being slightly favoured. |
warning/0001/hep-ph0001232.html | ar5iv | text | # Generalized Jacobi Identities, Curvature Relation, Schouten’s Identity, a Phase Rule and Derivation of O(𝑞⁴) Effective Lagrangian in the Presence of External Fields Directly Within Heavy Baryon Chiral Perturbation Theory
## 1 Introduction
Heavy Baryon Chiral Perturbation Theory (HBChPT) is a nonrelativistic (with respect to the “heavy” baryons) effective field theory used for studying meson-baryon interactions at low energies, typically below the mass of the first non-Goldstone resonance (See ). The degrees of freedom of SU(2) ($``$ isospin) HBChPT (which will be considered in this paper) are the (derivatives of) pion-triplet, the nucleon fields and the external fields.
Recently, a method was developed to generate HBChPT Lagrangian ($`_{\mathrm{HBChPT}}`$) for off-shell nucleons directly within HBChPT, which as stated in , can prove useful when applying HBChPT to nuclear processes in which the nucleons are bound, and hence off-shell. This method has the advantage of not having to bother to start with the relativistic BChPT Lagrangian and then carry out the nonrelativistic reduction. It is thus shorter than the standard approach to HBChPT as given in , showed explicitly up to O$`(q^3)`$ in . In the context of off-shell nucleons, the upshot of the method developed is a phase rule (See (6),) to implement charge conjugation invariance (along with Lorentz invariance, parity and hermiticity) directly within HBChPT. The phase rule, along with additional reductions from a variety of algebraic identities, was used to construct, directly within HBChPT, a complete list of off-shell O($`q^3`$) terms (in the isospin-conserving approximation and in the abscence of external fields). We also showed that the on-shell limit of the list of terms obtained matches the corresponding list in (in which the HBChPT Lagrangian up to O$`(q^3`$) was constructed starting from the relativistic BChPT Lagrangian).
For a complete and precise calculation in the single-nucleon sector to one loop, e.g., 1-loop corrections to pion production off a single (on-shell) nucleon, because of convergence problems (assocated with the amplitude “$`D_2`$” for pion production off a single nucleon), one needs to go up to O$`(q^4)`$ (See ). An overcomplete list of the divergent O($`q^4`$) $`\pi `$-nucleon interaction terms in the presence of external fields was constructed in , but again starting from the relativistic BChPT Lagrangian. In , we constructed a complete list of O$`(q^4)`$ terms working entirely within the framework of HBChPT in the absence of external fields and in the isospin-conserving approximation. In this paper, we extend the list to include external fields and do not assume isospin symmetry.
Section 2 has the basics and sets up the notations. In Section 3, reductions obtained in addition to (6) due to algebraic identities such as generalized Jacobi identities, Schouten’s identitiy and curvature relation, etc. are discussed. In Section 4, the complete lists of O$`(q^4)`$ terms is given. In Section 5 we discuss the derivation of the external field-dependent on-shell O$`(q^4`$) Lagrangian, again within HBChPT using the techniques of . Section 6 has the conclusion and discussion on comparison of the results of this paper with those of a recent paper by Fettes et al , in which the O$`(q^4)`$ list in the presence of external fields is derived, but starting from the relativistic O$`(q^4`$) BChPT Lagrangian.
## 2 The Basics
The HBChPT Lagrangian is written in terms of the “upper component” H (and its hermitian adjoint $`\overline{\mathrm{H}}`$), exponentially parameterized matrix-valued meson fields $`U,u\sqrt{U}`$, baryon (“$`v_\mu ,\mathrm{S}_\nu `$”) and pion-field-dependent (“$`\mathrm{D}_\mu ,u_\nu ,\chi _\pm `$, $`F_{\mu \nu }^\pm ,v_{\mu \nu }^{(s)}`$”) building blocks defined below:
$$\mathrm{H}e^{i\mathrm{m}vx}\frac{1}{2}(1+\text{/}v)\psi ,$$
(1)
where $`\psi `$ Dirac spinor and $`m`$ the nucleon mass;
$`v_\mu \mathrm{nucleon}\mathrm{veclocity},`$
$`\mathrm{S}_\nu {\displaystyle \frac{i}{2}}\gamma ^5\sigma _{\nu \rho }v^\rho \mathrm{Pauli}\mathrm{Lubanski}\mathrm{spin}\mathrm{operator};`$ (2)
$$U=\mathrm{exp}\left(i\frac{\varphi }{F_\pi }\right),\mathrm{where}\varphi \stackrel{}{\pi }\stackrel{}{\tau },$$
(3)
where $`\stackrel{}{\tau }`$ nucleon isospin generators; $`\mathrm{D}_\mu =_\mu +\mathrm{\Gamma }_\mu iv_\mu ^{(s)}`$ where $`\mathrm{\Gamma }_\mu \frac{1}{2}((u^{}(_\mu ir_\mu )u+u(_\mu il_\mu )u^{})`$ ($`v_\mu ^{(s)}`$ is the isosinglet vector field needed to generate the electromagnetic current (See )); $`u_\mu i\left(u^{}(_\mu ir_\mu )uu(_\mu il_\mu )u^{}\right)`$; $`\chi _\pm u^{}\chi u^{}\pm u\chi ^{}u`$, where $`\chi 2B(s+ip)`$,$`s(`$ quark mass matrix) and $`p`$ being the external scalar and pseudo-scalar fields; $`F_{\mu \nu }^\pm u^{}F_{\mu \nu }^Ru\pm uF_{\mu \nu }^Lu^{}`$ where $`F_{\mu \nu }^R_{[\mu }r_{\nu ]}i.[r_\mu ,u_\nu ]`$, and $`F_{\mu \nu }^L_{[\mu }l_{\nu ]}i[l_\mu ,l_\nu ]`$ and $`v_{\mu \nu }^{(s)}_{[\mu }v_{\nu ]}^{(s)}`$ in which $`r_\mu V_\mu +A_\mu ,l_\mu V_\mu A_\mu `$, where $`V_\mu ,A_\mu `$ are external vector and axial-vector fields.
Terms of the $`_{(\mathrm{H})\mathrm{BChPT}}`$ constructed from products of building blocks will automatically be chiral invariant. Symbolically, a term in $`_{\mathrm{HBChPT}}`$ can be written as just a product of the building blocks to various powers (omitting $`\mathrm{H},\overline{\mathrm{H}}`$ as will be done in the rest of the paper except for Section 5):
$$\mathrm{D}_\alpha ^mu_\beta ^n\chi _+^p\chi _{}^qv_\sigma ^l(v_{\sigma \omega }^{(s)})^k(F_{\rho \lambda }^+)^t(F_{\mu \nu }^{})^u\mathrm{S}_\kappa ^\mathrm{r}(\mathrm{m},\mathrm{n},\mathrm{p},\mathrm{q},\mathrm{t},\mathrm{u},\mathrm{k})\mathrm{O}(\mathrm{q}^{\mathrm{m}+\mathrm{n}+2\mathrm{p}+2\mathrm{q}+2\mathrm{k}+2\mathrm{t}+2\mathrm{u}}).$$
(4)
A systematic path integral derivation for $`_{\mathrm{HBChPT}}`$ based on a paper by Mannel et al , starting from $`_{\mathrm{BChPT}}`$ was first given by Bernard et al . As shown by them, after integrating out h from the generating functional, one arrives at $`_{\mathrm{HBChPT}}`$ :
$$_{\mathrm{HBChPT}}=\overline{\mathrm{H}}\left(𝒜+\gamma ^0^{}\gamma ^0𝒞^1\right)\mathrm{H},$$
(5)
an expression in the upper components only i.e. for non-relativistic nucleons. So the terms of $`_{\mathrm{HBChPT}}`$ in this paper are given as operators on the H-spinors. For off-shell nucleons, $`\gamma ^0^{}\gamma ^0𝒞^1𝒜`$, and hence, listing $`𝒜`$-type terms will suffice.
The phase rule derived in can be modified to include external fields. After doing so, one gets: HBChPT terms (that are Lorentz scalar - isoscalars of even parity) made hermitian using a prescription for constructing hermitian (anti-)commutators discussed in , consisting of $`q\chi _{}`$’s, $`P[,]`$’s, $`j`$ (which can take only the values 0 or 1) $`ϵ^{\mu \nu \rho \lambda }`$’s, $`kv_{\mu \nu }^{(s)}`$’s, $`tF_{\rho \lambda }^+`$’s and $`uF_{\mu \nu }^{}`$’s, for which the following phase rule is satisfied, are the only terms allowed:
$$(1)^{q+P+j+k+t+u}=1.$$
(6)
In for $`k=t=u=0`$, (6) was used to generate complete lists up to O$`(q^3)`$ in the absence of external vector and axial-vector fields. In this paper, the same phase rule is used to construct complete lists of O$`(q^4)`$ including external fields.
Let $`A,B,C,D`$ be operators chosen from the pion-field dependent building blocks of (4). In what follows, and especially in Section 3, use will be made of a notation of : $`(A,B)[A,B]`$ or $`[A,B]_+`$. One can then show that apart from the (0,0,0,2,0,0,0)-, (0,0,2,0,0,0,0)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type terms (using the notation of (4)), the following is the complete list of O$`(q^4,\varphi ^{2n})`$ terms (using (6)):
$`(i)`$ $`(A,(B,(C,D)))(a)[A,[B,[C,D]_+]];(b)[A,[B,[C,D]]_+];`$ (7)
$`(c)[A,[B,[C,D]]]_+;(d)[A,[B,[C,D]_+]_+]_+;`$
$`(ii)`$ $`((A,B),(C,D))(a)[[A,B],[C,D]_+];(b)[[A,B]_+,[C,D]];`$
$`(c)[[A,B],[C,D]]_+;(d)[[A,B]_+,[C,D]_+]_+;`$
$`(iii)`$ $`i(A,(B,(C,D)))(a)i[A,[B,[C,D]_+]_+];(b)i[A,[B,[C,D]_+]]_+;`$
$`(c)i[A,[B,[C,D]]_+]_+;(d)i[A,[B,[C,D]]];`$
$`(iv)`$ $`i((A,B),(C,D))(a)i[[A,B],[C,D]];(b)i[[A,B]_+,[C,D]_+];`$
$`(c)i[[A,B],[C,D]_+]_+;(d)i[[A,B]_+,[C,D]]_+;`$
$`(v)`$ $`i(A,(B,C))(a)i[A,[B,C]]_+;(b)i[A,[B,C]_+];(c)AB;`$
$`(d)i[[A,B],C]_+;(e)i[[A,B]_+,C];`$
$`(vi)`$ $`(A,(B,C))(a)[A,[B,C]];(b)[A,[B,C]_+]_+;(c)AB`$
$`(d)[[A,B],C];(e)[[A,B]_+,C]_+,`$
where it is understood that of all the possible terms implied by $`(i)(A,(B,(C,D)))`$, $`(i)((A,B),(C,D))`$ and $`(i)(A,(B,C))`$, only those that are allowed by (6) are to be included. For (0,0,0,2,0,0,0)-, (0,0,2,0,0,0,0)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type terms, one needs to
include <sup>1</sup><sup>1</sup>1One need not consider (0,0,0,0,0,2,0). See (10) and the discussion thereafter.
$`(A,A)(\chi _+,\chi _+);(\chi _{},\chi _{});(F_{\mu \nu }^+,F_{\rho \lambda }^+);(v_{\mu \nu }^{(s)},v_{\rho \lambda }^{(s)});`$
$`(A,B)[F_{\mu \nu }^+,v_{\rho \lambda }^{(s)}]_+.`$ (8)
The list (7) holds good for O$`(q^4,\varphi ^{2n+1})`$ terms with the difference that there is an additional factor of $`i`$ multiplying the terms in $`(i),(ii)`$ and $`(vi)`$, and the $`i`$ in $`(iii),(iv)`$ and $`(v)`$, is absent. The reason for including $`i`$ only in some combinations of terms has to do with imposing charge conjugation invariance along with other symmetries directly within HBChPT (See ). The terms of (7) and their analogs for O$`(q^4,\varphi ^{2n+1})`$ are not all independent since they can be related by a number of linear relations: see next section (and for O($`q^3)`$).
## 3 Further Reduction due to Algebraic identities
In this section, we discuss further reduction in addition to the ones obtained from (6). The main result from is that one need not consider trace-dependent terms in SU(2) HBChPT if one assumes isospin conservation. Given that isospin violation enters only via $`\chi _\pm `$, thus, $`\chi _\pm `$-independent trace-dependent O$`(q^4)`$ terms can be eliminated in preference for $`\chi _\pm `$-independent trace-independent terms. We discuss reduction due to algebraic identities in the various categories of (7). Some of the algebraic reductions require one to consider more than one category at a time, e.g., for O$`(q^4,\varphi ^{2n})`$ terms, the generalized Jacobi identities in (3.1) require one to consider $`(i),(i)(AB),(ii)`$.
One can show that ():
$`[\mathrm{D}_\mu ,\mathrm{D}_\nu ]={\displaystyle \frac{1}{4}}[u_\mu ,u_\nu ]{\displaystyle \frac{i}{2}}F_{\mu \nu }^+iv_{\mu \nu }^{(s)},`$ (9)
$`[\mathrm{D}_\mu ,u_\nu ][\mathrm{D}_\nu ,u_\mu ]=F_{\mu \nu }^{}.`$ (10)
The first relation, referred to as the curvature relation, will be used extensively in conjunction with some generalized Jacobi identities and Schouten’s identity discussed below. As a consequence of (9), we will choose to always write $`[\mathrm{D}_\mu ,\mathrm{D}_\nu ]`$ in terms of $`[u_\mu ,u_\nu ],F_{\mu \nu }^+,v_{\mu \nu }^{(s)}`$. As a consequence of (10), we see that $`F_{\mu \nu }^{}`$ can be eliminated in preference for $`[\mathrm{D}_{[\mu },u_{\nu ]}]`$. However, the relative coefficients of $`[u_\mu ,u_\nu ],F_{\mu \nu }^+,v_{\mu \nu }^{(s)}`$, as well as the relative coefficients of $`[\mathrm{D}_\mu ,u_\nu ],[\mathrm{D}_\nu ,u_\mu ]`$ can be made arbitrary because each can be obtained from the nonrelativistic reduction of linearly independent terms. One interesting consequence of (10) is that in the absence of external fields, the commutator of the covariant derivative and the axial-vector building block is symmetric in the Lorentz indeces - something missed in .
Further, another source of major reduction in number of terms is the Schouten’s identity:
$$ϵ^{\mu _1\mu _2\mu _3\mu _4}X^{\mu _5}+\mathrm{cyclic}=0,$$
(11)
where $`\mathrm{X}_\mu `$ is any arbitrary (axial)vector.
It is because of the curvature relation that one requires to consider, e.g., some (4,0,0,0,0,0,0)-, (0,4,0,0,0,0,0)- (2,2,0,0,0,0,0)-, (2,0,0,0,1,0,0)-, (2,0,0,0,0,0,1)-, (0,2,0,0,1,0,0)-, (0,2,0,0,0,0,1)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type terms together in (3.1). Due to Schouten’s identity, e.g., the aforementioned 7-tuples with $`ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda `$ and $`ϵ^{\mu \nu \rho \lambda }\mathrm{S}_\lambda v`$ (where $`v_\rho `$ is contracted with a building block) are required to be considered together.
Finally, the O($`q^2`$) pion eom will be used for reduction in the number of linearly independent terms:
$$[\mathrm{D}_\mu ,u^\mu ]=\frac{i}{2}(\chi _{}\frac{1}{2}\chi _{})\frac{i}{2}\stackrel{~}{\chi }_{}.$$
(12)
### 3.1 O($`q^4,\varphi ^{2n})`$ Terms
In this subsection, we consider reduction in the number of independent O$`(q^4,\varphi ^{2n})`$ terms due to various algebraic identities. The following are the algebraic identities responsible for reduction in number of O$`(q^4,\varphi ^{2n})`$ terms: (13), (14), (3.1), (3.1), (3.1) (3.1) and (3.1). For (3.1) and (3.1), there are two sets each of terms (one $`ϵ^{\mu \nu \rho \lambda }`$-dependent and the other $`ϵ^{\mu \nu \rho \lambda }`$-independent), that need to be considered.
$`p=q=0`$ in (4): $`(A,(B,(C,D)));((A,B),(C,D))`$
This includes $`(i)(iv)`$ of (7). All terms in each of the first four types (of terms) in (7) \[$`(i)(iv)`$\] are linearly independent for unequal field operators A, B, C, D. However for (4,0,0,0,0,0,0), (0,4,0,0,0,0,0) and (2,2,0,0,0,0,0), L.C.-independent terms, one needs to consider A=C, B=D in (i) in equation (7). Using
$$[A,[B,[A,B]_+]]=[A,[B,[A,B]]_+]$$
(13)
only three of the four terms in (i) of equation (7), are linearly independent. Similarly, using
$$[[A,B],[A,B]_+]=[[A,B]_+,[A,B]],$$
(14)
only three of the four terms in (ii) of equation (7), are linearly independent.
There are some reductions possible due to some generalized Jacobi identities by considering : $`(i),(i)(AB),(ii)`$ of (7)($`ϵ^{\mu \nu \rho \lambda }`$ -independent terms), and $`(iii),(iii)(AB),(iv)`$ of (7) ($`ϵ^{\mu \nu \rho \lambda }`$-dependent terms). The reason why one can not hope to get reductions by considering any other pairs of types of terms in $`(i)(iv)`$ (in (7)), is because one can get (linear) algebraic relationships only between those terms which are (both) independent of (have) an overall factor of $`i`$.
$`(i),(i)(AB),(ii)`$ of (7)
One can show the following 6 generalized Jacobi identities:
$`[A,[B,[C,D]_+]][[A,B],[C,D]_+]=(i)(a)(AB)`$
$`[A,[B,[C,D]_+]][[A,B]_+,[C,D]_+]_+=(i)(d)(AB)`$
$`[A,[B,[C,D]]_+][[A,B]_+,[C,D]]=(i)(b)(AB)`$
$`[A,[B,[C,D]]_+][[A,B],[C,D]]_+=(i)(c)(AB)`$
$`[A,[B,[C,D]]]_+[[A,B]_+,[C,D]]=(i)(c)(AB)`$
$`[A,[B,[C,D]_+]_+]_+[[A,B],[C,D]_+]=(i)(d)(AB).`$
Further, one can apply the following to $`(B,(C,D))`$ contained in $`(A,(B,(C,D)))`$:
$`[B,[C,D]][[B,C],D]=[C,[B,D]]`$
$`[B,[C,D]_+]_+[[B,C]_+,D]_+=[C,[B,D]]`$
$`[B,[C,D]][[B,C]_+,D]_+=[C,[B,D]_+]_+,`$ (16)
and:
$`i[B,[C,D]]_+i[[B,C],D]_+=i[C,[B,D]_+]`$
$`i[B,[C,D]]_+i[[B,C]_+,D]=i[C,[B,D]]_+`$
$`i[B,[C,D]_+]i[[B,C]_+,D]=i[C,[B,D]_+].`$ (17)
The three identities in (3.1) are similar to the ones that occur in SUSY graded Lie algebra for $`B,D`$ fermionic and $`C`$ bosonic fields, $`B,C`$ fermionic and $`D`$ bosonic fields, and $`B,C,D`$ fermionic fields, respectively.
(1) Using (3.1), (9), (3.1) and (3.1), one needs to consider the following (4,0,0,0,0,0,0)-, (0,4,0,0,0,0,0)- (2,2,0,0,0,0,0)-, (2,0,0,0,1,0,0)-, (2,0,0,0,0,0,1)-, (0,2,0,0,1,0,0)-, (0,2,0,0,0,0,1)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type $`ϵ^{\mu \nu \rho \lambda }`$-independent terms together:
$`(v\mathrm{D},(\mathrm{D}_\mu ,(v\mathrm{D},\mathrm{D}^\mu ))),(\mathrm{D}_\mu ,(v\mathrm{D},(v\mathrm{D},\mathrm{D}^\mu ))),((v\mathrm{D},\mathrm{D}_\mu ),(v\mathrm{D},\mathrm{D}^\mu ))`$
$`(vu,(u_\mu ,(vu,u^\mu ))),(u_\mu ,(vu,(vu,u^\mu ))),((vu,u_\mu ),(vu,u^\mu )),`$
$`(v\mathrm{D},(\mathrm{D}_\mu ,(vu,u^\mu ))),(\mathrm{D}_\mu ,(v\mathrm{D},(u_\mu ,vu))),((v\mathrm{D},\mathrm{D}_\mu ),(vu,u^\mu ))),`$
$`(vu,(u_\mu ,(v\mathrm{D},\mathrm{D}^\mu ))),(u_\mu ,(vu,(v\mathrm{D},\mathrm{D}^\mu ))),`$
$`(v\mathrm{D},((\mathrm{D}_\mu ,vu),u^\mu )),(u^\mu ,((\mathrm{D}_\mu ,vu),v\mathrm{D}),((v\mathrm{D},u^\mu ),(\mathrm{D}_\mu ,vu)),`$
$`(\mathrm{D}_\mu ,((v\mathrm{D},u^\mu ),vu)),(vu,((v\mathrm{D},u^\mu ),\mathrm{D}_\mu )),`$
$`v^\nu (v\mathrm{D},(\mathrm{D}_\mu ,F^{+\mu \nu })),v^\nu (\mathrm{D}_\mu ,(v\mathrm{D},F^{+\mu \nu })),v^\nu ((v\mathrm{D},\mathrm{D}_\mu ),F^{+\mu \nu }),`$
$`v^\nu (v\mathrm{D},(\mathrm{D}_\mu ,v^{(s)\mu \nu })),v^\nu (\mathrm{D}_\mu ,(v\mathrm{D},v^{(s)\mu \nu })),v^\nu ((v\mathrm{D},\mathrm{D}_\mu ),v^{(s)\mu \nu }),`$
$`v^\nu (vu,(u_\mu ,F^{+\mu \nu })),v^\nu (u_\mu ,(vu,F^{+\mu \nu })),v^\nu ((vu,u_\mu ),F^{\mu \nu }),`$
$`v^\nu (vu,(u_\mu ,v^{(s)\mu \nu })),v^\nu (u_\mu ,(vu,v^{(s)\mu \nu })),v^\nu ((vu,u_\mu ),v^{(s)\mu \nu }),`$
$`v^\kappa F_{\kappa \mu }^+v_\sigma F^{+\sigma \mu },v^\kappa v_{\kappa \mu }^{(s)}v_\sigma F^{+\sigma \mu },v^\kappa v_{\kappa \mu }^{(s)}v_\sigma v^{(s)\sigma \mu }.`$
One needs to do a careful counting of the total number of identities that one can write down using (3.1) and (9), and the total number of terms in those identities.
(2) Similarly, one will need to consider (4,0,0,0,0,0,0)-, (0,4,0,0,0,0,0)- (2,2,0,0,0,0,0)-, (2,0,0,0,1,0,0)-, (2,0,0,0,0,0,1)-, (0,2,0,0,1,0,0)-, (0,2,0,0,0,0,1)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type terms together:
$`(\mathrm{D}_\nu ,(\mathrm{D}_\mu ,(\mathrm{D}^\nu ,\mathrm{D}^\mu ))),((\mathrm{D}_\nu ,\mathrm{D}_\mu ),(\mathrm{D}^\nu ,\mathrm{D}^\mu ))`$
$`(u_\nu ,(u_\mu ,(u^\nu ,u^\mu ))),((u_\nu ,u_\mu ),(u^\nu ,u^\mu )),`$
$`(\mathrm{D}_\nu ,(\mathrm{D}_\mu ,(u^\nu ,u^\mu ))),((\mathrm{D}_\nu ,\mathrm{D}_\mu ),(u^\nu ,u^\mu )))`$
$`(u_\nu ,(u_\mu ,(\mathrm{D}^\nu ,\mathrm{D}^\mu )))`$
$`(\mathrm{D}^\mu ,(\mathrm{D}^\nu ,F_{\mu \nu }^+)),([\mathrm{D}^\mu ,\mathrm{D}^\nu ],F_{\mu \nu }^+)),`$
$`(\mathrm{D}^\mu ,(\mathrm{D}^\nu ,v_{\mu \nu }^{(s)})),([\mathrm{D}^\mu ,\mathrm{D}^\nu ],v_{\mu \nu }^{(s)})),`$
$`(u^\mu ,(u^\nu ,F_{\mu \nu }^+)),([u^\mu ,u^\nu ],F_{\mu \nu }^+)),`$
$`(u^\mu ,(u^\nu ,v_{\mu \nu }^{(s)})),([u^\mu ,u^\nu ],v_{\mu \nu }^{(s)}))`$
$`(F_{\mu \nu }^+)^2,v_{\mu \nu }^{(s)}F^{+\mu \nu },(v_{\mu \nu }^{(s)})^2`$ (19)
$`(iii),(iii)(AB)`$ and $`(iv)`$; (vi) of (7)
One can show the following generalized Jacobi identities to be true:
$`i[A,[B,[C,D]_+]_+]i[[A,B]_+,[C,D]_+]=(iii)(a)(AB)`$
$`i[A,[B,[C,D]_+]_+]i[[A,B],[C,D]_+]_+=(iii)(b)(AB)`$
$`i[A,[B,[C,D]_+]]_+i[[A,B]_+,[C,D]_+]=(iii)(b)(AB)`$
$`i[A,[B,[C,D]]_+]_+i[[A,B],[C,D]]=(iii)(c)(AB)`$
$`i[A,[B,[C,D]]_+]_+i[[A,B]_+,[C,D]_+]=(iii)(d)(AB)`$
$`i[A,[B,[C,D]]]i[[A,B],[C,D]]=(iii)(d)(AB)`$
Again one can apply the generalized Jacobi-iieuntities (3.1) and (3.1) to $`(B,(C,D))`$ contained in $`(A,(B,(C,D)))`$.
(3) The identities (3.1), (3.1) and (3.1) along with (9) require one to consider the following category of $`ϵ^{\mu \nu \rho \lambda }`$-dependent (4,0,0,0,0,0,0)-, (0,4,0,0,0,0,0)- (2,2,0,0,0,0,0)-, (2,0,0,0,1,0,0)-, (2,0,0,0,0,0,1)-, (0,2,0,0,1,0,0)-, (0,2,0,0,0,0,1)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- and (0,0,0,0,1,0,1)-type terms together:
$`iϵ^{\mu \nu \rho \lambda }v_\rho [(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,(\mathrm{D}_\lambda ,\mathrm{S}\mathrm{D}))),(u_\mu ,(u_\nu ,(u_\lambda ,\mathrm{S}u))),`$
$`(\mathrm{S}\mathrm{D},(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,[\mathrm{D}_\nu ,\mathrm{D}_\lambda ]))(\mathrm{D}_\mu ,(\mathrm{S}\mathrm{D},[\mathrm{D}_\nu ,\mathrm{D}_\lambda ])),`$
$`(\mathrm{S}u,(u_\mu ,([u_\nu ,u_\lambda ])),(u_\mu ,(\mathrm{S}u,[u_\nu ,u_\lambda ])),`$
$`([u_\mu ,u_\nu ],(\mathrm{D}_\lambda ,\mathrm{S}\mathrm{D})),(u_\mu ,(u_\nu ,(\mathrm{D}_\lambda ,\mathrm{S}\mathrm{D}))),`$
$`(\mathrm{D}_\mu ,(\mathrm{S}\mathrm{D},[u_\nu ,u_\lambda ])),(\mathrm{S}\mathrm{D},(\mathrm{D}_\mu ,[u_\nu ,u_\lambda ])),`$
$`(u_\mu ,(\mathrm{S}u,[\mathrm{D}_\nu ,\mathrm{D}_\lambda ])),(\mathrm{S}u,(u_\mu ,[\mathrm{D}_\nu ,\mathrm{D}_\lambda ])),`$
$`(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,(u_\lambda ,\mathrm{S}u))),([\mathrm{D}_\mu ,\mathrm{D}_\nu ],(u_\lambda ,\mathrm{S}u))`$
$`(\mathrm{D}_\mu ,((\mathrm{D}_\nu ,u_\lambda ),\mathrm{S}u)),(\mathrm{S}u,((\mathrm{D}_\nu ,u_\lambda ),\mathrm{D}_\mu ));`$
$`((\mathrm{D}_\mu ,\mathrm{S}u),(\mathrm{D}_\nu ,u_\lambda ));`$
$`(\mathrm{D}_\mu ,((\mathrm{D}_\nu ,\mathrm{S}u),u_\lambda )),(u_\lambda ,((\mathrm{D}_\nu ,\mathrm{S}u),\mathrm{D}_\mu )),`$
$`i(F_{\mu \nu }^+,(\mathrm{D}_\lambda ,\mathrm{S}\mathrm{D})),i(\mathrm{D}_\lambda ,(\mathrm{S}\mathrm{D},F_{\mu \nu }^+)),i(\mathrm{S}\mathrm{D},(\mathrm{D}_\lambda ,F_{\mu \nu }^+)),`$
$`i(v_{\mu \nu }^{(s)},(\mathrm{D}_\lambda ,\mathrm{S}\mathrm{D})),i(\mathrm{D}_\lambda ,(\mathrm{S}\mathrm{D},v_{\mu \nu }^{(s)})),i(\mathrm{S}\mathrm{D},(\mathrm{D}_\lambda ,v_{\mu \nu }^{(s)})),`$
$`i(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,F_{\kappa \lambda }^+))S^\kappa ,i([\mathrm{D}_\mu ,\mathrm{D}_\nu ],F_{\kappa \lambda }^+)S^\kappa ,i(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,v_{\kappa \lambda }^{(s)}))S^\kappa ,`$
$`i(u_\mu ,(u_\nu ,F_{\kappa \lambda }^+))S^\kappa ,i([u_\mu ,u_\nu ],F_{\kappa \lambda }^+)S^\kappa ,i(u_\mu ,(u_\nu ,v_{\kappa \lambda }^{(s)}))S^\kappa ,`$
$`i(F_{\mu \nu }^+,(u_\lambda ,\mathrm{S}u)),i(u_\lambda ,(\mathrm{S}u,F_{\mu \nu }^+)),i(\mathrm{S}u,(u_\lambda ,F_{\mu \nu }^+)),`$
$`i[u_\lambda ,\mathrm{S}u]_+v_{\mu \nu }^{(s)},[F_{\mu \nu }^+,F_{\kappa \lambda }^+]S^\kappa ].`$ (21)
Analogous to (3.1) and (3.1), one needs to do a careful counting of the total number of identities that one can write down using (3.1), and (9), and the total number of terms in those identities.
A similar analysis can be carried out for terms with $`v\mathrm{S}`$ in (3.1). However, due to (11), this set of terms has to be considered in conjunction with (3.1) $`(u_\kappa ,\mathrm{D}^\kappa vu,v\mathrm{D})`$
The identities in (3.1) are also used in, e.g., $`ϵ^{\mu \nu \rho \lambda }`$-dependent (1,1,0,1,0,0,0)-type terms.
(4) The identities (3.1), (3.1), (3.1) along with the curvature relation, require one to consider the following category of $`ϵ^{\mu \nu \rho \lambda }`$\- dependent (4,0,0,0,0,0,0)-, (0,4,0,0,0,0,0)- (2,2,0,0,0,0,0)-, (2,0,0,0,1,0,0)-, (2,0,0,0,0,0,1)-, (0,2,0,0,1,0,0)-, (0,2,0,0,0,0,1)-, (0,0,0,0,2,0,0)-, (0,0,0,0,0,0,2)- type terms together:
$`iϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda ((u_\kappa ,(u_\mu ,(u^\kappa ,u_\nu ))),(u_\mu ,(u_\kappa ,(u^\kappa ,u_\nu ))),((u_\mu ,u_\kappa ),(u_\nu ,u^\kappa )),`$
$`(\mathrm{D}_\kappa ,(\mathrm{D}_\mu ,(\mathrm{D}^\kappa ,\mathrm{D}_\nu ))),((\mathrm{D}_\mu ,(\mathrm{D}_\kappa ,(\mathrm{D}^\kappa ,\mathrm{D}_\nu ))),((\mathrm{D}_\kappa ,\mathrm{D}_\mu ),(\mathrm{D}^\kappa ,\mathrm{D}_\nu )),`$
$`((\mathrm{D}_\kappa ,(\mathrm{D}_\mu ,(u^\kappa ,u_\nu ))),((\mathrm{D}_\mu ,(\mathrm{D}_\kappa ,(u^\kappa ,u_\nu ))),((\mathrm{D}_\kappa ,\mathrm{D}_\mu ),(u^\kappa ,u_\nu )),`$
$`(u_\mu ,(u_\kappa ,(\mathrm{D}^\kappa ,\mathrm{D}_\nu ))),(u_\mu ,(u_\kappa ,(\mathrm{D}^\kappa ,\mathrm{D}_\nu ))),`$
$`(\mathrm{D}_\mu ,((\mathrm{D}_\nu ,u_\kappa ),u^\kappa )),((\mathrm{D}_\mu ,u^\kappa ),(\mathrm{D}_\nu ,u^\kappa )),`$
$`(u^\kappa ,((\mathrm{D}_\mu ,u^\kappa ),\mathrm{D}_\nu )),`$
$`i(\mathrm{D}^\kappa ,(\mathrm{D}_\mu ,F_{\kappa \nu }^+)),i(\mathrm{D}_\mu ,(\mathrm{D}^\kappa ,F_{\kappa \nu }^+)),i((\mathrm{D}^\kappa ,\mathrm{D}_\mu ),F_{\kappa \nu }^+)`$
$`i(\mathrm{D}^\kappa ,(\mathrm{D}_\mu ,v_{\kappa \nu }^{(s)})),i(\mathrm{D}_\mu ,(\mathrm{D}^\kappa ,v_{\kappa \nu }^{(s)})),i((\mathrm{D}^\kappa ,\mathrm{D}_\mu ),v_{\kappa \nu }^{(s)})`$
$`i(u^\kappa ,(u_\mu ,F_{\kappa \nu }^+)),i(u_\mu ,(u^\kappa ,F_{\kappa \nu }^+)),i((u^\kappa ,u_\mu ),F_{\kappa \nu }^+)`$
$`i(u^\kappa ,(u_\mu ,v_{\kappa \nu }^{(s)})),i(u_\mu ,(u^\kappa ,v_{\kappa \nu }^{(s)})),i((u^\kappa ,u_\mu ),v_{\kappa \nu }^{(s)})`$
$`i[F_\mu ^{+\kappa },F_{\kappa \nu }^+].).`$ (22)
A similar analysis can be carried out for terms with $`(u_\kappa ,\mathrm{D}^\kappa )(vu,v\mathrm{D})`$ in (3.1).
For (0,4,0,0,0,0,0)- and (1,3,0,0,0,0,0)-type terms, writing $`u_\mu =u_\mu ^a\tau ^a,[\mathrm{D}_\mu ,u_\nu ]=[\mathrm{D}_\mu ,u_\nu ]^a\tau ^a`$ (where $`[\mathrm{D}_\mu ,u_\nu ]^a_\mu u_\nu ^a+iϵ^{abc}\mathrm{\Gamma }_\mu ^bu_\nu ^c`$), one will need to consider the following relations:
$`(a)`$ $`[\tau ^b,[\tau ^c,\tau ^d]_+]=[\tau ^a,[\tau ^b,[\tau ^c,\tau ^d]]_+]=0;`$
$`[[\tau ^a,\tau ^b],[\tau ^c,\tau ^d]_+]=0;`$
$`i[[\tau ^a,\tau ^b]_+,[\tau ^c,\tau ^d]_+]=0;`$
$`(b)`$ $`[u_\mu ,u_\nu ]_+^2[u_\mu ,u_\nu ]^2=4u^4;`$
$`(c)`$ $`vuu_\mu vuu^\mu +h.c.=2u^2(vu)^2+[vu,u_\mu ]_+^2;`$
$`(d)`$ $`iϵ^{\mu \nu \rho \lambda }[[u_\mu ,u_\kappa ],u_\nu ,u^\kappa ]_+]_+=iϵ^{\mu \nu \rho \lambda }(u^2[u_\mu ,u_\nu ]u^\kappa [u_\mu ,u_\nu ]u_\kappa );`$
$`(e)`$ $`iϵ^{\mu \nu \rho \lambda }[[u_\mu ,vu],[u_\nu ,vu]_+]_+=iϵ^{\mu \nu \rho \lambda }\left((vu)^2[u_\mu ,u_\nu ]v[u_\mu ,u_\nu ]vu\right).`$
At least one of $`k,p,q,t,u`$ is $`0`$ in (4):$`(A,(B,C))`$
This includes $`(v)`$ and $`(vi)`$ of (7).
($`v`$) of (7)
The identities in (3.1) are used in, e.g., $`ϵ^{\mu \nu \rho \lambda }`$-independent (1,1,0,1,0,0,0)-type terms. When applying (3.1) to (2,0,1,0,0,0,0), because of (9), one will need to consider the following terms together:
$`iϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda ((\mathrm{D}_\mu ,(\mathrm{D}_\nu ,\chi _+)),([\mathrm{D}_\mu ,\mathrm{D}_\nu ],\chi _+);`$
$`(u_\mu ,(u_\nu ,\chi _+)),([u_\mu ,u_\nu ],\chi _+),i(F_{\mu \nu }^+,\chi _+),iv_{\mu \nu }^{(s)},\chi _+)).`$ (24)
### 3.2 O$`(q^4,\varphi ^{2n+1}`$) Terms
In this subsection, we consider the reduction in the number of O$`(q^4,\varphi ^{2n+1})`$ terms because of algebraic identities. The discussion in this subsection will be much briefer than the preceding (subsection).
$`(i)(iv)`$ of (7)
(1) The identities (3.1) are the same for O$`(q^4,\varphi ^{2n+1})`$ except for an overall factor of $`i`$. We will denote the analogue of (3.1) for O$`(q^4,\varphi ^{2n+1})`$ terms as (3.1).<sup>2</sup><sup>2</sup>2Similarly, the analogs of (7), (3.1), (3.1) and (3.1) will be denoted by (7), (3.1), (3.1) and (3.1). Using it together with and (9), one sees that one needs to consider the following set of terms together:
$`(a)`$ $`i(v\mathrm{D},(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,\mathrm{S}u))),i(\mathrm{D}_\mu ,(v\mathrm{D},(\mathrm{D}^\mu ,\mathrm{S}u))),i((v\mathrm{D},\mathrm{D}_\mu ),(\mathrm{D}^\mu ,\mathrm{S}u)),`$ (25)
$`i(\mathrm{S}u,(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,v\mathrm{D}))),i(\mathrm{D}_\mu ,(\mathrm{S}u,(\mathrm{D}^\mu ,v\mathrm{D})));`$
$`(\mathrm{D}^\mu ,(\mathrm{S}u,F_{\mu \nu }^+))v^\nu ,(\mathrm{S}u,(\mathrm{D}^\mu ,F_{\mu \nu }^+))v^\nu ,((\mathrm{D}^\mu ,\mathrm{S}u),F_{\mu \nu }^+)v^\nu ,`$
$`(\mathrm{D}^\mu ,(\mathrm{S}u,v_{\mu \nu }^{(s)}))v^\nu ,(\mathrm{S}u,(\mathrm{D}^\mu ,v_{\mu \nu }^{(s)}))v^\nu ,((\mathrm{D}^\mu ,\mathrm{S}u),v_{\mu \nu }^{(s)})v^\nu ;`$
$`i(vu,(u_\mu ,(\mathrm{S}u,\mathrm{D}^\mu ))),i(u_\mu ,(vu,(\mathrm{S}u,\mathrm{D}^\mu ))),i((vu,u_\mu ),(\mathrm{S}u,\mathrm{D}^\mu )),`$
$`i(\mathrm{D}_\mu ,(\mathrm{S}u,(u^\mu ,vu)))i(\mathrm{S}u,(\mathrm{D}_\mu ,(u^\mu ,vu)));`$
$`i(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,(v\mathrm{D},\mathrm{S}u))),i(\mathrm{D}^2,(v\mathrm{D},\mathrm{S}u)),`$
$`i(\mathrm{S}u,(v\mathrm{D},\mathrm{D}^2)),i(v\mathrm{D},(\mathrm{S}u,\mathrm{D}^2)).`$
$`(b)`$ $`i(v\mathrm{D},(\mathrm{D}_\mu ,(\mathrm{S}\mathrm{D},u^\mu ))),i(\mathrm{D}_\mu ,(v\mathrm{D},(\mathrm{S}\mathrm{D},u^\mu ))),i((v\mathrm{D},\mathrm{D}_\mu ),(\mathrm{S}\mathrm{D},u^\mu )),`$
$`i(u_\mu ,(\mathrm{S}\mathrm{D},(\mathrm{D}^\mu ,v\mathrm{D})))i(\mathrm{S}\mathrm{D},(u_\mu ,(\mathrm{D}^\mu ,v\mathrm{D})));`$
$`(u^\mu ,(\mathrm{S}\mathrm{D},F_{\mu \nu }^+))v^\nu ,(\mathrm{S}\mathrm{D},(u^\mu ,F_{\mu \nu }^+))v^\nu ,((u^\mu ,\mathrm{S}\mathrm{D}),F_{\mu \nu }^+)v^\nu ,`$
$`(u^\mu ,(\mathrm{S}\mathrm{D},v_{\mu \nu }^{(s)}))v^\nu ,(\mathrm{S}\mathrm{D},(u^\mu ,v_{\mu \nu }^{(s)}))v^\nu ,(u^\mu ,\mathrm{S}\mathrm{D}),v_{\mu \nu }^{(s)})v^\nu `$
$`i(vu,(u_\mu ,(u^\mu ,\mathrm{S}\mathrm{D}))),i(u_\mu ,(vu,(u^\mu ,\mathrm{S}\mathrm{D}))),i((vu,u_\mu ),(u^\mu ,\mathrm{S}\mathrm{D})),`$
$`i(\mathrm{S}\mathrm{D},(u_\mu ,(u^\mu ,vu))),i(u_\mu ,(\mathrm{S}\mathrm{D},(u^\mu ,vu)))`$
$`i(u^\mu ,(u_\mu ,(vu,\mathrm{D}^\mu ))),i(\mathrm{S}u,(u_\mu ,vu,\mathrm{D}^\mu )),`$
$`i((u_\mu ,\mathrm{S}u),(vu,\mathrm{D}^\mu )),i(\mathrm{D}_\mu ,(vu,(u^\mu ,\mathrm{S}u))),`$
$`i(vu,(\mathrm{D}_\mu ,(u^\mu ,\mathrm{S}u)));`$
$`i(u^\mu ,(u_\mu ,(vu,\mathrm{S}\mathrm{D}))),i(u^2,(vu,\mathrm{S}\mathrm{D})),`$
$`i(\mathrm{S}\mathrm{D},(vu,u^2)),i(vu,(\mathrm{S}\mathrm{D},u^2));`$
$`(c)`$ $`i(v\mathrm{D},(\mathrm{S}\mathrm{D},(\mathrm{D}^\mu ,u_\mu ))),i(\mathrm{S}\mathrm{D},(v\mathrm{D},(\mathrm{D}^\mu ,u_\mu ))),i((v\mathrm{D},\mathrm{S}\mathrm{D}),(\mathrm{D}^\mu ,u_\mu )),`$
$`i(u_\mu ,(\mathrm{D}^\mu ,(v\mathrm{D},\mathrm{S}\mathrm{D}))),i(\mathrm{D}_\mu ,(u^\mu ,(v\mathrm{DS}\mathrm{D})));`$
$`v^{[\rho }\mathrm{S}^{\lambda ]}(\mathrm{D}_\mu ,(u^\mu ,F_{\rho \lambda }^+)),v^{[\rho }\mathrm{S}^{\lambda ]}(u^\mu ,(\mathrm{D}_\mu ,F_{\rho \lambda }^+)),v^{[\rho }\mathrm{S}^{\lambda ]}((\mathrm{D}_\mu ,u^\mu ),F_{\rho \lambda }^+)`$
$`v^{[\rho }\mathrm{S}^{\lambda ]}(\mathrm{D}_\mu ,(u^\mu ,v_{\rho \lambda }^{(s)})),v^{[\rho }\mathrm{S}^{\lambda ]}(u^\mu ,(\mathrm{D}_\mu ,v_{\rho \lambda }^{(s)})),v^{[\rho }\mathrm{S}^{\lambda ]}((\mathrm{D}_\mu ,u^\mu ),v_{\rho \lambda }^{(s)})`$
$`i(vu,(\mathrm{S}u,(u^\mu ,\mathrm{D}_\mu ))),i(\mathrm{S}u,(vu,(\mathrm{D}^\mu ,u_\mu ))),i((vu,\mathrm{S}u),(\mathrm{D}^\mu ,u_\mu )),`$
$`i(\mathrm{D}_\mu ,(u^\mu ,(vu,\mathrm{S}u))),i(u_\mu ,(\mathrm{D}^\mu ,(vu,\mathrm{S}u)));`$
$`(d)`$ $`i(\mathrm{D}_\mu ,(\mathrm{S}\mathrm{D},(v\mathrm{D},u^\mu ))),i(\mathrm{S}\mathrm{D},(\mathrm{D}_\mu ,(v\mathrm{D},u^\mu ))),i((\mathrm{S}\mathrm{D},\mathrm{D}_\mu ),(v\mathrm{D},u^\mu )),`$
$`i(u_\mu ,(v\mathrm{D},(\mathrm{D}^\mu ,\mathrm{S}\mathrm{D}))),i(v\mathrm{D},(u_\mu ,(\mathrm{D}^\mu ,\mathrm{S}\mathrm{D}))),`$
$`(u^\mu ,(v\mathrm{D},F_{\mu \nu }^+))\mathrm{S}^\nu ,(v\mathrm{D},(u^\mu ,F_{\mu \nu }^+))\mathrm{S}^\nu ,((u^\mu ,v\mathrm{D}),F_{\mu \nu }^+)\mathrm{S}^\nu ,`$
$`(u^\mu ,(v\mathrm{D},v_{\mu \nu }^{(s)}))\mathrm{S}^\nu ,(v\mathrm{D},(u^\mu ,v_{\mu \nu }^{(s)}))\mathrm{S}^\nu ,((u^\mu ,v\mathrm{D}),v_{\mu \nu }^{(s)})\mathrm{S}^\nu ,`$
$`i(u_\mu ,(\mathrm{S}u,(u^\mu ,v\mathrm{D}))),i(\mathrm{S}u,(u_\mu ,(u^\mu ,v\mathrm{D}))),i((u_\mu ,\mathrm{S}u),(u^\mu ,v\mathrm{D}))),`$
$`i(v\mathrm{D},(u_\mu ,(u^\mu ,\mathrm{S}u)),i(u_\mu ,(v\mathrm{D},(u^\mu ,\mathrm{S}u));`$
$`i(u_\mu ,(u^\mu ,(\mathrm{S}u,v\mathrm{D}))),i(u^2,(\mathrm{S}u,v\mathrm{D})),`$
$`i(v\mathrm{D},(\mathrm{S}u,u^2)),i(\mathrm{S}u,(v\mathrm{D},u^2));`$
$`(e)`$ $`i(\mathrm{D}_\mu ,(\mathrm{S}\mathrm{D},(\mathrm{D}^\mu ,vu))),i(\mathrm{S}\mathrm{D},(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,vu))),i((\mathrm{D}_\mu ,\mathrm{S}\mathrm{D}),(\mathrm{D}^\mu ,vu))),`$
$`i(vu,(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,\mathrm{S}\mathrm{D})),i(\mathrm{D}_\mu ,(vu,(\mathrm{D}^\mu ,\mathrm{S}\mathrm{D}));`$
$`(vu,(\mathrm{D}^\mu ,F_{\mu \nu }^+))\mathrm{S}^\nu ,(\mathrm{D}^\mu ,(vu,F_{\mu \nu }^+))\mathrm{S}^\nu ,((vu,\mathrm{D}^\mu ),F_{\mu \nu }^+)\mathrm{S}^\nu ,`$
$`(vu,(\mathrm{D}^\mu ,v_{\mu \nu }^{(s)}))\mathrm{S}^\nu ,(\mathrm{D}^\mu ,(vu,v_{\mu \nu }^{(s)}))\mathrm{S}^\nu ,((vu,\mathrm{D}^\mu ),v_{\mu \nu }^{(s)})\mathrm{S}^\nu ,`$
$`i(u_\mu ,(\mathrm{S}u,(vu,\mathrm{D}^\mu ))),i(\mathrm{S}u,(u_\mu ,(vu,\mathrm{D}^\mu ))),i((\mathrm{S}u,u_\mu ),(vu,\mathrm{D}^\mu )),`$
$`i(\mathrm{D}_\mu ,(vu,(u^\mu ,\mathrm{S}u))),i(vu,(\mathrm{D}_\mu ,(u^\mu ,\mathrm{S}u)));`$
$`i(\mathrm{D}_\mu ,(\mathrm{D}^\mu ,(\mathrm{S}\mathrm{D},vu))),i(\mathrm{D}^2,(\mathrm{S}\mathrm{D},vu)),`$
$`i(vu,(\mathrm{S}\mathrm{D},\mathrm{D}^2)),i(\mathrm{S}\mathrm{D},(vu,\mathrm{D}^2)).`$
Using (3.1) and (3.1), one needs to consider $`(a)(e)`$ together.
Similarly, using (3.1)$``$ and (9), one needs to consider simultaneously
$`(a)`$ $`i(v\mathrm{D},(\mathrm{S}\mathrm{D},(v\mathrm{D},vu))),i(\mathrm{S}\mathrm{D},(v\mathrm{D},(v\mathrm{D},vu))),i((v\mathrm{D},\mathrm{S}\mathrm{D}),(v\mathrm{D},vu))),`$ (26)
$`i(vu,(v\mathrm{D},(v\mathrm{D},\mathrm{S}\mathrm{D})),i(v\mathrm{D},(vu,(v\mathrm{D},\mathrm{S}\mathrm{D}));`$
$`v^{[\rho }\mathrm{S}^{\lambda ]}(v\mathrm{D},(vu,F_{\rho \lambda }^+)),v^{[\rho }\mathrm{S}^{\lambda ]}(vu,(v\mathrm{D},F_{\rho \lambda }^+)),v^{[\rho }\mathrm{S}^{\lambda ]}((v\mathrm{D},vu),F_{\rho \lambda }^+);`$
$`v^{[\rho }\mathrm{S}^{\lambda ]}(v\mathrm{D},(vu,v_{\rho \lambda }^{(s)})),v^{[\rho }\mathrm{S}^{\lambda ]}(vu,(v\mathrm{D},v_{\rho \lambda }^{(s)})),v^{[\rho }\mathrm{S}^{\lambda ]}((v\mathrm{D},vu),v_{\rho \lambda }^{(s)});`$
$`i(vu,(\mathrm{S}u,(vu,v\mathrm{D}))),i(\mathrm{S}u,(vu,(vu,v\mathrm{D}))),i((\mathrm{S}u,vu),(vu,v\mathrm{D})),`$
$`i(vu,(v\mathrm{D},(vu,\mathrm{S}u)),i(v\mathrm{D},(vu,(vu,\mathrm{S}u));`$
$`i(vu,(\mathrm{S}\mathrm{D},(v\mathrm{D})^2)),i(\mathrm{S}\mathrm{D},(vu,(v\mathrm{D})^2)),`$
$`i((v\mathrm{D})^2,(\mathrm{S}\mathrm{D},vu)),i(v\mathrm{D},(v\mathrm{D},(\mathrm{S}\mathrm{D},vu)));`$
$`i(v\mathrm{D},\mathrm{S}u,(vu)^2)),i(\mathrm{S}u,(v\mathrm{D},(vu)^2)),`$
$`i((vu)^2,(\mathrm{S}u,v\mathrm{D})),i(vu,(vu,(\mathrm{S}u,v\mathrm{D}));`$
$`(b)`$ $`i(v\mathrm{D},(v\mathrm{D},(v\mathrm{D},\mathrm{S}u))),i((v\mathrm{D})^2,(v\mathrm{D},\mathrm{S}u)),`$
$`i[\mathrm{S}u,(v\mathrm{D})^3]_+,i(v\mathrm{D},(\mathrm{S}u,(v\mathrm{D})^2)).`$
Using (3.1) and (3.1), one needs to consider $`(a)`$ and $`(b)`$ together.
(2) Using (3.1), (3.1), (3.1), (9) and (11), one sees that one has to consider the following sets of terms together:
$`(a)`$ $`ϵ^{\mu \nu \rho \lambda }((\mathrm{D}_\mu ,(\mathrm{D}_\nu ,(\mathrm{D}_\rho ,u_\lambda ))),([\mathrm{D}_\mu ,\mathrm{D}_\nu ],(\mathrm{D}_\rho ,u_\lambda ))),`$ (27)
$`(u_\mu ,(\mathrm{D}_\nu ,[\mathrm{D}_\rho ,\mathrm{D}_\lambda ])),(\mathrm{D}_\mu ,(u_\nu ,[\mathrm{D}_\rho ,\mathrm{D}_\lambda ]));`$
$`i(\mathrm{D}_\mu ,(u_\nu ,F_{\rho \lambda }^+)),i(u_\nu ,(\mathrm{D}_\mu ,F_{\rho \lambda }^+)),i((u_\nu ,\mathrm{D}_\mu ),F_{\rho \lambda }^+);`$
$`i(u_\mu ,(\mathrm{D},v_{\rho \lambda }^{(s)})),i(\mathrm{D}_\mu ,u_\nu )v_{\rho \lambda }^{(s)};`$
$`(u_\mu ,(u_\nu ,(u_\rho ,\mathrm{D}_\lambda ))),([u_\mu ,u_\nu ],(u_\rho ,\mathrm{D}_\lambda )),`$
$`(\mathrm{D}_\mu ,(u_\nu ,[u_\rho ,u_\lambda ])),(u_\mu ,(\mathrm{D}_\nu ,[u_\rho ,u_\lambda ])));`$
$`(b)`$ $`ϵ^{\mu \nu \rho \lambda }v_\rho ((v\mathrm{D},(\mathrm{D}_\mu ,(\mathrm{D}_\nu ,u_\lambda ))),((v\mathrm{D},\mathrm{D}_\mu ),(\mathrm{D}_\nu ,u_\lambda ))),(\mathrm{D}_\mu ,(v\mathrm{D},(\mathrm{D}_\nu ,u_\lambda ))),`$
$`(u_\mu ,(\mathrm{D}_\nu ,(v\mathrm{D},\mathrm{D}_\lambda ))),(\mathrm{D}_\mu ,(u_\nu ,(v\mathrm{D},\mathrm{D}_\lambda )));`$
$`(vu,(u_\mu ,(u_\nu ,\mathrm{D}_\lambda ))),(u_\mu ,(vu,(u_\nu ,\mathrm{D}_\lambda ))),((u_\mu ,vu),(u_\nu ,\mathrm{D}_\lambda )),`$
$`(\mathrm{D}_\mu ,(u_\nu ,(vu,u_\lambda ))),(u_\mu ,(\mathrm{D}_\nu ,(vu,u_\lambda )))),`$
$`i(\mathrm{D}_\mu ,(u_\nu ,F_{\kappa \lambda }^+))v^\kappa ,i(u_\nu ,(\mathrm{D}_\mu ,F_{\kappa \lambda }^+))v^\kappa ,i((u_\nu ,\mathrm{D}_\mu ),F_{\kappa \lambda }^+)v^\kappa ;`$
$`i(u_\mu ,(\mathrm{D},v_{\kappa \lambda }^{(s)}))v^\kappa ,i(\mathrm{D}_\mu ,u_\nu )v_{\kappa \lambda }^{(s)}v^\kappa ;`$
$`(c)`$ $`ϵ^{\mu \nu \rho \lambda }v_\rho ((\mathrm{D}_\mu ,(\mathrm{D}_\nu ,(v\mathrm{D},u_\lambda ))),([\mathrm{D}_\mu ,\mathrm{D}_\nu ],(v\mathrm{D},u_\lambda ));`$
$`(v\mathrm{D},(u_\lambda ,[\mathrm{D}_\mu ,\mathrm{D}_\nu ])),(u_\mu ,(v\mathrm{D},[\mathrm{D}_\mu ,\mathrm{D}_\nu ]));`$
$`(u_\mu ,(u_\nu ,(v\mathrm{D},u_\lambda ))),([u_\mu ,u_\nu ],(v\mathrm{D},u_\lambda )));`$
$`(u_\lambda ,(v\mathrm{D},[u_\mu ,u_\nu ])),(v\mathrm{D},(u_\lambda ,[u_\mu ,u_\nu ]))`$
$`(v\mathrm{D},(u_\lambda ,F_{\mu \nu }^+)),(u_\lambda ,(v\mathrm{D},F_{\mu \nu }^+)),((u_\lambda ,v\mathrm{D}),F_{\mu \nu }^+));`$
$`(v\mathrm{D},(u_\lambda ,v_{\mu \nu }^{(s)})),(u_\lambda ,(v\mathrm{D},v_{\mu \nu }^{(s)})),((u_\lambda ,v\mathrm{D}),v_{\mu \nu }^{(s)}));`$
$`(d)`$ $`ϵ^{\mu \nu \rho \lambda }v_\rho ((\mathrm{D}_\mu ,(\mathrm{D}_\nu ,(\mathrm{D}_\lambda ,vu))),([\mathrm{D}_\mu ,\mathrm{D}_\nu ],(vu,\mathrm{D}_\lambda ));`$
$`(\mathrm{D}_\lambda ,(vu,[\mathrm{D}_\mu ,\mathrm{D}_\nu ])),(vu,(\mathrm{D}_\lambda ,[\mathrm{D}_\mu ,\mathrm{D}_\nu ]));`$
$`(u_\mu ,(u_\nu ,(vu,\mathrm{D}_\lambda ))),([u_\mu ,u_\nu ],(vu,\mathrm{D}_\lambda )));`$
$`(vu,(\mathrm{D}_\lambda ,[u_\mu ,u_\nu ])),(\mathrm{D}_\lambda ,(vu,[u_\mu ,u_\nu ]));`$
$`(\mathrm{D}_\lambda ,(vu,F_{\mu \nu }^+)),(vu,(\mathrm{D}_\lambda ,F_{\mu \nu }^+)),((vu,\mathrm{D}_\lambda ),F_{\mu \nu }^+));`$
$`(\mathrm{D}_\lambda ,(vu,v_{\mu \nu }^{(s)})),(vu,(\mathrm{D}_\lambda ,v_{\mu \nu }^{(s)})),((vu,\mathrm{D}_\lambda ),v_{\mu \nu }^{(s)})).`$
At least one of $`k,p,q,t,u`$ is $`0`$ in (4):$`(A,(B,C))`$
$`(v)`$ of (7)
Using (9) and (3.1), one sees that one will have to consider the following set of terms together:
$`(\mathrm{S}\mathrm{D},(v\mathrm{D},\chi _{})),(v\mathrm{D},(\mathrm{S}\mathrm{D},\chi _{})),((v\mathrm{D},\mathrm{S}\mathrm{D}),\chi _{});`$
$`i[F_{\mu \nu }^+,\chi _{}]_+v^{[\mu }\mathrm{S}^{\nu ]},iv_{\mu \nu }^{(s)}\chi _{}v^{[\mu }\mathrm{S}^{\nu ]};`$
$`(vu,(\mathrm{S}u,\chi _{})),(\mathrm{S}u,(vu,\chi _{})),((\mathrm{S}u,vu),\chi _{}).`$ (28)
$`(vi)`$ of (7)
Using (3.1), one sees that one will have to consider the following set of terms together:
$$i(\mathrm{S}\mathrm{D},(vu,\chi _+)),i(vu,(\mathrm{S}\mathrm{D},\chi _+)),i((\mathrm{S}\mathrm{D},vu),\chi _+);v\mathrm{S}.$$
(29)
Note that because of parity constraints and the algebra of the $`\mathrm{S}_\mu `$s (See ), there are no Levi Civita-dependent (2,0,0,1,0,0,0)-, (0,2,0,1,0,0,0)- and (1,1,1,0,0,0,0)-type terms.
### 3.3 Isospin Violation
As noted in Section 3, isospin violation enters via $`\chi _\pm `$, we thus need to reconsider the $`(m,n,p,q,t,u,k)`$ with $`p`$ or $`q0`$. We will also have to include trace-dependent terms for these type of terms. The following identities are used in arriving at terms in Table 2:
$`u^\mu \chi _+u_\mu +u^2\chi _++\mathrm{h}.\mathrm{c}.=2u^2\chi _++u^\mu [u_\mu ,\chi _+]_+,`$
$`vu\chi _+vu+u^2\chi _++\mathrm{h}.\mathrm{c}.=2(vu)^2\chi _++vu[vu,\chi _+]_+,`$
$`[\chi _\pm ,\chi _\pm ]_+=2\chi _\pm \chi _\pm +{\displaystyle \frac{1}{2}}[\chi _\pm ,\chi _\pm ]_+{\displaystyle \frac{1}{2}}\chi _\pm ^2,`$
$`[F_{\mu \nu }^+,\chi _\pm ]_+={\displaystyle \frac{1}{2}}[F_{\mu \nu }^+,\chi _\pm ]_++{\displaystyle \frac{1}{2}}[F_{\mu \nu }^+,\chi _\pm ]_+,`$
$`i[v\mathrm{D},[\mathrm{S}u,\chi _+]_+]_+={\displaystyle \frac{i}{2}}[v\mathrm{D},[\mathrm{S}u,\chi _+]_+]_++{\displaystyle \frac{i}{2}}[v\mathrm{D},[\mathrm{S}u,\chi _+]_+,`$
$`i[v\mathrm{D},[\mathrm{S}u,\chi _+]_+]_+={\displaystyle \frac{i}{2}}[v\mathrm{D},[\mathrm{S}u,\chi _+]_+]_++{\displaystyle \frac{i}{2}}[v\mathrm{D},[\mathrm{S}u,\chi _+]_+,`$
$`+{\displaystyle \frac{i}{2}}[v\mathrm{D},[\mathrm{S}u,\chi _+]_+,`$
$`ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda \left([\mathrm{D}_\mu ,[u_\nu ,\chi _{}]_+]_+=[\mathrm{D}_\mu ,[u_\nu ,\chi _{}]_+]_++{\displaystyle \frac{1}{2}}[\mathrm{D}_\mu ,[u_\nu ,\chi _{}]_+]_+\right),`$
$`iϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda ([[u_\mu ,u_\nu ],\chi _+]_+=[[u_\mu ,u_\nu ],\chi _+]_++{\displaystyle \frac{1}{2}}[[u_\mu ,u_\nu ],^H\chi _+]_+),`$
$`i[\mathrm{D}_\mu ,[u^\mu ,\chi _{}]_+]={\displaystyle \frac{i}{2}}[\mathrm{D}_\mu ,[u^\mu ,\chi _{}]_+]+{\displaystyle \frac{i}{2}}[\mathrm{D}_\mu ,[u^\mu ,\chi _{}]_+],`$
$`i[[\mathrm{D}_\mu ,u^\mu ],\chi _{}]_+={\displaystyle \frac{i}{2}}[[\mathrm{D}_\mu ,u^\mu ],\chi _{}]_++{\displaystyle \frac{i}{2}}[[\mathrm{D}_\mu ,u^\mu ],\chi _{}]_+,`$
$`i[u_\mu ,[\mathrm{D}^\mu ,\chi _{}]]_+={\displaystyle \frac{i}{2}}[u_\mu ,[\mathrm{D}^\mu ,\chi _{}]]_++{\displaystyle \frac{i}{2}}[u_\mu ,[\mathrm{D}^\mu ,\chi _{}]_+.`$ (30)
## 4 The Lists of Independent Terms in $`_{\mathrm{HBChPT}}`$ (off-shell nucleons)
In this section, using (6), and the algebraic reductions of Section 3, we list all possible $`𝒜`$-type terms of O$`(q^4,\varphi ^{2n})`$, and O$`(q^4,\varphi ^{2n+1})`$ in Tables 1 and 2, that are allowed by (6) and have not been eliminated in Section 3. As noted in Section 2 (and ), for off-shell nucleons, $`\gamma ^0^{}\gamma ^0𝒞^1𝒜`$. Hence, it is sufficient to list only $`𝒜`$-type terms (for off-shell nucleons).
Using the algebraic identities of Section 3, if we end up with $`m`$ independent identities in $`n(>m)`$ terms, then we can take $`(nm)`$ linearly independent terms. Even though the phase rule (6) and linear independence of terms are sufficient for listing terms in the O$`(q^4`$) HBChPT Lagrangian for off-shell nucleons, however, if for a given choice of terms and group of terms in (7), we find similar group of terms in , then while listing the $`(nm)`$ terms, preference is given to including terms that also figure in Table 1 of . The reason for doing the same is that this allows for an easy identification of the finite terms, given that the divergent (counter) terms have been worked out in .
In tables 1 (that one gets if one assumes isospin symmetry) and 2 (that one gets if one includes isospin violation), the allowed 7-tuples $`(m,n,p,q,t,u,k)`$ are listed along with the corresponding terms. The main aim is to find the number of finite O($`q^4)`$ terms, given that the UV divergent terms have already been worked out in . For this purpose, the terms in tables 1 and 2 are labeled as F denoting the finite terms and D denoting the divergent terms. For the purpose of comparison with , we have also indicated which terms in table 1 of the D-type terms correspond to. The LECs of O$`(q^4)`$ terms in are denoted by $`d_i,i=1`$ to 199. Further, the $`i=188`$ term in Table 1 of should have $`\mathrm{S}_\rho `$ instead of $`v_\rho `$.
Overall, one gets 27 finite and 79 divergent (counter) terms at O$`(q^4)`$.
## 5 On-shell reduction
In this section, we discuss the derivation of the on-shell O($`q^4`$) $`_{\mathrm{HBChPT}}`$, directly within HBChPT using the techniques of .
The main result obtained in extended to include external fields in the context of complete on-shell reduction within HBChPT was the following rule:
$`𝒜\mathrm{type}\mathrm{terms}\mathrm{of}\mathrm{the}\mathrm{form}\overline{\mathrm{H}}\mathrm{S}\mathrm{D}𝒪\mathrm{H}+\mathrm{h}.\mathrm{c}.`$
$`\mathrm{or}\overline{\mathrm{H}}\mathrm{v}\mathrm{D}𝒪\mathrm{H}+\mathrm{h}.\mathrm{c}.`$
$`\mathrm{or}\overline{\mathrm{H}}𝒪^\mu \mathrm{D}_\mu \mathrm{H}+\mathrm{h}.\mathrm{c}.\mathrm{can}\mathrm{be}\mathrm{eliminated}`$
$`\mathrm{except}\mathrm{for}𝒪_\mu (i^{m_1+l_5+l_7+t+u+1},\mathrm{or}ϵ^{\nu \lambda \kappa \rho }\times \mathrm{\Omega })\times u_\mu \mathrm{\Lambda }`$
$`\mathrm{with}l_11,\mathrm{\Omega }1(i)\mathrm{for}()^{m_1+l_5+l_7+t+u+1}=1(1),`$
$`\mathrm{or}`$
$`𝒪_\mu (i^{m_1+l_5+l_7+k+t+u},\mathrm{or}ϵ^{\nu \lambda \kappa \rho }\times \mathrm{\Omega }^{})\times \mathrm{D}_\mu \mathrm{\Lambda },`$
$`\mathrm{with}l_11,\mathrm{\Omega }^{}1(i)\mathrm{for}()^{m_1+l_5+l_7+k+t+u}=1(1).`$ (31)
In (5)
$$\mathrm{\Lambda }\underset{i=1}{\overset{M_1}{}}𝒱_{\nu _i}\underset{j=1}{\overset{M_2}{}}u_{\rho _j}(vu)^{l_1}u^{2l_2}\chi _+^{l_3}\chi _{}^{l_4}([v\mathrm{D},])^{l_5}(\mathrm{D}_\beta \mathrm{D}^\beta )^{l_6}(u_\alpha \mathrm{D}^\alpha )^{l_7}(v_{\sigma \omega }^{(s)})^k(F_{\rho \lambda }^+)^t(F_{\mu \nu }^+)^u\mathrm{S}_\kappa ^r,$$
(32)
where $`𝒱_{\nu _i}v_{\nu _i}\mathrm{or}\mathrm{D}_{\nu _\mathrm{i}}`$. where $`𝒱_{\nu _i}v_{\nu _i}\mathrm{or}\mathrm{D}_{\nu _\mathrm{i}}`$. The number of $`\mathrm{D}_{\nu _i}`$s in (32) equals $`m_1(M_1)`$. Assuming that Lorentz invariance, parity and hermiticity have been implemented, the choice of the factors of $`i`$ in (5) automatically incorporates the phase rule (6). In (5), it is only the contractions of the building blocks that has been indicated. Also, traces have not been indicated. It is understood that all (anti-)commutators in the HBChPT Lagrangian are to be expanded out until one hits the first $`\mathrm{D}_\mu `$, so that the $`𝒜`$-type HBChPT term can be put in the form $`\overline{\mathrm{H}}𝒪^\mu \mathrm{D}_\mu \mathrm{H}`$ +h.c.
The complete on-shell O($`q^4`$) HBChPT Lagrangian can be shown to be given by:
$`_{\mathrm{HBChPT}}^{(4)}=A^{(4)}+{\displaystyle \frac{1}{2\mathrm{m}}}\gamma ^0B^{(2)}{}_{}{}^{}\gamma _{}^{0}B^{(2)}`$
$`+{\displaystyle \frac{1}{2\mathrm{m}}}\left[\gamma ^0B^{(3)}{}_{}{}^{}\gamma _{}^{0}B^{(1)}+\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}B^{(3)}\right]`$
$`{\displaystyle \frac{1}{4\mathrm{m}^2}}\left[\gamma ^0B^{(2)}{}_{}{}^{}\gamma _{}^{0}C^{(1)}B^{(1)}+\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}C^{(1)}B^{(2)}\right]`$
$`{\displaystyle \frac{1}{4\mathrm{m}^2}}\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}C^{(2)}B^{(1)}+{\displaystyle \frac{1}{8\mathrm{m}^3}}\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}(C^{(1)})^2B^{(1)}.`$ (33)
Using
$$B_{\mathrm{OS}}^{(2)}[v^{(s)},F^+,\chi _{}]=\alpha _3\gamma ^5\chi _3i\alpha _7\gamma ^5v^{[\mu }\mathrm{S}^{\nu ]}v_{\mu \nu }^{(s)}+i\alpha _8\gamma ^5v^{[\mu }\mathrm{S}^{\nu ]}F_{\mu \nu }^++\alpha _9\chi _{},$$
(34)
(OS$``$on-shell) one gets:
$`{\displaystyle \frac{1}{2\mathrm{m}}}\gamma ^0B^{(2)}{}_{}{}^{}\gamma _{}^{0}B^{(2)}[v^{(s)},F^+,\chi _{}]`$
$`={\displaystyle \frac{1}{2\mathrm{m}}}[+(\alpha _3\chi _{}+\alpha _9\chi _{})^2\alpha _1[[vu,\mathrm{S}u],(\alpha _{}+\alpha _9\chi _{})]_+`$
$`v^{[\mu }\mathrm{S}^{\nu ]}[(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+,(\alpha _3\chi _{}+\alpha _9\chi _{})]_++i\alpha _4[(\alpha _3\chi _{}+\alpha _9\chi _{}),[v\mathrm{D},vu]]_+`$
$`+(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+(\alpha _7v_\rho ^{(s)\nu }+\alpha _8F_\rho ^{+\nu }){\displaystyle \frac{i}{2}}ϵ^{\nu \lambda \alpha \beta }v_\alpha v^\mu v^\rho \mathrm{S}_\beta (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)(\alpha _7v_{\rho \lambda }^{(s)}+\alpha _8F_{\rho \lambda }^+)`$
$`i\alpha _1\alpha _7\left({\displaystyle \frac{1}{4}}v^\mu [[vu,u^\nu ],v_{\mu \nu }^{(s)}]_+\right)+i\alpha _1\alpha _8\left({\displaystyle \frac{1}{4}}v^\mu [[vu,u^\nu ],F_{\mu \nu }^+]_++{\displaystyle \frac{1}{2}}ϵ^{\nu \rho \alpha \beta }v_\alpha \mathrm{S}_\beta v^\mu [[vu,u_\rho ],F_{\mu \nu }^+]\right)`$
$`+i\alpha _2\alpha _7\left({\displaystyle \frac{1}{2}}ϵ^{\nu \rho \alpha \beta }v_\alpha \mathrm{S}_\beta v^\mu [[vu,u^\rho ]_+,v_{\mu \nu }^{(s)}]_+\right)`$
$`+i\alpha _2\alpha _8\left({\displaystyle \frac{1}{4}}v^\mu [[vu,u^\nu ]_+,F_{\mu \nu }^+]{\displaystyle \frac{1}{2}}ϵ^{\nu \rho \alpha \beta }v_\alpha \mathrm{S}_\beta v^\mu [[vu,u^\rho ]_+,F_{\mu \nu }^+]_+\right)`$
$`\alpha _3v^{[\mu }\mathrm{S}^{\nu ]}[(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+),\chi _{}]_+\alpha _4v^{[\mu }\mathrm{S}^{\nu ]}[(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+),[v\mathrm{D},vu]]_+`$
$`i\alpha _2\alpha _4[[vu,\mathrm{S}u]_+,[v\mathrm{D},vu]]+i\alpha _3\alpha _4[\chi _{},[v\mathrm{D},vu]]_+].`$ (35)
Using (34) and $`C_{\mathrm{OS}}^{(1)}`$, and eliminating all terms proportional to the nonrelativistic eom by field redefinition of H, one gets:
$`{\displaystyle \frac{1}{4\mathrm{m}^2}}\gamma ^0B^{(2)}{}_{}{}^{}\gamma _{}^{0}C^{(1)}B^{(1)}[v^{(s)},F^+]+\mathrm{h}.\mathrm{c}.`$
$`=(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)[2iv^\mu [v\mathrm{D},\mathrm{D}^\nu ]+g_A^0\mathrm{D}^\nu \mathrm{S}u+2ϵ^{\nu \rho \lambda \sigma }v^\mu v_\lambda \mathrm{S}_\sigma [v\mathrm{D},\mathrm{D}_\rho ]`$
$`{\displaystyle \frac{i}{2}}g_A^0ϵ^{\nu \rho \lambda \sigma }v^\mu v_\lambda \mathrm{D}_\rho u_\sigma g_A^0v^\mu \mathrm{S}^\nu [v\mathrm{D},vu]{\displaystyle \frac{g_{A}^{0}{}_{}{}^{2}}{4}}vuu^\nu v^\mu `$
$`+{\displaystyle \frac{ig_{A}^{0}{}_{}{}^{2}}{2}}ϵ^{\alpha \nu \rho \omega }v^\mu v_\rho \mathrm{S}_\omega vuu_\alpha +g_A^0\mathrm{S}^\nu v^\mu \mathrm{D}ug_A^0v^\mu \mathrm{S}\mathrm{D}u^\nu `$
$`g_A^0v^\mu \mathrm{S}^\nu [v\mathrm{D},vu]+{\displaystyle \frac{i}{2}}ϵ^{\rho \lambda \alpha \nu }v^\mu v_\alpha u_\rho \mathrm{D}_\lambda +g_A^0v^\mu u^\nu \mathrm{S}\mathrm{D}g_A^0v^\mu \mathrm{S}u\mathrm{D}^\nu `$
$`+{\displaystyle \frac{ig_{A}^{0}{}_{}{}^{2}}{2}}({\displaystyle \frac{1}{2}}v^\mu u^\nu +iϵ^{\nu \rho \lambda \sigma }v^\mu v_\lambda \mathrm{S}_\sigma u_\rho )]`$
$`+(2(\alpha _3\chi _{}+\alpha _9\chi _{})v\mathrm{DS}\mathrm{D}i{\displaystyle \frac{g_A^0}{2}}(\alpha _3\chi _{}+\alpha _9\chi _{})[v\mathrm{D},vu]+{\displaystyle \frac{g_{A}^{0}{}_{}{}^{2}}{2}}(\alpha _3\chi _{}+\alpha _9\chi _{})vu\mathrm{S}u`$
$`+g_A^0(\alpha _3\chi _{}+\alpha _9\chi _{})[{\displaystyle \frac{1}{4}}vu\mathrm{S}u+{\displaystyle \frac{i}{4}}u_\mu \mathrm{D}^\mu +{\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda u_\mu \mathrm{D}_\nu ]+\mathrm{h}.\mathrm{c}.]`$ (36)
Similarly, using:
$$C_\mathrm{S}^{(2)}[v^{(s)},F^+,\chi _+]=\alpha _6\chi _+\alpha _7ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda v_{\mu \nu }^{(s)}+\alpha _8ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda F_{\mu \nu }^+\alpha _{10}\chi _+$$
(37)
and
$$B^{(1)}=2i\gamma ^5\mathrm{S}\mathrm{D}\frac{g_A^0}{2}\gamma ^5vu,$$
(38)
and eliminating all terms proportional to the nonrelativistic eom by field redefinition of H, one sees that:
$`{\displaystyle \frac{1}{4\mathrm{m}^2}}\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}C^{(2)}B^{(1)}[v^{(s)},F^+,\chi _+]`$
$`={\displaystyle \frac{1}{4\mathrm{m}^2}}[ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}\mathrm{D}(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)\mathrm{D}_\lambda +(i\mathrm{D}^\mu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)\mathrm{D}^\nu `$
$`+g_A^0\mathrm{S}u(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)v^\mu \mathrm{D}^\nu +g_A^0v^\mu \mathrm{D}^\nu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)\mathrm{S}u)`$
$`{\displaystyle \frac{g_{A}^{0}{}_{}{}^{2}}{4}}ϵ^{\mu \nu \rho \lambda }v_\rho vu(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)\mathrm{S}_\lambda vu`$
$`+2ϵ^{\mu \nu \rho \lambda }v_\rho \left({\displaystyle \frac{ig_A^0}{4}}v^\omega \mathrm{D}_\mu (\alpha _7v_{\omega \nu }^{(s)}+\alpha _8F_{\omega \nu }^+)u_\lambda \mathrm{D}_\mu (\alpha _7v_{\omega \nu }^{(s)}+\alpha _8F_{\omega \nu }^+)\mathrm{D}^\omega \mathrm{S}_\lambda \right)`$
$`g_A^0(\mathrm{S}\mathrm{D}(\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)u^\nu v^\sigma \mathrm{D}^\rho v^\sigma (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)u_\rho `$
$`+{\displaystyle \frac{ig_A^0}{4}}u^\nu v^\sigma (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)vu{\displaystyle \frac{g_A^0}{2}}ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda u_\mu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)v^\sigma vu`$
$`+({\displaystyle \frac{ig_A^0}{4}}ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{D}_\lambda (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)vu{\displaystyle \frac{ig_{A}^{0}{}_{}{}^{2}}{4}}v^\mu u^\nu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)vu`$
$`{\displaystyle \frac{g_{A}^{0}{}_{}{}^{2}}{2}}ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda v^\kappa u_\mu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)vu`$
$`g_A^0\mathrm{D}^\nu (\alpha _7v_{\mu \nu }^{(s)}+\alpha _8F_{\mu \nu }^+)\mathrm{S}^\nu +h.c.)`$
$`+(g_{A}^{0}{}_{}{}^{2}[{\displaystyle \frac{1}{4}}vu(\alpha _6\chi _++\alpha _{10}\chi _+)vu{\displaystyle \frac{1}{4}}u_\mu (\alpha _6\chi _++\alpha _{10}\chi _+)u^\mu +{\displaystyle \frac{i}{2}}ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda u_\mu (\alpha _6\chi _++\alpha _{10}\chi _+)u_\nu ]`$
$`+\mathrm{D}_\mu \chi _+\mathrm{D}^\mu 2iϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda \mathrm{D}_\mu \chi _+\mathrm{D}_\nu `$
$`+(ig_A^0vu\chi _+\mathrm{S}\mathrm{D}+\mathrm{h}.\mathrm{c}.)+{\displaystyle \frac{g_{A}^{0}{}_{}{}^{2}}{4}}vu\chi _+vu)]`$ (39)
Using:
$`B_{\mathrm{OS}}^{(3)}[v^{(s)},F^+,\chi _+,\chi _{}]=\beta _4\gamma ^5[vu,\chi _+]_++\beta _{14}\gamma ^5[\chi _{},\mathrm{S}u]`$
$`+\beta _{17}\gamma ^5\mathrm{S}^\nu [\mathrm{D}^\mu ,v_{\mu \nu }^{(s)}]i\beta _{18}\gamma ^5v^\mu [F_{\mu \nu }^+,u^\nu ]+\beta _{19}ϵ^{\mu \nu \rho \lambda }\gamma ^5\mathrm{S}_\lambda [F_{\mu \nu }^+,u_\rho ]_+`$
$`+\beta _{20}ϵ^{\mu \nu \rho \lambda }\gamma ^5v_{\mu \nu }^{(s)}u_\rho \mathrm{S}_\lambda +\beta _{21}ϵ^{\mu \nu \rho \lambda }\gamma ^5v_\rho \mathrm{S}_\lambda [F_{\mu \nu }^+,vu]`$
$`+\beta _{22}\gamma ^5\mathrm{S}^\nu [\mathrm{D}^\mu ,F_{\mu \nu }^+]+\beta _{23}\gamma ^5[vu,\chi _+]_+`$
$`+\beta _2\gamma ^5vu\chi _++i\beta _{25}\gamma ^5[v\mathrm{D},\chi _{}],`$ (40)
$`\frac{1}{2\mathrm{m}}\left[\gamma ^0B^{(3)}{}_{}{}^{}\gamma _{}^{0}B^{(1)}+\gamma ^0B^{(1)}{}_{}{}^{}\gamma _{}^{0}B^{(3)}\right][v^{(s)},F^+,\chi _{},\chi _+]`$ and eliminating all terms proportional to the nonrelativistic eom by field redefinition of H, one gets:
$`{\displaystyle \frac{1}{2\mathrm{m}}}[\beta _4(2i[\mathrm{S}\mathrm{D},[vu,\chi _+]_+]_+{\displaystyle \frac{g_A^0}{2}}[vu,[vu,\chi _+]_+]_+)`$
$`+\beta _{14}(i[\mathrm{D}_\mu ,[\chi _{},u^\mu ]]_++g_A^0[\mathrm{S}u,[\chi _{},vu]]_+ϵ^{\mu \nu \rho \lambda }v_\rho \mathrm{S}_\lambda [\mathrm{D}_\nu ,[\chi _{},u_\mu ]]`$
$`{\displaystyle \frac{g_A^0}{2}}[vu,[\chi _{},\mathrm{S}u]]_+)`$
$`+({\displaystyle \frac{ig_A^0}{2}}[\mathrm{D}^\mu ,(i\beta _{17}v_{\mu \nu }^{(s)}+i\beta _{22}F_{\mu \nu }^+)]v^\nu \mathrm{S}uu`$
$`{\displaystyle \frac{1}{2}}[\mathrm{D}^\mu ,(i\beta _{17}v_{\mu \nu }^{(s)}+i\beta _{22}F_{\mu \nu }^+)]\mathrm{D}_\nu `$
$`+iϵ^{\nu \rho \lambda \sigma }v_\lambda \mathrm{S}_\sigma [\mathrm{D}^\mu ,(i\beta _{17}v_{\mu \nu }^{(s)}+i\beta _{22}F_{\mu \nu }^+)]\mathrm{D}^\rho `$
$`i{\displaystyle \frac{g_A^0}{2}}\mathrm{S}^\nu [\mathrm{D}^\mu ,(i\beta _{17}v_{\mu \nu }^{(s)}+i\beta _{22}F_{\mu \nu }^+)]vu)`$
$`+\beta _{18}(2v^\mu [F_{\mu \nu }^+,u^\nu ]\mathrm{S}\mathrm{D}+{\displaystyle \frac{ig_A^0}{2}}v^\mu [F_{\mu \nu }^+,u^\nu ]vu)+\beta _{19}({\displaystyle \frac{i}{2}}ϵ^{\lambda \alpha \rho \beta }[F_{\lambda \alpha }^+,u_\rho ]_+(ig_A^0v_\beta \mathrm{S}u\mathrm{D}_\beta )`$
$`2([F_{\lambda \alpha }^+,u_\nu ]_+v^\lambda \mathrm{S}^\alpha \mathrm{D}^\nu \mathrm{S}^\alpha [F_{\lambda \alpha }^+,vu]_+\mathrm{D}^\lambda +v^\alpha [F_{\lambda \alpha }^+,\mathrm{S}u]_+\mathrm{D}^\nu )`$
$`+{\displaystyle \frac{g_A^0}{2}}ϵ^{\lambda \alpha \rho \beta }\mathrm{S}_\beta [F_{\lambda \alpha }^+,u_\rho ]_+vu)+\beta _{20}({\displaystyle \frac{i}{2}}ϵ^{\lambda \alpha \rho \beta }v_{\lambda \alpha }^{(s)}u_\rho (ig_A^0\mathrm{S}u\mathrm{D}_\beta )+{\displaystyle \frac{g_A^0}{2}}ϵ^{\lambda \alpha \rho \beta }v_{\lambda \alpha }^{(s)}u_\rho vu`$
$`2(v_{\lambda \alpha }^{(s)}v^\lambda \mathrm{S}^\alpha u_\nu \mathrm{D}^\nu \mathrm{S}^\lambda v_{\lambda \alpha }^{(s)}vu\mathrm{D}^\lambda +v^\lambda v_{\lambda \alpha }^{(s)}\mathrm{S}u\mathrm{D}^\nu ))`$
$`+\beta _{21}({\displaystyle \frac{i}{2}}ϵ^{\lambda \alpha \rho \beta }v_\rho [F_{\lambda \alpha }^+,vu]\mathrm{D}_\beta +2[F_{\lambda \alpha }^+,vu]v^\lambda ({\displaystyle \frac{1}{4}}u^\alpha +{\displaystyle \frac{i}{2}}ϵ^{\mu \nu \rho \kappa }v_\rho \mathrm{S}_\kappa u_\mu )`$
$`2\mathrm{S}^\alpha [F_{\lambda \alpha }^+,vu]\mathrm{D}^\lambda {\displaystyle \frac{g_A^0}{2}}ϵ^{\lambda \alpha \rho \beta }v_\rho \mathrm{S}_\beta [F_{\lambda \alpha }^+,vu]vu)`$
$`+\beta _{23}(2i[vu,\chi _+]_+\mathrm{S}\mathrm{D}+{\displaystyle \frac{g_A^0}{2}}vu[vu,\chi _+]_+`$
$`+\beta _{24}(vu\chi _+\mathrm{S}\mathrm{D}+{\displaystyle \frac{g_A^0}{2}}(vu)^2\chi _+)`$
$`+\beta _{25}(2[v\mathrm{D},\chi _{}\mathrm{S}\mathrm{D}+i{\displaystyle \frac{g_A^0}{2}}[v\mathrm{D},\chi _{}]vu)].`$ (41)
The set $`\{\beta _i\}`$ can be related to the set $`\{b_i\}`$ of .
## 6 Conclusion
A complete list of O$`(q^4)`$ terms for off-shell nucleons was obtained working within HBChPT using a phase rule obtained in , along with reductions from algebraic identities. We also obtain the on-shell O($`q^4`$) terms, again within the framework of HBChPT. For off-shell nucleons, one gets a total of 106 $`𝒪(q^4)`$ terms (given in Tables 1 and 2). Of these 27 are finite. Contrary to what is claimed in , the earlier version of this paper itself contained (overcomplete) list of terms including external fields. Besides, the whole point of this paper (and that of and ) is to carry out the $`1/\mathrm{m}`$-reduction without actually doing it; this is carried out by developing a method of imposing charge conjugation invariance directly within the nonrelativistic framework in terms of a phase rule (6) that can be used directly within HBChPT. Also, once having obtained the list of nonrelativistic terms, constructing their relativistic counterparts (for reasons given in ) can easily be done by following . We are getting fewer terms than the ones given in . Also, instead of working with symmetrized and anti-symmetrised commutators of $`\mathrm{D}_\mu `$ and $`u_\nu `$, we use (10) to eliminate $`F_{\mu \nu }^{}`$ altogether. |
warning/0001/math-ph0001042.html | ar5iv | text | # Semiclassical Limit for the Schrödinger Equation with a Short Scale Periodic Potential
## 1. Introduction
A basic problem of solid state physics is to understand the motion of electrons in the periodic potential which is generated by the ionic cores. While this problem is quantum mechanical, many electronic properties of solids can be understood already in the semiclassical approximation . One argues that if the wave packet spreads over many lattice spacings, the kinetic energy $`(\mathrm{}k)^2/2m`$ is modified to the $`n`$-th band energy $`E_n(k)`$. Otherwise the electron responds to external fields, $`E_{\mathrm{ex}}`$, $`B_{\mathrm{ex}}`$, as in the case of vanishing periodic potential. Thus the semiclassical equations of motion are
(1.1)
$$\begin{array}{c}\dot{r}=v_n(k)=_kE_n(k)\hfill \\ \\ \mathrm{}\dot{k}=e(E_{\mathrm{ex}}(r)+v_n(k)B_{\mathrm{ex}}(r)),\hfill \end{array}$$
where $`r`$ is the position and $`k`$ the quasimomentum of the electron. Note that there is a semiclassical evolution for each band separately.
The goal of our paper is to understand on a mathematical level how these semiclassical equations arise from the underlying Schrödinger equation. We consider only the case where $`B_{\mathrm{ex}}=0`$.
The setup is rather obvious. We start from the Schrödinger equation
(1.2)
$$i\frac{}{t}\psi =H\psi $$
with Hamiltonian
(1.3)
$$H=\frac{1}{2}\mathrm{\Delta }+V(x)+W(\epsilon x).$$
The electron moves in $`^d`$ and the solution to (1.2) defines the unitary time evolution $`U^\epsilon (t)\psi (x)=\mathrm{e}^{itH}\psi (x)=\psi (x,t)`$ in $`L^2(^d)`$. We have chosen units such that $`\mathrm{}=1`$ and the mass of the particle $`m=1`$. $`V(x)`$ is a periodic potential with average lattice spacing $`a`$. The precise conditions on $`V`$ will be spelled out in the following section, where we also describe the direct fiber integral decomposition for periodic Schrödinger operators. The lattice spacing $`a`$ defines the microscopic spatial scale. $`W(\epsilon x)`$ is an external electrostatic potential with dimensionless scale parameter $`\epsilon `$, $`\epsilon 1`$, which means that $`W`$ is slowly varying on the scale of the lattice. For real metals the condition of slow variation is satisfied even for the strongest external electrostatic fields available, cf. , Chapter 13.
The external forces due to $`W`$ are of order $`\epsilon `$ and therefore have to act over a time of order $`\epsilon ^1`$ to produce finite changes, which defines the macroscopic time scale. We will mostly work in the microscopic coordinates $`(x,t)`$ of (1.2). For sake of comparison we note that the macroscopic space-time scale $`(x^{},t^{})`$ is defined through $`x=\epsilon ^1x^{}`$ and $`t=\epsilon ^1t^{}`$. With this scale change Eqs. (1.2), (1.3) read
(1.4)
$$\begin{array}{c}i\epsilon \frac{}{t^{}}\psi =H\psi ,\hfill \\ \\ H=\left(\epsilon ^2\frac{1}{2}\mathrm{\Delta }^{}+V(x^{}/\epsilon )+W(x^{})\right)\hfill \end{array}$$
with initial conditions $`\psi ^\epsilon (x^{})=\epsilon ^{d/2}\psi (x^{}/\epsilon )`$. If $`V=0`$, Eq. (1.4) is the usual semiclassical limit with $`\epsilon `$ set equal to $`\mathrm{}`$. Thus our problem is to understand how an additional periodic, but rapidly oscillating potential modifies the standard picture.
The two scale problem (1.2), (1.3) can be attacked along several routes. A first choice would be time dependent WKB . In the limit $`\epsilon 0`$, for each energy band separately, one obtains a Hamilton-Jacobi equation for the phase and a transport equation for the amplitude of the wave function $`\psi (x,t)`$. As a main draw-back of this method, generically, the solution to the Hamilton-Jacobi equation develops singularities after some finite macroscopic time. If $`V=0`$, it is well understood how to go beyond such caustics by introducing new coordinates on the Lagrangian manifold. For (1.2), (1.3) a corresponding program has not yet been attempted. The results are valid only over a finite macroscopic time span with a duration depending on the initial wave function.
Another variant is to establish the semiclassical limit through the convergence of Wigner functions. In our context one defines a band Wigner function $`W_n^\epsilon (r,k,t)`$ depending on the band index $`n`$ and as a function of the position and quasimomentum. One then wants to prove that in the limit $`\epsilon 0`$ $`W_n^\epsilon (t)`$ converges to $`\overline{W}_n(t)`$, which is the initial band Wigner function $`\overline{W}_n(0)`$ evolved according to the semiclassical flow (1.1). Such a result is established in for the case of zero external potential, the general case being left open as a challenging problem.
A third approach to the semiclassical limit for $`V=0`$ is the strong convergence of Heisenberg operators . We briefly recall its main features. We define, as unbounded operators on $`L^2(^d)`$,
$`x(t):=\mathrm{e}^{itH}x\mathrm{e}^{itH},`$
$`p(t):=\mathrm{e}^{itH}p\mathrm{e}^{itH},p=i_x,`$
where $`H`$ is the Hamiltonian in (1.3) with $`V=0`$. The goal is to establish the strong limit of
$$x^\epsilon (t)\psi =\epsilon x(\epsilon ^1t)\psi ,p^\epsilon (t)\psi =p(\epsilon ^1t)\psi $$
as $`\epsilon 0`$ with $`\psi `$ in a suitable domain. In the trivial case of free motion, $`W=0`$, this amounts to the strong convergence of $`x^\epsilon (t)\psi =(\epsilon x+pt)\psi `$, $`p^\epsilon (t)\psi =p\psi `$, which yields $`lim_{\epsilon 0}x^\epsilon (t)=pt`$, $`lim_{\epsilon 0}p^\epsilon (t)=p`$. The general case requires more work . One obtains the strong limits
(1.5)
$$\begin{array}{c}\underset{\epsilon 0}{lim}x^\epsilon (t)=r(p,t),\hfill \\ \\ \underset{\epsilon 0}{lim}p^\epsilon (t)=u(p,t).\hfill \end{array}$$
Here $`r(p,t)`$, $`u(p,t)`$ are solutions of
(1.6)
$$\dot{r}=u,\dot{u}=W(r)$$
with initial conditions $`r_0=0`$, $`u_0=p`$. The initial condition $`r_0=0`$ reflects that $`|\psi |^2`$ looks like $`\delta (r)`$ on the macroscopic scale, provided that $`\psi _2=1`$. For general initial conditions, $`r_00`$, we would have to shift the initial $`\psi `$ by $`\epsilon ^1r_0`$.
The strong operator convergence may look slightly abstract, but all the desired physical information can be deduced. E.g., the initial $`\psi `$ defines the momentum distribution $`|\widehat{\psi }(k)|^2`$ independent of $`\epsilon `$ and the $`\delta (r)`$ spatial distribution in the limit $`\epsilon 0`$. Then, according to (1.5), for small $`\epsilon `$ the position distribution at time $`t`$ is given by
$`{\displaystyle _^d}f(x)|\psi ^\epsilon (x,t)|^2𝑑x=(\psi ,f(x^\epsilon (t))\psi )`$
$`(\psi ,f(r(p,t))\psi )={\displaystyle |\widehat{\psi }(k)|^2f(r(k,t))𝑑k},`$
which means that the phase space distribution $`\delta (r)|\widehat{\psi }(k)|^2drdk`$ is transported according to the semiclassical flow (1.6). The spatial marginal of this distribution at time $`t`$ is the desired approximation to the true position distribution $`|\psi ^\epsilon (x,t)|^2`$. $`|\psi ^\epsilon (x,t)|^2`$ may oscillate rapidly on small scales and some averaging, as embodied by the test function $`f`$, is needed.
In this paper we investigate the semiclassical limit (1.2), (1.3) through the strong convergence of the position operator $`x^\epsilon (t)`$. We will show that, in the limit $`\epsilon 0`$, $`x^\epsilon (t)`$ is diagonal with respect to the band index and in each band the structure is analogous to (1.5) with $`p`$ replaced by the quasimomentum $`k`$ and (1.6) replaced by (1.1). More generally we will consider the semiclassical limit of the Weyl quantized operators $`a^W(\epsilon x,p)`$, whose classical symbol is periodic in $`p`$.
To give a short outline: In the following section we collect some properties of periodic Schrödinger operators. In Section 3 we state our main results, which are proved in Sections 5, 6, and 7, respectively. In Section 4 we discuss some implications for the position and quasimomentum distributions, and, more generally, for the band Wigner functions. The difficulties arising from band crossings are explained in Section 9.
## 2. Periodic Schrödinger operators
For the periodic potential $`V`$ we will need only some rather minimal assumptions, which we state as
###### Condition ($`\mathrm{C}_{\mathrm{per}}`$).
Let $`\mathrm{\Gamma }^d`$ be the lattice generated by the basis $`\{\gamma _1,\mathrm{},\gamma _d\}`$, $`\gamma _i^d`$. Then $`V(x+\gamma )=V(x)`$ for all $`x^d`$, $`\gamma \mathrm{\Gamma }`$. Furthermore, we assume $`V`$ to be infinitesimally operator bounded with respect to $`H_0`$. The last condition is satisfied, e.g., if $`VL^p(M)`$, where $`M`$ is the fundamental domain of $`\mathrm{\Gamma }`$, and $`p=2`$ for $`d3`$ and $`p>d/2`$ for $`d>3`$, respectively.
($`\mathrm{C}_{\mathrm{per}}`$) will be assumed throughout.
We recall the Bloch-Floquet theory for the spectral representation of
(2.1)
$$H_{\mathrm{per}}=\frac{1}{2}p^2+V(x).$$
The reciprocal lattice $`\mathrm{\Gamma }^{}`$ is defined as the lattice generated by the dual basis $`\{\gamma _1^{},\mathrm{},\gamma _d^{}\}`$ determined by $`\gamma _i\gamma _j^{}=2\pi \delta _{ij}`$, $`i,j=1,\mathrm{},d`$. The fundamental domain of $`\mathrm{\Gamma }`$ is denoted by $`M`$, the one of $`\mathrm{\Gamma }^{}`$ by $`M^{}`$. $`M^{}`$ is usually referred to as first Brillouin zone. If we identify opposite edges of $`M`$, resp. $`M^{}`$, then it becomes a flat $`d`$-torus denoted by $`𝕋=^d/\mathrm{\Gamma }`$, resp. $`𝕋^{}=^d/\mathrm{\Gamma }^{}`$.
Let us introduce the Bloch-Floquet transformation, which should be viewed as a discrete Fourier transform, through
$$(𝒰\psi )(k,x):=\underset{\gamma \mathrm{\Gamma }}{}\mathrm{e}^{i(x+\gamma )k}\psi (x+\gamma ),(k,x)^{2d},$$
for $`\psi 𝒮(^d)`$. Clearly,
(2.2)
$$\begin{array}{c}(𝒰\psi )(k,x^{}+\gamma )=(𝒰\psi )(k,x^{}),\hfill \\ \\ (𝒰\psi )(k^{}+\gamma ^{},x)=\mathrm{e}^{ix\gamma ^{}}(𝒰\psi )(k^{},x).\hfill \end{array}$$
Therefore it suffices to specify $`𝒰\psi `$ on the set $`M^{}\times M`$ and, if needed, extend it to all of $`^{2d}`$ by (2.2). The linear map $`𝒰:L^2(^d)𝒮(^d):=_M^{}^{}L^2(M)𝑑k,`$ with $`dk`$ the normalized Lebesgue measure on $`M^{}`$, has norm one and can thus be extended to all of $`L^2(^d)`$ by continuity. $`𝒰`$ is surjective as can be seen from the inverse mapping
$$(𝒰^1\varphi )(x):=_M^{}\mathrm{e}^{ixk}\varphi (k,x)𝑑k,$$
which has norm one. Thus $`𝒰:L^2(^d)`$ is unitary.
To transform $`H_{\mathrm{per}}`$ under $`𝒰`$, we first note that $`\stackrel{~}{p}=𝒰p𝒰^1=D_x+k`$, with $`D_x=i_x`$. Therefore
$$\stackrel{~}{H}_{\mathrm{per}}:=𝒰H_{\mathrm{per}}𝒰^1=_M^{}^{}H_{\mathrm{per}}(k)𝑑k,$$
and
$$H_{\mathrm{per}}(k)=\frac{1}{2}(D_x+k)^2+V(x),k^d.$$
$`H_{\mathrm{per}}(k)`$ acts on $`L^2(M)`$ with $`k`$-independent domain $`D:=H^2(𝕋)`$. $`\psi D`$ is periodic in $`x`$. $`H_{\mathrm{per}}(k)`$ is a semi-bounded self-adjoint operator, since by condition ($`\mathrm{C}_{\mathrm{per}}`$) $`V`$ is infinitesimally operator bounded with respect to $`\mathrm{\Delta }`$ . In particular, $`H_{\mathrm{per}}(k)`$ is an entire analytic family of type (B) in the sense of Kato for $`k^d`$. Since the resolvent of $`H_0(k)=\frac{1}{2}(D_x+k)^2`$ is compact, the resolvent $`R_\lambda (H_{\mathrm{per}}(k)):=(H_{\mathrm{per}}(k)\lambda )^1`$, $`\lambda \sigma (H_{\mathrm{per}}(k))`$, is also compact, and $`H_{\mathrm{per}}(k)`$ has a complete set of (normalized) eigenfunctions $`\phi _n(k)H^2(𝕋)`$, $`n`$, called Bloch functions. The corresponding eigenvalues $`E_n(k)`$, $`n`$, accumulate at infinity and we enumerate them according to their magnitude and multiplicity, $`E_1(k)E_2(k)\mathrm{}`$ . $`E_n(k)`$ is called the $`n`$-th band function. We note that $`H_{\mathrm{per}}(k)=\mathrm{e}^{ix\gamma ^{}}H_{\mathrm{per}}(k+\gamma ^{})\mathrm{e}^{ix\gamma ^{}}`$. Therefore $`E_n(k)`$ is periodic with respect to $`\mathrm{\Gamma }^{}`$. If $`E_{n1}(k)<E_n(k)<E_{n+1}(k)`$ for all $`kM^{}`$ (in particular $`E_n(k)`$ is nondegenerate), then the $`n`$-th band is isolated. In this case $`E_n`$ and the corresponding projection operator are real analytic functions as a consequence of analytic perturbation theory . We denote by $``$ the set of indices of isolated bands.
It will be convenient to have also a notation for the spectral subspaces. Let $`P_n(k):L^2(M)L^2(M)`$ denote the orthogonal projection onto the $`n`$-th eigenspace of $`H_{\mathrm{per}}(k)`$. Similarly, we set $`Q_n(k)=\mathrm{𝟏}P_n(k)`$. Their direct fiber integral is denoted by
$$\stackrel{~}{P}_n=_M^{}^{}P_n(k)𝑑k.$$
$`\stackrel{~}{P}_n`$ projects onto the $`n`$-th band subspace in $``$ and $`P_n=𝒰^1\stackrel{~}{P}_n𝒰`$ projects onto the $`n`$-th band subspace in $`L^2(^d)`$. We have
(2.3) $`(\stackrel{~}{P}_n\psi )(k,)`$ $`=`$ $`P_n(k)\psi (k,)=(\phi _n(k),\psi (k))_{L^2(M)}\phi _n(k,)`$
$`=`$ $`\psi _n(k)\phi _n(k,).`$
The coefficient functions $`\psi _nL^2(M^{})`$ and are called the Bloch coefficients in the $`n`$-th band subspace. For the index set $``$ of isolated bands we set $`P_{}=_nP_n`$.
###### Remark.
To have a concise notation, we will use a tilde for operators acting on $``$. Thus if $`A`$ is an operator on $`L^2(^d)`$, then $`\stackrel{~}{A}=𝒰A𝒰^1`$. If $`A`$ has a direct fiber decomposition, then $`\stackrel{~}{A}=_M^{}^{}A(k)𝑑k`$ with $`A(k)`$ acting on the fiber $`L^2(M)`$ of $``$.
## 3. Main results
For the potentials we assume $`(\mathrm{C}_{\mathrm{per}})`$ for $`V`$ and in addition
###### Condition ($`\mathrm{C}_{\mathrm{ex}}`$).
The external potential $`W𝒮(^d)`$.
To state the semiclassical limit, we first have to explain the classical dynamics which will serve as a comparison. For each $`n`$ the classical phase space is $`^d\times 𝕋^{}`$, where $`𝕋^{}=^d/\mathrm{\Gamma }^{}`$. As $`n`$-th band Hamiltonian we have
$$h_n(r,k)=E_n(k)+W(r),(r,k)^d\times 𝕋^{},$$
and the classical dynamics in the $`n`$-th band is governed by
(3.1)
$$\dot{r}_n=_kE_n(k_n),\dot{k}_n=_rW(r_n).$$
Since we want to prove the strong convergence of the position operator, as in the case $`V0`$, we have to lift (3.1) to operators on $``$. For this purpose we solve (3.1) with initial condition $`r_n(0)=0`$, $`k_n(0)=k`$. We denote the solution by $`(r_n(t;k),k_n(t;k))`$, regarded as functions of $`k𝕋^{}`$. For $`\psi `$, we define
$$(R(t)\psi )(k,x)=\underset{n}{}r_n(t;k)\stackrel{~}{P}_n\psi (k,x),$$
and analogously, for later use,
$$(K(t)\psi )(k,x)=\underset{n}{}k_n(t;k)\stackrel{~}{P}_n\psi (k,x).$$
###### Theorem 3.1.
Let the conditions $`(\mathrm{C}_{\mathrm{per}})`$, $`(\mathrm{C}_{\mathrm{ex}})`$ be satisfied. Let
$$x^\epsilon (t)=\epsilon U^\epsilon (t/\epsilon )xU^\epsilon (t/\epsilon ).$$
Then for every $`\psi \text{Ran}P_{}D(|x|)H^2`$, with $`H^2`$ the second Sobolev space,
$$\underset{\epsilon 0}{lim}x^\epsilon (t)\psi =𝒰^1R(t)𝒰\psi $$
strongly.
Theorem 3.1 will be proved in several steps. First we show that in the semiclassical limit transitions from and to isolated band subspaces are suppressed on the level of the unitary groups. We define $`H_{\mathrm{diag}}^n=P_nHP_n+Q_nHQ_n`$ and $`U_{\mathrm{diag}}^{\epsilon ,n}(t):=\mathrm{exp}(itH_{\mathrm{diag}}^n)`$. In Section 5 we will prove
###### Theorem 3.2.
For any $`n`$ we have
$$\underset{\epsilon 0}{lim}\left(U^\epsilon (t/\epsilon )U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )\right)=0$$
in $`B(H^1,L^2)`$, where $`H^1`$ is the first Sobolev space.
The position operator is not diagonal with respect to the $`n`$-th band subspace and we define its diagonal part by $`x_{\mathrm{diag}}^n=P_nxP_n+Q_nxQ_n`$ with the time evolution
$$x_{\mathrm{diag}}^{\epsilon ,n}(t):=\epsilon U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x_{\mathrm{diag}}^nU_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon ).$$
Our second step is to prove that the off-diagonal part of $`x^\epsilon (t)`$ vanishes in the limit $`\epsilon 0`$.
###### Theorem 3.3.
For $`n`$
(3.2)
$$\underset{\epsilon 0}{lim}\left(x^\epsilon (t)x_{\mathrm{diag}}^{\epsilon ,n}(t)\right)=0$$
in $`B(H^2,L^2)`$.
By construction we have $`[x_{\mathrm{diag}}^{\epsilon ,n}(t),P_n]=0`$ and it suffices to study the dynamics in the $`n`$-th band subspace. This subspace is isomorphic to $`L^2(𝕋^{})`$ and, up to errors of higher order, $`x_{\mathrm{diag}}^{\epsilon ,n}(t)`$ can be replaced by $`x_{\mathrm{sc}}^{\epsilon ,n}(t)`$ whose time evolution is governed by a Hamiltonian of the form
$$\stackrel{~}{H}_{\mathrm{sc}}^{\epsilon ,n}=E_n(k)+W(i\epsilon _k).$$
At this stage we can apply the standard machinery of semiclassics, except that formally the roles of position and momentum have been interchanged and the new position space is the flat torus rather than $`^d`$.
So far we focused on the position operator, since the electronic density is the most accessible quantity experimentally and it corresponds in essence to a suitable function of the position. On more general grounds one would like to characterize a wider class of semiclassical observables. One further obvious candidate is the momentum $`p`$. In the Bloch-Floquet basis we have $`\stackrel{~}{p}=k+D_x`$. $`k`$ is semiclassical, as being canonically conjugate to $`i_k`$:
###### Theorem 3.4.
Let
$$k^\epsilon (t)=U^\epsilon (t/\epsilon )𝒰^1k𝒰U^\epsilon (t/\epsilon ).$$
Then for every $`\psi \text{Ran}P_{}`$
(3.3)
$$\underset{\epsilon 0}{lim}k^\epsilon (t)\psi =𝒰^1K(t)𝒰\psi $$
strongly.
On the other hand, $`D_x`$ is unbalanced because there is no extra factor of $`\epsilon `$. Thus $`p(t/\epsilon )`$ has a limit only when averaged over time (compare with Section 6).
It is relatively easy to see (cf. Section 8) that Theorems 3.1 and 3.4 imply the semiclassical limit also for bounded functions of $`x^\epsilon (t)`$ resp. of $`k^\epsilon (t)`$ (cf. Lemma 8.1). Next note that for $`\mathrm{\Gamma }^{}`$-periodic functions $`g`$, $`g(+\gamma ^{})=g()`$ for all $`\gamma ^{}\mathrm{\Gamma }^{}`$, we have $`𝒰g(p)𝒰^1=g(k)`$ and hence, by the functional calculus for self-adjoint operators, $`g(p^\epsilon (t))=g(k^\epsilon (t))`$. Therefore we introduce the set $`𝒪(0)C(^d\times ^d,)`$ of bounded and continuous semiclassical symbols which vanish if the first argument approaches infinity and are $`\mathrm{\Gamma }^{}`$-periodic in their second argument. For $`a𝒪(0)`$ we introduce its Weyl quantization
(3.4)
$$(a^\mathrm{W}\psi )(x)=\frac{1}{(2\pi )^d}a(\frac{x+y}{2},\xi )e^{i(xy)\xi }\psi (y)𝑑\xi 𝑑y$$
as a bounded operator on $`L^2(^d)`$. The operator corresponding to the symbol $`a(\epsilon x,\xi )`$ will be denoted by $`a^{\mathrm{W},\epsilon }`$ and we set, as before,
(3.5)
$$a^{\mathrm{W},\epsilon }(t)=U^\epsilon (t/\epsilon )a^{\mathrm{W},\epsilon }U^\epsilon (t/\epsilon ).$$
###### Theorem 3.5.
Let the conditions $`(\mathrm{C}_{\mathrm{per}})`$, $`(\mathrm{C}_{\mathrm{ex}})`$ be satisfied and $`a𝒪(0)`$. Then for every $`\psi P_{}L^2`$ we have
$$\underset{\epsilon 0}{lim}a^{\mathrm{W},\epsilon }(t)\psi =𝒰^1a(R(t),K(t))𝒰\psi .$$
## 4. Semiclassical distributions
Theorems 3.1 and 3.5 tell us how the quantum distributions behave in the semiclassical limit. Let us first consider the initial $`\psi P_{}`$. Its scaled position distribution is $`\epsilon ^d|\psi (x/\epsilon )|^2`$ which converges to $`\delta (x)`$ as a measure. The quasimomentum distribution $`_n|\psi _n(k)|^2`$ is independent of $`\epsilon `$. Thus it is natural to chose
(4.1)
$$\rho (drdk)=\underset{n}{}\delta (r)|\psi _n(k)|^2drdk=\underset{n}{}\rho _n(drdk)$$
as initial distribution for the semiclassical flow (3.1). We could consider more general initial measures at the expense of making $`\psi `$ itself $`\epsilon `$-dependent. For example the shifted initial measure $`_n\delta (rr_0)|\psi _n(k)|drdk`$ is approximated by $`\psi (x\epsilon ^1r_0)`$. Under (3.1) $`\rho (drdk)`$ evolves to $`\rho (drdk,t)=_n\rho _n(drdk,t)`$. Each $`\rho _n`$ satisfies weakly the transport equation
(4.2)
$$\frac{}{t}\rho _n=E_n(k)_r\rho _n+V(r)_k\rho _n$$
with initial condition $`\rho _n(drdk,0)=\rho _n(drdk)`$. We define the position and quasimomentum marginals through
(4.3)
$$\rho (dr,t)=_M^{}\rho (drdk,t),\rho (dk,t)=_^d\rho (drdk,t).$$
To connect with the quantum evolution we consider the quantum mechanical position distribution
(4.4)
$$\rho ^\epsilon (dx,t)=\epsilon ^d|\psi (x/\epsilon ,t/\epsilon )|^2dx$$
as a probability measure on $`^d`$. From Theorem 3.1 and Lemma 8.1 we conclude that
(4.5)
$$\underset{\epsilon 0}{lim}f(x^\epsilon (t))\psi =𝒰^1f(R(t))𝒰\psi $$
for $`fC_{\mathrm{}}(^d)`$. In particular,
(4.6)
$$\underset{\epsilon 0}{lim}\rho ^\epsilon (dx,t)f(x)=\underset{\epsilon 0}{lim}(\psi ,f(x^\epsilon (t))\psi )=(𝒰\psi ,f(R(t))𝒰\psi )$$
and we only have to compute the expression on the right hand side. Using that
$$(𝒰\psi )(x,k)=\underset{n}{}\psi _n(k)\rho _n(x,k)$$
we have
(4.7)
$$(𝒰\psi ,f(R(t))𝒰\psi )=\underset{n}{}_M^{}|\psi _n(k)|^2f(r_n(t;k))𝑑k=\underset{n}{}_M\rho _n(dr,t)f(r).$$
Thus the positional distribution $`\rho ^\epsilon (dx,t)`$ converges weakly as a measure to the incoherent sum $`_n\rho _n(dr,t)`$.
By the same reasoning, if $`g`$ is a $`\mathrm{\Gamma }^{}`$-periodic function, then by Theorem 3.4 and Lemma 8.1
(4.8)
$$\underset{\epsilon 0}{lim}g(p(t/\epsilon ))\psi =𝒰^1g(K(t))𝒰\psi .$$
Therefore, if $`\rho ^\epsilon (k,t)dk`$ denotes the spectral measure for the quasimomentum operator at time $`t/\epsilon `$, we have
(4.9)
$$\underset{\epsilon 0}{lim}\rho ^\epsilon (k,t)dk=\underset{n}{}\rho _n(dk,t)$$
weakly as measures.
More generally for $`\psi L^2`$ we define the scaled Wigner function by
(4.10)
$$W^\epsilon (x,k,t)=\underset{\gamma \mathrm{\Gamma }}{}\epsilon ^d\psi (\epsilon ^1x\frac{1}{2}\gamma ,\epsilon ^1t)\psi ^{}(\epsilon ^1x+\frac{1}{2}\gamma ,\epsilon ^1t)e^{ik\gamma }$$
with $`x^d`$, $`kM^{}`$. We think of $`W^\epsilon `$ as a signed, bounded measure over $`^d\times M^{}`$. The Wigner function yields expectations of Weyl quantized operators through
(4.11)
$$(\psi ,e^{iH^\epsilon t/\epsilon }a^{\mathrm{W},\epsilon }e^{iH^\epsilon t/\epsilon }\psi )=_{^d\times M^{}}W^\epsilon (x,k,t)a(x,k)𝑑x𝑑k$$
with $`a`$ $`\mathrm{\Gamma }^{}`$-periodic in its second argument. From Theorem 3.5 we therefore deduce that
(4.12)
$$\underset{\epsilon 0}{lim}W^\epsilon (r,k,t)drdk=\rho (drdk,t)$$
weakly as measures. The limits (4.6) and (4.9) are the particular cases, where either $`a(x,k)=f(x)`$ or $`a(x,k)=g(k)`$.
## 5. Convergence of the unitary groups
By definition, the time evolution generated by $`H_{\mathrm{per}}`$ leaves invariant the band subspaces Ran($`P_n`$) for all $`n`$. However, $`W^\epsilon (x)=W(\epsilon x)`$ does not respect the Bloch decomposition and it will induce transitions between different bands. Since $`W^\epsilon `$ is of slow variation, we expect such transitions to have a small amplitude as stated in Theorem 3.2.
$`W^\epsilon `$ transforms under $`𝒰`$ as
(5.1) $`(𝒰W^\epsilon \psi )(k,x)`$ $`=`$ $`{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\mathrm{e}^{i(x+\gamma )k}W(\epsilon (x+\gamma ))\psi (x+\gamma )`$
$`=`$ $`{\displaystyle \underset{\gamma \mathrm{\Gamma }}{}}\mathrm{e}^{i(x+\gamma )k}(2\pi )^{d/2}{\displaystyle _^d}\widehat{W}(p)\mathrm{e}^{i\epsilon (x+\gamma )p}𝑑p\psi (x+\gamma )`$
$`=`$ $`(2\pi )^{d/2}{\displaystyle _^d}\widehat{W}(p)(𝒰\psi )(k\epsilon p,x)𝑑p`$
$`=`$ $`(2\pi )^{d/2}{\displaystyle _^d}\widehat{W}^\epsilon (p)(𝒰\psi )(kp,x)dp=:(\stackrel{~}{W}^\epsilon 𝒰\psi )(k,x),`$
where $`\widehat{W}^\epsilon (p)=\epsilon ^d\widehat{W}(p/\epsilon )`$ and we adopt the quasiperiodic extension (2.2). Since $`\widehat{W}𝒮(^d)`$, the integral (5.1) is well-defined and $`\stackrel{~}{W}^\epsilon =𝒰W^\epsilon 𝒰^1`$ acts on $``$ as convolution with $`\widehat{W}^\epsilon `$ in the fiber parameter $`k`$. $`\widehat{W}^\epsilon `$ approximates a Dirac delta in the limit $`\epsilon 0`$ and the shift in (5.1) becomes the identity operator.
In the Bloch-Floquet representation the full Hamiltonian (1.3) becomes
$$(\stackrel{~}{H}\psi )(k,)=H_{\mathrm{per}}(k)\psi (k,)+(\stackrel{~}{W}^\epsilon \psi )(k,).$$
We expect the diagonal part of $`W^\epsilon `$ to be dominant with the off-diagonal piece as a small correction. For such a decomposition it turns out to be convenient to fix the index $`n`$ of an isolated band and to project along $`P_n`$ and its complement $`Q_n=\mathrm{𝟏}P_n`$. For $`n`$ we define the diagonal part $`H_{\mathrm{diag}}^n`$ of $`H`$ as
$$H_{\mathrm{diag}}^n=P_nHP_n+Q_nHQ_n,$$
and the off-diagonal part of the external potential as
$$W_{\mathrm{od}}^{\epsilon ,n}=Q_nW^\epsilon P_n+P_nW^\epsilon Q_n.$$
Then
$$H=H_{\mathrm{diag}}^n+W_{\mathrm{od}}^{\epsilon ,n}=(H_{\mathrm{per}}+W_{\mathrm{diag}}^{\epsilon ,n})+W_{\mathrm{od}}^{\epsilon ,n}.$$
We note that $`W_{\mathrm{diag}}^{\epsilon ,n}`$ and $`W_{\mathrm{od}}^{\epsilon ,n}`$ are bounded operators and set
$$U^\epsilon (t)=\mathrm{e}^{itH},U_{\mathrm{diag}}^{\epsilon ,n}(t)=\mathrm{e}^{itH_{\mathrm{diag}}^n}.$$
To prove Theorem 3.2 we start by writing the difference of the two unitary groups in the Bloch representation as
(5.2)
$$\stackrel{~}{U}^\epsilon (t/\epsilon )\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )=i\epsilon _0^{t/\epsilon }\stackrel{~}{U}^\epsilon (\epsilon ^1ts)\left(\epsilon ^1\stackrel{~}{W}_{\mathrm{od}}^{\epsilon ,n}\right)\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)𝑑s$$
and we have to investigate the operator $`\stackrel{~}{W}_{\mathrm{od}}^{\epsilon ,n}`$. By definition, for $`\psi `$, we have
$$(\stackrel{~}{Q}_n\stackrel{~}{W}^\epsilon \stackrel{~}{P}_n\psi )(k)=(2\pi )^{d/2}_^d\widehat{W}^\epsilon (p)Q_n(k)P_n(kp)\psi (kp)𝑑p,$$
which vanishes strongly in the limit $`\epsilon 0`$, since $`\widehat{W}^\epsilon `$ localizes around $`p=0`$. To control the long times in (5.2) we need uniform convergence of order $`o(\epsilon )`$, however. To have a more detailed information on $`W_{\mathrm{od}}^{\epsilon ,n}`$ we Taylor expand of $`P_n(kp)`$ around $`P_n(k)`$, leading, as we will show, to
(5.3)
$$(\stackrel{~}{Q}_n\stackrel{~}{W}^\epsilon \stackrel{~}{P}_n\psi )(k)=\epsilon (2\pi )^{d/2}_^d\widehat{F}^\epsilon (p)Q_n(k)_kP_n(k)\psi (kp)𝑑p+o(\epsilon ).$$
Here $`\widehat{F}^\epsilon (p):=\widehat{W}^\epsilon (p)\frac{p}{\epsilon }`$ is the Fourier transform of $`F^\epsilon (x)=(D_xW)(\epsilon x)`$ and we will associate to $`\widehat{F}^\epsilon `$ the operator $`\stackrel{~}{F}^\epsilon `$ as in the case of $`\widehat{W}^\epsilon `$,
$$(\stackrel{~}{F}^\epsilon \psi )(k)=(2\pi )^{d/2}_^d\widehat{F}^\epsilon (p)\psi (kp)𝑑p.$$
To justify (5.3) we first need to calculate $`_kP_n(k)`$.
###### Lemma 5.1.
Let $`n`$. Then
(5.4) $`_kP_n(k)`$ $`=`$ $`Q_n(k)R_{E_n(k)}(H_{\mathrm{per}}(k))(D_x+k)P_n(k)`$
$`P_n(k)(D_x+k)R_{E_n(k)}(H_{\mathrm{per}}(k))Q_n(k),`$
where $`R_\lambda (H)=(H\lambda )^1`$ is the resolvent of $`H`$. Thus $`P_n()C^{\mathrm{}}(M^{};B(L^2(M)))`$.
###### Proof.
Using contour integrals we write
$$_kP_n(k)=\frac{1}{2\pi i}_{c_n(k)}_kR_\lambda (H_{\mathrm{per}}(k))𝑑\lambda ,$$
where $`c_n(k)`$ is a closed rectifiable curve in the complex spectral plane which encircles $`E_n(k)`$ only. From
$`0`$ $`=`$ $`_k\mathrm{𝟏}=_k(H_{\mathrm{per}}(k)\lambda )R_\lambda (H_{\mathrm{per}}(k))`$
$`=`$ $`(D_x+k)R_\lambda (H_{\mathrm{per}}(k))+(H_{\mathrm{per}}(k)\lambda )_kR_\lambda (H_{\mathrm{per}}(k)),`$
we infer
$$_kR_\lambda (H_{\mathrm{per}}(k))=R_\lambda (H_{\mathrm{per}}(k))(D_x+k)R_\lambda (H_{\mathrm{per}}(k)).$$
Hence we get
(5.5) $`Q_n(k)_kP_n(k)=Q_n(k)_kP_n(k)(P_n(k)+Q_n(k))`$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{c_n(k)}}Q_n(k)R_\lambda (H_{\mathrm{per}}(k))(D_x+k)R_\lambda (H_{\mathrm{per}}(k))P_n(k)𝑑\lambda `$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{c_n(k)}}R_\lambda (H_{\mathrm{per}}(k))Q_n(k){\displaystyle \frac{1}{E_n(k)\lambda }}𝑑\lambda (D_x+k)P_n(k)`$
$`=`$ $`R_{E_n(k)}(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k),`$
where the term $`Q_n(k)_kP_n(k)Q_n(k)`$ vanishes, since in this case the integrand is an analytic function on the whole interior of $`c_n(k)`$. Note that $`P_n(k)`$ projects onto a subspace of finite energy, on which $`D_x+k`$ is bounded. The statement about continuity for this term then follows from the continuity of $`P_n(k)`$, $`E_n(k)`$ and the assumption that $`E_n(k)`$ is isolated from the remainder of the spectrum. An analogous computation for $`P_n(k)_kP_n(k)`$ leads to the second term in (5.4).
Finally, $`P_n()C^{\mathrm{}}(M^{};B(L^2(M)))`$ follows by induction. ∎
From $`Q_n(k)+P_n(k)=\mathrm{𝟏}`$ we conclude that $`Q_n(k)`$ is differentiable as well and that $`_kQ_n(k)=_kP_n(k)`$.
###### Lemma 5.2.
Let $`n`$. Then
$$\stackrel{~}{W}_{\mathrm{od}}^{\epsilon ,n}=\epsilon \left(\stackrel{~}{Q}_n_k\stackrel{~}{P}_n+\stackrel{~}{P}_n_k\stackrel{~}{Q}_n\right)\stackrel{~}{F}^\epsilon +o(\epsilon )$$
in $`B(,)`$, where $`_k\stackrel{~}{P}_n:=_M^{}^{}_kP_n(k)𝑑k`$.
###### Proof.
We will treat only the $`\stackrel{~}{Q}_n\stackrel{~}{W}^\epsilon \stackrel{~}{P}_n`$ part of $`\stackrel{~}{W}_{\mathrm{od}}^{\epsilon ,n}`$ explicitly, since the argument for the second part is analogous.
Let $`\psi `$. By Lemma 5.1 we are allowed to write the following well-defined identity, setting $`e_p=p/|p|`$,
(5.7) $`{\displaystyle \widehat{W}^\epsilon (p)P_n(kp)\psi (kp)𝑑p}`$
$`=`$ $`{\displaystyle \widehat{W}^\epsilon (p)|p|\left(|p|^1(P_n(kp)P_n(k))+e_p_kP_n(k)\right)\psi (kp)𝑑p}`$
$`+P_n(k){\displaystyle \widehat{W}^\epsilon (p)\psi (kp)𝑑p}{\displaystyle \widehat{W}^\epsilon (p)p_kP_n(k)\psi (kp)𝑑p}`$
$`=`$ $`\epsilon {\displaystyle \widehat{W}^\epsilon (p)\frac{|p|}{\epsilon }\left(\frac{P_n(kp)P_n(k)}{|p|}+e_p_kP_n(k)\right)\psi (kp)𝑑p}`$
$`+(2\pi )^{d/2}P_n(k)(\stackrel{~}{W}^\epsilon \psi )(k)\epsilon (2\pi )^{d/2}_kP_n(k)(\stackrel{~}{F}^\epsilon \psi )(k).`$
If in (5.7), (5.7) we apply $`Q_n(k)`$ from the left, the first term of (5.7) vanishes and it remains to show that (5.7), divided by $`\epsilon `$, tends to zero uniformly for all $`\psi `$.
We split the integral into two parts. Let $`R>0`$ be arbitrary, $`B_R=\{p|p|R\}`$. We start with
$`{\displaystyle \underset{M^{}}{}}{\displaystyle \underset{B_R}{}}\widehat{W}^\epsilon (p){\displaystyle \frac{|p|}{\epsilon }}({\displaystyle \frac{P_n(kp)P_n(k)}{|p|}}`$
$`+e_p_kP_n(k))\psi (kp)dp_{L^2(M)}dk`$
$``$ $`\underset{kM^{}}{sup}\underset{pB_R}{sup}|p|^1\left(P_n(kp)P_n(k)\right)+e_p_kP_n(k)`$
$`\times {\displaystyle _{B_R}}\left|\widehat{W}^\epsilon (p)\right|{\displaystyle \frac{|p|}{\epsilon }}{\displaystyle _M^{}}\psi (kp)_{L^2(M)}dkdp`$
$``$ $`\psi _{}\widehat{F}^\epsilon _{L^1}`$
$`\times \underset{kM^{},pB_R}{sup}|p|^1(P_n(kp)P_n(k))+e_p_kP_n(k).`$
Since $`\widehat{F}^\epsilon _{L^1}`$ does not depend on $`\epsilon `$ and since the difference quotient approaches the derivative uniformly on the compact domain $`M^{}`$, the $``$-norm of the first part tends to zero uniformly. For the remaining part we have
$`{\displaystyle \underset{M^{}}{}}{\displaystyle \underset{|p|>R}{}}\widehat{W}^\epsilon (p){\displaystyle \frac{|p|}{\epsilon }}({\displaystyle \frac{P_n(kp)P_n(k)}{|p|}}`$
$`+e_p_kP_n(k))\psi (kp)dp_{L^2(M)}dk`$
$``$ $`\psi _{}\widehat{F}^\epsilon _{L^1(B_R^c)}`$
$`\times \underset{kM^{},p^d}{sup}|p|^1(P_n(kp)P_n(k))+e_p_kP_n(k),`$
which tends to zero uniformly as $`\epsilon 0`$, since $`\widehat{F}^\epsilon _{L^1(B_R^c)}0`$ for any fixed $`R>0`$. ∎
As a consequence of Lemma 5.2 the difference of the two unitary groups in Eq. (5.2) can be written as
(5.8) $`\stackrel{~}{U}^\epsilon (t/\epsilon )\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )`$
$`=i\epsilon {\displaystyle _0^{t/\epsilon }}\stackrel{~}{U}^\epsilon (\epsilon ^1ts)\left(\stackrel{~}{Q}_n_k\stackrel{~}{P}_n+\stackrel{~}{P}_n_k\stackrel{~}{Q}_n\right)\stackrel{~}{F}^\epsilon \stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)𝑑s+o(1).`$
We have to estimate the integral without losing one order of $`\epsilon `$ from the integration over time. As in the proof in of the adiabatic theorem the idea is to rewrite the integrand as a time derivative, i.e. as a commutator of $`\stackrel{~}{H}_{\mathrm{diag}}^n`$ with an appropriately chosen operator $`A`$, at least up to an unavoidable error $`o(1)`$.
Let us define for $`n`$
$$B_n(k)=R_{E_n(k)}^2(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k).$$
###### Lemma 5.3.
For $`n`$ we have
$$\stackrel{~}{Q}_n_k\stackrel{~}{P}_n+\stackrel{~}{P}_n_k\stackrel{~}{Q}_n=[\stackrel{~}{B}_n+\stackrel{~}{B}_n^{},\stackrel{~}{H}_{\mathrm{per}}].$$
###### Proof.
Using the spectral decomposition and recalling
$$Q_n(k)_kP_n(k)=R_{E_n(k)}(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k)$$
from Lemma 5.1, one directly computes
$`B_n(k)H_{\mathrm{per}}(k)H_{\mathrm{per}}(k)B_n(k)`$
$`=`$ $`(H_{\mathrm{per}}(k)E_n(k))R_{E_n(k)}^2(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k)`$
$`=`$ $`R_{E_n(k)}(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k)`$
$`=`$ $`Q_n(k)_kP_n(k).`$
The lemma then follows from $`\stackrel{~}{P}_n_k\stackrel{~}{Q}_n=(\stackrel{~}{Q}_n_k\stackrel{~}{P}_n)^{}`$. ∎
###### Lemma 5.4.
$`[B_n+B_n^{},\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n}]0`$ in $`B(,)`$ as $`\epsilon `$ tends to zero.
###### Proof.
To have a concise notation in the following, expressions like $`\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n}P_n(k)`$ are understood in the sense that $`\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n}`$ acts on all $`k`$-depending objects on its right hand side. We recall that $`\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n}=\stackrel{~}{P}_n\stackrel{~}{W}^\epsilon \stackrel{~}{P}_n+\stackrel{~}{Q}_n\stackrel{~}{W}^\epsilon \stackrel{~}{Q}_n`$. Hence
$`[B_n(k),\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n}]`$
$`=`$ $`Q_n(k)[R_{E_n(k)}^2(H_{\mathrm{per}}(k))Q_n(k)(D_x+k)P_n(k),\stackrel{~}{W}^\epsilon ]P_n(k).`$
We now examine the commutators $`[P_n(k),\stackrel{~}{W}^\epsilon ]`$, $`[D_x+k,\stackrel{~}{W}^\epsilon ]`$ and $`[R_{E_n(k)}^2Q_n(k),\stackrel{~}{W}^\epsilon ]`$ one by one. It follows from the proof of Lemma 5.2 that $`[P_n(k),\stackrel{~}{W}^\epsilon ]`$ vanishes as $`\epsilon 0`$ and the analogous statement for $`[R_{E_n(k)}^2Q_n(k),\stackrel{~}{W}^\epsilon ]`$ can be shown to hold by a similar argument. Thus it remains to discuss the commutator $`[D_x+k,\stackrel{~}{W}^\epsilon ]`$. For $`\psi H^1(^d)`$ we compute
$`(2\pi )^{d/2}([D_x+k,\stackrel{~}{W}^\epsilon ]𝒰\psi )(k)`$
$`=`$ $`{\displaystyle _^d}\widehat{W}^\epsilon (p)(((D_x+k)(D_x+kp))𝒰\psi )(kp)𝑑p`$
$`=`$ $`\epsilon {\displaystyle _^d}\widehat{W}^\epsilon (p)\epsilon ^1p(𝒰\psi )(kp)𝑑p`$
$`=`$ $`\epsilon (\stackrel{~}{F}^\epsilon 𝒰\psi )(k),`$
which clearly vanishes uniformly for $`\psi L^2`$ as $`\epsilon 0`$, since $`F𝒮(^d,^d)`$. ∎
In summary we have shown that
$$\left(\stackrel{~}{Q}_n_k\stackrel{~}{P}_n+\stackrel{~}{P}_n_k\stackrel{~}{Q}_n\right)\stackrel{~}{F}^\epsilon =\left([\stackrel{~}{B}_n+\stackrel{~}{B}_n^{},\stackrel{~}{H}_{\mathrm{diag}}^n]+o(1)\right)\stackrel{~}{F}^\epsilon ,$$
and it remains to check
###### Lemma 5.5.
$`[\stackrel{~}{H}_{\mathrm{diag}}^n,\stackrel{~}{F}^\epsilon ]0`$ in $`B(𝒰H^1,)`$ as $`\epsilon `$ tends to zero.
###### Proof.
The commutator
$$[H_{\mathrm{per}},F^\epsilon ]=\frac{1}{2}\epsilon ^2(\mathrm{\Delta }F^\epsilon )\frac{1}{2}\epsilon (F^\epsilon )\frac{1}{2}\epsilon (F^\epsilon )$$
vanishes in $`B(H^1,L^2)`$ as $`\epsilon 0`$. The commutator $`[\stackrel{~}{W}_{\mathrm{diag}}^{\epsilon ,n},\stackrel{~}{F}^\epsilon ]`$ vanishes in $`B(,)`$, since the commutator of $`\stackrel{~}{P}_n`$ and $`\stackrel{~}{Q}_n`$ with $`\stackrel{~}{F}^\epsilon `$ are both of uniform order $`o(1)`$ (in $`B(,)`$) and $`[\stackrel{~}{W}^\epsilon ,\stackrel{~}{F}^\epsilon ]`$ vanishes identically. ∎
Defining
$$\stackrel{~}{A}_n=\left(\stackrel{~}{B}_n+\stackrel{~}{B}_n^{}\right)\stackrel{~}{F}^\epsilon ,$$
it follows that the integrand in (5.8) can be written as
$$\left(\stackrel{~}{Q}_n_k\stackrel{~}{P}_n+\stackrel{~}{P}_n_k\stackrel{~}{Q}_n\right)\stackrel{~}{F}^\epsilon =[\stackrel{~}{A}_n,\stackrel{~}{H}_{\mathrm{diag}}^n]+o(1),$$
where $`o(1)`$ is in the norm of $`B(𝒰H^1,)`$. (Note that for $`A^\epsilon B(L^2,L^2)`$, $`lim_{\epsilon 0}`$ $`A^\epsilon =0`$ in $`B(L^2,L^2)`$ implies, in particular, that also $`lim_{\epsilon 0}A^\epsilon =0`$ in $`B(H^1,L^2)`$).
We are now ready for the
###### Proof of Theorem 3.2.
Since $`U_{\mathrm{diag}}^{\epsilon ,n}(t):H^1H^1`$ is bounded uniformly in $`t`$ and $`\epsilon `$ (cf. Section 6), we obtain for the difference (5.8) of the unitary groups,
(5.9) $`\left(\stackrel{~}{U}^\epsilon (t/\epsilon )\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )\right)`$
$`=i\epsilon {\displaystyle _0^{t/\epsilon }}\stackrel{~}{U}^\epsilon (\epsilon ^1ts)[\stackrel{~}{A}_n,\stackrel{~}{H}_{\mathrm{diag}}^n]\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)𝑑s+o(1).`$
Abbreviating $`X^n(s)=\stackrel{~}{U}^\epsilon (s)\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)`$ and $`\stackrel{~}{A}_n(s)=\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)\stackrel{~}{A}_n\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)`$, we get, using partial integration in (5.9),
$`i\epsilon {\displaystyle _0^{t/\epsilon }}\stackrel{~}{U}^\epsilon (t/\epsilon )X^n(s)\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)[\stackrel{~}{A}_n,\stackrel{~}{H}_{\mathrm{diag}}^n]\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)𝑑s`$
$`=`$ $`\epsilon \stackrel{~}{U}^\epsilon (t/\epsilon ){\displaystyle _0^{t/\epsilon }}X^n(s)\left({\displaystyle \frac{d}{ds}}\stackrel{~}{A}_n(s)\right)𝑑s`$
$`=`$ $`\epsilon \left(\stackrel{~}{A}_n\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )\stackrel{~}{U}^\epsilon (t/\epsilon )\stackrel{~}{A}_n\right)`$
$`\epsilon \stackrel{~}{U}^\epsilon (t/\epsilon ){\displaystyle _0^{t/\epsilon }}\left({\displaystyle \frac{d}{ds}}X^n(s)\right)\stackrel{~}{A}_n(s)𝑑s`$
$`=`$ $`\epsilon \left(\stackrel{~}{A}_n\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )\stackrel{~}{U}^\epsilon (t/\epsilon )\stackrel{~}{A}_n\right)`$
$`i\epsilon \stackrel{~}{U}^\epsilon (t/\epsilon ){\displaystyle _0^{t/\epsilon }}\stackrel{~}{U}^\epsilon (s)\stackrel{~}{W}_{\mathrm{od}}^{\epsilon ,n}\stackrel{~}{A}_n\stackrel{~}{U}_{\mathrm{diag}}^{\epsilon ,n}(s)𝑑s.`$
For $`\epsilon 0`$ the first term vanishes since $`\stackrel{~}{A}_n`$ is bounded and the second term vanishes, since $`W_{\mathrm{od}}^{\epsilon ,n}`$ tends to zero uniformly according to Lemma 5.2. ∎
## 6. Convergence of the position operator
In this section we will study the asymptotics of the position operator $`x^\epsilon (t)`$. As in the case of the unitaries we have to establish that the off-diagonal contributions to $`x^\epsilon (t)`$ vanish in the limit $`\epsilon 0`$.
###### Proof of Theorem 3.3.
Let $`\psi D(|x|)H^2`$ and $`n`$. Then
(6.2) $`\left(x^\epsilon (t)x_{\mathrm{diag}}^{\epsilon ,n}(t)\right)\psi `$
$``$ $`\left(x^\epsilon (t)U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x^\epsilon U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )\right)\psi `$
$`+\left(U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x^\epsilon U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x_{\mathrm{diag}}^{\epsilon ,n}(t)\right)\psi .`$
In order to estimate (6.2), note that we have
(6.3)
$$x^\epsilon (t)\psi =\epsilon x\psi +\epsilon _0^{t/\epsilon }U^\epsilon (s)D_xU^\epsilon (s)\psi 𝑑s$$
and
$`U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x^\epsilon U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )`$
$`=`$ $`\epsilon x\psi +\epsilon {\displaystyle _0^{t/\epsilon }}U_{\mathrm{diag}}^{\epsilon ,n}(s)\left(D_x+i[W_{\mathrm{diag}}^{\epsilon ,n},x]\right)U_{\mathrm{diag}}^{\epsilon ,n}(s)\psi 𝑑s`$
$`=`$ $`\epsilon x\psi +\epsilon {\displaystyle _0^{t/\epsilon }}U_{\mathrm{diag}}^{\epsilon ,n}(s)D_xU_{\mathrm{diag}}^{\epsilon ,n}(s)\psi 𝑑s+o(1).`$
The last equality holds, since $`[W_{\mathrm{diag}}^{\epsilon ,n},x]=o(1)`$ in $`B(L^2)`$, as follows immediately from the fact that $`[W^\epsilon ,P_n]=o(1)`$ and $`[W^\epsilon ,Q_n]=o(1)`$, cf. proof of Lemma 5.2. Hence, using (6.3), the remaining term from (6.2) is
(6.5) $`\epsilon {\displaystyle _0^{t/\epsilon }}\left(U^\epsilon (s)D_xU^\epsilon (s)U_{\mathrm{diag}}^{\epsilon ,n}(s)D_xU_{\mathrm{diag}}^{\epsilon ,n}(s)\right)\psi 𝑑s`$
$`=`$ $`{\displaystyle _0^t}\left(U^\epsilon (s/\epsilon )U_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon )\right)D_xU_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon )\psi 𝑑s`$
$`+{\displaystyle _0^t}U^\epsilon (s/\epsilon )D_x\left(U^\epsilon (s/\epsilon )U_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon )\right)\psi 𝑑s.`$
Using the fact that $`V`$ and $`W`$ are infinitesimally operator bounded with respect to $`\frac{1}{2}\mathrm{\Delta }`$ and that $`\psi H^2`$, we get for $`\psi (s):=U_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon )\psi `$
$`D_x^2\psi (s)`$ $``$ $`H_{\mathrm{diag}}^{\epsilon ,n}\psi (s)+(V+W_{\mathrm{diag}}^{\epsilon ,n})\psi (s)`$
$``$ $`H_{\mathrm{diag}}^{\epsilon ,n}\psi +c_1D_x^2\psi (s)+c_2\psi ,`$
with $`c_1<\frac{1}{2}`$ and $`c_2<\mathrm{}`$. Hence $`D_xU_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon )\psi _{H^1}c\psi _{H^2}`$ with $`c`$ independent of $`s`$ and $`\epsilon `$ and we can apply Theorem 3.2 to conclude that the operator acting on $`\psi `$ in (6.5) vanishes in $`B(H^2,L^2)`$ as $`\epsilon 0`$.
We come to (6.5). Let $`\psi (s)=(U^\epsilon (s/\epsilon )U_{\mathrm{diag}}^{\epsilon ,n}(s/\epsilon ))\psi `$, then, by Cauchy-Schwarz,
$$D_x\psi (s)^2=(\psi (s),D_x^2\psi (s))\psi (s)D_x^2\psi (s).$$
The first factor tends to zero by Theorem 3.2 whereas the second is uniformly bounded by the same argument as in the treatment of (6.5) a few lines above.
Next we rewrite (6.2) as
$$\epsilon U_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )x_{\mathrm{od}}^nU_{\mathrm{diag}}^{\epsilon ,n}(t/\epsilon )$$
with $`x_{\mathrm{od}}^n:=Q_nxP_n+P_nxQ_n`$. This certainly vanishes as $`\epsilon 0`$ if $`x_{\mathrm{od}}^n`$ can be shown to be a bounded operator. To see this, note that in Bloch representation $`x`$ acts as $`i_k`$. Hence
$$(𝒰Q_nxP_n\psi )(k)=iQ_n(k)_kP_n(k)(𝒰\psi )(k)=iQ_n(k)(_kP_n(k))(𝒰\psi )(k)$$
and thus $`Q_nxP_n=\stackrel{~}{Q}_n_k\stackrel{~}{P}_n`$. Finally also $`P_nxQ_n`$ is bounded, since it is the adjoint of $`Q_nxP_n`$. ∎
## 7. Semiclassical equations of motion for the position operator
As we have shown, on the macroscopic scale the position and quasimomentum operators commute with the projection on isolated bands. Thus it remains to investigate the semiclassical limit for each isolated band separately. For this purpose we note that any $`\psi \stackrel{~}{P}_n`$ is of the form $`\psi _n(k)\phi _n(x,k)`$ with $`\psi _nL^2(M^{})`$. Since $`\phi _n`$ already satisfies (2.2), we have to extend the Bloch coefficients periodically. We determine now how $`H_{\mathrm{diag}}^{\epsilon ,n}`$ acts on $`L^2(M^{})`$. We have $`[H_{\mathrm{per}},\stackrel{~}{P}_n]=0`$ and therefore $`H_{\mathrm{per}}`$ acts as multiplication by $`E_n(k)`$. For $`W_{\mathrm{diag}}^{\epsilon ,n}`$ we have
(7.1) $`\left(\stackrel{~}{P}_n\stackrel{~}{W}^\epsilon \stackrel{~}{P}_n𝒰\psi \right)(k,x)`$
$`=`$ $`(2\pi )^{d/2}{\displaystyle _^d}\widehat{W}^\epsilon (p)(\phi _n(k),\phi _n(kp))_{L^2(M)}\psi _n(kp)𝑑p\phi _n(k,x)`$
$`=:`$ $`(\stackrel{~}{W}^{\epsilon ,n}\psi _n)(k)\phi _n(x,k).`$
Thus $`H_{\mathrm{diag}}^{\epsilon ,n}`$ restricted to $`\stackrel{~}{P}_n`$ is unitarily equivalent to $`H^{\epsilon ,n}:=E_n(k)+\stackrel{~}{W}^{\epsilon ,n}`$.
To be able to use techniques from semiclassics we next approximate $`\stackrel{~}{W}^{\epsilon ,n}`$ by the operator $`\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}=W(i\epsilon _k)`$, where $`_k`$ is understood with periodic boundary conditions on $`^d/\mathrm{\Gamma }^{}`$.
###### Lemma 7.1.
For any $`n`$
(7.2)
$$\stackrel{~}{W}^{\epsilon ,n}=\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}+o(\epsilon )$$
in $`B(L^2(M^{}))`$.
###### Proof.
By definition we have
$$\left(\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}\psi \right)(k)=(2\pi )^{d/2}_^d\widehat{W}^\epsilon (p)\psi _n(kp)𝑑p,$$
and therefore
(7.3) $`\left(\left(\stackrel{~}{W}^{\epsilon ,n}\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}\right)\psi \right)(k)=`$
$`=(2\pi )^{d/2}{\displaystyle _^d}\widehat{W}^{\epsilon ,n}(p)\left((\phi _n(k),\phi _n(kp))_{L^2(M)}1\right)\psi (kp)𝑑p.`$
As to be shown, there exists a constant $`c`$ such that
(7.4)
$$\left|(\phi _n(k),\phi _n(kp))_{L^2(M)}1\right|c|p|^2$$
for Lebesgue almost all $`k`$. Therefore we conclude
$`\left(\stackrel{~}{W}^{\epsilon ,n}\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}\right)\psi _{L^2(M^{})}`$ $``$ $`c\epsilon ^2{\displaystyle \left|\widehat{W}^\epsilon (p)\frac{|p|^2}{\epsilon ^2}\right||\psi (kp)|𝑑p}_{L^2(M^{})}`$
$``$ $`c^{}\epsilon ^2\psi _{L^2(M^{})}.`$
To show (7.4) note that one can chose $`\phi _n(k)`$ such that the map $`k\phi _n(k)L^2(M)`$ is smooth Lebesgue almost everywhere. This is because according to Lemma 5.1 the projections $`P_n(k)`$ depend smoothly on $`k`$ and hence one can locally define $`\phi _n(k)=P_n(k)\phi _n(k_0)/P_n(k)\phi _n(k_0)`$. Now we can cover $`M^{}`$ by finitely many open disjoint sets $`U_i`$ such that $`M^{}_iU_i`$ is a set of Lebesgue measure zero and $`\phi _n(k)`$ can be defined on the closure of each $`U_i`$ in the way described above. One obtains a family $`\phi _n(k)`$ of eigenfunctions which is smooth except at the boundaries between the sets, where we pick $`\phi _n(k)`$ with an arbitrary phase. Wherever $`\phi _n(k)`$ is smooth, Taylor expansion yields $`\phi _n(kp)=\phi _n(k)p_k\phi _n(k)+\frac{1}{2}p(\phi _n)(k^{})p`$, where $`(\phi _n)`$ denotes the Hessian and $`\frac{1}{2}p(\phi _n)(k^{})p`$ is the Lagrangian remainder. In view of $`(\phi _n(k),_k\phi _n(k))_{L^2(M)}=0`$, which follows from comparing (5.4) with
$$(_kP_n\psi )(k)=(\phi _n(k),\psi (,k))_k\phi _n(k)+(_k\phi _n(k),\psi (,k))\phi _n(k),$$
we obtain
$$\left|(\phi _n(k),\phi _n(kp))_{L^2(M)}1\right|c(k)|p|^2.$$
Here $`c(k)=\frac{1}{2}_{i,j}|(\phi _n(k^{}),_{k_i}_{k_j}\phi _n(k^{}))|`$. However, $`c(k)`$ is bounded uniformly in $`k`$, since $`\phi _n(k)`$ is smooth on each compact $`\overline{U_i}`$. ∎
We define now the semiclassical Hamiltonian $`H_{\mathrm{sc}}^{\epsilon ,n}`$
(7.5)
$$H_{\mathrm{sc}}^{\epsilon ,n}=E_n(k)+W(i\epsilon _k)$$
acting on $`L^2(M^{})`$. Then Lemma 7.1 shows that the difference $`H^{\epsilon ,n}H_{\mathrm{sc}}^{\epsilon ,n}`$ is of order $`o(\epsilon )`$ uniformly in $`B(L^2(M^{}))`$ and hence (cf. Section 5) the difference of the corresponding unitary groups approaches zero as $`\epsilon 0`$.
###### Corollary 7.2.
Let $`U_{\mathrm{sc}}^{\epsilon ,n}(t)=\mathrm{e}^{itH_{\mathrm{sc}}^{\epsilon ,n}}`$ and $`U^{\epsilon ,n}(t)=\mathrm{e}^{itH^{\epsilon ,n}}`$, then
$$\underset{\epsilon 0}{lim}\left(U^{\epsilon ,n}(t/\epsilon )U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )\right)=0$$
in $`B(L^2(M^{}))`$.
The semiclassical limit for $`U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )`$ on $`L^2(𝕋^{})`$ is well studied. We refer to . As a consequence the strong limits
(7.6) $`\underset{\epsilon 0}{lim}U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )(i\epsilon _k)U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )`$ $`=`$ $`r_n(t;k),`$
(7.7) $`\underset{\epsilon 0}{lim}U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )kU_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )`$ $`=`$ $`k_n(t;k)`$
exist on $`H^1(𝕋^{})`$. $`r_n`$ and $`k_n`$ act as multiplication operators and are defined as in (3.1) with initial conditions $`(r_n(0),k_n(0))=(0,k)`$.
Since the restriction of $`\epsilon x_{\mathrm{diag}}^n`$ to the $`n`$-th band subspace is unitarily equivalent to $`i\epsilon _k`$ on $`L^2(𝕋^{})`$, we can, in view of Theorem 3.3, conclude the proof of Theorem 3.1 by showing
###### Lemma 7.3.
In $`B(L^2(M^{}))`$ we have
(7.8)
$$\underset{\epsilon 0}{lim}\left(U^{\epsilon ,n}(t/\epsilon )(i\epsilon _k)U^{\epsilon ,n}(t/\epsilon )U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )(i\epsilon _k)U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )\right)=0.$$
###### Proof.
The proof of (7.8) is analogous to the proof of Theorem 3.3 in Section 5, however, simpler. As in (6.3) we have
$$U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )(i\epsilon _k)U_{\mathrm{sc}}^{\epsilon ,n}(t/\epsilon )=i\epsilon _k+\epsilon _0^{t/\epsilon }U_{\mathrm{sc}}^{\epsilon ,n}(s)[i_k,H_{\mathrm{sc}}^{\epsilon ,n}]U_{\mathrm{sc}}^{\epsilon ,n}(s)𝑑s$$
and
$`U^{\epsilon ,n}(t/\epsilon )(i\epsilon _k)U^{\epsilon ,n}(t/\epsilon )=`$
$`=`$ $`i\epsilon _k+\epsilon {\displaystyle _0^{t/\epsilon }}U^{\epsilon ,n}(s)[i_k,H^{\epsilon ,n}]U^{\epsilon ,n}(s)𝑑s`$
$`=`$ $`i\epsilon _k+\epsilon {\displaystyle _0^{t/\epsilon }}U^{\epsilon ,n}(s)\left([i_k,H_{\mathrm{sc}}^{\epsilon ,n}]+[i_k,\mathrm{\Delta }\stackrel{~}{W}^{\epsilon ,n}]\right)U^{\epsilon ,n}(s)𝑑s,`$
where $`\mathrm{\Delta }\stackrel{~}{W}^{\epsilon ,n}:=\stackrel{~}{W}^{\epsilon ,n}\stackrel{~}{W}_{\mathrm{sc}}^{\epsilon ,n}`$. Now $`[i_k,H_{\mathrm{sc}}^{\epsilon ,n}]=i_kE_n(k)`$ is bounded, and (7.8) follows from Corollary 7.2 if we can show that $`[i_k,\mathrm{\Delta }\stackrel{~}{W}^{\epsilon ,n}]=o(1)`$ in $`B(L^2(M^{}))`$. Noting that $`(\mathrm{\Delta }\stackrel{~}{W}^{\epsilon ,n}\psi )(k)`$ is given by (7.3), this can be shown by an argument similar to the one in Lemma 7.1.
## 8. Semiclassical equations of motion for general observables
We proceed to more general semiclassical observables. First note that Theorem 3.4 follows immediately from the results obtained so far (Theorem 3.2, Corollary 7.2 and (7.7)), since multiplication with $`k`$ in Bloch representation is bounded. Hence we now have that
(8.1)
$$\underset{\epsilon 0}{lim}x^\epsilon (t)\psi 𝒰^1R(t)𝒰\psi =0$$
for all $`\psi \text{Ran}P_{}D(|x|)H^2`$ and that
(8.2)
$$\underset{\epsilon 0}{lim}k^\epsilon (t)\psi 𝒰^1K(t)𝒰\psi =0$$
for all $`\psi \text{Ran}P_{}`$. We next consider bounded continuous functions of $`x^\epsilon (t)`$ and $`k^\epsilon (t)`$:
###### Lemma 8.1.
Let $`fC_{\mathrm{}}(^d)`$ and $`gC(M^{})`$. Then for all $`\psi \text{Ran}P_{}`$ we have
(8.3)
$$\underset{\epsilon 0}{lim}\left(f(x^\epsilon (t))𝒰^1f(R(t))𝒰\right)\psi =0$$
and
(8.4)
$$\underset{\epsilon 0}{lim}\left(g(k^\epsilon (t))𝒰^1g(K(t))𝒰\right)\psi =0.$$
###### Proof.
We will sketch the proof for $`x^\epsilon (t)`$. First note that $`\overline{R}(t):=𝒰^1R(t)𝒰`$ is a bounded self-adjoint operator and commutes with $`P_{}`$. Hence the sets $`D_\pm :=(\overline{R}(t)\pm i)(\text{Ran}P_{}D(|x|)H^2)`$ are dense in $`P_{}`$ (Since $`R`$ and $`x^\epsilon `$ are vectors of operators in $`^d`$, note that this and the following statements hold component wise). For $`\psi D_\pm `$ we have
(8.5)
$$\left[(x^\epsilon (t)\pm i)^1(\overline{R}(t)\pm i)^1\right]\psi =(x^\epsilon (t)\pm i)^1(\overline{R}(t)x^\epsilon (t))\phi $$
for $`\phi =(\overline{R}(t)\pm i)^1\psi \text{Ran}P_{}D(|x|)H^2`$. Thus, by Theorem 3.2, (8.5) strongly approaches zero as $`\epsilon 0`$ and, since $`D_\pm `$ are dense in $`P_{}`$, $`(x^\epsilon (t)\pm i)^1`$ strongly approach $`(\overline{R}(t)\pm i)^1`$ on $`P_{}`$.
Using the fact that polynomials in $`(x_j\pm i)^1`$, $`j=1,\mathrm{},d`$, are dense in $`C_{\mathrm{}}(^d)`$ one concludes that the convergence $`x^\epsilon (t)\overline{R}(t)`$ on Ran$`P_{}`$ in the “strong resolvent sense” implies
$$\underset{\epsilon 0}{lim}\left(f(x^\epsilon (t))f(\overline{R}(t))\right)\psi =0$$
for all $`fC_{\mathrm{}}(^d)`$ and $`\psi \text{Ran}P_{}`$ (cf. Theorem VIII.20 in ). However, by the functional calculus for self-adjoint operators we have $`f(𝒰^1R(t)𝒰)=𝒰^1f(R(t))𝒰`$ and (8.3) follows.
Clearly (8.4) follows analogously. ∎
###### Proof of Theorem 3.5.
Let $`a𝒪(0)`$. Referring again to the general Stone-Weierstraß theorem we can uniformly approximate $`a(x,\xi )`$ by a sum of products, i.e. $`a(x,\xi )=_{i=0}^{\mathrm{}}a_if_i(x)g_i(\xi )`$ with $`f_iC_{\mathrm{}}(^d)`$, $`g_iC(M^{})`$, $`|a_i|<\mathrm{}`$ and $`sup_{i,x^d,\xi M^{}}`$ $`|f_i(x)g_i(\xi )|<\mathrm{}`$. Hence in order to prove Theorem 3.5 we are left to show that for arbitrary $`fC_{\mathrm{}}(^d)`$ and $`gC(M^{})`$ we have
(8.6)
$$\left(f(x)g(\xi )\right)^{W,\epsilon }(t)𝒰^1f(R(t))g(K(t))𝒰$$
strongly on Ran$`P_{}`$. To see this recall the so called product rule for quantum observables (cf. ). It states, in particular, that for two symbols $`A,B𝒪(0)`$
$$\underset{\epsilon 0}{lim}\left((AB)^{W,\epsilon }A^{W,\epsilon }B^{W,\epsilon }\right)\psi =0.$$
Applied to our case this yields
$$\left(f(x)g(\xi )\right)^{W,\epsilon }(t)\left(f(x)^{W,\epsilon }g(\xi )^{W,\epsilon }\right)(t)=f(x^\epsilon (t))g(k^\epsilon (t)).$$
Finally, since $`f`$ and $`g`$ are bounded, Lemma 8.1 implies (8.6) and thus Theorem 3.5. ∎
## 9. Band crossings
We proved the semiclassical limit for isolated bands only. In principle, there are two distinct mechanisms of how this assumption could be violated. First of all a band could be isolated but have a constant multiplicity larger than one. This occurs, e.g., for the Dirac equation where because of spin the electron and positron bands are both two-fold degenerate. A systematic study is only recent and leads to a matrix valued symplectic structure for the semiclassical dynamics. For periodic potentials degeneracies are the exception. They form a real analytic subvariety of the Bloch variety $`B=\{(k,\lambda )^d\times fL^2(M):H_{\mathrm{per}}(k)f=\lambda f\}`$ and have a dimension at least one less than the dimension of $`B`$ . Thus points of band crossings have a $`k`$-Lebesgue measure zero. From the study of band structures in solids one knows that band crossings indeed occur. Thus it is of interest to understand the extra complications coming from band crossings.
There are two types of band crossings. The first one is removable through a proper analytic continuation of the bands. In a way, removable band crossings correspond to a wrong choice of the fundamental domain. E.g. for $`V=0`$ we may artificially introduce a lattice $`\mathrm{\Gamma }`$. The bands touch then at the boundary of $`M^{}`$. Upon analytic continuation we recover the single band $`E_1(k)=k^2/2`$ with $`M^{}=^d`$. In one dimension all band crossings can be removed . Thus, with the adjustment discussed, our result fully covers the case $`d=1`$. For $`d2`$ generically band crossings cannot be removed.
It is then of great physical interest to understand how a wave packet tunnels into a neighboring band through points of degeneracy (or almost degeneracy). For a careful asymptotic analysis in particular model systems we refer to the monumental work of G. Hagedorn . Gerard considers a model system with two bands in two dimensions, i.e., the role of $`\frac{1}{2}\mathrm{\Delta }+V`$ is taken by $`\left(\begin{array}{cc}k_1& k_2\\ k_2& k_1\end{array}\right)`$. He investigates the semiclassical limit and proves that the particle may tunnel to the other band with a probability which depends on how well the initial wave packet is concentrated near a semiclassical orbit hitting the singularity.
## Acknowledgments
FH gratefully acknowledges the financial support by the Deutsche Forschungsgemeinschaft via the Graduiertenkolleg Mathematik im Bereich ihrer Wechselwirkung mit der Physik at the LMU München. |
warning/0001/astro-ph0001174.html | ar5iv | text | # ORIGIN OF COLOR GRADIENTS IN ELLIPTICAL GALAXIES
## 1 INTRODUCTION
Stellar populations in an elliptical galaxy are not uniform. Stars at a galaxy center are redder than those in the outer region and colors in a galaxy are progressively getting bluer with an increasing radius (e.g., Vader et al. 1988; Franx, Illingworth, & Heckman 1989; Peletier et al. 1990a; Peletier, Valentijn, & Jameson 1990b). Peletier et al. (1990a) made surface photometry in the $`U`$-, $`B`$-, and $`R`$-bands for a sample of 39 nearby elliptical galaxies and examined their color gradients in $`UR`$ and $`BR`$. They found that typical color gradients $`\mathrm{\Delta }(UR)/\mathrm{\Delta }\mathrm{log}r`$ and $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}r`$ are $`0.20`$ mag$`/`$dex and $`0.09`$ mag$`/`$dex, respectively, and demonstrated that a dispersion of the color gradients is small, i.e., only 0.02 mag $`/`$dex in both colors. Since many of elliptical galaxies show radial gradients in line strengths such as Mg<sub>2</sub>, Fe<sub>1</sub>(5270 Å) and Fe<sub>2</sub>(5335 Å) (e.g., Carollo, Danziger, & Buson 1993; Davies, Sadler, & Peletier 1993; Gonzalez 1993; Kobayashi & Arimoto 1999), the color gradients have been naively assumed to originate from a metallicity gradient inside a galaxy.
However, such an interpretation for the origin of the color gradient is premature, because stellar populations of either higher metallicity or older age can make a galaxy redder. This problem, which is called the age-metallicity degeneracy, was first pointed out by Worthey, Trager, & Faber (1996) and then discussed by Arimoto (1996). For example, the degeneracy makes it difficult to interpret the origin of the tight correlation between colors and magnitudes of elliptical galaxies; brighter elliptical galaxies tend to have redder colors. This correlation called color-magnitude (CM) relation can be excellently reproduced by a metallicity sequence with a galactic wind model based on a monolithic collapse scenario (e.g., Arimoto & Yoshii 1987), where more massive elliptical galaxies should be enriched more in metals and thus become redder. However, Worthey et al. (1996) claimed that an age sequence of elliptical galaxies can equivalently reproduce the CM relation if brighter elliptical galaxies are older and thus redder. To break this degeneracy, Kodama & Arimoto (1997) built up two model sequences (metallicity and age sequences) which are normalized to reproduce the CM relation of ellipticals in Coma cluster by using their evolutionary synthesis model. They compared the evolution of the model CM relation with the observed relations of ellipticals in distant clusters. The CM relations produced by a metallicity sequence agree with the observed relations out to $`z1`$, while those produced by an age sequence deviates from the observed ones significantly even at $`z0.20.3`$, showing the origin of the CM relation is primarily a metallicity variation with a galaxy mass. The CM relation can also be reproduced in a hierarchical galaxy formation scenario (Kauffmann & Charlot 1998). Although a model of Kauffmann & Charlot (1998) allows a more extended period of star formation in elliptical galaxies, it shows that the CM relation is produced by a metallicity variation. It is thus worth emphasizing that the interpretation of CM relation with a metallicity sequence is robust and independent of detailed assumptions on galaxy formation processes.
Recently, the origin of color and line strength gradients in elliptical galaxies has been discussed with elaborate models. Martinelli, Matteucci, & Colafrancesco (1998) tried to reproduce the color gradients by assuming that a galactic wind blows later in the inner part of a galaxy due to a deeper potential well defined mainly by dark matter. Adopting a multi-zone model that takes into account gas dynamics, local star formation, and chemical evolution, Tantalo et al. (1998) reproduced radial gradients of colors and line strengths to some extent. Nevertheless, these detailed modelings are not fully successful; the inner part of a galaxy becomes too iron enriched due to an extended period of star formation there.
In this paper, adopting a much simpler approach without entering into details of physical processes of galaxy formation and evolution, we try to depict essential aspect of the origin of color gradients. Our approach is similar to that adopted by Kodama & Arimoto (1997) for studying the origin of the CM relation. By using a population synthesis model, we first make two different model galaxies, both of which can reproduce a typical color gradient of elliptical galaxies at $`z=0`$, by changing either mean stellar metallicity or age, and let evolve them back in time. The evolution of color gradients thus predicted are then confronted with the observed ones in distant ellipticals extracted from the Hubble Deep Field North (HDF-N; Williams et al. 1996). As was demonstrated by Kodama & Arimoto (1997), the best way to disentangle the age and metallicity effects on galaxy colors is to look back galaxies at high redshift.
It should be noted that dust extinction in elliptical galaxies may have some effects on the color gradients (Goudfrooij & de Jong 1994; Wise & Silva 1996). Witt, Thronson, & Capuano (1992) calculated a radiative transfer within elliptical galaxies by assuming a diffuse distribution of dust and suggested that surface brightness profiles and color gradients could be well reproduced by dust effects. In fact, it has been believed that the IRAS detected FIR emission for about half of the elliptical galaxies observed. However, since many of the detections were made with a $`3\sigma `$ threshold, a significant fraction of the previously claimed detections may be spurious. Only $`1217\%`$ of the observed elliptical galaxies are detected with a sufficiently high confidence level (Bregman et al. 1998). Therefore, in this paper, we have chosen to focus on age and metallicity effects only. Effects of dust extinction on the color gradients are still open and will be studied in our subsequent paper.
This paper is organized as follows. The model prescription is given in § 2. The sample selection and data reduction of elliptical galaxies in the HDF-N are described in § 3. The observed color gradients are compared with the models in § 4. Discussion and conclusions are given in § 5 and § 6, respectively. The cosmological parameters adopted throughout this paper are $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_0=0.2`$ and $`\mathrm{\Lambda }=0`$ unless otherwise noted.
## 2 MODELS
The color gradients can be reproduced by either a metallicity gradient or an age gradient, but these effects are degenerate at $`z=0`$. We attempt to break up such degeneracy by confronting model color gradients with the observed ones of ellipticals at high redshifts in the HDF-N. We have built two sequences of evolutionary models under the alternative assumptions: 1) the color gradients originate from the metallicity gradient of old stellar populations, or 2) the color gradients arise from the age gradient of stars which have the same metallicity within the galaxy.
The metallicity sequence is constructed as follows: We assume that the star formation lasted progressively longer towards galaxy center; i.e., the galactic wind blew later in the inner region, so that the mean stellar metallicity becomes higher and the colors become redder. Even at the galaxy center, however, the star formation stopped at early times ($`t_{\mathrm{gw}}=0.83`$ Gyr) and the stellar ages are almost the same at everywhere within a galaxy; $``$ 15 Gyr at $`z=0`$ (see Tables 1 and 2). The assumption seems to be justified, because the duration of star forming activity is determined from the local dynamical potential (Larson 1974) and should be longer in the inner region. To build a metallicity gradient, however, this hypothesis is not unique; other hypotheses such as higher star formation efficiency, a larger yield (i.e., a flatter initial mass function (IMF) slope), or higher metallicity of infalling gas towards the galaxy center can also produce the metallicity gradient. However, such details are not essential in the present study, because so far as the mean stellar metallicity increases in the inner part of a galaxy while stars are uniformly old, the resulting evolution of color gradient is essentially the same. We should stress this, because this paper aims for showing how color gradients of stellar populations behave as a function of lookback time, not for seeking the best model that can explain the observed color gradients of sampled galaxies in the HDF-N. The galactic wind is conventionally introduced to build up the metallicity gradient and it might be possible that such gradients can also be established if elliptical galaxies formed via hierarchical clustering.
The age sequence is constructed as follows: The star formation started earlier in the inner regions, but lasted equally long ($`t_{\mathrm{gw}}=0.83`$ Gyr) everywhere in a galaxy. The period chosen here is the same as that assumed for the center of metallicity sequence model. In this way, the mean stellar ages $`t_{\mathrm{age}}`$ become older and colors redder in the inner regions, while the mean stellar metallicities are almost the same everywhere (see Tables 1 and 2). It is hard to comprehend the physical process to produce such age gradient, although it may occur if a series of young dwarf galaxies accrete onto a massive galaxy and are tidally disrupted at the outer region of a galaxy before falling into the center. However, it is not our aim anyway to construct a physically motivated age sequence model.
Although we do not discuss in this paper in detail, we have tried to build an alternative age sequence model in which star formation began everywhere 15 Gyr ago and lasted longer at the outer parts of the galaxy, keeping the same metallicity everywhere as that at the galaxy center. As a result, the mean stellar ages become progressively younger from the galaxy center towards the outer region. However, this model fails to reproduce the observed color gradients in nearby ellipticals. If old stellar populations are contained at the outer regions and if the mean metallicity is as rich as that at the center, there is no way by making their mean ages young to build up the color gradients as steep as the observed ones.
Both metallicity and age sequences are constructed to reproduce typical color gradients observed for nearby elliptical galaxies; $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}r=0.09\pm 0.02`$ mag/dex (i.e., $`0.07`$ and $`0.11`$ mag/dex) and $`BR=1.633`$ mag at $`r=r_\mathrm{e}/10`$ ($`r_\mathrm{e}`$ refers to an effective radius) which is derived from the mean $`BR`$ color at $`r_\mathrm{e}/2`$ and $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}r=0.09`$ mag/dex (Peletier et al. 1990a). We also construct the models which reproduce the color gradients of $`0.09\pm 0.04`$ mag/dex (i.e., $`0.05`$ and $`0.13`$ mag/dex) within which the gradients of almost all the sample ellipticals in Peletier et al. (1990a) are included. Thus we have a set of four model sequences each for the age sequence and the metallicity sequence.
To calculate gradients of photo-chemical properties of a galaxy, we construct a galaxy model consisting of a ’shell’ at each radius, and assume that each ’shell’ evolves independently. The star formation rate (SFR) in a ’shell’ is proportional to the gas fraction (Schmidt law) with a time scale of 0.1 Gyr and a primordial gas is supplied to the shell with a rate of $`\mathrm{exp}(t/0.1\mathrm{Gyr})`$. An infall of a primordial gas onto the shell may sound rather strange, since it is more likely that the inner shells suffer from the infall of enriched gas from the outer shells. However, a proper modeling of such effect is beyond our scope. Since the infall rate we employed in the present study is significantly high, the resulting evolutionary behavior of the color gradient remains almost the same even if we adopt the ’simple’ model prescription instead of the infall model. For an IMF, a power-law mass spectrum with a slope of $`x=1.10`$ is assumed in the range of $`0.05M_{}M50M_{}`$. We note that this IMF slope is the same as in Kodama et al. (1998), and was introduced to increase chemical yields to reproduce the reddest end of the Coma CM relation. The nucleosynthesis yields of SNe Ia and SNe II are taken from Tsujimoto et al. (1995). The chemical enrichment by SNe Ia is calculated with a metallicity dependent SN Ia rate, of which detail formulation is given by Kobayashi, Tsujimoto & Nomoto (1999). The spectral evolution is calculated using a spectral synthesis database of Kodama & Arimoto (1997). The adopted parameter values for galactic wind epoch $`t_{\mathrm{gw}}`$ and stellar age $`t_{\mathrm{age}}`$ at each radius for the four color gradients are listed in Table 1. The resulting photo-chemical properties of a galaxy at $`z=0`$ are summarized in Table 2, including gas abundances and colors at each radius. It should be noted that our models have an SN II-like abundance pattern, though the value \[Mg/Fe\] $`+0.4`$ is slightly larger than that suggested by the observational estimates (Worthey, Faber & Gonzalez 1992; Kobayashi & Arimoto 1999); a larger \[Mg/Fe\] ratio is a result of the assumed flat IMF ($`x=1.10`$).
## 3 GALAXIES FOR COMPARISON
### 3.1 Sample Elliptical Galaxies
To compare theoretical color gradients with those of elliptical galaxies at high redshifts, the archival data of the HDF-N (Williams et al. 1996) are used. Our sample galaxies are selected from those brighter than $`I_{814,AB}=22`$ mag, in such a way that we can derive reliable surface brightness profiles and color gradients. All these galaxies have spectroscopic redshifts. Elliptical galaxies are identified by using a bulge-to-total ($`B/T`$) luminosity ratio derived by Marleau & Simard (1998) who obtained the $`B/T`$ ratios in the $`I_{814}`$-band of HDF galaxies brighter than $`I_{814}=26`$ mag by decomposing quantitatively the surface brightness profile into the bulge and the disk components. We define galaxies with $`B/T>0.5`$ as ellipticals according to Marleau & Simard (1998). Our resulting sample consists of ten elliptical galaxies with redshifts spanning from $`z=0.089`$ to $`1.015`$ as listed in Table 3. It is important to note here that our sample selection does not rely on any color information. Franceschini et al. (1998) also made a sample of elliptical galaxies in the HDF using surface brightness profiles. Their selection, however, is not based on the bulge - disk decomposition. Consequently, within $`I_{814,AB}<22`$ mag, our sample is smaller than that of Franceschini et al. (1998), except for the galaxy 4-241.1 at $`z=0.318`$. The Franceschini et al’s sample includes several galaxies having the $`B/T`$ ratio smaller than 0.5, which are presumably disk dominant galaxies and thus are not included in our sample (see also Kodama, Bower & Bell 1999). It should be noted that $`B/T>0.5`$ in the $`I_{814}`$-band may be slightly loose to isolate elliptical galaxies. $`B/T>0.5`$ in the $`B`$-band roughly corresponds to local early-type galaxies (e.g., Simien & de Vaucouleurs 1986). By using bulge and disk models of Kodama et al. (1999), we found that the $`I_{814}`$-band $`B/T`$ ratios of ellipticals are larger than $``$ 0.7 at $`z=0`$ and change little from $`z=0`$ to 1 in the observer’s frame. Thus our sample may include galaxies as late as Sa, though they must be minority because most of the $`B/T`$ ratios of our sample galaxies are larger than 0.7 (Table 3).
We measured $`V_{606}I_{814}`$ colors of our sample galaxies and present them in Table 3 and Figure 1. These colors are obtained for a 10 kpc aperture. Seven of the sample galaxies have red colors consistent with those for ellipticals with passive evolution as shown in Figure 1 (hereafter we call these galaxies ‘red ellipticals’), and the other three galaxies (4-241.1, 2-251.0, and 4-928.0) have blue colors (see also Figure 2; hereafter referred to as ‘blue ellipticals’ for convenience). Azimuthally averaged surface brightness profiles of the sample ellipticals in the $`I_{814}`$-band are shown in Figure 2. The average is taken along an ellipse fitted to an isophote of the $`I_{814}`$-band image. These profiles are well represented by an $`r^{1/4}`$ law for both the red ellipticals and the blue ellipticals. Effective radii ($`r_\mathrm{e}`$) of the sample ellipticals are obtained by the $`r^{1/4}`$ fitting and results are shown in Table 3. The fitting is done by removing data points in the inner and outer regions, since they are unreliable due to an effect of a point spread function (PSF) and low signal-to-noise ratios (S/Ns), respectively. The obtained effective radii range from $`r_e=`$1.4 to 7.9 kpc, which are typical for nearby giant ellipticals (e.g., Bender, Burstein, & Faber 1992).
### 3.2 Color Gradients
To examine color gradients of our sample galaxies, the images in F606W ($`V_{606}`$)- and F814W ($`I_{814}`$)-bands were used, because S/Ns of the images are better in these bands than in the other bands (F300W and F450W). First, the sky value around each object was determined by using the ”phot” task in the IRAF apphot package (”mode” in an annulus with an inner radius of mostly $`3^{\prime \prime }4^{\prime \prime }`$ and a width of 0.$`{}_{}{}^{\prime \prime }4`$ is adopted) and was subtracted. Next, the angular resolutions of the blue images and the red images were adjusted and a $`V_{606}I_{814}`$ color map of each galaxy was made. Figure 3 shows the resulting color maps. Most of the color maps for the red ellipticals show symmetric structure and ordinary color gradient. In two cases (2-456.0 and 2-121.0), slight asymmetry in the color distribution is seen. This does not seem to be caused by a misalignment of the images. For the blue ellipticals, variety of the color distribution is seen, which we will discuss in §5.1. Next, a $`V_{606}I_{814}`$ color profile is derived with the azimuthally averaged radial profiles as a function of semi-major radius from the galaxy center. Resulting color profiles are shown in Figures 4 and 5 for the red ellipticals and in Figure 7 for the blue ellipticals. It should be noted that in plotting observed color profiles, a zero-point of the observed color for each galaxy is shifted to be compared with those by the models. The amounts of these shifts are given in the last column of Table 3. Most of the shifts are less than 0.2 mag; for intrinsically luminous objects the colors tend to be shifted to bluer colors and vice versa. Thus the color differences are presumably caused by an effect of the color-magnitude relation of elliptical galaxies. The zero-point offsets could also come from photometric errors in HST data (e.g., Holtzman et al. 1995; Ellis et al. 1997; Kodama et al. 1998) and from uncertainty of the population synthesis model. Needless to say, zero-point offsets do not affect the color gradients. An error bar attached to each data point in the color profile includes an photometric error, a local sky subtraction error, and a dispersion of colors along each elliptical isophote. The data points of 2-456.0 and 2-121.0 have relatively large error bars in the middle region of the profiles; these objects have asymmetric color distributions which results in the relatively large dispersions of colors in the regions.
Finally, slopes of the color profiles $`\mathrm{\Delta }(V_{606}I_{814})/\mathrm{\Delta }\mathrm{log}(r/r_\mathrm{e})`$ of the sample galaxies are derived by applying a least square fit to the color profiles. This fit was done after removing the unreliable data points in the outer region ($`r>r_\mathrm{e}`$) and in the innermost region. The removed data points in the inner most regions are shown by crosses in the color profiles (Figures 4 and 5); they clearly deviate from the $`r^{1/4}`$ fits presumably due to the effect of a PSF as seen in Figure 2. The slopes of the color profiles for the blue ellipticals presented in Figure 7 are not derived.
## 4 ORIGIN OF COLOR GRADIENTS
Once the spectra are calculated at various radii of a model galaxy in such a way that they reproduce the color gradient $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}(r/r_\mathrm{e})=0.09`$ mag/dex at $`z=0`$, the model gradient in any colors at any redshifts can be predicted to compare with the observed gradients. The model color gradients at the redshift of each sample elliptical galaxy are over-plotted in each panel of Figures 4 and 5 for the age gradient and the metallicity gradient models, respectively. Solid and dotted lines in each panel show the gradients in $`V_{606}I_{814}`$ color in the observer’s frame corresponding to the gradients of $`0.09\pm 0.02`$ mag/dex and $`0.09\pm 0.04`$ mag/dex in the $`BR`$ color at $`z=0`$, respectively.
It is clearly shown that the model gradient made by the age gradient begins to deviate from the observed ones at a redshift of $`z0.3`$ and the deviation is getting worse as a redshift increases (Figure 4). On the contrary, the model gradient made by the metallicity gradient agrees well with the observed color gradients within the effective radius at any redshift from $`z0.1`$ to $`1`$ (Figure 5). Considering that the model gradients are calibrated only at $`z=0`$, we insist that the agreement between the model and the observed gradients is excellent. Such agreement is more clearly seen in Figure 6 which shows both the observed color gradients and the model gradients as a function of redshift. Solid and dotted lines show the evolution of model color gradients in $`V_{606}I_{814}`$ in the observer’s frame for the metallicity gradient and the age gradient, respectively. These model gradients are calculated in a region of $`1\mathrm{log}r/r_\mathrm{e}0`$ for the metallicity gradient and for the age gradient at $`z<0.5`$, and $`1\mathrm{log}r/r_\mathrm{e}0.5`$ for the age gradient at $`z>0.5`$. Dots show the observed gradients of the sample galaxies within $`r_\mathrm{e}`$ calculated by a least square fitting to the data points shown as filled circles in Figures 4 and 5. Error bars indicate the fitting errors of 1$`\sigma `$.
## 5 DISCUSSION
### 5.1 Blue Galaxies in Our Sample
Since our sample is selected independently of galaxy colors, relatively blue elliptical-like galaxies are included in our sample (3/10). Among these three, 4-241.1 ($`z=0.318`$) has a nearly flat color profile with a slightly bluer color in the inner part (Figure 3 and the top panel of Figure 7). The spectrum of this galaxy (from the Hawaii Active Catalog; Cowie 1997) shows Balmer absorption lines shortward of 4000Å in the rest frame as well as emission lines such as \[OII\]$`\lambda `$3727, H$`\beta `$, \[OIII\]$`\lambda `$4959, $`\lambda `$ 5007, and H$`\alpha `$. Thus this galaxy is expected to have some young stellar populations. The galaxy 2-251.0 ($`z=0.960`$) has a color gradient in the opposite sense to those of ordinary ellipticals; i.e., the inner region of the galaxy is bluer than the outer region (Figure 3 and the middle panel of Figure 7). Since this galaxy has been suggested as an AGN (Franceschini et al. 1998), this opposite gradient is probably caused by AGN and/or associating starburst activity in the central part of the galaxy. Although both galaxies seem to have the younger ages, our model by the age gradient cannot reproduce their color gradients. Accordingly, neither the metallicity gradient nor the age gradient can explain their observed color gradients. These galaxies presumably have on-going star formation (or post-starburst) in the inner part of the galaxies, which may be caused by a galaxy interaction as expected from their rather disturbed morphology. Our model cannot apply to these cases. The other galaxy 4-928.0 ($`z=1.015`$) has a rather steep color gradient as seen in Figure 3 and the bottom panel of Figure 7. The observed color profile of the galaxy is shifted by 0.4 mag to be compared with the model color gradients (solid and dotted lines). The observed color gradient of this galaxy may be explained by the age gradient if we tune ages of stellar populations in the model galaxy. However, the mean age is set to be 0.2 Gyr at $`r_\mathrm{e}`$ and 9 Gyr at the innermost point. Such a large age at the center cannot be realized at this redshift ($`z=1.015`$) unless a different set of cosmological parameters is adopted. Finally, note that colors of these galaxies are not consistent with those of passively evolving ellipticals formed at high $`z`$. Much larger sample is needed in future works to find whether the population of the blue galaxies is really the minority in the whole population of elliptical galaxies.
### 5.2 Gradients of Metal Absorption Line Strengths
Kobayashi & Arimoto (1999) showed that central velocity dispersions of elliptical galaxies correlate well with their mean metallicities. They found, however, that this correlation has a significant scatter. They also showed that Mg<sub>2</sub> gradients do not correlate with any physical parameters including a galaxy mass. Such a large dispersion and an absence of the correlation are contrary to what monolithic collapse simulations predicted. As the common origin for the large dispersion and the absence of the correlation, they suggested the following possibilities: (1) a significant age spread among elliptical galaxies, (2) age gradients inside elliptical galaxies, (3) dust obscuration, (4) effects of merging on metallicity gradients inside elliptical galaxies, and (5) intrinsic scatter which elliptical galaxies have had since they formed. Since we confirm in this paper that the color gradients are caused by the metallicity gradient, we can reject the possibility (2).
There seems to be a problem that color gradients in elliptical galaxies do not correlate with their Mg<sub>2</sub> gradients, though only a small number of elliptical galaxies are examined both for color gradient and Mg<sub>2</sub> gradient (Peletier 1989). However, the absence of this correlation may be a spurious result caused by a different way to derive the gradients between the color and Mg<sub>2</sub> gradients. Some effects of dust extinction and/or, possibly, observational errors on both gradients might suppress the true correlation. In any case, this problem will be investigated in detail in our forthcoming paper.
Worthey et al. (1992) suggested that the color-magnitude relation of elliptical galaxies may originate from the fact that more luminous ellipticals tend to contain more Mg relative to Fe. One can claim that it could be possible that the same holds for the internal gradients. However, Kobayashi & Arimoto (1999) have recently carried out a careful study of line strength gradients of 80 elliptical galaxies having the best quality of spectra but find no evidence supporting for the \[Mg/Fe\] gradient. Therefore, we conclude that the color gradients are not originating from the gradient of the \[Mg/Fe\] ratio.
## 6 CONCLUSIONS
The origin of color gradients in elliptical galaxies is examined. A typical color gradient of nearby elliptical galaxies is reproduced by two alternative model sequences, metallicity gradient and age gradient. This age-metallicity degeneracy is broken by looking back the evolution of color gradient towards high redshifts; the predicted color gradients at high redshifts are compared with the observed color gradients of ten elliptical galaxies in the HDF-N out to $`z1`$. We find that the observed color gradients of the seven red galaxies, whose colors are consistent with those of passively evolving galaxies, are in excellent agreement with the metallicity gradient at any redshift. These conclusions are independent of cosmological parameters and parameters for an evolutionary model of galaxy (IMF and SFR) within a reasonable range. Thus we conclude that the primary origin for the color gradient in elliptical galaxies is the metallicity gradient in the old stellar populations.
This work was financially supported in part by Grant-in-Aids for the Scientific Research (No. 0940311 and 09740173) by the Japanese Ministry of Education, Culture, Sports and Science. CK and TK thank Research Fellowships of the Japan Society for the Promotion of Science (JSPS) for Young Scientists.
Figure Caption
Fig. 1 — $`V_{606}I_{814}`$ colors of our sample galaxies in a 10 kpc aperture as a function of redshift. Solid line shows a color track for a passively evolving elliptical and dotted line for a non-evolving elliptical. Filled circles indicate red galaxies whose colors are nearly consistent with that of the passively evolving or non-evolving ellipticals. Open circles indicate the remaining bluer galaxies.
Fig. 2 — Azimuthally averaged radial surface brightness profiles of our sample galaxies in the $`I_{814}`$-band. (a) Red galaxies of which colors are consistent with that of a passively evolving galaxy. (b) Galaxies of which colors are bluer than that of a passively evolving galaxy (see also Figure 1). Filled circles show the data points and solid lines show the fitted lines with the $`r^{1/4}`$ law.
Fig. 3 — $`V_{606}I_{814}`$ color maps for (a) the red galaxies whose colors are consistent with that of a passively evolving galaxy, and (b) the “blue” ellipticals. Plus sign of each object denotes the centroid of the galaxy in the $`I_{814}`$-band image. The circle represents an effective radius centered on the centroid. Object ID, redshift, and effective radius are also shown. Note that a color bar is different from galaxy to galaxy, to show the color distribution clearly.
Fig. 4 — Observed color gradients of the red elliptical galaxies together with the model gradients by the age gradient. Filled circles represent the data points. Crosses refer to the data points which are not used in deriving slopes of color profiles. The models show the predicted color gradients seen at each object’s redshift. Solid lines correspond to $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}(r/\mathrm{r}_\mathrm{e})=0.09\pm 0.02`$ mag/dex at $`z=0`$, and the dotted lines correspond to $`\mathrm{\Delta }(BR)/\mathrm{\Delta }\mathrm{log}(r/\mathrm{r}_\mathrm{e})=0.09\pm 0.04`$ mag/dex at $`z=0`$. Zero-points of the observed colors are slightly shifted by $`\mathrm{\Delta }(V_{606}I_{814})`$ in Table 3 to compare the observed gradients with the models.
Fig. 5 — Same as Fig. 4, but for the case of the metallicity gradient.
Fig. 6 — $`V_{606}I_{814}`$ color gradient versus redshift for our red sample galaxies. Filled circles represent color gradients of the sample ellipticals, which are determined within effective radii without the innermost regions of the galaxies. Solid and dotted lines indicate the evolutions of color gradients caused by the metallicity gradient and the age gradient, respectively. (See text for the model.)
Fig. 7 — Color gradients of the galaxies which have relatively blue colors. The color gradient of 4-928.0, which has a rather steep color gradient, is plotted together with a model color gradient made by age gradient fitted to the observed one by tuning ages of stellar populations in the model galaxy. (See text for the model.) The observed color for 4-928.0 is shifted by 0.4 mag to compare the model gradients. |
warning/0001/hep-ph0001218.html | ar5iv | text | # On the structure of large 𝑁_𝑐 cancellations in baryon chiral perturbation theory
## I Introduction
Baryon chiral perturbation theory can be used to systematically compute the properties of baryons as a function of the light quark masses $`m_q`$. A physical quantity’s non-analytic dependence on $`m_q`$ is calculable from pion-loop graphs; its analytic dependence has contributions both from pion-loop graphs and from low-energy constants that are present in the chiral Lagrangian.
It is convenient to formulate baryon chiral perturbation theory in terms of velocity-dependent baryon fields, so that the expansion of the baryon chiral Lagrangian in powers of $`m_q`$ and $`1/m_B`$ (where $`m_B`$ is the baryon mass) is manifest . This formulation is called heavy baryon chiral perturbation theory. The earliest application of heavy baryon chiral perturbation theory was to baryon axial vector currents . Two important results were obtained from this analysis. First, the baryon axial coupling ratios were found to be close to their $`SU(6)`$ values with an $`F/D`$ ratio close to $`2/3`$, the value predicted by the non-relativistic quark model. Second, there were large cancellations in the one-loop corrections to the baryon axial vector currents between loop graphs with intermediate spin-$`1/2`$ octet and spin-$`3/2`$ decuplet baryon states. It was later proven using the $`1/N_c`$ expansion that the baryon axial couplings ratios should have $`SU(6)`$ values with $`F/D=2/3`$, up to corrections of order $`1/N_c^2`$ for pions .<sup>*</sup><sup>*</sup>*The ratio $`F/D`$ is equal to $`2/3`$ in the quark representation, and to $`5/8`$ in the Skyrme representation. The quark and Skyrme representations are equivalent in large $`N_c`$ up to corrections of relative order $`1/N_c^2`$ for pions . We will use the quark representation in this paper. In addition, it was shown that axial vector current loop graphs with octet and decuplet intermediate states cancel to various orders in $`N_c`$. For nucleon and Delta intermediate states, there is a cancellation of the one-loop graphs to two orders in $`N_c`$; each individual one-loop diagram is of order $`N_c`$, but the sum of all one-loop diagrams is of order $`1/N_c`$ . Similar large-$`N_c`$ cancellations also occur for other baryonic quantities .
We would like to find a calculational scheme that simultaneously exhibits both the $`m_q`$ and $`1/N_c`$ expansions. In the chiral limit $`m_q0`$, pions become massless Goldstone boson states. There is an expansion about the chiral limit in powers of $`m_q/\mathrm{\Lambda }_\chi `$, or equivalently, in powers of $`m_\pi ^2/\mathrm{\Lambda }_\chi ^2`$, where $`\mathrm{\Lambda }_\chi 1\mathrm{GeV}`$ is the scale of chiral symmetry breaking and $`m_\pi `$ is the pion mass.For $`SU(3)_L\times SU(3)_R`$ chiral symmetry, there are loop corrections involving the pions, kaons and $`\eta `$ which depend on the pion, kaon and $`\eta `$ masses, respectively. Chiral perturbation theory depends on the expansion parameters $`m_\pi ^2/\mathrm{\Lambda }_\chi ^2`$, $`m_K^2/\mathrm{\Lambda }_\chi ^2`$ and $`m_\eta ^2/\mathrm{\Lambda }_\chi ^2`$. Large-$`N_c`$ chiral perturbation theory also depends on the $`\eta ^{}`$ mass. In the large-$`N_c`$ limit, the nucleon and Delta become degenerate, $`\mathrm{\Delta }M_\mathrm{\Delta }M_N1/N_c0`$, and form a single irreducible representation of the contracted spin-flavor symmetry of baryons in large-$`N_c`$ QCD . There is an expansion in powers of $`1/N_c`$ about the large-$`N_c`$ limit. We will consider a combined expansion in $`m_q/\mathrm{\Lambda }_\chi `$ and $`1/N_c`$ about the double limit $`m_q0`$ and $`N_c\mathrm{}`$.
Loop graphs in heavy baryon chiral perturbation theory have a calculable dependence on the ratio $`m_\pi /\mathrm{\Delta }`$. In general, this dependence is described by a function $`F(m_\pi ,\mathrm{\Delta })`$. In the chiral limit $`m_q0`$ with $`\mathrm{\Delta }`$ held fixed, the function can be expanded in powers of $`m_\pi /\mathrm{\Delta }`$,
$$F(m_\pi ,\mathrm{\Delta })=F_0+\left(\frac{m_\pi }{\mathrm{\Delta }}\right)F_1+\left(\frac{m_\pi }{\mathrm{\Delta }}\right)^2F_2+\mathrm{},$$
(1)
whereas in the $`1/N_c0`$ limit with $`m_\pi `$ held fixed, the function can be expanded in powers of $`\mathrm{\Delta }/m_\pi `$,
$$F(m_\pi ,\mathrm{\Delta })=\overline{F}_0+\left(\frac{\mathrm{\Delta }}{m_\pi }\right)\overline{F}_1+\left(\frac{\mathrm{\Delta }}{m_\pi }\right)^2\overline{F}_2+\mathrm{}.$$
(2)
The difference between the two expansions in Eqs. (1) and (2) is commonly referred to as the non-commutativity of the chiral and large-$`N_c`$ limits .
It is important to remember, however, that the conditions for heavy baryon chiral perturbation theory (including Delta states) to be valid are that $`m_\pi \mathrm{\Lambda }_\chi `$ and $`\mathrm{\Delta }\mathrm{\Lambda }_\chi `$. The ratio $`m_\pi /\mathrm{\Delta }`$ is not constrained and can take any value. The entire dependence of a physical quantity on $`m_\pi /\mathrm{\Delta }`$ is calculable in heavy baryon chiral perturbation theory , so the ratio $`m_\pi /\mathrm{\Delta }`$ need not be small or large for calculations. In the real world, $`m_\pi /\mathrm{\Delta }0.5`$, so it is useful to have a calculational scheme that retains the full functional dependence of $`F(m_\pi ,\mathrm{\Delta })`$ on the ratio $`m_\pi /\mathrm{\Delta }`$. A straightforward approach is to simply calculate the full dependence on $`m_\pi /\mathrm{\Delta }`$ of the loop graphs, and evaluate the loop correction at the physical value $`m_\pi /\mathrm{\Delta }0.5`$ . Another common procedure advocated in the literature is to not include intermediate Delta particles explicitly in loops, but to incorporate their effects into the low-energy constants of the effective Lagrangian . The disadvantage of this second approach is that one finds large numerical cancellations between loop diagrams with intermediate nucleon states and low-energy constants containing the effects of Delta states. These cancellations are guaranteed to occur as a consequence of the contracted spin-flavor symmetry which is present in the $`N_c\mathrm{}`$ limit. The large-$`N_c`$ spin-flavor symmetry responsible for the cancellations is hidden in this approach because including only the spin-$`1/2`$ baryons in the chiral Lagrangian breaks the large-$`N_c`$ spin-flavor symmetry explicitly, since the spin-$`1/2`$ and spin-$`3/2`$ baryons together form an irreducible representation of spin-flavor symmetry. Because the sum of the loop contributions with intermediate octet and decuplet states respects spin-flavor symmetry and is much smaller (by powers of $`1/N_c`$) than each individual loop contribution separately, it is important to keep the large-$`N_c`$ spin-flavor symmetry of the baryon chiral Lagrangian and the large-$`N_c`$ cancellations manifest.
In this paper, we will show how one can combine heavy baryon chiral perturbation theory with the $`1/N_c`$ expansion so that the full-dependence on $`m_\pi /\mathrm{\Delta }`$ is retained and the $`1/N_c`$ cancellations are explicit. This method has the advantage that the loop correction to the baryon axial isovector current, which is order $`1/N_c`$, is automatically obtained to be of this order, instead of as the sum of two contributions (loop correction and counterterm) of order $`N_c`$ which cancel to two powers in $`1/N_c`$. Note that at higher orders the cancellations become more severe, and it is even more important to keep the $`1/N_c`$ cancellations manifest. For example, at two loops, each loop diagram is naively of order $`N_c^2`$, whereas the sum of all two-loop diagrams is order $`1/N_c^2`$. Not including the $`1/N_c`$ cancellations in a systematic way gives a misleading picture of the baryon chiral expansion—one finds higher order corrections that grow with $`N_c`$, which is incorrect. Including the $`1/N_c`$ cancellations restores the $`1/N_c`$ power counting so that the loop corrections are suppressed by the factor $`1/N_c^L`$, where $`L`$ is the number of loops.
The organization of this paper is as follows. In Sec. II, we begin with a brief overview of the $`1/N_c`$ cancellations occurring in the one-loop correction to the baryon axial vector currents. In Sec. III, we derive the formula for the one-loop correction to the baryon axial vector currents for arbitrary $`\mathrm{\Delta }/m_\pi `$, in a form that is convenient for later use. The structure of large-$`N_c`$ cancellations for $`\mathrm{\Delta }/m_\pi =0`$ and $`\mathrm{\Delta }/m_\pi 0`$ are discussed in Secs. IV and V, respectively. The general power counting for large-$`N_c`$ cancellations is derived to all orders in the baryon hyperfine mass splitting $`\mathrm{\Delta }`$, and it is determined that the dominant large-$`N_c`$ cancellations are present only in terms that are of low and finite order in $`\mathrm{\Delta }`$. A procedure for subtracting and isolating these large-$`N_c`$ cancellations is given. Other contributions to axial vector current renormalization are briefly presented in Sec. VI. Our conclusions are summarized in Sec. VII.
## II Overview
A brief review of heavy baryon chiral perturbation theory and the $`1/N_c`$ baryon chiral Lagrangian can be found in Ref. , so only a few salient facts will be repeated here. The pion-baryon vertex is proportional to $`g_A/f`$, where $`f`$ is the decay constant of the $`\pi `$ meson. In the large-$`N_c`$ limit, $`g_AN_c`$ and $`f\sqrt{N_c}`$, so that the pion-baryon vertex is of order $`\sqrt{N_c}`$ and grows with $`N_c`$. The baryon propagator is $`i/(kv)`$ and is $`N_c`$-independent, as is the pion propagator. In the $`\overline{\mathrm{MS}}`$ scheme, all loop integrals are given by the pole structure of the propagators, so loop integrals do not depend on $`N_c`$.
The tree-level matrix element of the baryon axial vector current is of order $`N_c`$, since $`g_A`$ is of order $`N_c`$. The one-loop diagrams that renormalize the baryon axial vector current are shown in Fig. 1.
Each of the one-loop corrections in Fig. 1(a,b,c) involves two pion-baryon vertices, and is order $`N_c`$ times the tree-level graph.
The matrix elements of the space components of the baryon axial vector current between initial and final baryon states $`B`$ and $`B^{}`$ will be denoted by
$$B^{}\left|\overline{\psi }\gamma ^i\gamma _5T^a\psi \right|B=\left[A^{ia}\right]_{B^{}B},$$
(3)
where $`B`$ and $`B^{}`$ are baryons in the lowest-lying irreducible representation of contracted-$`SU(6)`$ spin-flavor symmetry, i.e. the spin-$`1/2`$ octet and spin-$`3/2`$ decuplet baryons. The Feynman diagram amplitude for $`BB^{}+\pi (k)`$ is $`\left[A^{ia}\right]_{B^{}B}𝐤^i/f`$, where $`𝐤`$ is the three-momentum of the emitted pion. The time-component of the axial current has zero matrix element between static baryons, and is represented in the heavy baryon formulation by a higher dimension operator in the effective Lagrangian. The matrix elements $`\left[A^{ia}\right]_{B^{}B}`$ of the spatial components of the axial vector current can be written in terms of the octet and decuplet pion coupling constants $`F`$, $`D`$, $`𝒞`$, and $``$ , each of which is of order $`N_c`$.
The one-loop correction to the baryon axial vector current, in the limit that the Delta-nucleon mass difference is neglected, is proportional to the double commutator
$$\delta A^{ia}\frac{1}{f^2}[A^{jb},[A^{jb},A^{ia}]],$$
(4)
where the sum over intermediate baryon states is given by matrix multiplication of the $`A^{ia}`$ matrices. Naively, the double commutator is of order $`N_c^3`$, and $`f\sqrt{N_c}`$ so that $`\delta A^{ia}`$ is of order $`N_c^2`$. One of the results of the $`1/N_c`$ analysis for baryons is that the double commutator is of order $`N_c`$, rather than $`N_c^3`$ . Each individual term in the sum Eq. (4) is of order $`N_c^3`$, but there is a cancellation in the sum over intermediate baryons, which is guaranteed by the spin-flavor symmetry of large-$`N_c`$ QCD . The cancellation only occurs when the ratios of $`F`$, $`D`$, $`𝒞`$, and $``$ are close to their $`SU(6)`$ values.An important point to note is that large-$`N_c`$ QCD predicts only the ratios of $`F/D`$, $`𝒞/D`$, and $`/D`$, the overall normalization of the coupling constants is not fixed by the symmetry. The large-$`N_c`$ cancellations depend on the coupling ratios being close to their $`SU(6)`$ values, and do not depend on the overall normalization of the couplings. The large-$`N_c`$ cancellation implies that the one-loop correction to the axial current is $`1/N_c`$ times the tree-level value, instead of $`N_c`$ times the tree-level value. Similarly, the two-loop correction is $`1/N_c^2`$ times the tree-level value, instead of $`N_c^2`$ times the tree-level value (see Ref. for an explicit calculation in the degeneracy limit). The one-loop large-$`N_c`$ cancellations will be discussed more fully in Secs. IV and V. The formalism for making large-$`N_c`$ cancellations manifest is provided in Sec. V.
The large-$`N_c`$ cancellation in the one-loop correction to the baryon axial vector current can be seen numerically from explicit computation in heavy baryon chiral perturbation theory. The baryon axial vector current matrix element at one-loop has the form
$$A=\alpha +\left(\overline{\beta }\overline{\lambda }\alpha \right)\frac{m^2}{16\pi ^2f^2}\mathrm{ln}\frac{m^2}{\mu ^2}+\mathrm{}$$
(5)
where $`\alpha `$ is the tree-level contribution, $`\overline{\beta }`$ is the vertex correction, $`\overline{\lambda }`$ is the wavefunction renormalization, $`m`$ is the $`\pi `$, $`K`$ or $`\eta `$ mass, and the Delta-nucleon mass difference has been neglected for simplicity so there is only a chiral-logarithmic contribution. \[The full one-loop correction will be discussed in the next section.\] For the case of $`p\left|\overline{u}\gamma ^\mu \gamma _5d\right|n`$, the coefficients are
$`\alpha `$ $`=`$ $`D+F,`$ (6)
$`\overline{\lambda }_\pi `$ $`=`$ $`{\displaystyle \frac{9}{4}}(F+D)^2+2𝒞^2,`$ (7)
$`\overline{\lambda }_K`$ $`=`$ $`{\displaystyle \frac{1}{2}}(9F^26FD+5D^2+𝒞^2),`$ (8)
$`\overline{\lambda }_\eta `$ $`=`$ $`{\displaystyle \frac{1}{4}}(3FD)^2,`$ (9)
$`\overline{\lambda }_\eta ^{}`$ $`=`$ $`2D^2,`$ (10)
$`\overline{\beta }_\pi `$ $`=`$ $`{\displaystyle \frac{1}{4}}(F+D)^3+{\displaystyle \frac{16}{9}}(F+D)𝒞^2{\displaystyle \frac{50}{81}}𝒞^2FD,`$ (11)
$`\overline{\beta }_K`$ $`=`$ $`{\displaystyle \frac{1}{3}}(3F^3+3F^2DFD^2+D^3)+{\displaystyle \frac{2}{9}}(F+3D)𝒞^2{\displaystyle \frac{10}{81}}𝒞^2{\displaystyle \frac{1}{2}}(F+D),`$ (12)
$`\overline{\beta }_\eta `$ $`=`$ $`{\displaystyle \frac{1}{12}}(F+D)(3FD)^2,`$ (13)
$`\overline{\beta }_\eta ^{}`$ $`=`$ $`{\displaystyle \frac{1}{12}}(F+D)(3FD)^2.`$ (14)
The coefficients for the other matrix elements can be found in the literature . The subscripts $`\pi `$, $`K`$ and $`\eta `$ denote the contributions from $`\pi `$, $`K`$ and $`\eta `$ loops. To illustrate the cancellation, we have plotted the one-loop coefficients $`\left(\overline{\beta }\overline{\lambda }\alpha \right)`$
for the axial currents (or equivalently, the couplings) $`NN\pi `$, $`\mathrm{\Sigma }\mathrm{\Lambda }\pi `$, $`\mathrm{\Sigma }\mathrm{\Sigma }\pi `$, $`\mathrm{\Xi }\mathrm{\Xi }\pi `$, $`\mathrm{\Lambda }N\overline{K}`$, $`\mathrm{\Sigma }N\overline{K}`$, $`\mathrm{\Xi }\mathrm{\Lambda }\overline{K}`$, and $`\mathrm{\Xi }\mathrm{\Sigma }\overline{K}`$ in Figs. 2, 3, and 4. For simplicity, the coefficients are plotted as a function of $`F/D`$ only — the other coupling ratios have been fixed at their $`SU(6)`$ values $`𝒞/D=2`$ and $`/D=3`$. The best fit to the baryon axial currents has the axial coupling ratios close to their $`SU(6)`$ values , so this is a reasonable approximation. The large-$`N_c`$ analysis indicates that there should be some cancellation in the loop correction when $`F/D`$ is close to the $`SU(6)`$ value of $`2/3`$. This suppression is evident separately for the $`\pi `$, $`K`$ and $`\eta `$ loops for all eight processes. This is the cancellation pointed out phenomenologically in Ref. , and later proved in Refs. . We will study this cancellation quantitatively in terms of the $`1/N_c`$ expansion in this work.
## III One-loop correction to the axial current
The one-loop diagrams that contribute to the baryon axial vector current are shown in Fig. 1. Figs. 1(a,b,c) are of order $`N_c`$ times the tree-level vertex, and Fig. 1(d) is of order $`1/N_c`$ times the tree-level vertex. The large-$`N_c`$ cancellations occur between Figs. 1(a,b,c), so we will concentrate on these three diagrams in this section. The contribution from Fig. 1(d) is considered briefly in Sec. VI. Both contributions can be found in Ref. .<sup>§</sup><sup>§</sup>§Fig. 1(d) is linear in the pion-baryon coupling constants $`F`$, $`D`$ and $`𝒞`$, whereas Figs. 1(a,b,c) are cubic, so it is easy to identify the two pieces in existing calculations.
All the loop graphs we need can be written in terms of the basic loop integral
$$\delta ^{ij}F(m,\mathrm{\Delta },\mu )=\frac{i}{f^2}\frac{d^4k}{(2\pi )^4}\frac{(𝐤^i)(𝐤^j)}{(k^2m^2)(kv\mathrm{\Delta }+iϵ)},$$
(15)
where $`\mu `$ is the scale parameter of dimensional regularization. Evaluating the integral gives
$`24\pi ^2f^2F(m,\mathrm{\Delta },\mu )`$ (18)
$`=\mathrm{\Delta }\left(\mathrm{\Delta }^2{\displaystyle \frac{3}{2}}m^2\right)\mathrm{ln}{\displaystyle \frac{m^2}{\mu ^2}}{\displaystyle \frac{8}{3}}\mathrm{\Delta }^3{\displaystyle \frac{7}{2}}\mathrm{\Delta }m^2`$
$`+\{\begin{array}{cc}2\left(m^2\mathrm{\Delta }^2\right)^{3/2}\left[\frac{\pi }{2}\mathrm{tan}^1\left(\frac{\mathrm{\Delta }}{\sqrt{m^2\mathrm{\Delta }^2}}\right)\right],\hfill & \left|\mathrm{\Delta }\right|m\text{,}\hfill \\ & \\ \left(\mathrm{\Delta }^2m^2\right)^{3/2}\mathrm{ln}\left(\frac{\mathrm{\Delta }\sqrt{\mathrm{\Delta }^2m^2}}{\mathrm{\Delta }+\sqrt{\mathrm{\Delta }^2m^2}}\right),\hfill & \left|\mathrm{\Delta }\right|>m\text{ .}\hfill \end{array}`$
### A Wavefunction Renormalization
The wavefunction renormalization graph for baryon $`B`$ is shown in Fig. 5, where one sums over all possible intermediate baryons $`B_I`$. The loop graph is equal to
$`iG_B`$ $`=`$ $`{\displaystyle \underset{j,k,b,B_I}{}}{\displaystyle \frac{i^2}{f^2}}\left[A^{kb}\right]_{BB_I}\left[A^{jb}\right]_{B_IB}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{(𝐤^k)(𝐤^j)}{\left(k^2m_b^2\right)\left(\left(k+p\right)v\left(M_IM\right)+iϵ\right)}},`$ (19)
where $`b=1,\mathrm{},9`$ or $`\pi ,K,\eta ,\eta ^{}`$ labels the intermediate meson.The $`\eta ^{}`$ is a ninth Goldstone boson in the large $`N_c`$ limit . Our formulae apply to the $`\eta ^{}`$ corrections, with flavor matrix $`\lambda ^9=\sqrt{2/3}`$. The formalism for including the $`\eta ^{}`$ is described in detail in Ref. . The wavefunction renormalization correction for baryon $`B`$ is
$$Z_B=\frac{G_B}{(pv)}|_{pv=0}.$$
(20)
The wavefunction correction to the axial vector current matrix element $`B^{}\left|A^{ia}\right|B`$ depends on
$`Z_{B^{}B}={\displaystyle \frac{1}{2}}\left(Z_B^{}+Z_B\right),`$ (21)
which can be written in terms of the function $`F(m,\mathrm{\Delta },\mu )`$ defined in Eq. (15) as
$$Z_{B^{}B}=\underset{j,b,B_I}{}\left[A^{jb}\right]_{B^{}B_I}\left[A^{jb}\right]_{B_IB}\frac{F(m_b,\mathrm{\Delta }_{B_IB},\mu )}{\mathrm{\Delta }_{B_IB}},$$
(22)
where
$$\mathrm{\Delta }_{B_1B_2}M_{B_1}M_{B_2}.$$
(23)
The wavefunction renormalization correction $`Z_{B^{}B}`$ is diagonal in flavor and spin.
### B Vertex Correction
The one-loop correction to the matrix element $`B^{}\left|A^{ia}\right|B`$ from the vertex graph Fig. 6 can be written as
$`\left[\delta A^{ia}\right]_{B^{}B}^{\mathrm{vertex}}={\displaystyle \underset{j,k,b,B_1,B_2}{}}{\displaystyle \frac{i}{f^2}}\left[A^{kb}\right]_{B^{}B_2}\left[A^{ia}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}`$ (24)
$`\times {\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{(𝐤^k)(𝐤^j)}{\left(k^2m_b^2\right)\left(kv\left(M_1M\right)+iϵ\right)\left(\left(kq\right)v\left(M_2M\right)+iϵ\right)}},`$ (25)
where $`q`$ is the outgoing momentum transfer at the axial vertex. For octet-octet matrix elements, $`qv=0`$, whereas for decuplet-octet transition matrix elements, $`qv=MM^{}`$, the average decuplet-octet mass difference. One can rewrite the denominator of Eq. (25) using the identity
$`{\displaystyle \frac{1}{(k^0\mathrm{\Delta }_1+iϵ)(k^0\mathrm{\Delta }_2+iϵ)}}={\displaystyle \frac{1}{(\mathrm{\Delta }_1\mathrm{\Delta }_2)}}\left[{\displaystyle \frac{1}{(k^0\mathrm{\Delta }_1+iϵ)}}{\displaystyle \frac{1}{(k^0\mathrm{\Delta }_2+iϵ)}}\right]`$ (26)
so that
$`\left[\delta A^{ia}\right]_{B^{}B}^{\mathrm{vertex}}`$ $`=`$ $`{\displaystyle \underset{j,b,B_1,B_2}{}}\left[A^{jb}\right]_{B^{}B_2}\left[A^{ia}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}`$ (28)
$`\times {\displaystyle \frac{1}{\mathrm{\Delta }_{B_1B}\mathrm{\Delta }_{B_2B^{}}}}\left[F(m_b,\mathrm{\Delta }_{B_1B},\mu )F(m_b,\mathrm{\Delta }_{B_2B^{}},\mu )\right]`$
where $`\mathrm{\Delta }_{B_1B_2}`$ is defined in Eq. (23).
### C Total correction
The total correction to the baryon axial vector current matrix element from Figs. 1(a,b,c) is
$`\left[\delta A^{ia}\right]_{B^{}B}`$ $`=`$ $`\left[\delta A^{ia}\right]_{B^{}B}^{\mathrm{vertex}}+{\displaystyle \frac{1}{2}}\left\{{\displaystyle \underset{B_1}{}}Z_{B^{}B_1}\left[A^{ia}\right]_{B_1B}+{\displaystyle \underset{B_2}{}}\left[A^{ia}\right]_{B^{}B_2}Z_{B_2B}\right\}`$ (29)
$`=`$ $`{\displaystyle \underset{j,b,B_1,B_2}{}}\left[A^{jb}\right]_{B^{}B_2}\left[A^{ia}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}{\displaystyle \frac{F(m_b,\mathrm{\Delta }_{B_1B},\mu )F(m_b,\mathrm{\Delta }_{B_2B^{}},\mu )}{\mathrm{\Delta }_{B_1B}\mathrm{\Delta }_{B_2B^{}}}}`$ (32)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,b,B_1,B_2}{}}\left[A^{jb}\right]_{B^{}B_2}\left[A^{jb}\right]_{B^{}B_2}\left[A^{ia}\right]_{B_2B_1}{\displaystyle \frac{F(m_b,\mathrm{\Delta }_{B_2B_1},\mu )}{\mathrm{\Delta }_{B_2B_1}}}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,b,B_1,B_2}{}}\left[A^{ia}\right]_{B^{}B_2}\left[A^{jb}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}{\displaystyle \frac{F(m_b,\mathrm{\Delta }_{B_1B},\mu )}{\mathrm{\Delta }_{B_1B}}}.`$
(In addition to $`\delta A^{ia}`$ of Eq. (32), there are also the contributions of Fig. 1(d) and the low-energy constants, which are considered in Sec. VI.) Eq. (32) includes the full dependence on $`\mathrm{\Delta }/m`$ of the one-loop correction. We want to rewrite this expression so that the large-$`N_c`$ cancellations are manifest.
In the limit that the octet and decuplet baryons are degenerate, all the mass differences $`\mathrm{\Delta }_{AB}0`$, and
$$\frac{1}{\mathrm{\Delta }_{B_1B}\mathrm{\Delta }_{B_2B^{}}}\left[F(m_b,\mathrm{\Delta }_{B_1B},\mu )F(m_b,\mathrm{\Delta }_{B_2B^{}},\mu )\right]F^{(1)}(m_b,0,\mu ),$$
(33)
where $`F^{(n)}`$ is defined by
$$F^{(n)}(m_b,\mathrm{\Delta },\mu )\frac{^nF(m_b,\mathrm{\Delta },\mu )}{\mathrm{\Delta }^n}.$$
(34)
In this limit, the correction to the axial current Eq. (32) reduces to
$`\left[\delta A^{ia}\right]_{B^{}B}`$ $`=`$ $`{\displaystyle \underset{j,b,B_1,B_2}{}}F^{(1)}(m_b,0,\mu )`$ (37)
$`\{\left[A^{jb}\right]_{B^{}B_2}\left[A^{ia}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}`$
$`+{\displaystyle \frac{1}{2}}\left[A^{jb}\right]_{B^{}B_2}\left[A^{jb}\right]_{B_2B_1}\left[A^{ia}\right]_{B_1B}+{\displaystyle \frac{1}{2}}\left[A^{ia}\right]_{B^{}B_2}\left[A^{jb}\right]_{B_2B_1}\left[A^{jb}\right]_{B_1B}\}`$
Let us adopt the more compact notation that $`A^{ia}`$ represents a matrix with matrix elements $`\left[A^{ia}\right]_{B^{}B}`$, and summation over intermediate baryon states is denoted by matrix multiplication. Then Eq. (37) can be written as
$`\delta A^{ia}`$ $`=`$ $`{\displaystyle \underset{b,j}{}}F^{(1)}(m_b,0,\mu )\left\{A^{jb}A^{ia}A^{jb}+{\displaystyle \frac{1}{2}}A^{jb}A^{jb}A^{ia}+{\displaystyle \frac{1}{2}}A^{ia}A^{jb}A^{jb}\right\}`$ (38)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{b,j}{}}F^{(1)}(m_b,0,\mu )[A^{jb},[A^{jb},A^{ia}]],`$ (39)
which is the double-commutator form originally derived in Ref. .
The loop integral in the degeneracy limit $`\mathrm{\Delta }0`$ reduces to
$$F^{(1)}(m_b,0,\mu )=\frac{1}{16\pi ^2f^2}m_b^2\left(\frac{11}{3}+\mathrm{ln}\frac{m_b^2}{\mu ^2}\right).$$
(40)
The $`\mathrm{ln}m_b/\mu `$ term is non-analytic in the quark mass, and is called a “chiral log.” The constant ($`11/3`$) piece is analytic in the quark masses, and has the same form as higher dimension terms in the chiral Lagrangian. The constant term is scheme-dependent, but the chiral logarithm is universal.
We discuss the structure of the large-$`N_c`$ cancellations for the baryon axial vector currents in the next two sections. First, in Sec. IV, the cancellations are studied in the degeneracy limit. The generalization to non-degenerate baryons is given in Sec. V.
## IV Large-$`N_c`$ cancellations: $`\mathrm{\Delta }/m_\pi =0`$
The large-$`N_c`$ cancellations in the degeneracy limit for the one-loop correction to the baryon axial vector current follow from the double-commutator form of Eq. (39). The pion decay constant $`f\sqrt{N_c}`$, so the function $`F^{(1)}(m_b,0,\mu )`$ is of order $`1/N_c`$. Each axial vector current matrix element is of order $`N_c`$ (recall that $`g_A`$ is of order $`N_c`$), so the correction $`\delta A^{ia}`$ is naively of absolute order $`N_c^2`$, i.e. of order $`N_c`$ relative to the tree-level value $`A^{ia}`$. The large-$`N_c`$ consistency conditions derived in Ref. imply that the double commutator $`[A^{jb},[A^{jb},A^{ia}]]`$ is of order $`N_c`$ rather the naive order $`N_c^3`$, *provided one sums over all baryon states in a complete multiplet of the large-$`N_c`$ $`SU(6)`$ spin-flavor symmetry, i.e. over both the octet and decuplet, and uses axial coupling ratios given by the large-$`N_c`$ spin-flavor symmetry.*
Before discussing the cancellation in the double commutator, we first review some necessary large-$`N_c`$ formalism. The baryon matrix element of the axial vector current in QCD can be expanded in a $`1/N_c`$ expansion in terms of $`SU(6)`$ spin-flavor operators , For recent reviews of the large-$`N_c`$ spin-flavor symmetry, see Refs. .
$$G^{ia}=q^{}\frac{\sigma ^i}{2}\frac{\lambda ^a}{2}q,T^a=q^{}\frac{\lambda ^a}{2}q,J^i=q^{}\frac{\sigma ^i}{2}q,$$
(41)
where $`q`$ and $`q^{}`$ are $`SU(6)`$ operators that create and annihilate states in the fundamental representation of $`SU(6)`$, and $`\sigma ^i`$ and $`\lambda ^a`$ are the Pauli spin and Gell-Mann flavor matrices. The lowest mass baryon multiplet transforms under $`SU(6)`$ as a completely symmetric tensor with $`N_c`$ indices. For $`N_c=3`$, this representation decomposes under spin and flavor into a spin-1/2 octet and a spin-3/2 decuplet. The baryon axial vector current $`A^{ia}`$ in the large-$`N_c`$ limit has the form
$`A^{ia}=a_1G^{ia}+{\displaystyle \underset{n=2,3}{\overset{N_c}{}}}b_n{\displaystyle \frac{1}{N_c^{n1}}}𝒟_n^{ia}+{\displaystyle \underset{n=3,5}{\overset{N_c}{}}}c_n{\displaystyle \frac{1}{N_c^{n1}}}𝒪_n^{ia},`$ (42)
where the coefficients are of order one. The operators $`𝒟_n^{ia}`$ are diagonal operators with nonzero matrix elements only between states with the same spin, and the operators $`𝒪_n^{ia}`$ are purely off-diagonal operators with nonzero matrix elements only between states of different spin. The explicit forms for these operators can be found in Ref. . At the physical value $`N_c=3`$, Eq. (42) reduces to
$`A^{ia}`$ $`=`$ $`a_1G^{ia}+b_2{\displaystyle \frac{1}{N_c}}J^iT^a+b_3{\displaystyle \frac{1}{N_c^2}}𝒟_3^{ia}+c_3{\displaystyle \frac{1}{N_c^2}}𝒪_3^{ia},`$ (43)
where
$`𝒟_3^{ia}`$ $`=`$ $`\{J^i,\{J^j,G^{ja}\}\},`$ (44)
$`𝒪_3^{ia}`$ $`=`$ $`\{J^2,G^{ia}\}{\displaystyle \frac{1}{2}}\{J^i,\{J^j,G^{ja}\}\}.`$ (45)
The four conventional $`SU(3)`$ baryon axial couplings $`F`$, $`D`$, $`𝒞`$ and $``$ for the baryon octet and decuplet can be written as linear combinations of the coefficients $`a_1`$, $`b_2`$, $`b_3`$ and $`c_3`$ of the $`1/N_c`$ expansion,
$`\begin{array}{c}D={\displaystyle \frac{1}{2}}a_1+{\displaystyle \frac{1}{6}}b_3,\hfill \\ F={\displaystyle \frac{1}{3}}a_1+{\displaystyle \frac{1}{6}}b_2+{\displaystyle \frac{1}{9}}b_3,\hfill \\ 𝒞=a_1{\displaystyle \frac{1}{2}}c_3,\hfill \\ ={\displaystyle \frac{3}{2}}a_1{\displaystyle \frac{3}{2}}b_2{\displaystyle \frac{5}{2}}b_3.\hfill \end{array}`$ (50)
The leading order prediction of large-$`N_c`$ QCD is obtained by dropping the $`1/N_c`$ suppressed terms in Eq. (42), i.e. the 3-body operators $`𝒟_3^{ia}`$ and $`𝒪_3^{ia}`$. The $`G^{ia}`$ operator gives $`F=2D/3`$, $`𝒞=2D`$ and $`=3D`$, so that the coupling ratios, but not necessarily their absolute normalization, are those predicted by $`SU(6)`$ symmetry. The 2-body operator $`J^iT^a`$ corrects these relations. The correction is of relative order $`1/N_c^2`$ for pions.
The baryon matrix elements of $`J^i`$ for the low-lying baryons in the $`SU(6)`$ representation are of order unity. The $`N_c`$ dependence of matrix elements of $`G^{ia}`$ and $`T^a`$ is more subtle, and depends on the particular component $`a`$ chosen, as well as on the initial and final state baryon . For the purposes of this paper, we will use the naive estimate that matrix elements of $`G^{ia}`$ and $`T^a`$ are both of order $`N_c`$, which is the largest they can be. We focus upon baryons with spins of order unity. The $`N_c`$ counting rules are summarized as:
$$G^{ia}N_c,T^aN_c,J^i1.$$
(51)
The $`1/N_c`$ expansion of a baryonic matrix element can be written as an expansion in powers of $`G^{ia}/N_c`$, $`T^a/N_c`$ and $`J^i/N_c`$. The counting rules Eq. (51) show that each factor $`J`$ leads to a $`1/N_c`$ suppression factor.
We can now understand the origin of the large-$`N_c`$ cancellations in Eq. (39). At leading order in $`N_c`$, the axial current operator $`A^{ia}`$ can be replaced by $`a_1G^{ia}`$, and has matrix elements of order $`N_c`$. The commutator $`[A^{ia},A^{jb}]=[a_1G^{ia},a_1G^{jb}]`$ is naively of order $`N_c^2`$, since each $`G^{ia}`$ is of order $`N_c`$. However, the commutation relation
$$[G^{ia},G^{jb}]=\frac{i}{4}\delta ^{ij}f^{abc}T^c+\frac{i}{6}\delta ^{ab}ϵ^{ijk}J^k+\frac{i}{2}d^{abc}ϵ^{ijk}G^{kc},$$
(52)
shows that matrix elements of the commutator $`[G^{ia},G^{jb}]`$ are at most of order $`N_c`$, since the right-hand side of Eq. (52) is at most of order $`N_c`$. Thus there is a factor of $`N_c`$ cancellation between the various terms in the commutator $`[G^{ia},G^{jb}]`$ from the summation over intermediate baryon states. Similarly, one finds that there is a factor of $`N_c`$ cancellation in the sum over intermediate states for the commutators
$$[T^a,G^{ib}]=if^{abc}G^{ic},$$
(53)
and
$$[T^a,T^b]=if^{abc}T^c,$$
(54)
where the naive counting rule Eq. (51) has been used to estimate the order in $`N_c`$ of both sides of these equations. The basic reason for the cancellation is that the maximum order in $`N_c`$ an $`r`$-body operator matrix element can be is $`N_c^r`$, (an $`r`$-body operator is one with $`r`$ $`q`$’s and $`r`$ $`q^{}`$’s, i.e. can be written as a polynomial of order $`r`$ in $`J^i`$, $`G^{ia}`$ and $`T^a`$), but the commutator of an $`r`$-body and $`s`$-body operator is at most an $`r+s1`$ body operator. Thus, every commutator potentially leads to a cancellation by one factor of $`N_c`$. However, not every commutator gives a factor of $`N_c`$ cancellation. The commutators
$$[J^i,J^j]=iϵ^{ijk}J^k,$$
(55)
and
$$[J^i,G^{ja}]=iϵ^{ijk}G^{ka},$$
(56)
have no cancellations, since both sides are of order one and order $`N_c`$, respectively. The reason that there is no cancellation in Eqs. (55) and (56) is that $`J^i`$ is a one-body operator whose matrix elements are of order unity, rather than of order $`N_c`$.
Equations (53)–(56) lead to the conclusion that each commutator produces a $`N_c`$ cancellation, unless a factor of $`J^i`$ is eliminated. The double-commutator in Eq. (39) has a cancellation of $`N_c^2`$, because $`[G^{jb},[G^{jb},G^{ia}]]J+G+T`$, so that the double-commutator is order $`N_c`$, rather than $`N_c^3`$. This was the cancellation observed numerically in Ref. , and later proven in the $`1/N_c`$ expansion in Refs. .
## V Large-$`N_c`$ cancellations: $`\mathrm{\Delta }/m_\pi 0`$
In this section, we analyze the large-$`N_c`$ cancellations in the renormalization of the baryon axial vector current for finite $`\mathrm{\Delta }/m_\pi `$.
Equation (32) can be expanded in a power series in $`\mathrm{\Delta }`$. Expanding the function $`F(m,\mathrm{\Delta },\mu )`$ in a power series, and collecting terms gives
$`\delta A^{ia}={\displaystyle \underset{j,b}{}}\{{\displaystyle \frac{1}{2}}F^{(1)}(m_b,0,\mu )[A^{jb},[A^{jb},A^{ia}]]+{\displaystyle \frac{1}{2}}F^{(2)}(m_b,0,\mu )\{A^{jb},[A^{ia},[,A^{jb}]]\}`$ (57)
$`+{\displaystyle \frac{1}{6}}F^{(3)}(m_b,0,\mu )([A^{jb},[[,[,A^{jb}]],A^{ia}]]+{\displaystyle \frac{1}{2}}[[,A^{jb}],[[,A^{jb}],A^{ia}]])+\mathrm{}\}`$ (58)
where $``$ is the baryon mass matrix. In deriving this result we have converted explicit sums over intermediate baryons to implicit sums in the matrix multiplications. One can use either the baryon mass matrix $``$, or the baryon mass-splitting matrix $`\mathrm{\Delta }`$ in Eq. (34), since $``$ differs from $`\mathrm{\Delta }`$ by the average baryon mass times the unit matrix, which commutes and drops out of Eq. (57). To evaluate Eq. (57) to all orders in $`\mathrm{\Delta }/m_\pi `$ would be extremely difficult, since one would have to sum an infinite series, with each term having a coefficient which is a complicated commutator/anticommutator of $``$’s and $`A^{ia}`$’s.
We would like to evaluate graphs in heavy baryon chiral perturbation theory so that the $`1/N_c`$ cancellations are manifest, and do not occur as numerical cancellations at the end of the calculation. We will now show that the large $`N_c`$ cancellations only occur in the first few terms of Eq. (57), so that the remaining terms can be summed using conventional heavy baryon chiral perturbation theory in the usual manner.
The expansion Eq. (57) has a different structure depending on whether one has an even or odd number of insertions of the baryon mass operator $``$. Terms with $`2r`$ insertions of $``$ have $`2r+2`$ commutators, whereas terms with $`2r+1`$ insertions of $``$ have $`2r+2`$ commutators and one anticommutator.
The general form of the baryon mass operator in the $`1/N_c`$ expansion in the $`SU(3)`$ limit is
$$=N_c\left[m_0+m_1\frac{J^2}{N_c^2}+m_2\frac{(J^2)^2}{N_c^4}+\mathrm{}\right].$$
(59)
The importance of a given term in the $`1/N_c`$ expansion can be obtained by counting powers of $`J`$. Each factor of $`J/N_c`$ leads to a $`1/N_c`$ suppression, since $`J`$ is of order unity according to the counting rules Eq. (51).
We now have all the necessary ingredients to count the power in $`1/N_c`$ of a general term in Eq. (57). The operators $`A^{ia}`$ and $``$ are one-body operators, with naive order $`N_c`$, so the $`^r`$ term in the expansion in Eq. (57) is naively of order $`N^{3+r1}`$, including the factor of $`1/N_c`$ from the $`1/f^2`$ in the loop integral $`F`$, as shown in row (C) of Table I.
The number of commutators in each term is listed in the next line in this table. Every commutator (naively) leads to a decrease in the naive $`N_c`$ order by unity, since the commutator of an $`r`$-body and an $`s`$-body operator is at most an $`r+s1`$ body operator. This leads to the $`N_c`$ power given in row (E). Finally, we need to count the powers of $`J`$ in each term. Each factor of $``$ has at least two $`J`$’s, since the $`m_0N_c`$ term in Eq. (59) drops out of the expression Eq. (57). Each factor of $`A^{ia}`$ can have $`0`$ $`J`$’s, as is clear from Eqs. (42)–(45). The number of $`J`$’s in the original expression is listed in row (F), where $`p0`$ is the number of extra $`J`$’s from $``$ or $`A^{ia}`$, beyond the minimum values of two and zero, respectively. Finally, note that each commutator can be used to eliminate one power of $`J`$. Thus the net power of $`J`$ left is given by subtracting the number of commutators from the number of $`J`$’s. The minimum number of $`J`$’s is non-negative, and is listed in row (G). Thus, the final $`N_c`$ power (row (H)) is given by the net power in row (E) minus the minimum number of $`J`$’s in row (G) since there is an additional $`1/N_c`$ factor for each $`J`$. One can compare this with the “usual” $`N_c`$ counting rule listed in row (I). The usual counting rule is obtained by including a factor of $`N_c`$ for each $`A^{ia}`$ (i.e. for each factor of $`F`$, $`D`$, $`𝒞`$ or $``$), a factor of $`1/N_c`$ for the $`1/f^2`$, and a factor of $`1/N_c`$ for each power of $`\mathrm{\Delta }`$, with $`p0`$ representing $`1/N_c`$ suppressed terms.
One interesting point can be noted from Table I. The dominant $`1/N_c`$ corrections from the baryon mass splittings are due to multiple insertions of the $`J^2`$ term in the baryon mass matrix. Two insertions of the $`J^2`$ term (the $`p=0`$ term in the $`^2`$ column) is $`N_c`$ more important than one insertion of the $`J^4`$ term (the $`p=2`$ term in the $`^1`$ column).
There is an extra cancellation in the term linear in $``$ that is not apparent in Table I. We will discuss this new large-$`N_c`$ cancellation momentarily. Including this effect, one sees from Table I that all terms in the expansion of Eq. (32) with two or more powers of $``$ have the same $`N_c`$ behavior as one finds with the usual $`N_c`$ counting, i.e. these terms have no extra cancellations. One can therefore treat all terms with two or more powers of $``$ by conventional heavy baryon chiral perturbation theory—compute all the graphs, with vertices written in terms of $`F`$, $`D`$, $`𝒞`$ and $``$. The only terms that have to be treated specially are those with zero or one power of $``$. To compute graphs in the conventional way omitting the first two terms in Eq. (57) is trivial; one simply rewrites the loop integral Eq. (15) by explicitly extracting the first three terms in an expansion in $`\mathrm{\Delta }`$,
$$F(m_b,\mathrm{\Delta },\mu )=F(m_b,0,\mu )+F^{(1)}(m_b,0,\mu )\mathrm{\Delta }+\frac{1}{2}F^{(2)}(m_b,0,\mu )\mathrm{\Delta }^2+\stackrel{~}{F}(m_b,\mathrm{\Delta },\mu ).$$
(60)
and takes the standard expressions for the loop corrections written in terms of $`F`$, $`D`$, $`𝒞`$ and $``$, with $`F\stackrel{~}{F}`$. This procedure sums the entire series in $`\left(\mathrm{\Delta }/m_\pi \right)^r`$ starting with $`r=3`$. To this result is added the first two terms in Eq. (57). One needs to extract three terms from $`F`$ to obtain the first two terms in Eq. (57), since since $`F(m_b,0,\mu )`$ cancels out of the correction.
One can now analyze the first two terms in the expansion of Eq. (32), which are given in Eq. (57). The first term is the double commutator term discussed in the previous section. We see from Table I that this term is naively of order $`N_c^2`$, but actually is at most of order $`N_c^0`$. This is consistent with the loop expansion being an expansion in $`\mathrm{}/N_c`$, since the one-loop correction is of order $`1/N_c`$ relative to the tree-level contribution of order $`N_c`$. It is also apparent from Table I that all terms of order $`^0`$ with $`p=0,1,2`$ are equally important. Since there are no powers of the baryon mass operator, the $`p`$ factors of $`J`$ must all arise from $`1/N_c`$ corrections in the axial vertices $`A^{ia}`$. The expression for the axial current relevant for $`N_c=3`$ is given in Eq. (43), from which it follows that terms with $`p=0,1,2`$ in the product $`AAA`$ are of the form $`GGG`$, $`GGJT`$, $`GJTJT`$, $`GG𝒟_3`$, and $`GG𝒪_3`$. All these terms contribute at the same order to the double-commutator, whereas according to the usual counting one would have expected the $`p=0`$ product $`GGG`$ to be one power of $`N_c`$ more important that the $`p=1`$ product $`GGJT`$, which in turn would be more important by one power of $`N_c`$ than the $`p=2`$ products. This result has an important consequence: the one-loop correction is very sensitive to the deviations of the axial coupling ratios from their $`SU(6)`$ values. While the deviations are small corrections to the couplings themselves, their importance gets enhanced in the one-loop coefficient, because the leading term (proportional to $`a_1^3`$) is $`1/N_c^2`$ suppressed. Thus, for example, the $`a_1`$ term is the dominant contribution to $`F`$, $`D`$, $`𝒞`$ and $``$, and the $`c_3`$ term is a $`1/N_c^2`$ correction, but the $`a_1^3`$ and $`a_1^2c_3`$ terms are just as important in the one-loop correction. Explicit forms for the one-loop correction in terms of $`a_1`$, $`b_2`$, $`b_3`$ and $`c_3`$ will be given elsewhere.
The second term in Eq. (57) is
$$\frac{1}{2}F^{(2)}(m_b,0,\mu )\{A^{jb},[A^{ia},[,A^{jb}]]\}$$
(61)
and is at most of order $`N_c`$ (using the value $`(1p)`$ for $`p=0`$), the same order as the tree-level contribution to the axial current. This result is surprising, because the $`1/N_c`$ expansion is a semiclassical expansion in $`\mathrm{}/N_c`$. One should be able to obtain the leading in $`1/N_c`$ contributions from classical field theory. For example, it was shown that the $`N_cm_q^{3/2}`$ one-loop correction to the baryon mass could be obtained from the energy of the pion cloud coupled to a classical baryon source . The term in Eq. (61) involves the baryon mass splitting, which is a quantum effect. The order $`N_c`$ contribution to Eq. (61) comes from using $`A^{ia}=a_1G^{ia}`$, and $`\mathrm{\Delta }=m_2J^2/N_c`$,
$$\frac{a_1^3m_2}{2N_c}F^{(2)}(m_b,0,\mu )\{G^{jb},[G^{ia},[J^2,G^{jb}]]\}.$$
(62)
The operator factor $`\{G^{jb},[G^{ia},[J^2,G^{jb}]]\}`$ is naively of order $`N_c^3`$, which implies that the correction Eq. (62) is an order $`N_c`$ correction to the axial currents, since $`F^{(2)}(m_b,0,\mu )`$ is of order $`1/N_c`$. However, an explicit computation of the operator product using the identities in Ref. gives
$`\{G^{jb},[G^{ia},[J^2,G^{jb}]]\}`$ $`=`$ $`\{J^2,G^{ia}\}+{\displaystyle \frac{1}{2}}\left(N_f+N_c\right)J^iT^a{\displaystyle \frac{1}{2}}\left(N_f2\right)G^{ia},`$ (63)
which is only of order $`N_c^2`$, using the $`N_c`$-counting rules in Eq. (51). The order $`N_c^3`$ part of Eq. (63) vanishes, which is a new cancellation in the one-loop correction to the axial vector current. Consequently, Eq. (61) is of order $`N_c^0`$ rather than order $`N_c`$, and is consistent with being a quantum correction.
## VI Other contributions
We have computed the chiral logarithmic correction to the axial vector current from Figs. 1(a,b,c). There is also the contribution from Fig. 1(d), which is
$$\delta A^{ia}=\frac{1}{2}\underset{b}{}[T^b,[T^b,A^{ia}]]I(m_b),$$
(64)
where
$`I(m_b)`$ $`=`$ $`{\displaystyle \frac{i}{f^2}}{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{1}{k^2m_b^2}}={\displaystyle \frac{m_b^2}{16\pi ^2f^2}}\left(\mathrm{ln}m_b^2/\mu ^21\right).`$ (65)
This contribution is of order $`1/N_c`$ relative to the tree-level contribution, and does not involve any cancellations between Delta and nucleon states.
In addition to the loop corrections, one has the contribution from low-energy constants multiplying higher dimension operators in the heavy baryon chiral Lagrangian. These terms are analytic in the quark mass $`m_q`$. The analytic contributions from the chiral Lagrangian can be of order $`N_c`$, i.e. the same order in $`N_c`$ as the tree-level contribution.
## VII Conclusions
We have shown how to rewrite loop corrections in heavy baryon chiral perturbation theory so as to include the full dependence on the Delta-nucleon mass difference, while at the same time including the cancellations that follow from the large-$`N_c`$ spin-flavor symmetry of baryons. The treatment in this paper has included the decuplet-octet mass difference, but neglected the $`SU(3)`$ splittings of the octet and decuplet baryons. It is possible to generalize our analysis by including the $`SU(3)`$ mass splittings in the baryon mass operator $``$.
The one-loop correction to the baryon axial vector currents is very sensitive to deviations of the axial couplings from their $`SU(6)`$ symmetry ratios, since the correction that depends only on the $`SU(6)`$ coupling ratios (the $`GGG`$ term) is suppressed by $`1/N_c^2`$, and the first subleading correction (the $`GGJT`$ term) is suppressed by $`1/N_c`$. Thus, the normally second subleading terms with two powers of $`J`$ in the axial currents are as important as these two contributions. We also have found a new cancellation in the one-loop correction to the baryon axial vector current in the term linear in the baryon mass splittings.
The large-$`N_c`$ cancellations play an important role in the one-loop corrections to the axial vector current, and become more important at higher loops. They also play an important role in the one-loop corrections to other baryon properties, such as the baryon masses . At one loop, the $`m_q^{3/2}`$ correction to the baryon mass from Fig. 7
is of order $`N_c`$, the same order as the tree-level baryon mass term, and there is no cancellation between nucleon and Delta states. However, at two loops, the graphs in Fig. 8 produce $`m_q^{5/2}`$ corrections to the baryon mass that are formally of order $`N_c^2`$, but have cancellations which make the net correction of order one.
###### Acknowledgements.
This work was supported in part by the Department of Energy under Grant No. DOE-FG03-97ER40546. R.F.M. was supported by CONACYT (Mexico) under the UC-CONACYT agreement of cooperation and by CINVESTAV (Mexico). C.P.H. acknowledges support from the Schweizerischer Nationalfonds and Holderbank-Stiftung. E.J. was supported in part by the Alfred P. Sloan Foundation and by the National Young Investigator program through Grant No. PHY-9457911 from the National Science Foundation. |
warning/0001/astro-ph0001349.html | ar5iv | text | # Infrared Imaging of 𝑧= 2.43 Radio Galaxy B3 0731+438 with the Subaru Telescope — Detection of H𝛼 Ionization Cones of a Powerful Radio Galaxy
## 1. Introduction
It is well known that in high-redshift radio galaxies, the elongated morphologies of the emission lines and rest-frame UV continuum light are often aligned with the radio axes (McCarthy et al. 1987; Chambers et al. 1987). The cause of this “alignment-effect” is still unclear, although it is almost certain that this phenomenon is due to the anisotropic structures of AGNs at the centers of the galaxies. Several hypotheses have been suggested, such as induced star formation by the passage of radio jets (Rees 1989; Begelman, Cioffi 1989; de Young 1989), inverse-Compton scattering of cosmic-background radiation (Daly 1992a, 1992b), dust or electron scattering of anisotropic radiation from a hidden AGN (di Serego Alighieri et al. 1989; Fabian 1989), and nebular emission from surrounding clouds ionized by UV radiation from a central engine (Dickson et al. 1995).
To examine the structure and properties of the H$`\alpha `$ emission-line region in a high-redshift radio galaxy, we carried out high-resolution infrared imaging observations of a powerful radio galaxy, B3 0731+438, at $`z=`$ 2.429 in the $`K^{}`$-band and the 2.25 $`\mu \mathrm{m}`$ narrow-band. This object is a typical FR II radio galaxy with strong double-lobed radio-emitting hot spots and a central core (Carilli et al. 1997). The optical $`R`$-band image shows a diffuse morphology with a few knots, while the narrow-band image presents aligned Ly$`\alpha `$ emission clouds (McCarthy 1991). The optical spectrum displays an extremely strong L$`\alpha `$ emission line with a rest-frame equivalent width of 900 $`\mathrm{\AA }`$ (McCarthy 1991), which is one of the the largest values among known high-$`z`$ radio galaxies. Infrared spectra reveal a strong H$`\alpha `$+\[N ii\] emission line (McCarthy et al. 1992; Eales, Rawlings 1993; Evans 1998), contributing $`20`$$`30\%`$ flux in the $`K`$-band. Diagnostic emission-line ratios are consistent with the presence of a Seyfert 2 nucleus (Evans 1998).
We describe our observations and data reductions in section 2, and present the results in section 3. Discussions are presented in section 4 and we summarize them in section 5. Cosmological constants are assumed to be $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.1`$ throughout this paper. The scale at $`z`$ = 2.429 is thus 11.3 kpc/<sup>′′</sup>, and the look-back time is 13 Gyr, while the age of the universe is 16.5 Gyr.
## 2. Observations and Data Reduction
### 2.1. Observations
B3 0731+438 was observed using a 2.25 $`\mu `$m narrow-band filter ($`\lambda /\mathrm{\Delta }\lambda =100`$) on 1999 February 25 and 27, and using a $`K^{}`$-band filter on February 27, with a near-infrared camera CISCO (Motohara et al. 1998) mounted on the Cassegrain focus of the Subaru telescope. The total field of view of the camera was $`2^{}\times 2^{}`$ with a pixel scale of $`0.^{\prime \prime }116`$ pixel<sup>-1</sup>. The exposure time of a single frame was 60 s for the narrow-band and 20 s for the $`K^{}`$-band images. To subtract the background sky emission, we nodded the telescope slightly($`10^{\prime \prime }`$) and acquired 6 (narrow-band) or 12 ($`K^{}`$-band) frames at each of eight different positions, resulting in 48 or 96 frames per each dataset, respectively. The total exposure time was 5760 s for the narrow-band and 1920 s for the $`K^{}`$-band. The atmospheric conditions were photometric and the seeing was stable during the observation, which was 0$`.^{\prime \prime }`$6 on February 25 and 0$`.^{\prime \prime }`$4 on 27. The logs of the observations are summarized in table 1.
### 2.2. Data Reduction
We carried out data reduction of the $`K^{}`$-band frames as follows. First, a “standard $`K^{}`$ flat” frame was produced by median-averaging 250 available $`K^{}`$-band frames taken on February 23, 25, and 27, excluding the B3 0731+438 frames. Then, every B3 0731+438 frame was divided by the standard $`K^{}`$ flat frame. After applying cosmetic corrections for bad pixels, we made a “sky” frame by median-averaging the frames of which prominent objects were masked and replaced by the surrounding sky value with appropriate noise. We smoothed this sky frame with a 20$`\times `$20 pixels boxcar filter, normalized the average of the pixel value to unity and multiplied it by the standard $`K^{}`$ flat frame to make a “self-calibrated $`K^{}`$ flat” frame.
Using this self-calibrated $`K^{}`$ flat frame, a sky frame was generated again by following the above procedure from the beginning, and was subtracted from each frame. The 12 frames in each set were shifted by a sub-pixel offset according to a reference star in the frames and averaged to create a “minor” frame. The final $`K^{}`$-band image was produced by median-averaging these minor frames. The seeing size of the final image was 0$`.^{\prime \prime }`$4.
Narrow-band frames were processed in the same way as for the $`K^{}`$-band, except that no self-calibrated flat frame was produced, due to the small sky background flux in this band. Two versions of final narrow-band images were made: one was a “total” image produced from all the frames; the other was a “good-seeing” image produced only from the frames taken on February 27. The seeing size was 0$`.^{\prime \prime }`$6 for the total image and 0$`.^{\prime \prime }`$4 for the good-seeing image. We created a color image of B3 0731+438 from the $`K^{}`$-band and the “good-seeing” narrow-band image, which is shown in figure 1.
To investigate the emission-line and continuum properties of B3 0731+438, we made two post-processed images. One was a line-free continuum image ($`K`$-continuum image) produced by subtracting the scaled narrow-band image from the $`K^{}`$-band image. The other was a continuum-subtracted emission-line image (H$`\alpha `$+\[N ii\] image) produced by subtracting the scaled $`K`$-continuum image from the total narrow-band image. We present these images in figure 2 together with their contours.
## 3. Results
### 3.1. Photometry
The results of aperture photometry of B3 0731+438 are given in table 2 together with that of other observations. The aperture is $`9.^{\prime \prime }4`$ diameter. The flux calibration was done using images of UKIRT faint standard stars (FS 15, FS 21, and FS 23) taken before or after the observations. The $`K^{}`$\- and narrow-band flux of the standard stars were calculated by interpolating their $`K`$\- and $`H`$-band flux.
### 3.2. The $`K`$-Continuum Image
The $`K`$-continuum image of B3 0731+438 in figure 2b appears to comprise a very compact core and an extended diffuse nebula, which can be seen in the deconvolved image (figure 2c) more clearly. The compact core couldn’t be resolved even with a spatial resolution as high as 0$`.^{\prime \prime }`$4. Therefore, we modeled the $`K`$-continuum image with a two-component profile consisting of a stellar core and an exponential disk.
The profile of the core is assumed to be that of a field star fitted by a modified moffat function. The moffat function is written as
$$I(x,y)=I_c\left[1+\left(\frac{r}{\alpha }\right)^\beta \right],$$
(1)
where $`r`$ is the radial distance from the center. However, the stellar image is distorted elliptically because of imperfect operation of the telescope (mirror support, auto-guiding and so on) during the observations. We therefore had to introduce a PSF modification, defined as $`r^2=X^2+Y^2`$ and
$$\left(\begin{array}{c}X\\ Y\end{array}\right)=\left(\begin{array}{cc}\sqrt{1+e}& 0\\ 0& \sqrt{1e}\end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\mathrm{\Theta }& \mathrm{sin}\mathrm{\Theta }\\ \mathrm{sin}\mathrm{\Theta }& \mathrm{cos}\mathrm{\Theta }\end{array}\right)\left(\begin{array}{c}x\\ y\end{array}\right),$$
(2)
where $`\mathrm{\Theta }`$ and $`e`$ are the position angle and the ellipticity of a stellar image, respectively. The values of these parameters for the field stars are $`\alpha =`$ 0$`.^{\prime \prime }`$42, $`\beta =3.1`$, $`\mathrm{\Theta }=45^{}`$, and $`e=0.12`$.
We set five free parameters to reconstruct the $`K`$-continuum profile of B3 0731+438. They are the peak height, the position relative to the core profile, and the effective radius of the exponential disk and the peak height of the core. The results of the fitting are shown in figure 3, and the obtained photometric data are listed in table 2. The FWHM of the exponential disk profile is 1$`.^{\prime \prime }`$6, which corresponds to 18 kpc. The peak of the disk is located at 0$`.^{\prime \prime }`$24 south of the core.
### 3.3. The Emission-Line Image
The H$`\alpha `$+\[N ii\] image in figure 2a shows a unique morphology. Diffuse line emission extends out to 3$`.^{\prime \prime }`$3 from the center, corresponding to 37 kpc. They are aligned with the axis of the radio hot spots, and both ends of the clouds fork into two directions. Such a morphology suggests the existence of biconical clouds radiating H$`\alpha `$+\[N ii\] emission lines. The total flux of H$`\alpha `$+\[N ii\] emission is $`3.5\times 10^{18}`$ W m<sup>-2</sup>, which corresponds to a luminosity of $`3.2\times 10^{37}`$ W.
We used square aperture photometry of the northern cone, the southern cone and the central core, which are marked as square boxes in figure 2. The results are listed in table 3. The most striking result is the extremely large rest-frame equivalent width of more than 1000 $`\mathrm{\AA }`$ of the H$`\alpha `$+\[N ii\] line in the northern and southern cones. When the results of both cones are combined, the equivalent width at the cones is $`1098_{288}^{+516}\mathrm{\AA }`$.
## 4. Discussion
Few narrow-band imagings aimed at the rest-frame optical emission-line of powerful radio galaxies at a redshift of $`z>2`$ have been carried out, because the lines redshift into the infrared wavelength. Armus et al. (1998) imaged the \[O iii\] emission-line morphology of $`z=3.594`$ radio galaxy 4C 19.71 through a 2.3 $`\mu \mathrm{m}`$ narrow-band filter. They found a large ($`70`$ kpc) aligned nebular, whose length is the same as that of the separation of the radio hot spots. However, its morphology is like a “corridor” between the core and the radio hot spots, and does not show a conic structure. The total \[O iii\] luminosity is $`2\times 10^{37}`$ W, which is comparable to a H$`\alpha `$+\[N ii\] luminosity of $`3\times 10^{37}`$ W of B3 0731+438, if we assume \[O iii\]/H$`\alpha `$+\[N ii\] = 1. Egami et al. (1999) observed the H$`\alpha `$+\[N ii\] emission-line morphology of $`z=2.269`$ 4C 40.36, and found aligned H$`\alpha `$+\[N ii\] knots extended linearly over $``$ 20kpc, but no emission-line cone was seen. They also found an unresolved ($`<2`$ kpc) continuum core, as we found in B3 0731+438.
### 4.1. Properties of Extended Emission Line Clouds
Concerning the emission mechanism of the aligned H$`\alpha `$+\[N ii\] morphology of B3 0731+438, there are four major hypotheses, as described before. Among them, it is impossible to explain the observed emission-line spectrum by the inverse-Compton scattering of the cosmic-background radiation, from which no spectrum feature is expected.
Regarding the second possibility of induced star formation, it is known that the Wolf–Rayet galaxy NGC 4861 shows a H$`\alpha `$+\[N ii\] equivalent width larger than 900 $`\mathrm{\AA }`$ (McQuade et al. 1995). Model spectra of a few-Myr old galaxies also show equivalent widths larger than 1000 $`\mathrm{\AA }`$ (Calzetti 1997). On the other hand, the shape of the L$`\alpha `$ cloud with an equivalent width of 900 $`\mathrm{\AA }`$ reported by McCarthy (1991), is pinched at the peak of the $`K`$-continuum image, and matches the H$`\alpha `$+\[N ii\] morphology that we have observed. Such a morphology suggests that L$`\alpha `$ photons are radiated from the same region of the H$`\alpha `$+\[N ii\] cones, and that both emission lines have a common ionization source. Since a stellar system cannot account for an equivalent width of L$`\alpha `$ as large as 900 $`\mathrm{\AA }`$ (Charlot, Fall 1993), we infer that the alignment of the H$`\alpha `$+\[N ii\] image is not caused by a star-forming region.
We cannot rule out the third possibility of scattered light of anisotropic radiation from the central engine by dust or electrons, because the maximum H$`\alpha `$+\[N ii\] equivalent widths of high-redshift QSOs observed are 800 $`\mathrm{\AA }`$ (Baker et al. 1999; Hill et al. 1993; Espey et al. 1989) and our 1$`\sigma `$ lower limit is 800 $`\mathrm{\AA }`$. We therefore evaluated the L$`\alpha `$/H$`\alpha `$ ratio to examine the contribution of nebular emission to the H$`\alpha `$ flux, because the extended L$`\alpha `$ emission is radiated as nebular emission. The observed L$`\alpha `$/H$`\alpha `$ is 5.5, assuming that all of the L$`\alpha `$ flux ($`3.1\times 10^{18}\mathrm{W}\mathrm{m}^2`$; McCarthy 1991) is radiated from the cones and that H$`\alpha `$/\[N ii\]=4. On the other hand, L$`\alpha `$/H$`\alpha `$ under the case-B condition of the low-density limit is 8.75 (Binette et al. 1992), and only a small amount of dust, such as $`E(BV)=0.04`$, halves the ratio. We infer from this result that the majority of H$`\alpha `$ luminosity is not scattered light from the central engine, but nebular emission radiated from the cones, itself.
Accordingly, we examined whether a gas cloud ionized by anisotropic UV radiation from a hidden AGN can reproduce an equivalent width as large as 1100 $`\mathrm{\AA }`$. First, we extracted the physical properties of the cones and the ionizing source from the data in table 3, using the same method as that carried out by Baum and Heckman (1989). Assuming the gas to be fully ionized and in the case-B condition, the luminosity of the H$`\alpha `$ emission line is written as
$$L(\mathrm{H}\alpha )=n_\mathrm{e}^2\alpha _{\mathrm{H}\alpha }^{\mathrm{eff}}h\nu _{\mathrm{H}\alpha }Vf,$$
(3)
where $`n_\mathrm{e}`$ is the electron density, $`h`$ the Planck constant, $`\nu _{\mathrm{H}\alpha }`$ the frequency of the H$`\alpha `$ line, $`V`$ the volume of the line emitting region, $`f`$ the volume filling factor, and $`\alpha _{\mathrm{H}\beta }^{\mathrm{eff}}=6.04\times 10^{14}`$ (cm<sup>3</sup> s<sup>-1</sup>) (Osterbrock 1989), the H$`\alpha `$ recombination coefficient under case-B. We assumed H$`\alpha `$/\[N ii\] = 4 according to simulations by CLOUDY90 (see following) and took $`V`$ for a cylinder and a half-cut cylinder for the northern and southern cones, respectively.
Because we do not know the volume filling factor, $`f`$, we must estimate it from direct measurements of other radio galaxies. From measurements of the density of ionized gas using sulphur lines in low-redshift radio galaxies, their filling factor is estimated to be in range $`10^4`$$`10^5`$ (Heckman et al. 1982; van Breugel et al. 1985). While the same value is observed in a high-redshift radio galaxy (Rush et al. 1997), we adopted the value $`f=10^4`$. However, the reader should keep in mind that this value is uncertain and may differ by an order of magnitude.
The mass of the line emission gas was calculated using the relation
$$M_{\mathrm{gas}}=Vfn_\mathrm{e}m_\mathrm{H},$$
(4)
where $`m_\mathrm{H}`$ is the hydrogen mass. We next assumed that the line emission gas was filling the cones at the beginning, and was swept up by the passage of radio jets. We thus deduced the total mass of the gas surrounding the galaxy $`M_{\mathrm{tot}}=M_{\mathrm{gas}}{\displaystyle \frac{4\pi }{\mathrm{\Omega }}}`$, taking the opening angles of the cones $`\mathrm{\Omega }`$ to be 0.12$`\pi `$ str and 0.06$`\pi `$ str for the northern and southern cones, respectively.
The number of ionizing photons $`Q`$ can be written as (Osterbrock 1989)
$$Q=\frac{\alpha _\mathrm{B}}{\alpha _{\mathrm{H}\alpha }^{\mathrm{eff}}}\frac{L(\mathrm{H}\alpha )}{h\nu _{\mathrm{H}\alpha }},$$
(5)
where $`\alpha _\mathrm{B}=1.43\times 10^{13}`$ is the total recombination coefficient under the case-B condition. The total number of ionizing photons radiated by the central engine was calculated as $`Q_{\mathrm{tot}}=Q{\displaystyle \frac{4\pi }{\mathrm{\Omega }f_\mathrm{c}}}`$, where $`f_\mathrm{c}`$ is the covering factor of the cloud in the cone.
At last, the ionization parameter of the clouds is defined as
$$U=\frac{Q}{R^2\mathrm{\Omega }f_\mathrm{c}n_\mathrm{e}c},$$
(6)
where $`R`$ is the distance of the cloud from the central engine and $`c`$ the velocity of light. Because we do not know the covering factor, $`f_\mathrm{c}`$, we calculated the lower limits of $`Q_{\mathrm{tot}}`$ and $`U`$ assuming $`f_\mathrm{c}=1`$.
All of these calculated values are given in table 4. The electron density is on the order of 50 (cm<sup>-3</sup>), the ionization parameter 0.001, and the mass of the emission line cloud $`10^9M_{}`$. These values are similar to those deduced for the $``$100 kpc extended line emission clouds of other high-redshift radio galaxies (Rush et al. 1997; Heckman et al. 1991). The large luminosity of the ionizing photons ($`>10^{57}`$ photons $`\mathrm{s}^1`$) draws our attention, which is far larger than the values for typical low-$`z`$ radio galaxies ($`10^{5154}`$ photons $`\mathrm{s}^1`$) (Baum, Heckman 1989), even larger than that of radio loud QSOs ($`10^{5456}`$ photons $`\mathrm{s}^1`$) (Stockton, MacKenty 1987), and almost comparable to that of the brightest QSOs ($`10^{5758}`$ photons $`\mathrm{s}^1`$) (Hill et al. 1993). Here, we calculated the luminosity of ionizing photons for QSOs from their $`L(\mathrm{H}\alpha )`$ using equation (5), or from $`L(\mathrm{H}\beta )`$ with $`L(\mathrm{H}\alpha )/L(\mathrm{H}\beta )=4`$ taken from Kwan and Krolik (1981).
Next, we carried out a photoionization calculation using the code CLOUDY90 (Ferland et al. 1998) by assuming a wide range of hydrogen density, $`n_\mathrm{H}`$, and ionization parameter, $`U`$. The cloud distance from an ionizing source was set to 25 kpc and its thickness to 100 pc. We assumed the metal abundance to be $`Z=0.1Z_{}`$, according to the observations of $`z2`$ damped Ly$`\alpha `$ clouds (Pettini et al. 1994). The continuum spectrum of the ionizing source is set to a power law, $`f_\nu \nu ^\alpha `$, with $`\alpha =0.7`$ longward of $`912\mathrm{\AA }`$ and $`\alpha =2.5`$ shortward of $`912\mathrm{\AA }`$, taken from the composite spectrum of radio-loud QSOs (Cristiani, Rio 1990; Zheng et al. 1997). The resultant equivalent widths, calculated from the emission-line strength and the diffuse continuum radiation, are shown in table 5. Most of the simulated equivalent widths exceed 1000 $`\mathrm{\AA }`$, satisfying the observed value.
Thus, we suggest that the emission mechanism of the extended H$`\alpha `$+\[N ii\] cones of B3 0731+438 is nebular emission from clouds ionized by the strong UV radiation of the hidden AGN.
### 4.2. Spectral Energy Distribution
Assuming that the exponential disk component of the $`K`$-continuum image is a host galaxy of B3 0731+438 and the compact core a type-2 AGN, we reconstructed the SED of the B3 0731+438 with model spectra of a type-2 AGN and a galaxy. Because contamination of the \[O ii\] $`\lambda \lambda `$ 3727 emission line to the $`J`$-band flux is expected, we subtracted it, assuming an H$`\alpha `$/\[O ii\] ratio of 3 (McCarthy et al. 1995). The model spectrum of the galaxy was calculated by the spectrophotometric galaxy evolution model PEGASE (Fioc, Rocca-Volmerange 1997) under three variations of star-forming history, namely, an instantaneous burst model, and two exponential burst models with time scales of $`\tau =200`$ Myr and 2 Gyr, respectively. For the model spectrum of the type-2 AGN, we selected the dust-scattered AGN model calculated by Cimatti et al. (1994) with a scattering angle, $`\mathrm{\Theta }`$, of $`90^{}`$ and the power-law index of the incident continuum being $`\alpha =0.7`$, and assuming extinction of SMC dust (Prévot et al. 1984, Bouchet et al. 1985). We scaled these two spectra according to their observed $`K`$-continuum flux and fit the $`J`$\- and $`R`$-band flux by altering the age of the galaxy and the extinction of the type-2 AGN.
We show the results in figure 4 and table 6. The $`\tau =`$ 2 Gyr exponential burst model is not plausible, because the age of the best-fit galaxy model is 10 Gyr, which is larger than the cosmic age (3.5 Gyr for the assumed cosmological parameters). Both the $`\tau =`$ 200 Myr exponential burst model and the instantaneous burst model fit the SED. However, we prefer the instantaneous burst model because the observed flux of the $`R`$-band appears to be dominated by the diffuse galactic component, as can be seen in the $`R`$-band image of McCarthy (1991).
Consequently, we suppose that the age of the galaxy is 500 Myr, resulting in a formation redshift of $`z_{\mathrm{form}}3`$. The total stellar mass of the galaxy estimated from the model is 3$`\times 10^{11}M_{}`$, comparable to that of a typical 3C radio galaxy at $`z=1`$ (Best et al. 1997).
The H$`\alpha `$+\[N ii\] peak is located at the center of the galaxy and aligned to the radio axis. Infrared spectroscopic observations showed that the H$`\alpha `$ line width is narrow ($`<560\mathrm{km}\mathrm{s}^1`$ ; McCarthy et al. 1992) and that the emission-line ratios are similar to those of Seyfert 2 galaxies (Evans 1997). We therefore suggest that the H$`\alpha `$+\[N ii\] luminosity at the center of the galaxy is dominated by scattered radiation from the narrow-line regions of the hidden nucleus. However, the possibility of intense starburst activity cannot be ruled out entirely, and polarimetry is necessary for a confirmation.
## 5. Summary and Conclusions
We observed the powerful radio galaxy B3 0731+438 in the infrared $`K^{}`$-band and the 2.25 $`\mu `$m narrow-band, corresponding to the rest-frame 6000 $`\mathrm{\AA }`$ continuum and H$`\alpha `$+\[N ii\] emission line, respectively. We then produced a line-free $`K`$-continuum image and a continuum-subtracted H$`\alpha `$+\[N ii\] emission-line image.
Our observations are the first to show a cone-shaped H$`\alpha `$+\[N ii\] emission-line structure of a high-redshift radio galaxy at infrared wavelength. The radio-aligned cones suggest that the gas is ionized by the hidden AGN. The contribution of the scattered light from the hidden AGN to the H$`\alpha `$ luminosity is estimated to be small. On the other hand, we find that a gas cloud ionized by a power-law UV continuum can account for the large observed H$`\alpha `$+\[N ii\] equivalent width, using the values of the electron density and the ionization parameter estimated from the observed H$`\alpha `$ line luminosity. Together with their biconical structure, we infer that we have detected H$`\alpha `$ ionization cones of a high-redshift powerful radio galaxy for the first time. The estimated mass of the ionized gas cones is of the order of $`10^9M_{}`$, and the expected total mass of the gas surrounding the galaxy is $`10^{11}M_{}`$.
The H$`\alpha `$+\[N ii\] peak at the center of the galaxy is also aligned to the radio axis. We suppose that this peak is scattered light from the narrow-line regions of the hidden AGN, and that polarimetry is necessary for a confirmation.
The $`K`$-continuum image is separated into two components, assumed to be a type-2 AGN and an underlying host galaxy. The SED of the whole radio galaxy is modeled by a model spectrum of a 500 Myr-old instantaneous burst galaxy and dust-scattered power-law continuum with $`A_V=1.9`$ mag extinction. The stellar mass of the galaxy is $`3\times 10^{11}M_{}`$, which is comparable to that of a typical 3C radio galaxy at $`z=1`$.
We thank all staff of the Subaru Telescope, who supported us to set up our instrument, and helped us to make these observations. We would like to express our thanks to the engineering staff of Mitsubishi Electric Co. for their fine operation of the telescope during the test observation runs, and the staff of Fujitsu Co. for timely provision of control software.
We also appreciate M. Fioc and B. Rocca-Volmerange for generously offering their galaxy modeling code, PEGASE, and G. Ferland for the spectral synthesis code, CLOUDY90.
K.Motohara was financially supported by the Japan Society for the Promotion of Science.
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Figure 1. Color image of B3 0731+438, 11$`.^{\prime \prime }`$8 on a side. The blue color is assigned to the line-subtracted $`K^{}`$-band image and the red color to the 2.25 $`\mu `$m narrow-band. Both images were smoothed by a 1 pixel Gaussian filter. Green is assigned to the average of the red and the blue image.
Figure 2. Logarithmic gray-scale images of B3 0731+438. The squares are the apertures with which aperture photometry was carried out, while the crosses indicate the position of 4710 MHz radio hot spots, as mapped by Carilli et al. (1997). (a) The H$`\alpha `$+\[N ii\] image, smoothed by a Gaussian filter of 1 pixel. The first contour is at 1 $`\sigma `$ above the background level, while the subsequent contours are at levels of $`2n\sigma (n=1,2,3,\mathrm{})`$. $`\sigma `$ is taken from the non-smoothed image. (b) The $`K`$-continuum image, which was smoothed by a 3 pixel Gaussian filter to match the seeing size with the H$`\alpha `$+\[N ii\] image. The contours are at levels of $`n\sigma (n=1,2,3,\mathrm{})`$, where $`\sigma `$ is again taken from the non-smoothed image. The profile at the lower-right corner is the contours of the stellar image used for deconvolution, whose peak value is normalized to B3 0731+438. (c) The $`K`$-continuum image deconvolved by MEM deconvolution method.
Figure 3. (upper-left) The $`K`$-continuum image of B3 0731+438. (upper-right) The best-fit two-component model profile of the $`K`$-continuum image. (lower-left) The stellar-profile subtracted $`K`$-continuum image, which reveals the underlying galaxy. (lower-right) The residual image after the two-component profile is subtracted.
Figure 4. Spectral energy distribution of B3 0731+438. The open circle is from McCarthy (1991), the open square from Iwamuro (private communication), and the filled square from this work. The thick solid line is the best-fit models, and the thin solid line the spectra of (A) the type-2 AGN with dust extinction and (B) the galaxy-evolution model. The dotted line is the spectrum of the type-2 AGN without extinction. The star-formation models are: (a) instantaneous burst and (b) $`\tau =200`$ Myr exponential burst. |
warning/0001/math0001011.html | ar5iv | text | # Computing special values of partial zeta functions
## 1 Introduction
Let $`K/`$ be a totally real number field of degree $`n`$ with ring of integers $`𝒪_K`$, and let $`U𝒪_K^\times `$ be the subgroup of totally positive units. Let $`𝔣,𝔟𝒪_K`$ be relatively prime ideals. Then the *partial zeta function* associated to this data is defined by
$$\zeta _𝔣(𝔟,s):=\underset{𝔞𝔟}{}N(𝔞)^s,$$
where $`𝔞𝔟`$ means $`𝔞𝔟^1=(\alpha )`$, where $`\alpha `$ is a totally positive number in $`1+𝔣𝔟^1`$. According to a classical result of Klingen and Siegel, the special values $`\zeta _𝔣(𝔟,k)`$ are rational for nonpositive integers $`k`$. Moreover, the values $`\zeta _𝔣(𝔟,0)`$ are especially important because of their connection with the Brumer-Stark conjecture and the Leopoldt conjecture .
In , one of us (RS) gave a cohomological interpretation of these special values by showing that they can be computed in finite terms as periods of the *Eisenstein cocycle*. This is a cocycle $`\mathrm{\Psi }H^{n1}(GL_n();)`$, where $``$ is a certain $`GL_n()`$-module. Then two of us (PEG and RS) showed in that the Eisenstein cocycle is an effectively computable object. More precisely, using the cocycle one can express $`\zeta _𝔣(𝔟,k)`$ as a finite sum of *generalized Dedekind sums*, and that the latter can be effectively computed by a continued-fraction algorithm that uses a generalization of the classical Dedekind-Rademacher reciprocity law.
In this note we describe an ongoing project to build a database of $`\zeta _𝔣(𝔟,0)`$ for various fields $`K`$ and ideals $`𝔣,𝔟`$. We recall the definition of the Eisenstein cocycle and its relation to the special values (§2), and discuss the effective computation of Dedekind sums (§3). We conclude with examples of special values for some fields of degree 3 and 4 (§4).
## 2 Dedekind sums and the Eisenstein cocycle
### 2.1
Let $`\sigma `$ be a square matrix with integral columns $`\sigma _j^n`$ ($`j=1,\mathrm{},n`$), and let $`L^n`$ be a lattice of rank $`r1`$. Let $`v^n`$, and let $`e^n`$ with $`e_j1`$. Then the *Dedekind sum* $`S`$ associated to the data $`(L,\sigma ,e,v)`$ is defined by
$$S=S(L,\sigma ,e,v):=\underset{xL}{}{}_{}{}^{^{}}𝐞(x,v)\frac{det\sigma }{x,\sigma _1^{e_1}\mathrm{}x,\sigma _n^{e_n}}.$$
(1)
Here $`x,y:=x_iy_i`$ is the usual scalar product on $`^n`$, $`𝐞(t)`$ is the character $`\mathrm{exp}(2\pi it)`$, and the prime next to the summation means to omit terms for which the denominator vanishes. The series (1) converges absolutely if all $`e_j>1`$, but may only converge conditionally if $`e_j=1`$ for some $`j`$. In this latter case we can define the sum by the *$`Q`$-limit*
$$\underset{xL}{}{}_{}{}^{^{}}a(x)|_Q:=\underset{t\mathrm{}}{lim}\left(\underset{|Q(x)|<t}{}{}_{}{}^{^{}}a(x)\right),$$
(2)
where $`Q`$ is any finite product of real-valued linear forms on $`^n`$ that doesn’t vanish on $`^n\{0\}`$. One can precisely determine how the value of (1) depends on $`Q`$ (\[9, Thm. 7\]). The sum $`S`$ is always a rational number times a power of $`2\pi i`$.
### 2.2
We recall now the definition of the Eisenstein cocycle $`\mathrm{\Psi }`$ and its relationship with the special values $`\zeta _𝔣(𝔟,k)`$. For simplicity, we describe only material necessary to compute the special value at $`k=0`$, and refer to for other $`k`$.
Let $`𝒜=(A_1,\mathrm{},A_n)(GL_n())^n`$ be an $`n`$-tuple of matrices. For an $`n`$-tuple $`d=(d_1,\mathrm{},d_n)`$ of integers $`1d_in`$, let $`𝒜(d)^n`$ be the subspace generated by all columns $`A_{ij}`$ such that $`j<d_i`$. (Here $`A_{ij}`$ denotes the $`j`$th column of the matrix $`A_i`$.) Writing $`𝒜(d)^{}`$ for the orthogonal complement of $`𝒜(d)`$ in $`^n`$, we let
$$X(d)=𝒜(d)^{}\underset{i=1}{\overset{n}{}}\sigma _i^{},\text{where }\sigma _i=A_{id_i}\text{.}$$
(3)
The $`n`$-tuple $`𝒜`$ determines a decomposition of $`^n\{0\}`$ into linear strata
$$\underset{dD}{}X(d),$$
(4)
indexed by the finite set
$$D=D(𝒜)=\{dX(d)\mathrm{}\}.$$
Associated to this decomposition is a collection of rational functions $`\psi (𝒜)`$ on $`^n\{0\}`$, defined by
$$\psi (𝒜)(x)=\frac{det(\sigma _1,\mathrm{},\sigma _n)}{x,\sigma _1\mathrm{}x,\sigma _n},\text{if }xX(d)\text{.}$$
Note that $`\psi (𝒜)(x)`$ is well-defined by the construction of $`X(d)`$.
Let $`v^n`$, and let $`Q`$ be defined as in §2.1. Then the *Eisenstein cocycle* $`\mathrm{\Psi }`$ is defined as
$$\mathrm{\Psi }=\mathrm{\Psi }(𝒜)(Q,v):=(2\pi i)^n\underset{x^n}{}𝐞(x,v)\psi (𝒜)(x)|_Q.$$
One can show that $`\mathrm{\Psi }`$ is a homogeneous $`(n1)`$-cocycle for $`GL_n()`$. Furthermore, we can express $`\mathrm{\Psi }`$ in terms of Dedekind sums
$$\mathrm{\Psi }(𝒜)(Q,v)=(2\pi i)^n\underset{dD}{}S(L(d),\sigma ,\mathrm{𝟏},v)|_Q,$$
(5)
where $`\sigma `$ is the matrix with columns $`A_{id_i},`$ ($`i=1,\mathrm{},n`$), $`L(d)`$ is the lattice $`𝒜(d)^{}^n`$, and $`\mathrm{𝟏}`$ is the vector $`(1,\mathrm{},1)`$.
### 2.3
Now we describe how $`\mathrm{\Psi }`$ can be used to compute special values. Let $`W`$ be a $``$-basis for the fractional ideal $`𝔣𝔟^1=W_j`$, and let $`W^{}`$ be the dual basis with respect to the trace form. Via the $`n`$ real embeddings $`\tau _i`$, $`i=1,\mathrm{},n`$, any $`xK`$ determines a row vector $`(\tau _1(x),\mathrm{},\tau _n(x))`$. Hence we may identify $`W`$ with a matrix in $`GL_n()`$: the $`j`$th row of this matrix is the image of the $`j`$th basis element of $`W`$. Let
$$Q(X)=\underset{i}{}\underset{j}{}X_j(\tau _i(W_j^{})),$$
and let $`v^n`$ be defined by $`v_j=\mathrm{Tr}(W_j^{})`$.
Let $`\nu =n1`$, and let $`\epsilon _1,\mathrm{},\epsilon _\nu `$ be a basis for the totally positive units $`U`$. Using the regular representation $`\rho `$ with respect to the basis $`W`$, we identify the units $`\epsilon _j`$ with elements $`A_j=\rho (\epsilon _j)^tGL_n()`$. Using the bar notation
$$[A_1|\mathrm{}|A_\nu ]:=(1,A_1,A_1A_2,\mathrm{},A_1\mathrm{}A_\nu )(GL_n())^n,$$
we have the following proposition expressing the zeta values in terms of the Eisenstein cocycle:
###### Proposition 1
Let $`U_𝔣`$ be the subgroup $`U(1+𝔣)`$, and let $`\pi `$ run through all permutations of $`\{1,\mathrm{},\nu \}`$. Then
$$\zeta _𝔣(𝔟,0)=\eta \underset{\epsilon U/U_𝔣}{}\underset{\pi }{}\mathrm{sgn}(\pi )\mathrm{\Psi }([A_{\pi (1)}|\mathrm{}|A_{\pi (\nu )}])(Q,\rho (\epsilon )^tv).$$
Here $`\eta =\pm 1`$ is defined by
$$\eta =(1)^\nu \mathrm{sgn}(detW)\mathrm{sgn}(R),$$
where $`R=det(\mathrm{log}\tau _j(\epsilon _i))`$, $`1i,j\nu `$.
## 3 Diagonality and unimodularity
### 3.1
We define the *rank* of $`S=S(L,\sigma ,e,v)`$ to be the rank of the lattice $`L`$. It is easy to see that after a $`GL_n()`$ transformation, we may assume that $`L`$ is the sublattice $`Z^{\mathrm{}}`$ spanned by the first $`\mathrm{}`$ standard basis vectors, where $`\mathrm{}`$ is the rank of $`L.`$ Furthermore, by multiplying by an appropriate rational factor, permuting columns and repeating columns if necessary, we may assume the pair $`(Z^{\mathrm{}},\sigma )`$ satisfies the following conditions:
1. For each column $`\sigma _j`$, the vector of the first $`\mathrm{}`$ components of $`\sigma _j`$ is primitive and integral.
2. If two columns of $`\sigma `$ induce proportional linear forms on $`Z^{\mathrm{}}`$, then these two linear forms coincide on $`Z^{\mathrm{}}`$, and are adjacent columns of $`\sigma `$.
3. The vector $`e=\mathrm{𝟏}`$.
Let $`S(Z^{\mathrm{}},\sigma ,\mathrm{𝟏},v)`$ be a Dedekind sum satisfying the three conditions above. Let $`\pi :^N^{\mathrm{}}`$ be the projection on the first $`\mathrm{}`$ components, and let $`\pi (\sigma )`$ be the $`\mathrm{}\times n`$ matrix with columns $`\pi (\sigma _i)`$.
###### Definition 1
Let $`M(\sigma )`$ be the set of maximal minors of $`\pi (\sigma )`$. Then the *index* of $`S`$, denoted $`S`$, is defined to be
$$\underset{\tau M(\sigma )}{\mathrm{max}}|det\tau |.$$
A Dedekind sum is *unimodular* if $`S=1`$.
### 3.2
Now define a partition
$$[[n]]=\underset{k=1}{\overset{s}{}}I_k,\mathrm{}sn$$
(6)
as follows. Put
$$i,jI_k\text{if and only if}\pi (\sigma _i)=\pi (\sigma _j).$$
In other words, two elements of $`[[n]]`$ are in the same set of the partition if the corresponding columns of $`\sigma `$ induce the same linear form on $`Z^{\mathrm{}}`$.
Let $`p_k=\mathrm{\#}I_k`$.
###### Definition 2
The vector $`p(S)=(p_1,\mathrm{},p_s)`$ is called the *type* of $`S`$. A Dedekind sum is called *diagonal* if $`p(S)`$ has length $`\mathrm{}`$.
### 3.3
The virtue of diagonality is that a diagonal Dedekind sum $`S`$ may be evaluated as a finite sum of products of generalized Bernoulli polynomials. Furthermore, the number of terms in this finite sum is the index of $`S`$. Hence diagonal and unimodular Dedekind sums can be evaluated very rapidly.
In general, the Dedekind sums in (5) aren’t diagonal. However, we have the following theorem, which is the main result of :
###### Theorem 3.1
Every Dedekind sum $`S(L,\sigma ,e,v)`$ can be expressed as a finite rational linear combination of unimodular diagonal sums. If $`n`$, $`\mathrm{Rank}L`$, and $`e`$ are fixed, then this expression can be computed in time polynomial in $`\mathrm{log}S`$. Moreover, the number of terms in this expression is bounded by a polynomial in $`\mathrm{log}S`$.
The key ingredient in the proof of Theorem 3.1 is a “reciprocity law” for higher-dimensional Dedekind sums. For any nonzero point $`v^n`$, let $`v^{}`$ be the hyperplane $`\{xv,x=0\}`$. Let $`Q`$ be a finite product of real-valued linear forms on $`^n`$ that do not vanish on $`^n\{0\}`$.
###### Proposition 2
Let $`\sigma _0,\mathrm{},\sigma _n^n`$ be nonzero. For $`j=0,\mathrm{},n`$, let $`\sigma ^j`$ be the matrix with columns $`\sigma _0,\mathrm{},\widehat{\sigma }_j,\mathrm{},\sigma _n`$. Fix a lattice $`L^n`$, and assume $`e=\mathrm{𝟏}`$. Then for any $`v^n`$, we have the following identity among Dedekind sums:
$$\underset{j=0}{\overset{n}{}}(1)^jS(L,\sigma ^j,\mathrm{𝟏},v)|_Q=\underset{j=0}{\overset{n}{}}(1)^jS(L\sigma _j^{},\sigma ^j,\mathrm{𝟏},v)|_Q.$$
(7)
We refer to for proofs of the above statements. Here, in the following two sections, we show how Theorem 3.1 is applied with a rank $`2`$ example. For simplicity we ignore issues of convergence, and merely remark that all of our manipulations with sums are compatible with the $`Q`$-limit process (2).
### 3.4
Let $`L`$ be the lattice $`Z^2`$. Let
$$\sigma =\left(\begin{array}{ccc}1& 0& 1\\ 0& 1& 2\\ 0& 0& 1\end{array}\right),e=(1,1,2),\text{and }v=(0,0,0)\text{.}$$
Hence $`S(L,\sigma ,e,v)`$ denotes the absolutely convergent sum
$$\underset{(x,y)^2}{}{}_{}{}^{^{}}\frac{1}{xy(x+2y)^2},$$
where the prime on the summation indicates that we omit the terms $`(x,y)`$ for which $`x`$, $`y`$ or $`x+2y`$ vanish.
This sum isn’t diagonal, since $`\sigma `$ induces $`3`$ different linear forms on $`L`$ instead of $`2`$. To diagonalize $`S`$, we begin with the identity of rational functions
$$\frac{1}{xy(x+2y)^2}=\frac{1}{y(x+2y)^3}+\frac{2}{x(x+2y)^3}.$$
(8)
This is true provided none of the denominators vanishes. The numerators of the functions on the right come from expressing the third column of $`\sigma `$ as a linear combination of the first two:
$$(1,2)^t=1(1,0)^t+2(0,1)^t.$$
We want to sum both sides of (8) over pairs $`(x,y)^2`$ to obtain an identity among Dedekind sums of the form
$$\underset{(x,y)^2}{}{}_{}{}^{^{}}\frac{1}{xy(x+2y)^2}=\underset{(x,y)^2}{}{}_{}{}^{^{}}\frac{1}{y(x+2y)^3}+\underset{(x,y)^2}{}{}_{}{}^{^{}}\frac{2}{x(x+2y)^3}.$$
(9)
However, as written (9) is incorrect. The identity (8) only holds if none of $`x`$, $`y`$, or $`x+2y`$ vanish, but the sums on the right of (9) include some of these terms (for instance, the first sum on the right of (9) contains terms $`(x,y)`$ with $`x=0`$). We account for this by subtracting two rank $`1`$ Dedekind sums from the right of (9) as “correction terms”:
$`{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{xy(x+2y)^2}}`$ $`=`$ $`{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{y(x+2y)^3}}+{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{2}{x(x+2y)^3}}`$ (10)
$`{\displaystyle \underset{(x,y)^2}{}_{x=0}}{}_{}{}^{^{}}{\displaystyle \frac{1}{y(x+2y)^3}}{\displaystyle \underset{(x,y)^2}{}_{y=0}}{}_{}{}^{^{}}{\displaystyle \frac{1}{x(x+2y)^3}}`$
$`=`$ $`{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{y(x+2y)^3}}+{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{2}{x(x+2y)^3}}`$
$`{\displaystyle \underset{y}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{8y^4}}{\displaystyle \underset{x}{}}{}_{}{}^{^{}}{\displaystyle \frac{2}{x^4}}.`$
This equation is precisely an instance of the reciprocity law (Proposition 2). The three rank $`2`$ sums are the left of (7), and the two rank $`1`$ sums are the right of (7) (one rank $`1`$ sum in (7) vanishes identically). Note that all of the sums on the right of (LABEL:recip.law.ex) are now diagonal.
To diagonalize a general Dedekind sum $`S(L,\sigma ,\mathrm{𝟏},v)`$, one considers the configuration $`C^n`$ of linear subspaces consisting of $`(L)^{}`$ and the spaces generated by the points $`\sigma _1,\mathrm{},\sigma _n`$. One shows by investigating the geometry of $`C`$ that a point $`\sigma _0`$ can be found such that when Proposition 2 is applied with the tuple $`(\sigma _0,\mathrm{},\sigma _n)`$, the resulting Dedekind sums are “closer” to diagonality in a certain sense. It may take several applications of Proposition 2 to express a Dedekind sum as a linear combination of diagonal sums.
### 3.5
The second rank two sum on the right of (LABEL:recip.law.ex) has index two. We will show how to make this sum unimodular. Write
$$S(^2,\tau ,(1,3),v^{})=\underset{(x,y)^2}{}{}_{}{}^{^{}}\frac{1}{x(x+2y)^3}$$
where
$$\tau =(\tau _1,\tau _2)=\left(\begin{array}{cc}1& 1\\ 0& 2\end{array}\right)\text{and}v^{}=(0,0).$$
Let $`\rho `$ be the column vector $`(0,1)^t`$. We apply Proposition 2 to the triple $`(\rho ,\tau _1,\tau _2)`$:
$`{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{2}{x(x+2y)^3}}`$ $`=`$ $`{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{y(x+2y)^3}}+{\displaystyle \underset{(x,y)^2}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{x(x+2y)^2y}}`$ (11)
$`+{\displaystyle \underset{y}{}}{}_{}{}^{^{}}{\displaystyle \frac{1}{8y^4}}+{\displaystyle \underset{x}{}}{}_{}{}^{^{}}{\displaystyle \frac{2}{x^4}}.`$
Now all the terms on the right of (11) are diagonal and unimodular except for the second rank two sum. In fact, this sum is no longer diagonal. However, one further application of Proposition 2 as in §3.4 will make the third sum diagonal and unimodular. Hence we will have succeeded in expressing the original sum as a finite linear combination of diagonal, unimodular Dedekind sums.
In general, one must be able to construct the vector $`\rho `$ as above. An easy argument using Minkowski’s Theorem from the geometry of numbers guarantees the existence of $`\rho `$ . To construct $`\rho `$ in practice, one may use $`LLL`$-reduction of the lattice spanned by the *rows* of $`\sigma `$ and \[4, Conjecture 3.9\].
## 4 Examples
Here we present some numerical examples. For simplicity we compute $`\zeta =\zeta _𝔣(𝔟,0)`$, where $`𝔣=N𝒪_K`$ for various rational integers $`N`$, and $`𝔟=𝒪_K`$. These fields are the first entries in the tables of totally real fields with small discriminant, available from .
Cubic fields
* $`K=(\theta )`$, where $`\theta ^3+\theta ^22\theta 1=0`$ (discriminant $`49`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`1`$ | $`7`$ | $`2`$ | $`10`$ | $`16`$ | $`13`$ | $`1`$ | $`16`$ | $`9`$ | $`19`$ | $`18`$ | $`22`$ | $`8`$ |
| $`2`$ | $`0`$ | $`5`$ | $`2`$ | $`8`$ | $`3`$ | $`11`$ | $`10`$ | $`14`$ | $`24`$ | $`17`$ | $`26`$ | $`20`$ | $`19`$ | $`23`$ | $`10`$ |
| $`3`$ | $`2`$ | $`6`$ | $`8`$ | $`9`$ | $`2`$ | $`12`$ | $`5`$ | $`15`$ | $`14`$ | $`18`$ | $`8`$ | $`21`$ | $`6`$ | $`24`$ | $`23`$ |
* $`K=(\theta )`$, where $`\theta ^33\theta 1=0`$ (discriminant $`81`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`1`$ | $`7`$ | $`14`$ | $`10`$ | $`16`$ | $`13`$ | $`2`$ | $`16`$ | $`13`$ | $`19`$ | $`11`$ | $`22`$ | $`72`$ |
| $`2`$ | $`0`$ | $`5`$ | $`2`$ | $`8`$ | $`7`$ | $`11`$ | $`6`$ | $`14`$ | $`0`$ | $`17`$ | $`13`$ | $`20`$ | $`11`$ | $`23`$ | $`62`$ |
| $`3`$ | $`2/3`$ | $`6`$ | $`8/3`$ | $`9`$ | $`2/3`$ | $`12`$ | $`11/3`$ | $`15`$ | $`14/3`$ | $`18`$ | $`64/3`$ | $`21`$ | $`14/3`$ | $`24`$ | $`29/3`$ |
* $`K=(\theta )`$, where $`\theta ^3+\theta ^23\theta 1=0`$ (discriminant $`148`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`1`$ | $`7`$ | $`2`$ | $`10`$ | $`2`$ | $`13`$ | $`22`$ | $`16`$ | $`7`$ | $`19`$ | $`82`$ | $`22`$ | $`68`$ |
| $`2`$ | $`0`$ | $`5`$ | $`4`$ | $`8`$ | $`3`$ | $`11`$ | $`18`$ | $`14`$ | $`20`$ | $`17`$ | $`100`$ | $`20`$ | $`4`$ | $`23`$ | $`12`$ |
| $`3`$ | $`2`$ | $`6`$ | $`4`$ | $`9`$ | $`10`$ | $`12`$ | $`5`$ | $`15`$ | $`42`$ | $`18`$ | $`32`$ | $`21`$ | $`78`$ | $`24`$ | $`23`$ |
* $`K=(\theta )`$, where $`\theta ^3\theta ^24\theta 1=0`$ (discriminant $`169`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`3`$ | $`7`$ | $`6`$ | $`10`$ | $`8`$ | $`13`$ | $`2`$ | $`16`$ | $`25`$ | $`19`$ | $`238`$ | $`22`$ | $`160`$ |
| $`2`$ | $`0`$ | $`5`$ | $`1`$ | $`8`$ | $`5`$ | $`11`$ | $`6`$ | $`14`$ | $`56`$ | $`17`$ | $`6`$ | $`20`$ | $`8`$ | $`23`$ | $`386`$ |
| $`3`$ | $`2`$ | $`6`$ | $`8`$ | $`9`$ | $`26`$ | $`12`$ | $`11`$ | $`15`$ | $`17`$ | $`18`$ | $`52`$ | $`21`$ | $`10`$ | $`24`$ | $`89`$ |
* $`K=(\theta )`$, where $`\theta ^34\theta 1=0`$ (discriminant $`229`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`2`$ | $`7`$ | $`18`$ | $`10`$ | $`24`$ | $`13`$ | $`60`$ | $`16`$ | $`19`$ | $`19`$ | $`26`$ | $`22`$ | $`8`$ |
| $`2`$ | $`0`$ | $`5`$ | $`2`$ | $`8`$ | $`7`$ | $`11`$ | $`4`$ | $`14`$ | $`24`$ | $`17`$ | $`8`$ | $`20`$ | $`38`$ | $`23`$ | $`202`$ |
| $`3`$ | $`2`$ | $`6`$ | $`0`$ | $`9`$ | $`10`$ | $`12`$ | $`18`$ | $`15`$ | $`56`$ | $`18`$ | $`16`$ | $`21`$ | $`96`$ | $`24`$ | $`51`$ |
* $`K=(\theta )`$, where $`\theta ^3\theta ^24\theta +3=0`$ (discriminant $`257`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`6`$ | $`7`$ | $`16`$ | $`10`$ | $`4`$ | $`13`$ | $`10`$ | $`16`$ | $`5`$ | $`19`$ | $`42`$ | $`22`$ | $`72`$ |
| $`2`$ | $`0`$ | $`5`$ | $`3`$ | $`8`$ | $`9`$ | $`11`$ | $`6`$ | $`14`$ | $`0`$ | $`17`$ | $`26`$ | $`20`$ | $`12`$ | $`23`$ | $`112`$ |
| $`3`$ | $`2`$ | $`6`$ | $`8`$ | $`9`$ | $`4`$ | $`12`$ | $`8`$ | $`15`$ | $`21`$ | $`18`$ | $`16`$ | $`21`$ | $`12`$ | $`24`$ | $`58`$ |
Quartic fields
* $`K=(\theta )`$, where $`\theta ^4\theta ^33\theta ^2+\theta +1=0`$ (discriminant $`725`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`1`$ | $`7`$ | $`4`$ | $`10`$ | $`0`$ | $`13`$ | $`20`$ | $`16`$ | $`35`$ | $`19`$ | $`32`$ | $`22`$ | $`32`$ |
| $`2`$ | $`0`$ | $`5`$ | $`4`$ | $`8`$ | $`7`$ | $`11`$ | $`2`$ | $`14`$ | $`16`$ | $`17`$ | $`92`$ | $`20`$ | $`0`$ | $`23`$ | $`12`$ |
| $`3`$ | $`4`$ | $`6`$ | $`0`$ | $`9`$ | $`44`$ | $`12`$ | $`0`$ | $`15`$ | $`0`$ | $`18`$ | $`320`$ | $`21`$ | $`30`$ | $`24`$ | $`84`$ |
* $`K=(\theta )`$, where $`\theta ^4\theta ^34\theta ^2+4\theta +1=0`$ (discriminant $`1125`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`3`$ | $`4/5`$ | $`5`$ | $`4/3`$ | $`7`$ | $`124/3`$ | $`9`$ | $`116/5`$ | $`11`$ | $`72`$ | $`13`$ | $`676/3`$ | $`15`$ | $`4`$ |
| $`2`$ | $`0`$ | $`4`$ | $`2`$ | $`6`$ | $`0`$ | $`8`$ | $`7`$ | $`10`$ | $`128`$ | $`12`$ | $`0`$ | $`14`$ | $`144`$ | $`16`$ | $`51`$ |
* $`K=(\theta )`$, where $`\theta ^46\theta ^2+4=0`$ (discriminant $`1600`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`1/2`$ | $`7`$ | $`4`$ | $`10`$ | $`4`$ | $`13`$ | $`104`$ | $`16`$ | $`45/2`$ | $`19`$ | $`224`$ | $`22`$ | $`136`$ |
| $`2`$ | $`1`$ | $`5`$ | $`4`$ | $`8`$ | $`5/2`$ | $`11`$ | $`4`$ | $`14`$ | $`16`$ | $`17`$ | $`84`$ | $`20`$ | $`9`$ | $`23`$ | $`60`$ |
| $`3`$ | $`4`$ | $`6`$ | $`0`$ | $`9`$ | $`16`$ | $`12`$ | $`24`$ | $`15`$ | $`6`$ | $`18`$ | $`128`$ | $`21`$ | $`222`$ | $`24`$ | $`15`$ |
* $`K=(\theta )`$, where $`\theta ^44\theta ^2\theta +1=0`$ (discriminant $`1957`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`3`$ | $`4`$ | $`5`$ | $`20`$ | $`7`$ | $`8`$ | $`9`$ | $`52`$ | $`11`$ | $`4`$ | $`13`$ | $`500`$ |
| $`2`$ | $`0`$ | $`4`$ | $`1`$ | $`6`$ | $`0`$ | $`8`$ | $`11`$ | $`10`$ | $`0`$ | $`12`$ | $`3`$ | $`14`$ | $`8`$ |
* $`K=(\theta )`$, where $`\theta ^45\theta ^2+5=0`$ (discriminant $`2000`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`4`$ | $`11/5`$ | $`7`$ | $`52`$ | $`10`$ | $`0`$ | $`13`$ | $`32`$ | $`16`$ | $`31/10`$ | $`19`$ | $`412/5`$ | $`22`$ | $`568`$ |
| $`2`$ | $`2/5`$ | $`5`$ | $`0`$ | $`8`$ | $`71/10`$ | $`11`$ | $`4`$ | $`14`$ | $`16`$ | $`17`$ | $`296`$ | $`20`$ | $`30`$ | $`23`$ | $`148`$ |
| $`3`$ | $`4`$ | $`6`$ | $`0`$ | $`9`$ | $`4`$ | $`12`$ | $`12`$ | $`15`$ | $`340`$ | $`18`$ | $`112`$ | $`21`$ | $`428`$ | $`24`$ | $`117`$ |
* $`K=(\theta )`$, where $`\theta ^44\theta ^2+2=0`$ (discriminant $`2048`$).
| $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ | $`N`$ | $`N\zeta `$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | | $`3`$ | $`4`$ | $`5`$ | $`28`$ | $`7`$ | $`8`$ | $`9`$ | $`20`$ | $`11`$ | $`68`$ | $`13`$ | $`52`$ |
| $`2`$ | $`1/2`$ | $`4`$ | $`1/4`$ | $`6`$ | $`10`$ | $`8`$ | $`9/4`$ | $`10`$ | $`19`$ | $`12`$ | $`41/2`$ | $`14`$ | $`36`$ | |
warning/0001/math0001072.html | ar5iv | text | # A formula for Euler characteristic of line singularities on singular spaces
## 1. Introduction
For an analytic function germ $`f:(X,0)(,0)`$ with critical locus $`\mathrm{\Sigma }X`$, there is a local Milnor fibration induced by $`f`$. We are interested in the topology of the Milnor fibre $`F`$ of $`f`$ in the case when $`dim\mathrm{\Sigma }=1`$. It is well known that in this case the homotopy type of $`F`$ is not necessarily a bouquet of spheres in the middle dimension. The calculation of the Euler characteristic $`\chi (F)`$ of $`F`$ is of importance. There is a so called Iomdin-Lê formula which expresses the Euler characteristic of the Milnor fibre of $`f`$ by that of the series of $`f`$ with isolated singularities.
The question we are interested in is that if there is a way to express $`\chi (F)`$ by some “computable” invariants determined only by $`(f,\mathrm{\Sigma },X)`$. When $`X`$ is $`^m`$, the singular locus $`\mathrm{\Sigma }`$ of $`f`$ is a one dimensional complete intersection with isolated singularity, and the transversal singularity type of $`f`$ along $`\mathrm{\Sigma }`$ is Morse, Pellikaan answered this question positively. Pellikaan’s formula expresses the Euler characteristic in terms of the Jacobian number $`j(f)`$, $`\delta `$ and the Milnor number $`\mu (\mathrm{\Sigma })`$ of $`\mathrm{\Sigma }`$. These numbers can be computed directly by counting the dimensions of certain finite dimensional vector spaces. The development of computer algebra makes this kind of algebraic formulae more and more important and popular. In this article we answer the question by proving a similar formula for function germs with line singularities on a weighted homogeneous space $`X`$ with isolated complete intersection singularity (see Proposition 7). Remark that for a function germ $`f`$ with isolated singularity on a weighted homogeneous complete intersection with isolated singularity, Bruce and Robert have proved an algebraic formula for the Milnor number of $`f`$.
## 2. Non-isolated singularities on singular spaces
Let $`𝒪_^m`$ be the structure sheaf of $`^m`$. The stalk $`𝒪_{^m,0}`$ of $`𝒪_^m`$ at 0 is a local ring, consisting of germs at 0 of analytic functions on $`^m`$. The ring $`𝒪_{^m,0}`$ is often denoted by $`𝒪_m`$, or simply by $`𝒪`$ when no confusion can be caused. The unique maximal ideal of $`𝒪_m`$ is denoted by $`𝔪_m`$ or $`𝔪`$.
Let $`(X,0)(^m,0)`$ be the germ of a reduced analytic subspace $`X`$ of $`^m`$ defined by an ideal $`𝔥`$ of $`𝒪`$, generated by $`h_1,\mathrm{},h_p𝒪`$. Let $`𝔤`$ be the ideal generated by $`g_1,\mathrm{},g_n𝒪`$. The germ of the analytic space defined by $`𝔤`$ at $`0`$ is denoted by $`(\mathrm{\Sigma },0)`$. Write $`𝒪_X:=𝒪/𝔥`$ and $`𝒪_\mathrm{\Sigma }:=𝒪/𝔤.`$
Let $`\mathrm{Der}(𝒪)`$ denote the $`𝒪`$-module of germs of analytic vector fields on $`^m`$ at $`0`$. Then $`\mathrm{Der}(𝒪)`$ is a free $`𝒪`$-module with $`\frac{}{z_1},\mathrm{},\frac{}{z_m}`$ as basis, where $`z_1,\mathrm{},z_m`$ are the local coordinates of $`(^m,0).`$ $`\mathrm{Der}(𝒪)`$ is a Lie algebra with the bracket defined by $`[\xi ,\eta ]:=\xi \eta \eta \xi `$ for all $`\xi ,\eta \text{Der}.`$
Define $`\mathrm{Der}_𝔥(𝒪):=\{\xi \mathrm{Der}(𝒪)\xi (𝔥)𝔥\}`$, which is the $`𝒪`$-module of logarithmic vector fields along $`(X,0)`$ and a Lie subalgebra of $`\mathrm{Der}(𝒪)`$ . When $`𝔥`$ is a radical ideal defining the analytic space $`X`$, $`\mathrm{Der}_𝔥(𝒪)`$ is often denoted by $`D_X`$. Geometrically, $`D_X`$ consists of all the germs of vector fields that are tangent to the smooth part of $`X`$. When $`X`$ is a weighted homogeneous complete intersection with isolated singularity, one can write down precisely all the generators of $`D_X`$ (cf. ).
For $`f𝒪`$, the ideal $`J_X(f):=\{\xi (f)\xi D_X\}`$ is called the (relative) Jacobian ideal of $`f`$. Obviously, when $`X`$ is the whole space $`^m`$, namely, $`𝔥=\{0\},`$ then $`J_X(f)=J(f)`$, the Jacobian ideal of $`f`$.
Let $`𝒮=\{S_\alpha \}`$ be an analytic stratification of $`X`$, $`f:(X,0)(,0)`$ an analytic function germ. The critical locus $`\mathrm{\Sigma }_f^𝒮`$ of $`f`$ relative to the stratification $`𝒮`$ is the union of the critical loci of $`f`$ restricted to each of the strata $`S_\alpha `$, namely, $`\mathrm{\Sigma }_f^𝒮=_{S_\alpha 𝒮}\overline{\mathrm{\Sigma }_{f|S_\alpha }}`$. We denote $`\mathrm{\Sigma }_f^𝒮`$ by $`\mathrm{\Sigma }_f`$ when $`𝒮`$ is clear from the context. If the dimension of $`\mathrm{\Sigma }_f^𝒮`$ is not positive, we say that $`f`$ defines (or has) isolated singularities on $`X`$. If the dimension of $`\mathrm{\Sigma }_f^𝒮`$ is positive, we say that $`f`$ defines (or has) non-isolated singularities on $`X`$. If $`\mathrm{\Sigma }_f^𝒮`$ is one dimensional smooth complex manifold, we say that $`f`$ defines (or has) a line singularity on $`X`$.
For an analytic space $`X`$ embedded in a neighborhood $`U`$ of $`0^m`$, there is a logarithmic stratification $`𝒮_{\mathrm{log}}:=\{S_\alpha \}`$ of $`U`$ (see ). In general, $`𝒮_{\mathrm{log}}`$ is not locally finite. If $`𝒮_{\mathrm{log}}`$ is locally finite, then $`X`$ is said to be holonomic.
Let $`X`$ be of pure dimension. The collection $`𝒮_{\mathrm{log}}^{}=\{XS_\alpha S_\alpha 𝒮_{\mathrm{log}}\}`$ is a stratification of $`X`$ which will be called the logarithmic stratification of $`X`$ in this article. Especially, when $`X`$ has isolated singularity in $`0`$, then $`\{0\}`$ and the connected components of $`X\{0\}`$ form a holonomic logarithmic stratification of $`X`$. So 0 is always a critical point of any germ $`f:(X,0)(,0)`$ relative to this stratification. Hence for $`f𝔪`$, $`\mathrm{\Sigma }_f=\{pX\xi (f)(p)=0\text{ for all }\xi D_X\}`$ is the critical locus of $`f`$ relative to the logarithmic stratification. Obviously $`\mathrm{\Sigma }_f=X𝒱(J_X(f))`$.
###### Definition 1.
Let $`𝔥=(h_1,\mathrm{},h_p)𝔤=(g_1,\mathrm{},g_n)`$ be ideals of $`𝒪_{^m,0}`$. Define a subset of $`𝒪`$, called the primitive ideal of $`𝔤`$ relative to $`𝔥`$:
$$_𝔥𝔤:=\{f𝔤\xi (f)𝔤\text{ for all }\xi \mathrm{Der}_𝔥(𝒪)\}.$$
In the following we always assume that $`𝔥`$ and $`𝔤`$ are radical, $`X=𝒱(𝔥)`$ and $`\mathrm{\Sigma }=𝒱(𝔤).`$ In this case, $`_𝔥𝔤`$ is denoted by $`_X𝔤`$ or $`𝔤`$ when no confusion can be caused by this.
###### Remarks 2.
* When $`X`$ is smooth this definition was given by Pellikaan . It is straightaway to verify that $`_𝔥𝔤`$ is an ideal of $`𝒪`$, and $`𝔤^2+𝔥_𝔥𝔤𝔤`$ always holds. And for $`𝔤_i𝔥`$ $`(i=1,2)`$, we have $`_𝔥𝔤_1_𝔥𝔤_2=_𝔥\left(𝔤_1𝔤_2\right)`$;
* Geometrically, the relative primitive ideal collects all the functions whose zero level surfaces pass through $`\mathrm{\Sigma }`$ and are tangent to the regular part $`X_{\mathrm{reg}}`$ of $`X`$ along $`\mathrm{\Sigma }X_{\mathrm{reg}}`$.
* The singular locus (relative to $`𝒮_{\mathrm{log}}^{}`$) of $`f`$ is $`\mathrm{\Sigma }_f:=𝒱(J_X(f))X`$, and $`\mathrm{\Sigma }_ff^1(0)`$ if $`f(0)=0`$. If $`f𝔤`$, then $`\mathrm{\Sigma }X𝒱(J_X(f))=\mathrm{\Sigma }_f`$. Conversely, for $`f𝔪`$, we have $`f𝔤`$ when $`\mathrm{\Sigma }\mathrm{\Sigma }_f`$ and $`𝔤`$ is radical. The reason is: $`J_X(f)𝔤`$, and since $`f`$ takes finite values on $`\mathrm{\Sigma }_f`$ and $`0\mathrm{\Sigma }`$, $`f|_\mathrm{\Sigma }=0`$, so $`f^k𝔤`$ for some $`k`$. Hence $`f𝔤`$ since $`𝔤`$ is radical.
* The relative primitive ideals have been generalized to higher relative primitive ideals in . Under the assumption that $`𝔥`$ is pure dimension, $`𝔤`$ is radical, and the Jacobian ideal of $`𝔥`$ is not contained in any associated prime of $`𝔤`$, it was proved that the primitive ideal $`_𝔥𝔤`$ is the inverse image in $`𝒪`$ of the second symbolic power of the quotient ideal $`\overline{𝔤}:=𝔤/𝔥`$ of $`𝒪_X`$. Remark that the results in generalized the results of .
## 3. Transversal singularities
Let $`(X,0)(^{n+1},0)`$ be the germ of a reduced analytic space with isolated singularity in 0. Let $`\mathrm{\Sigma }`$ be a reduced curve germ on $`(X,0)`$ defined by $`𝔤`$ and have isolated singularity in 0. A germ $`f𝔤`$ is called a transversal $`A_1`$ singularity along $`\mathrm{\Sigma }`$ on $`X`$ if its singular locus $`\mathrm{\Sigma }_f=\mathrm{\Sigma }`$, and, for $`P\mathrm{\Sigma }0,f`$ has only $`A_1`$ singularity transversal to the branch of $`\mathrm{\Sigma }`$ containing $`P`$. It was proved in , that $`f`$ is a transversal $`A_1`$ singularity along $`\mathrm{\Sigma }`$ on $`X`$ if and only if the Jacobian number $`j(f):=dim(𝔤/(𝔥+J_X(f))<\mathrm{}`$.
There exist admissible linear forms $`l`$ (see ) such that $`\{l=0\}\mathrm{\Sigma }=\{0\}`$ and $`\{l=0\}`$ intersects both $`X`$ and $`\mathrm{\Sigma }`$ transversally, and $`\{l=0\}Xf^1(0)`$ has isolated singularity at the origin. We assume $`\{l=0\}`$ is the first coordinate hyperplane $`\{z_0=0\}`$ of $`^{n+1}`$.
Let $`D_X`$ be generated by $`\xi ^0,\xi ^1,\mathrm{}\xi ^s`$. Denote by $`D_X^0`$ the submodule of $`D_X`$ generated by those $`\xi ^i`$ such that if we write $`\xi ^i=_{j=0}^n\xi _j^i\frac{}{z_j}`$, then $`\xi _0^i𝔤`$, and by $`D_X^1`$ the submodule of $`D_X`$ generated by those $`\xi ^i`$ such that if we write $`\xi ^i=_{j=0}^n\xi _j^i\frac{}{z_j}`$, then $`\xi _0𝔤`$, thus $`D_X=D_X^0+D_X^1`$. Denote
$$J^0(f):=D_X^0(f)=\{\xi (f)\xi D_X^0\},J^1(f):=D_X^1(f)=\{\xi (f)\xi D_X^1\}.$$
If $`f`$ is clear from the context we just write $`J^0`$ and $`J^1`$.
###### Lemma 3.
Let $`f𝔤`$, and $`z_0,z_1,\mathrm{},z_n`$ be the coordinates of $`^{n+1}`$ such that $`z_0=0`$ is admissible. The transversal singularity type of $`f`$ along every branch of $`\mathrm{\Sigma }`$ is constant at all the points of $`\mathrm{\Sigma }\{0\}`$ if and only if
$$dim_{}\left(\frac{𝒪}{𝔤+((J^1+𝔥):J^0)}\right)<\mathrm{}.$$
###### Proof.
The inequality means that $`\mathrm{\Sigma }𝒱((J^1+𝔥):J^0)=\{0\}`$. For $`ϵ>0`$ small enough, let $`P\mathrm{\Sigma }\{z_0=t\}`$, $`0<|t|<ϵ`$. $`P𝒱((J^1+𝔥):J^0)`$ if and only if $`(J^0)_P(J^1)_P`$, where $`(J^i)_P`$ is the localization of the ideal at $`P`$. Since $`z_0=0`$ intersects both $`\mathrm{\Sigma }`$ and $`X`$ transversally and $`X`$ is smooth at $`P`$, we can choose the local coordinates such that locally the branch of $`\mathrm{\Sigma }`$ containing $`P`$ is the first axis. Furthermore, we can arrange the coordinate transformation such that under this transformation $`D_X^0`$ and $`D_X^1`$ are preserved. By this we mean that any derivation of $`X`$ with the first component non-zero at $`P`$ will remain non-zero at $`P`$ and any derivation of $`X`$ with first component zero at $`P`$ will remain zero at $`P`$. We let $`x=z_0,y_1,\mathrm{},y_{np}`$ be the new local coordinates in a neighborhood of $`P`$ in $`X`$, then at $`P`$, $`(J^0)_P=\left(\frac{f}{x}\right)𝒪_{X,P}`$ and $`(J^1)_P=(\frac{f}{y_1},\mathrm{},\frac{f}{y_{np}})𝒪_{X,P}`$. By , the inequality is equivalent to the constancy of the transversal singularity type of $`f`$. $`\mathrm{}`$
Let $`X`$, $`\mathrm{\Sigma }`$ and $`f𝔤`$ be the same as before. There is an integer $`k_0`$ such that for all $`kk_0`$, $`f_k=f+\frac{1}{k+1}x^{k+1}`$ defines an isolated singularity at $`O`$. Let $`\mu (f_k)`$ be the Milnor number of $`f_k`$. If $`X`$ is a weighted homogeneous complete intersection with isolated singularity, by 7.7, we have
$$\mu (f_k)=dim\left(\frac{𝒪}{𝔥+J_X(f_k)}\right).$$
By the exact sequence
$$0\frac{𝔤+J_X(f_k)}{𝔥+J_X(f_k)}\frac{𝒪}{𝔥+J_X(f_k)}\frac{𝒪}{𝔤+J_X(f_k)}0,$$
we know that
$`(\text{3}.1)`$
$$\mu (f_k)=dim\left(\frac{𝔤+J_X(f_k)}{𝔥+J_X(f_k)}\right)+dim\left(\frac{𝒪}{𝔤+J_X(f_k)}\right)$$
Define
$$e_k:=dim\frac{𝒪}{𝔤+J_X(f_k)},\sigma (\mathrm{\Sigma }X,0):=dim\left(\frac{𝒪_\mathrm{\Sigma }}{J_X(x)}\right),\text{multi}_x(\mathrm{\Sigma }):=dim\left(\frac{𝒪_\mathrm{\Sigma }}{(x)}\right).$$
Then obviously $`e_k=\sigma (\mathrm{\Sigma }X,0)+k\text{multi}_x(\mathrm{\Sigma }),`$ and
$$\frac{𝔤+J_X(f_k)}{𝔥+J_X(f_k)}=\frac{𝔤}{(𝔥+J_X(f_k))𝔤}.$$
Hence
$`(\text{3}.2)`$
$$\mu (f_k)=\sigma (\mathrm{\Sigma }X,0)+k\text{multi}_x(\mathrm{\Sigma })+dim\left(\frac{𝔤}{(𝔥+J_X(f_k))𝔤}\right)$$
## 4. Line singularities
In this section, we assume that $`\mathrm{\Sigma }`$ is a line in $`^{n+1}`$ defined by the ideal $`𝔤`$, and $`X`$ is a space with isolated complete intersection singularity defined by $`𝔥𝔤`$.
###### Lemma 4.
Let $`X`$ be a space with isolated complete intersection singularity, containing a line $`\mathrm{\Sigma }`$, which is chosen to be the first axis of a local coordinate system of $`^{n+1}`$. If $`X`$ is weighted homogeneous with respect to this coordinate system, then $`D_X=𝒪\xi _E+D_X^1`$ and
$`J_X(f_k)`$ $`=(\xi _E(f_k))𝒪+J^1(f_k)`$
$`(\xi _E(f_k))𝒪+D_X^1(f)+D_X^1(x^{k+1}),`$
where $`J^1(f_k)=\{\xi (f_k)\xi D_X^1\}𝔤`$.
###### Proof.
Let $`x,y_1,\mathrm{},y_n`$ be the local coordinates of $`(^{n+1},0)`$ in the statement of the lemma. We know by that the Euler derivation $`\xi _E=w_0x\frac{}{x}+\underset{k=1}{\overset{n}{}}w_ky_k\frac{}{y_k}D_X`$, where $`w_0,w_1,\mathrm{},w_n`$ are the weights of the coordinates $`x,y_1,\mathrm{},y_n`$ respectively. Let $`\xi =\xi ^0\frac{}{x}+\underset{k=1}{\overset{n}{}}\xi ^k\frac{}{y_k}D_X`$. Since $`𝒪_\mathrm{\Sigma }`$ is a principal ideal domain and $`X`$ has isolated singularity at $`O`$, $`\overline{\xi }^0(\overline{x})`$. Let $`\overline{\xi }^0=\overline{x}\overline{\xi }_1^0`$. Then $`\xi ^{}=\xi \frac{1}{w_0}\xi _1^0\xi _E`$ has the coefficient of $`\frac{}{x}`$ in $`𝔤`$. $`\mathrm{}`$
###### Lemma 5.
Let $`\mathrm{\Sigma }`$ and $`X`$ be the same as in Lemma 4. Let $`f𝔤`$ have transversal $`A_1`$ singularity along $`\mathrm{\Sigma }`$. For $`k>0`$ sufficiently large, we have
$$\stackrel{~}{J}:=\xi _E(f)𝔤+D_X^1(f)+𝔥=(𝔥+J_X(f_k))𝔤$$
###### Proof.
By Lemma 4, we have
$$(𝔥+J_X(f_k))𝔤=(\xi _E(f_k))𝔤+D_X^1(f_k)+𝔥.$$
For $`a(\xi _E(f_k))𝔤`$, $`a=a_0x^{k+1}+a_0\xi _E(f)𝔤`$. Since $`\xi _E(f)𝔤`$, $`a_0𝔤`$. Hence $`(\xi _E(f_k))𝔤=(\xi _E(f_k))𝔤`$ and
$`(\text{4}.1)`$
$$\begin{array}{cc}(𝔥+J_X(f_k))𝔤\hfill & =(x^{k+1}+\xi _E(f))𝔤+D_X^1(f_k)+𝔥\hfill \\ & (x^{k+1}+\xi _E(f))𝔤+x^kD_X^1(x)+D_X^1(f)+𝔥\hfill \end{array}$$
Since $`j(f)<\mathrm{}`$, there is a $`k_1`$ such that when $`k>k_1`$, $`x^k𝔤𝔥+J_X(f)=(\xi _E(f))+D_X^1(f)+𝔥`$.
By Lemma 3, there is a $`k_2`$ such that $`x^{k_2}𝔤+((D_X^1(f)+𝔥):(\xi _E(f)))`$, and $`x^{k_2}\xi _E(f)(\xi _E(f))𝔤+D_X^1(f)+𝔥`$. Hence when $`k>k_1+k_2`$,
$$a_0(x^{k+1}+\xi _E(f))+b_0x^k\stackrel{~}{J}\text{ and }(𝔥+J_X(f_k))𝔤\stackrel{~}{J}.$$
On the other hand, there is an integer $`n_1>>0`$ such that when $`kn_1`$ (see (3.2))
$$x^k𝔤(𝔥+J_X(f_k))𝔤\stackrel{~}{J}𝔤.$$
Then for $`k>n_1`$, by the equality in (4.1)
$`\stackrel{~}{J}`$ $`\xi _E\left(f+{\displaystyle \frac{x^{k+1}}{k+1}}\right)𝔤+x^{k+1}𝔤+D_X^1(f)+𝔥`$
$`=(\xi _E(f)+x^{k+1})𝔤+D_X^1(f)+𝔥+x^{k+1}𝔤`$
$`=(𝔥+J_X(f_k))𝔤+x^{k+1}𝔤`$
$`(𝔥+J_X(f_k))𝔤+x^k𝔤`$
$`(𝔥+J_X(f_k))𝔤+x^{kn_1}\stackrel{~}{J}`$
$`(𝔥+J_X(f_k))𝔤+𝔪\stackrel{~}{J}.`$
By Nakayama’s lemma, $`(𝔥+J_X(f_k))𝔤=\stackrel{~}{J}`$. $`\mathrm{}`$
###### Lemma 6.
Under the assumption of Lemma 5, we have
(0) $`1)`$ $`L:={\displaystyle \frac{𝔥+J_X(f)}{\stackrel{~}{J}}}{\displaystyle \frac{(\xi _E(f))}{(\xi _E(f))\stackrel{~}{J}}};`$
(0) $`2)`$ $`\mathrm{Ann}(L)=𝔤+((D_X^1(f)+𝔥):(\xi _E(f)));`$
(0) $`3)`$ $`L{\displaystyle \frac{𝒪}{\mathrm{Ann}(L)}}={\displaystyle \frac{𝒪}{𝔤+((D_X^1(f)+𝔥):(\xi _E(f)))}}.`$
###### Proof.
It is an easy exercise in commutative algebra, we omit it. $`\mathrm{}`$
From the exact sequence
$$0\frac{𝔥+J_X(f)}{(𝔥+J_X(f_k))𝔤}\frac{𝔤}{(𝔥+J_X(f_k))𝔤}\frac{𝔤}{𝔥+J_X(f)}0,$$
we have
$$\mu (f_k)=j(f)+e_k+dim\left(\frac{𝔥+J_X(f)}{(𝔥+J_X(f_k))𝔤}\right).$$
Notice that $`\sigma (\mathrm{\Sigma }X,0)=\text{multi}_x(\mathrm{\Sigma })=1`$, $`e_k=k+1`$. By Lemma 6, we have
$$\mu (f_k)=k+1+j(f)+dim\left(\frac{𝒪}{𝔤+((D_X^1(f)+𝔥):(\xi _E(f)))}\right)$$
Let $`F`$ and $`F_k`$ are the Milnor fibre of $`f`$ and $`f_k`$ respectively. Iomdin-Lê’s formula , says that
$$\chi (F)=\chi (F_k)+(1)^{dimX}(k+1).$$
But $`\chi (F_k)=1+(1)^{dimX1}\mu (f_k)`$. Since in our case $`\sigma (\mathrm{\Sigma }X,0)=1`$, we have proved
###### Proposition 7.
Let $`X`$ be a space with isolated complete intersection singularity containing a line $`\mathrm{\Sigma }`$, which is chosen to be the first axis of a local coordinate system of $`^{n+1}`$. Assume that $`X`$ is weighted homogeneous with respect to the coordinate system. For an analytic function $`f_X𝔤`$ with $`j(f)<\mathrm{}`$, the Euler characteristic of the Milnor fibre $`F`$ of $`f`$ is
$`(\text{4}.2)`$
$$\chi (F)=1+(1)^{dimX1}(j(f)+\nu )$$
where
$`\mathrm{}`$
$$\nu =dim\left(\frac{𝒪}{𝔤+((D_X^1(f)+𝔥):(\xi _E(f)))}\right).$$
###### Remark 8.
Formula (4.2) allows us to use a computer program to compute the Euler characteristic effectively. In fact, we have a small Singular program to calculate $`\chi (F)`$ for a function with critical locus a line on a hypersurface $`X`$. We use it to check the examples in Example 10.
Let $`\mathrm{\Sigma }`$ be a line in $`^{n+1}`$ defined by the ideal $`𝔤=(y_1,\mathrm{},y_n)`$, $`X`$ a space with isolated complete intersection singularity defined by $`𝔥=(h_1,\mathrm{},h_p)𝔤`$. By changing the generators of $`𝔤`$, we can write
$$h_i\underset{k=1}{\overset{p}{}}b_{ik}y_kmod𝔤^2$$
such that the determinant $`b`$ of the matrix $`B=(b_{ik})`$ is a non-zero divisor in $`𝒪_\mathrm{\Sigma }`$. Note that $`y_1,\mathrm{}y_p`$ are projected to zero or the generators of the torsion part $`T(M)`$ of the conormal module $`M=𝔤/(𝔤^2+𝔥)`$ of $`𝔤/𝔥`$, and $`y_{p+1},\mathrm{}y_n`$ form a basis of the free module $`N=M/T(M)`$. We call $`\lambda (\mathrm{\Sigma }X):=dim_{}T(M)=dim𝒪/(𝔤+(b))`$ the torsion number of the space pair $`(\mathrm{\Sigma },X)`$. For $`f=h_{kl}y_ky_l𝔤^2`$, define
$$\mathrm{\Delta }=det(h_{kl})_{p+1k,ln},\delta _f:=dim(\frac{𝒪}{𝔤+\mathrm{\Delta }}).$$
###### Question 9.
For $`f𝔤^2`$, does the equality $`\nu =2\lambda (\mathrm{\Sigma }X)+\delta _f1`$ always hold? Or under what conditions does it hold?
The equality in Question 9 and (4.2) show the geometric meaning of $`\chi (F)`$. By using Singular we have checked that it holds for all the examples we know.
###### Example 10.
Let $`X_{k,l}`$ be an $`A_{k,l}`$ singularity defined by $`h=x^ly+x^sz^2+yz`$. This is an $`A_k`$ singularity with $`k=2l+s1`$ ($`l1,s0`$). Since we take the line $`\mathrm{\Sigma }`$ with torsion number $`l`$ as $`x`$-axis, we have the definition equation (see ). By resolution of singularities, one can prove that for any function $`f`$ on $`X_{k,l}`$, if it has isolated line singularity and $`A_1`$ type transversal singularity, then the Milnor fibre $`F`$ of $`f`$ is a bouquet of circles (see ). Then $`\mu (f)=j(f)+\nu `$.
* We consider a function $`g:X_{k,l}.`$ For generic $`(a,b,c)^3`$, let $`g=ay^2+byz+cz^2`$, for example, take $`g=y^2yz+\frac{1}{2}z^2`$, a calculation shows that $`\mu (g)=6l3`$.
* Let $`f:X_{k,l}`$ be defined by $`f=y+\frac{1}{2}z^2`$. A calculation shows that
$$\mu (f)=\{\begin{array}{cc}l1\hfill & \text{ if }s=0,\hfill \\ l+3s2\hfill & \text{ if }1sl1,\hfill \\ 4l2\hfill & \text{ if }sl.\hfill \end{array}$$ |
warning/0001/astro-ph0001299.html | ar5iv | text | # The X-ray Properties of 𝑧>4 Quasars
## 1. Introduction
Quasars with redshifts larger than 4 were first discovered more than a decade ago (Warren et al. 1987). Recently many new $`z>4`$ quasars have been discovered (e.g., Fan et al. 1999, 2000) with expectations of an even greater increase in the near future due to new surveys (e.g., the Sloan Digital Sky Survey should identify $`1000`$ quasars with $`z>4.5`$ – Schneider 1999; York et al. 2000). Currently there are 85 $`z>4`$ quasars that have appeared in journals, and an additional number can be found on various World Wide Web pages. Quasars at $`z>4`$ provide us with direct information about the first 10% of cosmic time. They are among the most luminous objects known, and from the Eddington limit many require $``$ $`10^8`$$`10^9`$ M black holes. They have wide cosmological importance since they must be associated with deep potential wells in the earliest massive collapsed structures, and their strong evolution provides clues about the process by which the remarkably homogeneous $`z1000`$ Universe revealed by the cosmic microwave background is transformed into the inhomogeneous Universe seen today (e.g., Efstathiou & Rees 1988; Turner 1991).
To date $`z>4`$ quasars have been mainly studied at optical wavelengths (e.g., in order to determine their redshifts). There has also been some progress in studying their far-infrared and radio properties (e.g., Schmidt et al. 1995; Omont et al. 1996; McMahon et al. 1999). However, their properties as a whole at these and other wavelengths have not yet been fully explored. X-ray emission appears to be a universal property of quasars at $`z`$ 0–2, and X-rays have also been studied from many $`z`$ 2–4 quasars. However, at $`z>4`$ the X-ray properties of quasars are much less well understood; only six $`z>4`$ quasars have been detected in X-ray. The luminous X-ray emission from quasars reveals the physical conditions in the immediate vicinities of their black holes, and X-ray studies of high-redshift quasars can in principle discover if quasar central power sources and quasar environments evolve over cosmic time (e.g., Bechtold et al. 1994b; Blair et al. 1998; Elvis et al. 1998; Fiore et al. 1998).
We list the $`z>4`$ quasars previously detected in the X-ray band in Table 1. For each quasar we give, in the first seven columns, its coordinates, redshift, absorption-corrected X-ray flux in the detection band, and the reference to the paper where the X-ray detection was made. While the most basic X-ray properties (e.g., $`\alpha _{\mathrm{ox}}`$, the slope of a nominal power law between 2500 Å and 2 keV) of these quasars appear to be generally consistent with those of quasars at lower redshifts, the constraints are not tight and require substantial improvement. Comparisons of these properties with those of the majority of low-redshift quasars are difficult as most of the X-ray detected $`z>4`$ quasars were selected in different ways: three of them are X-ray selected objects, two are radio selected, and only one is optically selected.
All objects but one in Table 1 were detected using ROSAT, demonstrating the ability of this satellite to detect $`z>4`$ quasars. Encouraged by this we have systematically searched for detections of $`z>4`$ quasars in the ROSAT public database. In this paper we present our results, which double the number of $`z>4`$ quasars detected in X-rays. We increase the number of optically selected $`z>4`$ quasars from one to seven and provide limits for 15 others. In § 2 we present the database search and in § 3 we discuss our results.
Throughout this paper we use the cosmological parameters $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.5`$. We define the energy index $`\alpha `$ as $`f_\nu \nu ^\alpha `$ and likewise the photon index $`\mathrm{\Gamma }=\alpha +1`$ with photon flux density $`f(E)E^\mathrm{\Gamma }`$ in photons cm<sup>-2</sup> s<sup>-1</sup> keV<sup>-1</sup>. Unless otherwise noted, we use $`\alpha {}_{\mathrm{o}}{}^{}=0.5`$ in the UV-optical range and $`\alpha {}_{\mathrm{x}}{}^{}=1`$ in the X-ray range. These are representative values of these parameters for lower redshift quasars (e.g., Netzer 1990; Reeves et al. 1997).
## 2. Search and Analysis
We have searched the ROSAT public database<sup>1</sup><sup>1</sup>1Via: http://heasarc.gsfc.nasa.gov/W3Browse for all fields which include the optical positions of $`z>4`$ quasars in the literature. We have found 27 quasars’ positions (out of the total 85 $`z>4`$ quasars) to lie in ROSAT fields. For ten quasars we found only one observation, while for the others there were two or more. We retrieved from the ROSAT public database up to four observations (when available) for each quasar. We preferentially chose long observations where the quasar was close to the field’s center. Twenty-six quasars were observed by the Position Sensitive Proportional Counter (PSPC; Pfeffermann et al. 1987) with several of them also having High Resolution Imager (HRI; David et al. 1999) observations, and one quasar was observed only with the HRI. Out of the 27 quasars, 12 were the observation’s target and 15 were serendipitously in the detector’s field of view.
All observations were processed using the PROS software in IRAF.<sup>2</sup><sup>2</sup>2IRAF (Image Reduction and Analysis Facility) is distributed by the National
Optical Astronomy Observatories, which are operated by AURA, Inc., under coo-
perative agreement with the National Science Foundation. We have manually inspected the images (and smoothed versions thereof) and have measured the net counts around the optical positions of the quasars. Typically PSPC positions are good to $``$20–30″ (e.g., Voges et al. 1996), and indeed all our detections but one are $`20`$″ from the optical position. The aperture size for count extraction was scaled to take into account the size of the point spread function (PSF) at the off-axis angle of the quasar. We were able to detect, at above the 2.5$`\sigma `$ level, an X-ray source and extract the X-ray count rate for nine quasars. For the other 18 quasars we were only able to determine upper limits to their X-ray count rates. We have also searched for X-ray detections of 86 quasars listed on World Wide Web pages but not in the literature, and we find none.
In the upper part of Table 2 we present the quasars detected. Their names, coordinates and redshifts are given in the first four columns. The monochromatic AB<sub>1450(1+z)</sub> magnitude is listed in column (5). The AB<sub>1450(1+z)</sub> magnitudes are from Schneider, Schmidt, & Gunn (1991), Henry et al. (1994), Storrie-Lombardi et al. (1996), and Hook & McMahon (1998), with estimated errors of $`\pm 0.1`$ magnitudes. The Galactic column density, found using the H I map of Dickey & Lockman (1990), is given in column (6). We have used the AB<sub>1450(1+z)</sub> magnitude and a flux-density power law with $`\alpha _0=0.5`$ to compute the rest-frame 2500 Å flux density and luminosity which are listed in columns (7) and (8). The absolute $`B`$ magnitude, given in column (9), was found using the equation
$$M_B=\mathrm{AB}_{1450(1+z)}5\mathrm{log}(\sqrt{1+z}1)45.03$$
(1)
which is equation 5 of Schneider, Schmidt, & Gunn (1989) adapted for $`H_0=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>.
In columns (10)–(18) of Table 2 we list the X-ray observations and properties of the quasars. The sequence-ID and observing date are listed in columns (10) and (11). The angular distance of the quasar’s position from the center of the field is listed in column (12). The number of background-subtracted counts (in the broad band 0.1–2 keV) is given in column (13), and the vignetting, exposure-map and background corrected count-rate is listed in column (14). We use the PIMMS software (Mukai 1997) to define a power law with photon index $`\mathrm{\Gamma }=2`$, which is then used in the XSPEC software (Arnaud 1996) to evaluate the absorption-corrected 0.1–2 keV observed flux, which is listed in column (15). From this power law we also calculated the rest-frame 2-keV flux density and luminosity which are listed in columns (16) and (17). (Using PROS to calculate these quantities yields consistent results.) Finally, we list in column (18) the effective optical-to-X-ray power-law slope, $`\alpha _{\mathrm{ox}}`$, defined as:
$$\alpha _{\mathrm{ox}}=\frac{\mathrm{log}[(f_\nu (2\mathrm{keV})/f_\nu (2500\text{Å})]}{\mathrm{log}[\nu (2\mathrm{keV})/\nu (2500\text{Å})]},$$
(2)
where $`f_\nu `$’s are flux densities at the given wavelengths and $`\nu `$’s are the corresponding frequencies. The uncertainties for the fluxes, luminosities, and $`\alpha _{\mathrm{ox}}`$ values can be propagated from the relative error on the number of counts, which has a mean value of 30%. Another uncertainty involves our assumption of $`\mathrm{\Gamma }=2`$.
\[l\] Table 1 X-ray Detected Quasars at $`z>4`$ Known Prior to This Work Object RA (2000.0) Dec $`z`$ Flux<sup>a</sup><sup>a</sup>footnotemark: Band \[keV\] Ref. AB $`N_\mathrm{H}`$<sup>b</sup><sup>b</sup>footnotemark: $`f_\nu `$<sup>c</sup><sup>c</sup>footnotemark: $`\mathrm{log}(\nu L_\nu )`$ $`M_B`$ $`f_\nu `$<sup>d</sup><sup>d</sup>footnotemark: $`\mathrm{log}(\nu L_\nu )`$ $`\alpha _{\mathrm{ox}}`$ $`R`$ (1) (2) (3) (4) (5) (6) (7) (8) (9) (10) (11) (12) (13) (14) (15) (16) Q 0000$``$2619 00 03 22.9 $``$26 03 19 4.098 8.7$`\times 10^{14}`$<sup>e</sup><sup>e</sup>footnotemark: 0.1–2.4 1 17.5 1.67 4.767 46.9 $``$28.0 2.9 45.3 $``$1.62 $`<`$0.6 RX J1028.6$``$0844 10 28 37.7 $``$08 44 39 4.276 8.3$`\times 10^{13}`$ 0.1–2.4 2 20.6<sup>f</sup><sup>f</sup>footnotemark: 4.55 0.274 45.7 $``$25.0 28.6 46.3 $``$0.76 7400 RX J105225.9+571905 10 52 25.9 +57 19 07 4.45 2.3$`\times 10^{15}`$ 0.5–2.0 3 22.6 0.56 0.043 44.9 $``$23.1 0.19 44.2 $``$1.29 $`<`$190 GB 1428+4217 14 30 23.7 +42 04 36 4.715 $`1\times 10^{12}`$ 0.1–2.4 4 19.4 1.40 0.829 46.2 $``$26.4 37.2 46.5 $``$0.90 1900 GB 1508+5714 15 10 02.8 +57 02 44 4.301 $`1\times 10^{12}`$ 0.3–3.5 5, 6 19.8 1.47 0.573 46.0 $``$25.8 44.7 46.5 $``$0.81 2700 RX J1759.4+6638 17 59 27.9 +66 38 53 4.320 $`1.2\times 10^{14}`$ 0.5–2.0 7 19.3 4.23 0.908 46.2 $``$26.3 0.95 44.8 $``$1.53 30 <sup>a</sup>In units of erg s<sup>-1</sup> cm<sup>-2</sup>. <sup>b</sup>In units of $`10^{20}`$ cm<sup>-2</sup>. <sup>c</sup>At 2500 Å in units of $`10^{27}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup>. <sup>d</sup>At 2 keV in units of $`10^{31}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup>. <sup>e</sup>Computed using Bechtold et al. (1994a) model # 3. <sup>f</sup>Estimated from the 2500 Å flux given in Zickgraf et al. (1997). References.— (1) Bechtold et al. 1994a; (2) Zickgraf et al. 1997; (3) Schneider et al. 1998; (4) Fabian et al. 1997; (5) Mathur & Elvis 1995; (6) Moran et al. 1996; (7) Henry et al. 1994.
As radio-quiet quasars typically have $`\mathrm{\Gamma }`$ = 1.7–2.3, our assumption of $`\mathrm{\Gamma }=2`$ might introduce an additional uncertainty of $``$15% in the flux estimates. As the true $`\mathrm{\Gamma }`$ for each quasar is unknown (other than GB 1428+4217 and GB 1508+5714), we do not quote an error on the latter quantities. We estimate the total fractional uncertainties for the fluxes to be in the range of 30–50%.
X-ray images of all nine detected quasars are presented in Fig. 1 which was created using the adaptive smoothing method of Ebeling, White, & Rangarajan (2000) applied to the full-band images. For most of the images we used a 2.5$`\sigma `$ level of smoothing; for three low-$`\sigma `$ detected objects (BR 0951$``$0450, BRI0952$``$0115, and BR 1202$``$0725) we used a 2.0$`\sigma `$ level of smoothing.
For each of the PSPC observations listed in the upper part of Table 2 we also calculated the X-ray counts in the soft band (0.1–0.5 keV) and the hard band (0.5–2.0 keV). We used these to calculate the rest-frame 2-keV flux density, luminosity, and $`\alpha _{\mathrm{ox}}`$, in the same way as described above. The results were similar to the broad-band results. However, as the statistical uncertainties in the soft and hard bands were large due the small number of counts (note the small number of counts in the broad band – column of Table 2), we do not include these results in our analysis.
In the lower part of Table 2 we list the $`z>4`$ quasars which we were unable to detect in the X-ray observations. For these quasars we give 3$`\sigma `$ upper limits (listed in column ), where $`\sigma `$ is the square root of the counts in an aperture of size appropriate to the quasar’s off-axis angle and centered at the quasar’s optical position. For each quasar we list the one observation which gave the faintest upper limit.
For 27 objects the number of expected detections which are merely statistical fluctuations at the 2.5$`\sigma `$ level is 0.17; this suggests that none of our detections is likely to be a background fluctuation. We also estimated the probability that our X-ray quasar detections are merely of unrelated X-ray sources that happen to be coincident with the $`z>4`$ quasars’ optical positions. We shifted the 27 quasars’ positions by eight arcmin in eight different directions and looked again for X-ray detections (in the same manner as was done for the real optical positions). At the new positions we found seven X-ray sources which met our detection criteria regarding significant level and positional coincidence. This test shows that among our nine X-ray detections there might be one source which is not the counterpart of the quasar but an X-ray source which happened to be in that position by chance. Additional suggestive evidence that the number of false detections is small is that we have detected only the brighter quasars out of the total 27 (see Fig. 2 and § 3); if our detections had random contamination by foreground sources then their distribution would not be correlated with quasar luminosity.
For comparison purposes we present some properties of the X-ray detected $`z>4`$ quasars known prior to this work in Table 1. We list the AB<sub>1450(1+z)</sub> magnitudes in column (8) and the Galactic column densities in column (9). Based on the AB<sub>1450(1+z)</sub> magnitudes we calculate the rest-frame 2500 Å flux densities and luminosities which are listed in columns (10) and (11), and the absolute $`B`$ magnitudes in column (12). We used the X-ray fluxes and bands from columns (5) and (6) and the PIMMS software to define a simple power law with $`\mathrm{\Gamma }=2`$. This power law was used to estimate the rest-frame 2-keV flux den-
sity and luminosity which are listed in columns (13) and (14). We list the resulting $`\alpha _{\mathrm{ox}}`$ in column (15).
## 3. Discussion
In most cases where we detect X-ray sources at the quasars’ optical positions they are very close to the detection limit. Comparing the data for the quasars which are X-ray detected to those which are not, we notice several trends. The detected quasars are among the brighter in the optical band (see Fig. 2 for the $`M_B`$ distribution). If the undetected quasars have comparable X-ray luminosities to the detected ones, we suggest that they were mainly not detected since they were observed at the edges of the PSPC field where the PSF and vignetting are larger. In the two cases where high luminosity quasars were observed close to the PSPC field center (BRI 1050$``$0000, BR 1144$``$0723), we attribute the non-detections to short exposure times.
Three objects which we detected had already been reported as X-ray emitters in the past (Q 0000$``$2619, GB 1428+4217, and RX J1759.4+6638). The X-ray properties we have measured for them are in agreement with the previous reported properties (see Table 1). We have not detected the other three quasars in Table 1 since the data for RX J1028.6$``$0844 are not public (it was detected in the ROSAT All-Sky Survey), GB 1508+5714 was not observed by ROSAT, and RX J105225.9+571905 could be detected only in the “ultradeep” HRI survey (see Schneider et al. 1998 and references therein).
At present only a few radio-loud $`z>4`$ quasars are known. These include GB 1428+4217, GB 1508+5714, and RX J1028.6$``$0844 which have had their X-ray data published prior to this work (Table 1), PKS 1251$``$407, for which we present an X-ray upper limit, and GB 1713+2148, whose ROSAT observation is not yet public. Out of these objects, GB 1428+4217 and GB 1508+5714 show evidence for relativistic beaming (e.g., Moran & Helfand 1997; Fabian et al. 1999). We used the NRAO<sup>3</sup><sup>3</sup>3The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. VLA Sky Survey (NVSS; Condon et al. 1998) catalog and images at 1.4 GHz to estimate the quasars’ radio loudnesses, $`R`$, defined as the ratio of the radio flux (extra-
cted from the NVSS) to the optical flux at 4400 Å (estimated using the AB magnitude and a flux-density power law with $`\alpha _\mathrm{o}=0.5`$). Radio-loud quasars typically have $`R>100`$ and radio-quiet quasars have $`R<10`$ (e.g., Kellermann et al. 1989). We list $`R`$ in Table 1 column (16) and Table 2 column (19). Most of the quasars in this study are undetected by the NVSS, and thus we provide only upper limits<sup>4</sup><sup>4</sup>4The results reported here from the NVSS are consistent with those from the Faint Images of the Radio Sky at Twenty-cm (FIRST; Becker, White, & Helfand 1995) which currently covers only a third of our objects.. In addition to the above mentioned radio-loud quasars we detect PC 0027+0525 to be radio-loud and BRI 1050-0000, RX J1759.4+6638, and PC 2331+0216 to be intermediate between the two radio classes (see also McMahon et al. 1994). All other quasars (but two) have $`R`$ upper limits which designate them as being radio-quiet. This result is in agreement with other radio and far-IR studies of these $`z>4`$ quasars which find them to be radio-quiet (e.g., Schneider et al. 1992; Omont et al. 1996; McMahon et al. 1999). Our results are also consistent with the Schmidt et al. (1995) conclusion that only 5–10% of optically selected $`z>4`$ quasars are radio-loud.
In Fig. 3 we compare the $`\alpha _{\mathrm{ox}}`$ distribution of all the X-ray detected $`z>4`$ quasars with the $`\alpha _{\mathrm{ox}}`$ distribution of a sample of all 87 Palomar-Green (PG) quasars at $`z<0.5`$ from Brandt, Laor, & Wills (2000). We have translated the $`\alpha _{\mathrm{ox}}`$ given in Brandt et al. (2000) for a flux density at 3000 Å to that for a flux density at 2500 Å. To carry out this comparison we consider only the optically selected $`z>4`$ quasars<sup>5</sup><sup>5</sup>5The only objects which are not optically selected in our study are the five last objects listed in Table 1. since the PG sample is an optically selected sample. The $`\alpha _{\mathrm{ox}}`$ distribution for the seven optically-selected objects is consistent with the $`\alpha _{\mathrm{ox}}`$ distribution of the PG quasars and with the $`\alpha _{\mathrm{ox}}`$ distribution usually found for quasars (e.g., Wilkes et al. 1994; Green et al. 1995). We also use all optically selected objects in this study, including the ones with upper limits on their X-ray properties, to derive the mean $`\alpha _{\mathrm{ox}}`$. To that purpose we have used the ASURV software package Rev 1.2 (LaValley, Isobe & Feigelson 1992), which implements the survival analysis methods presented in Feigelson & Nelson (1985) and Isobe, Feigelson, & Nelson (1986).
We find the mean $`\alpha _{\mathrm{ox}}`$ to be $`1.49\pm 0.04`$ which is in agreement with past studies.
Three of the detected quasars have $`\alpha _{\mathrm{ox}}<1`$ (Fig. 3). This is not consistent with the PG quasars’ $`\alpha _{\mathrm{ox}}`$ distribution, and indeed these three quasars are the ones known to be radio-loud with two of them being blazar-type objects, while the PG quasars are mainly radio-quiet with no blazars among them.
The separation between the blazar-type and radio-quiet quasars can also be seen in Fig. 3, where we plot the rest-frame 2 keV versus 2500 Å luminosities. The X-ray fluxes for the three radio-loud quasars are about an order of magnitude higher than those of the radio-quiet ones. A line fit to the radio-quiet quasars’ data using ASURV yields
$$\frac{\nu L_\nu (2\mathrm{keV})}{10^{45}\mathrm{ergs}\mathrm{s}^1}=\left(0.175_{0.050}^{+0.070}\right)\left[\frac{\nu L_\nu (2500\text{Å})}{10^{45}\mathrm{ergs}\mathrm{s}^1}\right]^{0.69\pm 0.10}.$$
(3)
This relation is in agreement with the one found for lower-redshift quasars (e.g., Wilkes et al. 1994; Green et al. 1995), albeit within relatively large uncertainty owing to the small number of points and their distribution.
Three of our detected quasars have two ROSAT observations (Table 2). In all cases the fluxes from different observations agree to within the measurement uncertainty of $``$30%, and no flux variations over time are detected. The quasar GB 1428+4217 has six ROSAT observations and was found to vary by a factor of two over a timescale of two weeks (or less), which corresponds to less than 2.5 days in the source’s rest-frame, a result not unexpected in that this quasar is a flat-spectrum radio-loud blazar (Fabian et al. 1998, 1999).
The observed 0.1–2 keV band corresponds to a rest-frame band of 0.5–10 keV at $`z=4`$. Emission at the low end of this bandpass can originate from an accretion disk, but is mainly thought to arise from the surrounding corona (e.g., Fabian 1994). Detecting X-ray emission in this band from $`z>4`$ quasars suggests that similar processes are taking place in low-redshift and the highest-redshift quasars.
We have established that in almost all cases where $`z>4`$ quasars were observed at the center of the PSPC field and the exposure times were sufficiently long (as anticipated from their optical luminosities), an X-ray source was found at the optical position of the quasar. As two, and possibly three, of the objects detected before our paper are “peculiar” blazar-type objects, we have more than doubled the number of optically selected $`z>4`$ quasars detected in the X-ray band and determined that $`\alpha _{\mathrm{ox}}`$ in these quasars is similar to that of lower $`z`$ quasars. At present we have only been able to study in X-rays the most luminous $`z>4`$ quasars. New X-ray missions such as Chandra, XMM, Constellation-X, and XEUS should allow the study of the X-ray properties of considerably less luminous $`z>4`$ quasars. With thousands of $`z>4`$ quasars expected to be found in the next few years (e.g., by the Sloan Digital Sky Survey), there should be ample targets.
We are grateful for several valuable suggestions by David J. Helfand. We acknowledge the support of NASA LTSA grant NAG5-8107 (SK, WNB), the Alfred P. Sloan Foundation (WNB), and NSF grant 99-00703 (DPS). We thank Harald Ebeling for the use of his IDL software. |
warning/0001/hep-th0001111.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recently it has been realized that the noncommutative geometry has played a profound role in a specific compactification of the Matrix theory and also in superstring theory via D-branes with constant $`B`$ fields -. From deeper investigation of the noncommutative gauge theory via D-branes, Seiberg and Witten have proposed that the noncommutative gauge theories realized as effective theories on D-branes are equivalent to some ordinary gauge theories . In a single D-brane case, it has been known that the effective action on the brane is Dirac-Born-Infeld action if all derivative terms are neglected -. Thus the Dirac-Born-Infeld action should be consistent with the equivalence in this approximation. Indeed this has been shown in .
It is a very natural question whether the equivalence indeed holds beyond the approximation of neglecting all derivative terms, or not. If it holds without the approximation, the forms of derivative corrections have to be highly restricted. Moreover when this equivalence is strong enough, we can determine the effective action completely from it with the help of other requirements and study the dynamics of the D-branes using the action.
In this paper, we show that in the approximation of neglecting the fourth and higher order derivative terms the D-brane action computed in the superstring theory is consistent with the equivalence. Although, this is the Dirac-Born-Infeld action without two-derivative corrections, to show the equivalence we should take into account the orderings of the noncommutative field strength in the Dirac-Born-Infeld action. By taking appropriate ordering it becomes consistent and this is regarded as a non-trivial test of the equivalence.
With the mapping of the ordinary gauge field to noncommutative gauge field given in , we also explicitly construct general forms of the two-derivative corrections which satisfy the equivalence relation in the approximation of neglecting the four-derivative terms. Furthermore, we can construct the $`2n`$-derivative corrections which are consistent with the equivalence in the approximation of neglecting the $`(2n+2)`$-derivative terms. It should be emphasized that the results obtained in this paper are valid for arbitrary order of the field strength. In this paper, we regard $`F_{ij}𝒪(^0)`$ and $`A_i𝒪(^1)`$
On the other hand, in it has been shown that for the bosonic string case the known two-derivative correction is not consistent with the equivalence. In the consistent two-derivative corrections up to the quartic order of the field strength have been considered. As we will see later, the result obtained in this paper reproduce these corrections . This problem can be resolved by considering $`B`$-dependent field redefinition of the $`U(1)`$ gauge field . Therefore in order to constrain the effective action for the bosonic string case, we should include the $`B`$-dependent field redefinition. However two-derivative corrections allowed by the equivalence without the $`B`$-dependent field redefinition may also be allowed by the equivalence with it . Furthermore, the higher derivative corrections may capture some general structures of the effective action of the D-brane. Therefore the derivative corrections obtained in this paper are probably important.
This paper is organized as follows. In section 2, we briefly review the equivalence between noncommutative and ordinary gauge theories shown in . In section 3 it is shown that the certain noncommutative version of the DBI actions without two-derivative corrections are consistent with the equivalence in the approximation of neglecting the four-derivative terms. We also construct the consistent two-derivative corrections in this approximation. In section 4 we argue that the two-derivative corrections obtained in section 3 exhaust the consistent two-derivative corrections and also generalize these to the $`2n`$-derivative corrections. Finally section 5 is devoted to conclusion.
## 2 Noncommutative Gauge Theory
In this section we briefly review the equivalence between noncommutative and ordinary gauge theories shown in . We consider open strings in flat space, with metric $`g_{ij}`$, in the presence of a constant $`B_{ij}`$ and with a Dp-brane. Here we assume that $`B_{ij}`$ has rank $`p+1`$ and $`B_{ij}0`$ only for $`i,j=1,\mathrm{},p+1`$. The world-sheet action is
$$S=\frac{1}{4\pi \alpha ^{}}_\mathrm{\Sigma }g_{ij}_ax^i^ax^j\frac{i}{2}_\mathrm{\Sigma }B_{ij}x^i_\tau x^ji_\mathrm{\Sigma }A_i(x)_\tau x^i,$$
(2.1)
where $`\mathrm{\Sigma }`$ is the string world-sheet, $`_\tau `$ is the tangential derivative along the world-sheet boundary $`\mathrm{\Sigma }`$ and $`A_i`$ is a background gauge field. In the case that $`\mathrm{\Sigma }`$ is the upper half plane parameterized by $`\mathrm{}\tau \mathrm{}`$ and $`0\sigma \mathrm{}`$, the propagator evaluated at boundary points is -
$$x^i(\tau )x^j(\tau ^{})=\alpha ^{}(G^1)^{ij}\mathrm{log}(\tau \tau ^{})^2+\frac{i}{2}\theta ^{ij}ϵ(\tau \tau ^{}),$$
(2.2)
where $`G`$ and $`\theta `$ are the symmetric and antisymmetric tensors defined by
$$(G^1)^{ij}+\frac{1}{2\pi \alpha ^{}}\theta ^{ij}=\left(\frac{1}{g+2\pi \alpha ^{}B}\right)^{ij}.$$
(2.3)
From considerations of the string S-matrix, the $`B`$ dependence of the effective action for fixed $`G`$ can be obtained by replacing ordinary multiplication in the effective action for $`B=0`$ by the $``$ product defined by the formula
$$f(x)g(x)=e^{\frac{i}{2}\theta ^{ij}\frac{}{\xi ^i}\frac{}{\zeta ^j}}f(x+\xi )g(x+\zeta )|_{\xi =\zeta =0}.$$
(2.4)
It is likely that the gauge transformation also becomes noncommutative. In fact, using the point splitting regularization, $`S`$ is invariant under noncommutative gauge transformation
$$\widehat{\delta }\widehat{A}_i=\widehat{D}_i\lambda ,$$
(2.5)
where covariant derivative $`\widehat{D}_i`$ is defined as
$$\widehat{D}_iE(x)=_iE(x)+i\left(E(x)\widehat{A}_i\widehat{A}_iE(x)\right).$$
(2.6)
Conversely, using Pauli-Villars regularization, $`S`$ is invariant under ordinary gauge transformation
$$\delta A_i=_i\lambda .$$
(2.7)
Therefore, the effective Lagrangian obtained in this way becomes ordinary gauge theory. Thus this theory and the corresponding noncommutative gauge theory are equivalent under the field redefinition $`\widehat{A}=\widehat{A}(A)`$ since the coupling constants in the world-sheet theory are the spacetime fields. Because the two different gauge invariance should satisfy $`\widehat{A}(A)+\widehat{\delta }_{\widehat{\lambda }}\widehat{A}(A)=\widehat{A}(A+\delta _\lambda A)`$, the mapping of $`A`$ to $`\widehat{A}`$ for $`U(1)`$ case is obtained as a differential equation for $`\theta `$,
$`\delta \widehat{A}_i(\theta )=\delta \theta ^{kl}{\displaystyle \frac{}{\theta ^{kl}}}\widehat{A}_i(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{kl}[\widehat{A}_k(_l\widehat{A}_i+\widehat{F}_{li})+(_l\widehat{A}_i+\widehat{F}_{li})\widehat{A}_k]`$
$`\delta \widehat{F}_{ij}(\theta )=\delta \theta ^{kl}{\displaystyle \frac{}{\theta ^{kl}}}\widehat{F}_{ij}(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{kl}[2\widehat{F}_{ik}\widehat{F}_{jl}+2\widehat{F}_{jl}\widehat{F}_{ik}`$ (2.8)
$`\widehat{A}_k(\widehat{D}_l\widehat{F}_{ij}+_l\widehat{F}_{ij})(\widehat{D}_l\widehat{F}_{ij}+_l\widehat{F}_{ij})\widehat{A}_k],`$
where
$$\widehat{F}_{ij}=_i\widehat{A}_j_j\widehat{A}_ii\widehat{A}_i\widehat{A}_j+i\widehat{A}_j\widehat{A}_i.$$
(2.9)
In this map has been derived in a path integral form from D-brane world-volume perspective .<sup>§</sup><sup>§</sup>§ Although the differential equation has ambiguities , these ambiguities have no physical meaning because they correspond to the field redefinition.
In the approximation of neglecting the derivative terms, the effective Lagrangian is the Dirac-Born-Infeld Lagrangian
$$_{DBI}=\frac{1}{g_s(2\pi )^p(\alpha ^{})^{\frac{p+1}{2}}}\sqrt{det(g+2\pi \alpha ^{}(B+F))},$$
(2.10)
where $`F_{ij}=_iA_j_jA_i`$. Here $`g_s`$ is the closed string coupling and the normalization of the Lagrangian is same as the one taken in . Therefore the equivalent noncommutative gauge theory in the approximation has the following Lagrangian
$$\widehat{}_{DBI}=\frac{1}{G_s(2\pi )^p(\alpha ^{})^{\frac{p+1}{2}}}\sqrt{det(G+2\pi \alpha ^{}\widehat{F})}.$$
(2.11)
Note that all the multiplication entering the r.h.s of (2.11) can be regarded as the ordinary multiplication except those in the definition of $`\widehat{F}`$ because of the approximation. From the requirement $`_{DBI}=\widehat{}_{DBI}`$ for $`F=0`$, the overall normalization $`G_s`$ should be fixed as $`G_s=g_s\sqrt{det(G)/det(g+2\pi \alpha ^{}B)}`$.
Furthermore, in it has been proposed that the effective action can be written for arbitrary values of $`\theta `$. More precisely for given physical parameters $`g_s,g_{ij}`$ and $`B_{ij}`$ and an auxiliary parameter $`\theta `$, we define $`G_s,G_{ij}`$ and a two form $`\mathrm{\Phi }_{ij}`$ as
$`\left({\displaystyle \frac{1}{G+2\pi \alpha ^{}\mathrm{\Phi }}}\right)^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \alpha ^{}}}\theta ^{ij}+\left({\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}\right)^{ij}`$
$`G_s`$ $`=`$ $`g_s\left(det\left({\displaystyle \frac{1}{2\pi \alpha ^{}}}\theta +{\displaystyle \frac{1}{g+2\pi \alpha ^{}B}}\right)det(g+2\pi \alpha ^{}B)\right)^{\frac{1}{2}}.`$ (2.12)
Then the effective action $`\widehat{S}(G_s,G,\mathrm{\Phi },\theta ;\widehat{F})`$, in which the multiplication is the $`\theta `$-dependent $``$ product, is actually $`\theta `$-independent, i.e. $`\widehat{S}(G_s,G,\mathrm{\Phi },\theta ;\widehat{F})=S(g_s,g,B,\theta =0;F)`$. The effective action including $`\mathrm{\Phi }`$ may be obtained using a regularization which interpolates between Pauli-Villars and point splitting as in . In this paper, we simply assume this proposal.
In the approximation of neglecting the derivative of $`F`$, the equation
$$\delta _\mathrm{\Phi }=\delta \theta ^{kl}\frac{_\mathrm{\Phi }}{\theta ^{kl}}|_{g_s,g,B,A_ifixed}=\mathrm{total}\mathrm{derivative}+𝒪(^2),$$
(2.13)
should hold. Hereafter $`\delta `$ always denotes $`\delta \theta ^{kl}\frac{}{\theta ^{kl}}`$. Here $`_\mathrm{\Phi }`$ is the Lagrangian defined as
$$_\mathrm{\Phi }=\frac{1}{G_s(2\pi )^p(\alpha ^{})^{\frac{p+1}{2}}}\sqrt{det(G+2\pi \alpha ^{}(\widehat{F}+\mathrm{\Phi }))},$$
(2.14)
where the multiplication is the $``$ product except in the definition of $`\widehat{F}`$. Below for simplicity we set $`2\pi \alpha ^{}=1`$. The variation of $`G_s,G`$ and $`\mathrm{\Phi }`$ are
$`\delta G_s`$ $`=`$ $`{\displaystyle \frac{1}{2}}G_s\mathrm{Tr}(\mathrm{\Phi }\delta \theta ),`$
$`\delta G`$ $`=`$ $`G\delta \theta \mathrm{\Phi }+\mathrm{\Phi }\delta \theta G,`$
$`\delta \mathrm{\Phi }`$ $`=`$ $`\mathrm{\Phi }\delta \theta \mathrm{\Phi }+G\delta \theta G,`$ (2.15)
and the variation of $`\widehat{F}`$ is
$`\delta \widehat{F}_{ij}`$ $`=`$ $`(\widehat{F}\delta \theta \widehat{F})_{ij}\widehat{A}_k\delta \theta ^{kl}{\displaystyle \frac{1}{2}}(_l+\widehat{D}_l)\widehat{F}_{ij}+𝒪(^4)`$ (2.16)
$`=`$ $`(\widehat{F}\delta \theta \widehat{F})_{ij}\widehat{A}_k\delta \theta ^{kl}(_l{\displaystyle \frac{1}{2}}\theta ^{mn}_n\widehat{A}_l_m)\widehat{F}_{ij}+𝒪(^4).`$
Following , we get
$`\delta \left({\displaystyle \frac{1}{G_s}}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\right)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{G_s}}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\left(\mathrm{Tr}(\widehat{F}\delta \theta )+\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{ji}\widehat{A}_k\delta \theta ^{kl}{\displaystyle \frac{1}{2}}(_l+\widehat{D}_l)\widehat{F}_{ij}\right),`$ (2.17)
where the multiplication is the ordinary one except in $`\widehat{F}`$ and $`\widehat{D}_l`$. Now using
$$\frac{1}{2}(_l+\widehat{D}_l)\widehat{A}_k\frac{1}{2}(_k+\widehat{D}_k)\widehat{A}_l=\widehat{D}_l\widehat{A}_k_k\widehat{A}_l=\widehat{F}_{lk},$$
(2.18)
we see that
$`\delta \theta ^{kl}(_l+\widehat{D}_l)\left(\widehat{A}_kdet(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\right)`$
$`=\delta \theta ^{kl}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(\widehat{F}_{lk}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{ji}\widehat{A}_k(_l+\widehat{D}_l)\widehat{F}_{ij})+𝒪(^4),`$ (2.19)
is a total derivative. Thus we obtain the desired result
$$\delta \left(\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\right)=\mathrm{total}\mathrm{derivative}+𝒪(^4).$$
(2.20)
Note that the computation above shows that the Dirac-Born-Infeld Lagrangian (2.14) is $`\theta `$-independent even in the approximation of neglecting $`𝒪(^4)`$ terms.
## 3 Two-derivative terms
In this section, we will see that certain noncommutative Dirac-Born-Infeld actions are consistent with the equivalence in the approximation of neglecting four-derivative terms. We also give certain two-derivative corrections consistent with the equivalence in the approximation of neglecting $`𝒪(^4)`$ terms.
Because the multiplication in the Dirac-Born-Infeld Lagrangian $`_\mathrm{\Phi }`$ should be replaced by the $``$ product in the approximation, we first consider the ordering of the $`\widehat{F}_{ij}`$ in the Dirac-Born-Infeld Lagrangian in which the multiplication is the $``$ product. This Lagrangian is relevant in the approximation $`𝒪(^4)`$ and denoted as $`\widehat{}_\mathrm{\Phi }`$.
It seems that there are two natural ways of ordering. The first one is symmetrization of the $`(\widehat{F}+\mathrm{\Phi })_{ij}`$ in $`\widehat{}_\mathrm{\Phi }`$, as in the non-Abelian Dirac-Born-Infeld Lagrangian considered in . The another one is as follows. First we expand the square root of determinant using $`U_{2n}\mathrm{Tr}(G^1(\widehat{F}+\mathrm{\Phi }))^{2n}`$. Next keeping the order of $`(\widehat{F}+\mathrm{\Phi })_{ij}`$ in $`U_{2n}`$ indicated by the symbol $`\mathrm{Tr}`$, we symmetrize the polynomials of $`U_{2n}`$’s. Then replacing all the multiplication by $``$ product, we obtain the $`\widehat{}_\mathrm{\Phi }`$ from the $`_\mathrm{\Phi }`$.
To take the either one of these, we will show
$$\widehat{}_\mathrm{\Phi }=_\mathrm{\Phi }+𝒪(^4).$$
(3.1)
To show this, we first remember the fact $`fg+gf=2fg+𝒪(^4)`$. Thus taking the first way of ordering, we easily see that (3.1) is satisfied since $`(\widehat{F}+\mathrm{\Phi })_{ij}`$ is symmetrized. If we take the second way, using
$`\mathrm{Tr}\left[(G^1(\widehat{F}+\mathrm{\Phi }))\mathrm{}(G^1(\widehat{F}+\mathrm{\Phi }))\right]`$ (3.2)
$`=`$ $`\mathrm{Tr}\left[(G^1(\widehat{F}+\mathrm{\Phi }))^{2n}\right]`$
$`+i{\displaystyle \underset{m=0}{\overset{2n2}{}}}(m+1)\theta ^{kl}\mathrm{Tr}\left[(G^1_k\widehat{F})(G^1(\widehat{F}+\mathrm{\Phi }))^{2nm2}(G^1_l\widehat{F})(G^1(\widehat{F}+\mathrm{\Phi }))^m\right]+𝒪(^4)`$
$`=`$ $`\mathrm{Tr}\left[(G^1(\widehat{F}+\mathrm{\Phi }))^{2n}\right]+𝒪(^4),`$
we can also show (3.1). Note that if we take the other way of ordering, (3.1) is not necessarily satisfied.
From (2.20) and (3.1), we conclude that the noncommutative Dirac-Born-Infeld Lagrangian with one of these orderings, $`\widehat{}_\mathrm{\Phi }`$, satisfies the desired equation
$$\delta \widehat{}_\mathrm{\Phi }=\mathrm{total}\mathrm{derivative}+𝒪(^4).$$
(3.3)
Therefore this Lagrangian without two derivative corrections is allowed by the equivalence in the approximation of neglecting $`𝒪(^4)`$, but keeping an arbitrary order of $`\widehat{F}`$. This result is consistent with the calculations of the effective action for the superstring case because in this case there are no two-derivative terms in the effective action.
For the bosonic case, it has been known that the known two-derivative corrections derived from the string four-point amplitude and the $`\beta `$ function in the open string $`\sigma `$ model is not consistent with the equivalence with (2.8). However it can be shown that if the mapping of $`A`$ to $`\widehat{A}`$ (2.8) is modified by some field redefinition containing $`\theta ,F`$ and $`\mathrm{\Phi }`$, the equivalence is consistent with the result in and .
Although this modification should be applied for the bosonic case, we will consider the two-derivative corrections which are consistent with the equivalence using (2.8) in the rest of this section. This is because these corrections can be added consistently in the approximation even if we modify (2.8) and may capture some general structures of the effective action of the D-brane. We will also consider the $`2m`$-derivative corrections which are consistent with the equivalence using (2.8) in the approximation of neglecting $`(2m+2)`$-derivative terms in the next section. For the superstring case, it is possible that (2.8) is valid even if we do not neglect the higher-derivative terms. Hence it is important that the determination of these corrections.
Then we will show below that the two-derivative term
$$\widehat{}_2=\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}L_2,$$
(3.4)
satisfies
$$\delta \widehat{}_2=\mathrm{total}\mathrm{derivative}+𝒪(^4),$$
(3.5)
where
$`L_2`$ $`=`$ $`\left\{a_1(h_S)^{mn}(h_S)^{iq}(h_S)^{jp}+a_2(h_S)^{mi}(h_S)^{nq}(h_S)^{jp}\right\}\widehat{D}_m\widehat{F}_{ij}\widehat{D}_n\widehat{F}_{pq},`$
$`(h_S)^{ij}`$ $`=`$ $`\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{\mathrm{sym}}^{ij}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)^{ij}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{G\widehat{F}\mathrm{\Phi }}}\right)^{ij}`$ (3.6)
$`=`$ $`\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}G{\displaystyle \frac{1}{G\widehat{F}\mathrm{\Phi }}}\right)^{ij},`$
and $`a_1`$ and $`a_2`$ are some constants. Here it is not required to consider the ordering problem of (3.4) because the condition (3.5) means that the term is consistent with the equivalence in the approximation of neglecting $`𝒪(^4)`$. How these terms are derived is explained in the next section.
If we define a differential operator
$$\stackrel{~}{\delta }=\frac{1}{2}\delta \theta ^{kl}\widehat{A}_k(_l+\widehat{D}_l),$$
(3.7)
the differential of $`\widehat{}_2`$ is written as
$$\delta \widehat{}_2=\frac{1}{2}\delta \theta ^{kl}(_l+\widehat{D}_l)\left(\widehat{A}_k\widehat{}_2\right)+\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(\delta +\stackrel{~}{\delta })L_2,$$
(3.8)
where the first term of the r.h.s of (3.8) is a total derivative. Hence we consider the variations of $`h_S`$ and $`\widehat{D}\widehat{F}`$ under $`\delta +\stackrel{~}{\delta }`$.
From
$$(\delta +\stackrel{~}{\delta })\widehat{F}_{ij}=(\widehat{F}\delta \theta \widehat{F})_{ij}+𝒪(^4),$$
(3.9)
it obeys
$`(\delta +\stackrel{~}{\delta })h_{S}^{}{}_{}{}^{ij}`$ $`=`$ $`\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\left((G+\mathrm{\Phi })\delta \theta (G+\mathrm{\Phi })+(\delta +\stackrel{~}{\delta })\widehat{F}\right){\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{\mathrm{sym}}^{ij}`$ (3.10)
$`=`$ $`\left(\delta \theta {\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\widehat{F}\delta \theta \delta \theta \widehat{F}{\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{\mathrm{sym}}^{ij}+𝒪(^4)`$
$`=`$ $`\left(h_S(\widehat{F}\delta \theta )+(\delta \theta \widehat{F})h_S\right)^{ij}+𝒪(^4).`$
Remembering that $`[\delta ,_m]=0`$ and that $`\widehat{D}`$ explicitly depends on $`\theta `$ through $``$ product, we see that the commutation relation between $`\widehat{D}`$ and $`\delta `$ is
$`[\delta ,\widehat{D}_m]E`$ $`=`$ $`[\delta ,_m]E+i[E,\delta \widehat{A}_m]+\delta \theta ^{kl}_k\widehat{A}_m_lE+𝒪(^5E)`$ (3.11)
$`=`$ $`i[E,{\displaystyle \frac{1}{2}}\delta \theta ^{kl}\widehat{A}_k(_l\widehat{A}_m+\widehat{F}_{lm})]+\delta \theta ^{kl}_k\widehat{A}_m_lE+𝒪(^5E),`$
where $`E`$ is an arbitrary function. The computation of the commutation relation between $`\widehat{D}`$ and $`\stackrel{~}{\delta }`$ can be carried out straightforwardly
$`[\stackrel{~}{\delta },\widehat{D}_m]E`$ $`=`$ $`\stackrel{~}{\delta }\left(_mE+i[E,\widehat{A}_m]\right)\widehat{D}_m\left({\displaystyle \frac{1}{2}}\widehat{A}_k\delta \theta ^{kl}(_l+\widehat{D}_l)E\right)`$ (3.12)
$`=`$ $`{\displaystyle \frac{1}{2}}\delta \theta ^{kl}\left(\widehat{D}_m\widehat{A}_k(_l+\widehat{D}_l)E+i\widehat{A}_k[E,\widehat{F}_{ml}_l\widehat{A}_m]\right)+𝒪(^5E).`$
After some calculations using (3.11) and (3.12), we can find a simple result
$$[\delta +\stackrel{~}{\delta },\widehat{D}_m]E=(\widehat{F}\delta \theta )_m^l(\widehat{D}_lE)+𝒪(^5E),$$
(3.13)
and then we obtain
$$(\delta +\stackrel{~}{\delta })\widehat{D}_m\widehat{F}_{ij}=\left((\widehat{D}_m\widehat{F})\delta \theta \widehat{F}\right)_{ij}\left(\widehat{F}\delta \theta (\widehat{D}_m\widehat{F})\right)_{ij}(\widehat{F}\delta \theta )_m^l(\widehat{D}_l\widehat{F}_{ij})+𝒪(^5).$$
(3.14)
This and (3.10) imply that $`(\delta +\stackrel{~}{\delta })L_2=0`$. Thus we conclude that $`\delta \widehat{}_2=\mathrm{total}\mathrm{derivative}+𝒪(^6)`$. However $`𝒪(^4)`$ terms may exist if the ordinary multiplication in (3.4) is replaced by $``$ product. Thus only (3.5) is meaningful because we have not considered the ordering problems of (3.4).
Now we discuss the expansion about $`F`$ of two-derivative corrections of the effective Lagrangian (3.4) with $`B=\mathrm{\Phi }=\theta =0`$ and $`g_{ij}=\delta _{ij}`$. In this commutative description, the Dirac-Born-Infeld Lagrangian (2.11) becomes $`1/(g_s(2\pi )^{\frac{p+1}{2}})\sqrt{det(1+F)}`$. Using $`det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}=1\frac{1}{4}\mathrm{Tr}F^2+𝒪(F^4)`$ and $`h_S=\frac{1}{1F^2}=1+F^2+𝒪(F^4)`$, we see
$`\widehat{}_2`$ $`=`$ $`a_1{\displaystyle \frac{1}{g_s}}\left[\left(1{\displaystyle \frac{1}{4}}\mathrm{Tr}F^2\right)_mF_{ij}_mF_{ji}+(F^2)_{mn}_mF_{ij}_mF_{ji}+2(F^2)_{ik}_mF_{ij}_mF_{jk}\right]`$ (3.15)
$`+a_2{\displaystyle \frac{1}{g_s}}\left[\left(1{\displaystyle \frac{1}{4}}\mathrm{Tr}F^2\right)_mF_{im}_nF_{ni}+2(F^2)_{mj}_mF_{ij}_nF_{ni}+(F^2)_{iq}_mF_{im}_nF_{nq}\right]`$
$`+𝒪(F^4FF).`$
Here following , we define a basis of terms of order $`F^2FF`$ as
$`J_1=F_{kl}F_{lk}_nF_{ij}_nF_{ji},J_2=F_{kl}F_{li}_nF_{ij}_nF_{jk},`$
$`J_3=F_{ni}F_{im}_nF_{kl}_mF_{lk},J_4=F_{kl}F_{lk}_nF_{ni}_mF_{im},`$
$`J_5=F_{jk}F_{km}_nF_{ni}_mF_{ij},J_6=F_{kl}F_{lk}_m_mF_{ij}F_{ji},`$
$`J_7=_m_mF_{ij}F_{jk}F_{kl}F_{li}.`$ (3.16)
Therefore after some computations, we obtain
$$\widehat{}_2=a_1\frac{1}{g_s}\left(_mF_{ij}_mF_{ji}\frac{1}{4}J_1+2J_2+J_3\right)+a_2\frac{1}{g_s}\left(\frac{1}{2}_mF_{ij}_mF_{ji}J_5+\frac{1}{8}J_6\frac{1}{2}J_7\right)+\mathrm{},$$
(3.17)
where the ellipsis represent total derivative terms and $`𝒪(F^4FF)`$ terms. This is same as the effective Lagrangian obtained in . Thus the result obtained in this paper implies that the one obtained in is consistent even in the arbitrary order of $`F`$. It is noted that in only the equivalence without assuming the existence of $`\mathrm{\Phi }`$ is required.
## 4 General forms of the derivative corrections
In this section, we discuss the other two-derivative corrections which satisfy (3.5) and also the higher derivative corrections.
Requiring the noncommutative gauge invariance and the gauge invariance for the $`B`$ field, the most general two derivative terms are
$$\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\left[T^{ijklmn}(G^1,\widehat{F}+\mathrm{\Phi })\widehat{D}_m\widehat{F}_{ij}\widehat{D}_n\widehat{F}_{kl}+T^{ijmn}(G^1,\widehat{F}+\mathrm{\Phi })\widehat{D}_m\widehat{D}_n\widehat{F}_{ij}\right],$$
(4.1)
where $`T^{ijklmn}`$ and $`T^{ijmn}`$ are arbitrary tensors constructed from $`(G^1)^{pq}`$, $`M_{pq}(\widehat{F}+\mathrm{\Phi })_{pq}`$, $`detG`$ and $`detM`$. The tensors $`T`$ should not depend on $`(M^1)^{pq}`$ because it has a singularity at $`M=0`$.
We can easily generalize the invariance problem under $`\delta `$ of two-derivative terms to the one of higher derivative terms. Thus below we will look for the Lagrangian $`\widehat{}_m`$ with $`m`$-derivative which satisfies the invariant condition
$$\delta \widehat{}_m=\mathrm{total}\mathrm{derivative}+𝒪(^{m+2}),$$
(4.2)
which is the condition consistent with the equivalence in the approximation of neglecting $`𝒪(^{m+2})`$. To do this, let us consider a $`(2n,0)`$ tensor $`T^{p_1p_2\mathrm{}p_{2n}}`$ constructed from $`G^1`$ and $`M`$. It can be written as
$$T^{p_1p_2\mathrm{}p_{2n}}=\underset{k_1=0}{\overset{\mathrm{}}{}}\mathrm{}\underset{k_n=0}{\overset{\mathrm{}}{}}C_{\{k\}}(M^{k_1})^{p_1p_2}(M^{k_2})^{p_3p_4}\mathrm{}(M^{k_n})^{p_{2n1}p_{2n}},$$
(4.3)
where $`C_{\{k\}}`$ is some function of the scalars constructed from $`M`$ and $`G`$ and $`(M^{k_i})^{p_{2i1}p_{2i}}`$ means $`((G^1M)^{k_i}G^1)^{p_{2i1}p_{2i}}`$. We also consider a $`(0,2n)`$ tensor $`\widehat{J}_{p_1p_2\mathrm{}p_{2n}}`$ such that
$$\widehat{J}_{p_1p_2\mathrm{}p_{2n}}=\left\{(\widehat{D}\mathrm{}\widehat{D}\widehat{F})\mathrm{}(\widehat{D}\mathrm{}\widehat{D}\widehat{F})\right\}_{p_1p_2\mathrm{}p_{2n}}.$$
(4.4)
For given $`n`$ and $`m`$, where $`m`$ is the number of derivative $`\widehat{D}`$ in $`\widehat{J}`$, there are finite number $`s`$ of independent $`\widehat{J}_{p_1p_2\mathrm{}p_{2n}}`$ under the identification using Bianchi identity. We note that the total divergence terms should not be used for the identification. Taking a basis of these $`\widehat{J}^{(i)}`$, where $`1is`$, we will study the invariance of
$$\widehat{}_m=\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\underset{i=1}{\overset{s}{}}\underset{p_1,\mathrm{},p_{2n}}{}(T_{(i)})^{p_1p_2\mathrm{}p_{2n}}(\widehat{J}^{(i)})_{p_1p_2\mathrm{}p_{2n}},$$
(4.5)
where $`(T_{(i)})`$ is the tensor of the form (4.3) with the coefficients $`C_{\{k\}}^{(i)}`$. As like the derivation of (3.8), we can show that
$$\delta \widehat{}_m=\mathrm{total}\mathrm{derivative}+\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(\delta +\stackrel{~}{\delta })\left(\underset{i=1}{\overset{s}{}}\underset{p_1,\mathrm{},p_{2n}}{}(T_{(i)})^{p_1\mathrm{}p_{2n}}(\widehat{J}^{(i)})_{p_1\mathrm{}p_{2n}}\right).$$
(4.6)
However there is a possibility of cancellation between the variations under $`\delta `$ of the terms with different $`n`$’s. Note that this cancellation can occur only between the variations of the terms with $`n`$ and $`n+1`$. We will consider this later.
In order to proceed further, we require the invariance of $`\widehat{}_m`$ with $`\widehat{F}=0`$ first. From (4.6), we have conditions for the invariance as
$$\underset{i=1}{\overset{s}{}}\underset{p_1,\mathrm{},p_{2n}}{}(\delta T_{(i)})^{p_1p_2\mathrm{}p_{2n}}(\widehat{J}^{(i)})_{p_1\mathrm{}p_{2n}}=0.$$
(4.7)
It can be shown that $`\delta G^1|_{G=1}=\left(\mathrm{\Phi }\delta \theta +\delta \theta \mathrm{\Phi }\right)`$ and
$`\delta \left((G^1\mathrm{\Phi })^kG^1\right)|_{G=1}`$ $`=`$ $`\left(\delta \theta \mathrm{\Phi }^{k+1}+\mathrm{\Phi }\delta \theta \mathrm{\Phi }^k+\mathrm{}+\mathrm{\Phi }^{k+1}\delta \theta \right)`$ (4.8)
$`+\left(\delta \theta \mathrm{\Phi }^{k1}+\mathrm{\Phi }\delta \theta \mathrm{\Phi }^{k2}+\mathrm{}+\mathrm{\Phi }^{k1}\delta \theta \right),`$
where $`k1`$. Here we have set $`G=1`$ after operating $`\delta `$ for notational simplicity. Therefore to satisfy (4.7), we have to take
$$(T_{(i)})^{p_1\mathrm{}p_{2n}}=C^{(i)}(h_S)^{p_1p_2}(h_S)^{p_3p_4}\mathrm{}(h_S)^{p_{2n1}p_{2n}},$$
(4.9)
where $`C^{(i)}`$ is some function of the scalars constructed from $`M`$ and $`G`$. We also see that this $`C^{(i)}`$ should be some constant since $`\delta \mathrm{Tr}((G^1\mathrm{\Phi })^{2k})|_{G=1}=2k\mathrm{Tr}(\delta \theta (\mathrm{\Phi }^{2k1}\mathrm{\Phi }^{2k+1}))`$.
Now we require the condition (4.7) without taking $`\widehat{F}=0`$. Using (3.10) and (3.14) as in the previous section, one can easily show that the condition (4.7) is satisfied for
$$\widehat{}_m=\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\underset{i=1}{\overset{s_1}{}}C^{(i)}\underset{p_1,\mathrm{},p_{2n}}{}(h_S)^{p_1p_2}(h_S)^{p_3p_4}\mathrm{}(h_S)^{p_{2n1}p_{2n}}\overline{J}_{p_1p_2\mathrm{}p_{2n}}^{(i)},$$
(4.10)
where $`\{\overline{J}_{p_1p_2\mathrm{}p_{2n}}^{(i)}\}`$ is a basis of the form
$$((\widehat{D}\widehat{F})\mathrm{}(\widehat{D}\widehat{F}))_{p_1p_2\mathrm{}p_{2n}}.$$
(4.11)
Hence this term is allowed by the equivalence in the approximation. In particular, for two-derivative terms, (4.10) is equivalent to (3.4).
The terms containing $`(\widehat{D})^N\widehat{F}`$ with $`N2`$ are other candidates, but they can not satisfy the condition (4.7) generically because there are contributions from $`(\widehat{D})^N\widehat{F}`$ which can not be canceled by the ones from $`h_S`$ for the variation of $`\widehat{}_m`$ under $`\delta `$ as seen from (3.13) with $`E=(\widehat{D})^{N1}\widehat{F}`$. However, in some cases, these are absent because of the symmetry for the indices. For example, we consider
$$\widehat{_2}(h_S)^{ip}(h_S)^{jq}[\widehat{D}_p,\widehat{D}_q]\widehat{F}_{ij}.$$
(4.12)
Remember that
$$(\delta +\stackrel{~}{\delta })(\widehat{D}_p\widehat{D}_q\widehat{F}_{ij})=(\delta \theta )^{lk}\left((\widehat{D}_p\widehat{F}_{ql})(\widehat{D}_k\widehat{F}_{ij})+(\widehat{D}_q\widehat{F}_{il})(\widehat{D}_p\widehat{F}_{kj})+(\widehat{D}_p\widehat{F}_{il})(\widehat{D}_q\widehat{F}_{kj})\right)+\mathrm{},$$
(4.13)
where the ellipsis represents the terms which are canceled by the contribution from $`h_S`$ and $`𝒪(^6)`$. Thus we obtain $`(\delta +\stackrel{~}{\delta })((h_S)^{ip}(h_S)^{jq}[\widehat{D}_p,\widehat{D}_q]\widehat{F}_{ij})=𝒪(^6)`$ from the Bianchi identity and the symmetries of the indices. Therefore (4.12) is also allowed by the equivalence in the approximation though this vanishes at $`\theta =0`$. Note that (4.12) is the only allowed term with two-derivative and one $`\widehat{F}`$ because of the Bianchi identity and the symmetry for the indices of $`\widehat{F}_{ij}`$.
There are invariant combinations of the terms with different numbers of indices which implies different numbers of $`\widehat{F}`$. To illustrate these, we consider
$$\widehat{}_4^A=\frac{C}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(h_S)^{pq}(h_S)^{tu}(h_S)^{ic}(h_S)^{jd}(\widehat{D}_p\widehat{D}_q\widehat{F}_{ij})(\widehat{D}_t\widehat{D}_u\widehat{F}_{cd}),$$
(4.14)
where $`C`$ is some constant. In this case, using (4.13) the unnecessary terms are easily computed as
$$\delta \widehat{}_4^A2\frac{C}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(h_S)^{pq}(h_S)^{tu}(h_S)^{ic}(h_S)^{jd}(\delta \theta )^{lk}(\widehat{D}_p\widehat{F}_{ql})(\widehat{D}_k\widehat{F}_{ij})(\widehat{D}_t\widehat{D}_u\widehat{F}_{cd}),$$
(4.15)
where we neglect the terms which are canceled by the contribution from $`h_S`$ and $`𝒪(^8)`$. Let us define
$`(h_A)^{ij}`$ $`=`$ $`\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{\mathrm{anti}.\mathrm{sym}}^{ij}=\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}(\widehat{F}+\mathrm{\Phi }){\displaystyle \frac{1}{G\widehat{F}\mathrm{\Phi }}}\right)^{ij},`$ (4.16)
which obeys
$`(\delta +\stackrel{~}{\delta })h_{A}^{}{}_{}{}^{ij}`$ $`=`$ $`\left({\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\left((G+\mathrm{\Phi })\delta \theta (G+\mathrm{\Phi })+(\delta +\stackrel{~}{\delta })\widehat{F}\right){\displaystyle \frac{1}{G+\widehat{F}+\mathrm{\Phi }}}\right)_{\mathrm{anti}.\mathrm{sym}}^{ij}`$ (4.17)
$`=`$ $`(\delta \theta )^{ij}+\left(h_A(\widehat{F}\delta \theta )+(\delta \theta \widehat{F})h_A\right)^{ij}+𝒪(^4).`$
Then it can be seen that
$$\delta (\widehat{}_4^A+\widehat{}_4^B+\widehat{}_4^C)0,$$
(4.18)
where
$`\widehat{}_4^B`$ $`=`$ $`2{\displaystyle \frac{C}{G_s}}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(h_S)^{pq}(h_S)^{tu}(h_S)^{ic}(h_S)^{jd}(h_A)^{lk}(\widehat{D}_p\widehat{F}_{ql})(\widehat{D}_k\widehat{F}_{ij})(\widehat{D}_t\widehat{D}_u\widehat{F}_{cd}),`$
$`\widehat{}_4^C`$ $`=`$ $`{\displaystyle \frac{C}{G_s}}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}(h_S)^{pq}(h_S)^{tu}(h_S)^{ic}(h_S)^{jd}(h_A)^{lk}(h_A)^{ef}`$ (4.19)
$`\times (\widehat{D}_p\widehat{F}_{ql})(\widehat{D}_k\widehat{F}_{ij})(\widehat{D}_t\widehat{F}_{ue})(\widehat{D}_f\widehat{F}_{cd}).`$
As this example, for general terms of the form
$$\widehat{}_m=\frac{1}{G_s}det(G+\widehat{F}+\mathrm{\Phi })^{\frac{1}{2}}\underset{i=1}{\overset{s_1}{}}C^{(i)}\underset{p_1,\mathrm{},p_{2n}}{}(h_S)^{p_1p_2}(h_S)^{p_3p_4}\mathrm{}(h_S)^{p_{2n1}p_{2n}}\widehat{J}_{p_1p_2\mathrm{}p_{2n}}^{(i)},$$
(4.20)
we may construct the invariant combinations by adding certain terms like $`\widehat{}_4^B`$ and $`\widehat{}_4^C`$.
Therefore we conclude that the general forms of the allowed $`m`$-derivative corrections in the approximation of neglecting $`𝒪(^{m+2})`$ are given by (4.10) and (4.20) with certain terms like $`\widehat{}_4^B`$ and $`\widehat{}_4^C`$.
Finally we study the behavior of the derivative corrections at $`\theta =0`$ in the zero slope limit of , $`\alpha ^{}ϵ^{\frac{1}{2}}`$, $`g_{ij}=ϵ\delta _{ij}`$ with $`ϵ0`$. In it has been shown that
$$_{DBI}=\frac{1}{g_s(2\pi )^p(\alpha ^{})^{\frac{p+1}{2}}}\left(|\mathrm{Pf}(\mathrm{F}+\mathrm{B})|\frac{ϵ^2}{4}|\mathrm{Pf}(\mathrm{F}+\mathrm{B})|\mathrm{Tr}\frac{1}{(F+B)^2}+𝒪(ϵ^3)\right),$$
(4.21)
where the first term is a constant plus a total derivative. We can show that
$$h_S=\frac{ϵ}{(2\pi \alpha ^{})^2}\frac{1}{(F+B)^2}𝒪(ϵ^0),$$
(4.22)
and
$$h_A=\frac{1}{2\pi \alpha ^{}}\frac{1}{(F+B)}𝒪(ϵ^{\frac{1}{2}}).$$
(4.23)
From the dimensional analysis, the constants $`C^{(i)}`$ and $`C`$ in (4.10) and (4.20), respectively, are restricted. Indeed, we see that $`C`$ or $`C^{(i)}\alpha ^{(p+1)/2+n_S+n_A}`$, where $`n_S`$ and $`n_A`$ are the number of the $`h_S`$ and $`h_A`$ in the Lagrangians $`\widehat{}`$, respectively. Thus
$`\widehat{}`$ $``$ $`{\displaystyle \frac{1}{g_s(2\pi )^p(\alpha ^{})^{\frac{p+1}{2}}}}`$ (4.24)
$`\times \left(ϵ^{\frac{n_S}{2}}|\mathrm{Pf}(\mathrm{F}+\mathrm{B})|{\displaystyle \frac{1}{(F+B)^{2n_S}}}{\displaystyle \frac{1}{(F+B)^{n_A}}}J_{p_1p_2\mathrm{}p_{2(n_S+n_A)}}+𝒪(ϵ^{\frac{n_S}{2}+2})\right),`$
where $`J`$ is $`\widehat{J}^{(i)}`$ or $`\stackrel{~}{J}`$ with $`\widehat{D}=`$ and $`\widehat{F}=F`$. This is negligible compared with $`_{DBI}`$ if $`n_s>4`$. Therefore for the superstring case, the only remaining derivative corrections in the limit are the terms like $`\widehat{}_4^A+\widehat{}_4^B+\widehat{}_4^C`$. This result may have application for a deeper understanding of the relation between the instanton on the noncommutative space and the instanton solution in the Dirac-Born-Infeld Lagrangian with nonzero $`B`$ field .
## 5 Conclusion
We have considered the derivative corrections to the Dirac-Born-Infeld action consistent with the equivalence between the noncommutative gauge theories and the ordinary gauge theory. In particular, we have shown that in the approximation of neglecting the fourth and higher order derivative terms the D-brane action computed in the superstring theory is consistent with the equivalence.
We have also explicitly constructed the general forms of the $`2n`$-derivative corrections which satisfy this equivalence relation in the approximation of neglecting the $`(2n+2)`$-derivative terms. It may capture some general structures of the effective action of the D-brane.
It is interesting to generalize the results obtained in this paper to the effective theories on several D-branes. In this case, we should treat the non-Abelian gauge fields, so that the ordering problems exist even for the ordinary gauge fields which have not been solved yet. Thus the constraints using the equivalence are expected to be important for determination of the effective action on the several D-branes.
Acknowledgements
I would like to thank J. Hashiba, K. Hosomichi and K. Okuyama for useful discussions. I would also like to thank K. Hashimoto, T. Kawano and Y. Okawa for discussions on the derivative corrections. This work was supported in part by JSPS Research Fellowships for Young Scientists.
Note added:
As this article was being completed, we received the preprint which give the derivative corrections for the Dirac-Born-Infeld Lagrangian which is invariant under a simplified version of the Seiberg-Witten map. The two-derivative correction obtained in
$$g_{kl}g_{pq}\left(\frac{1}{g+F}\right)^{ij}_i\left(\frac{1}{g+F}\right)^{kp}_i\left(\frac{1}{g+F}\right)^{lq}=(h_S)^{ij}\mathrm{Tr}\left(h_S(_iF)h_S(_jF)\right),$$
(5.1)
coincides with the $`\widehat{L}_2`$ obtained in this paper with $`a_1=1`$ and $`a_2=0`$. Note that on the computation of the expansion about $`F`$, some terms are omitted in the eq.(4) in . |
warning/0001/quant-ph0001002.html | ar5iv | text | # Coherent states, displaced number states and Laguerre polynomial states for su(1,1) Lie algebra
## I Introduction
The harmonic oscillator is a fundamental exactly solvable physical system and the coherent state(CS) defined in this system is well studied. The generalization of the CS to multi-photon case and the extention to various systems have been made.
As for su(1,1) Lie algebra, the Perelomov’s coherent state(PCS) is well known . The su(1,1) Lie algebra is of great interest in quantum optics because it can characterize many kinds of quantum optical systems. In particular, the bosonic realization of su(1,1) describes the degenerate and non-degenerate parametric amplifiers.
The generators of su(1,1) Lie algebra, $`K_0`$ and $`K_\pm ,`$ satisfy the commutation relations
$$[K_+,K_{}]=2K_0,\text{ }[K_0,K_\pm ]=\pm K_\pm .$$
(1)
Its discrete representation is
$`K_+|n,k`$ $`=`$ $`\sqrt{(n+1)(2k+n)}|n+1,k,`$ (2)
$`K_{}|n,k`$ $`=`$ $`\sqrt{n(2k+n1)}|n1,k,`$ (3)
$`K_0|n,k`$ $`=`$ $`(n+k)|n,k.`$ (4)
Here $`|n,k(n=0,1,2,\mathrm{})`$ is the complete orthonormal basis and $`k=1/2,1,3/2,2,\mathrm{}`$ is the Bargmann index labeling the irreducible representation\[$`k(k1)`$ is the value of Casimir operator\]. We introduce the number operator $`𝒩`$ by
$$𝒩=K_0k,𝒩|n,k=n|n,k.$$
(5)
The PCS is defined as
$`|\alpha ,k_P`$ $`=`$ $`S(\xi )|0,k`$ (6)
$`=`$ $`(1|\alpha |^2)^k{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{{\displaystyle \frac{\mathrm{\Gamma }(2k+n)}{\mathrm{\Gamma }(2k)n!}}}\alpha ^n|n,k,`$ (7)
where $`\xi =r\mathrm{exp}(i\theta )`$ ,$`\alpha =\mathrm{exp}(i\theta )\mathrm{tanh}r`$ , $`\mathrm{\Gamma }(x)`$ is the gamma function, $`S(\xi )=\mathrm{exp}(\xi K_+\xi ^{}K_{})`$ is the su(1,1) displacement operator. There is another coherent state of su(1,1) which is known as the Barut-Girardello(BG) coherent state(BGCS). The BGCS is defined as the eigenstate of the lowering operator $`K_{}`$
$$K_{}|\alpha ,k_{BG}=\alpha |\alpha ,k_{BG},$$
(8)
and it can be expressed as
$$|\alpha ,k_{BG}=\sqrt{\frac{|\alpha |^{2k1}}{I_{2k1}(2|\alpha |)}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{\alpha ^n}{\sqrt{n!\mathrm{\Gamma }(n+2k)}}|n,k,$$
(9)
where $`I_\nu (x)`$ is the first kind modified Bessel function.
The PCS is defined as the displacement operator formalism, while the BGCS as the ladder operator formalism. We ask if the PCS admits the ladder operator formalism? The answer is affirmative. We will discuss it in the next section. We also give the ladder operator formalism of a general su(1,1) quantum state and find that the PCS is a su(1,1) nonlinear coherent state(NLCS). The complete expansion and exponential form of the su(1,1) NLCS are obtained.
There are three definitions of coherent states, that is, (1) the displacement operator acting on the vacuum states, (2) the eigenstates of the annihilation operator, (3) the minimum uncertainty states. These three definitions are identical only for the simplest harmonic oscillator. For su(1,1) system, the PCS is defined according to the first and the BGCS to the second. The minimum uncertainty states(MUSs) for su(1,1) are defined as
$$(\mu K_++\nu K_{})|\alpha ,k_{MUS}=\alpha |\alpha ,k_{MUS},$$
(10)
where $`\mu `$ and $`\nu `$ are complex constants satisfying $`|\mu /\nu |<1.`$ One type of the MUS is the Laguerre polynomial state(LPS). The LPS is only given formally in the literature. We will give the expansion of the LPS in terms of states $`|n,k`$ in section III. The PCS, BGCS and MUS cover the three definations of the coherent states. In section IV, we consider several interesting su(1,1) optical systems, namely, the density-dependent Holstein-Promakoff system, amplitude-squared system, and two-mode system. A conclusion is given in Sec.V.
## II su(1,1) coherent states and nonlinear coherent states
We consider a general state
$$|x,k_G=\underset{n=0}{\overset{\mathrm{}}{}}C(n,x,k)|n,k,$$
(11)
where $`x`$ denote parameter and all the coefficients $`C(n,x,k)`$ are non-zero. Now we try to give the ladder operator formalism of the above general state. The key point is to let the number operator $`𝒩`$ and the operator $`f(𝒩)K_+`$ act on Eq.(8), respectively. Here $`f(𝒩)`$ is a real function of $`𝒩.`$ The operations lead to
$`𝒩|x,k_G`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}C(n,x,k)n|n,k,`$ (12)
$`f(𝒩)K_+|x,k_G`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}f(n)C(n1,x,k)\sqrt{n(n+2k1)}|n,k.`$ (13)
If we choose
$$f(𝒩)=\frac{C(𝒩,k,x)\sqrt{𝒩}}{C(𝒩1,k,x)\sqrt{𝒩+2k1}},$$
(14)
the following equation is obtained
$$[𝒩f(𝒩)K_+]|x,k_G=0.$$
(15)
This is the ladder operator formalism of the general state $`|x,k_G.`$
Let us examine the algebraic structure involved in the general state. Define $`𝒜`$ as an associate algebra with generators
$$𝒩,A_+=\frac{C(𝒩,k,x)\sqrt{𝒩}}{C(𝒩1,k,x)\sqrt{𝒩+2k1}}K_+,A_{}=(A_+)^{}.$$
(16)
Then it is easy to verify that these operators satisfy the following relations
$$[𝒩,A_\pm ]=\pm A_\pm ,\text{ }A_+A_{}=S(𝒩),\text{ }A_{}A_+=S(𝒩+1),$$
(17)
where the function
$$S(𝒩)=\frac{𝒩^2C^2(𝒩,k,x)}{C^2(𝒩1,k,x)}$$
(18)
This algebra $`𝒜`$ is nothing but the generally deformed oscillator(GDO) algebra with the structure function $`S(𝒩)`$. So we see that the general state $`|x,k_G`$ bears generally deformed oscillator algebraic structure.
By acting the annihilation operator $`K_{}`$ on Eq.(11) from left, we get
$$[f(𝒩+1)(𝒩+2k)K_{}]|x,k_G=0.$$
(19)
In the derivation of the above equation, we have used the fact that the operator $`𝒩+1`$ is non-zero in the whole space. The function $`f(𝒩)`$ is completely determined by the coefficients of the state $`|x,k_G.`$ From the coefficients of the BGCS(Eq.(6)), the operator-valued function $`f(𝒩)=\alpha /(𝒩+2k1).`$ It is easily seen that Eq.(15) reduces to Eq.(5) as we expected. From the coefficients of the PCS(Eq.(4)), we obtain the corresponding operator-valued function $`f(𝒩)=\alpha .`$ This simple result leads to the ladder operator formalism of the PCS
$$\frac{1}{𝒩+2k}K_{}|\alpha ,k_P=\alpha |\alpha ,k_P$$
(20)
In fact, by direct verification, we have
$$[\frac{1}{𝒩+2k}K_{},K_+]=1.$$
(21)
Therefore, the exponential formalism of the PCS can be given by
$$|\alpha ,k_P=\mathrm{exp}(\alpha K_+)|0,k$$
(22)
up to a normalization constant.
Reminding the definition of nonlinear coherent states in Fock space, we can call the state $`|\alpha ,k_P`$ as a su(1,1) NLCS. The NLCS is defined as
$$G(𝒩)K_{}|\alpha ,k_{NL}=\alpha |\alpha ,k_{NL},$$
(23)
where $`G(𝒩)`$ is a real function of $`𝒩.`$ The PCS and the BGCS are recovered for the special choices of $`G(𝒩)=1/(𝒩+2k)`$ and $`G(𝒩)=1,`$ respectively. Thus the two coherent states, PCS and BGCS, are unified within the framework of the su(1,1) NLCS.
Assuming the expansion of the NLCS is
$$|\alpha ,k_{NL}=\underset{n=0}{\overset{\mathrm{}}{}}D(n,\alpha ,k)|n,k$$
(24)
and substituting it into Eq.(19), we get the following recursion relation
$$\frac{D(n+1,\alpha ,k)}{D(n,\alpha ,k)}=\frac{\alpha }{G(n)\sqrt{(n+1)(2k+n)}}.$$
(25)
Eq.(21) leads to
$$D(n,\alpha ,k)=\frac{\alpha ^nD(0,\alpha ,k)}{G(n1)G(n2)\mathrm{}G(0)\sqrt{n!\mathrm{\Gamma }(2k+n)/\mathrm{\Gamma }(2k)}}.$$
(26)
The combination of Eq.(20) and (22) gives the expansion of the su(1,1) NLCS,
$`|\alpha ,k_{NL}`$ $`=`$ $`D(0,\alpha ,k){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\alpha ^nD(0,\alpha ,k)\sqrt{\mathrm{\Gamma }(2k)}}{G(n1)G(n2)\mathrm{}G(0)\sqrt{n!\mathrm{\Gamma }(2k+n)}}}|n,k,`$ (27)
$`=`$ $`D(0,\alpha ,k){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\alpha ^n\mathrm{\Gamma }(2k)K_+^n}{G(n1)G(n2)\mathrm{}G(0)n!\mathrm{\Gamma }(2k+n)}}|0,k.`$ (28)
The coefficient $`D(0,\alpha ,k)`$ can be determined by normalization. Let $`G(𝒩)=1/(𝒩+2k),`$ we naturally reduce Eq.(23) to Eq.(4) up to a normalization constant.
One can show that
$`𝒩K_+`$ $`=`$ $`K_+(𝒩+1),f(𝒩)K_+=K_+f(𝒩+1),`$ (29)
$`[f(𝒩)K_+]^n`$ $`=`$ $`(K_+)^nf(𝒩+1)f(𝒩+2)\mathrm{}f(𝒩+n).`$ (30)
Then as a key step, by using Eq.(24) with
$$f(𝒩)=\frac{\alpha }{G(𝒩1)(𝒩+2k1)},$$
(31)
the NLCS is finally written in the exponential form
$`|\alpha ,k_{NL}`$ $`=`$ $`D(0,\alpha ,k){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}[{\displaystyle \frac{\alpha }{G(𝒩1)(𝒩+2k1)}}K_+]^n|0,k,`$ (32)
$`=`$ $`D(0,\alpha ,k)\mathrm{exp}[{\displaystyle \frac{\alpha }{G(𝒩1)(𝒩+2k1)}}K_+]|0,k.`$ (33)
From the above equation, the exponential form of the BGCS is easily obtained by setting $`G(𝒩)=1`$,
$$|\alpha ,k_{BG}=[\mathrm{exp}\frac{\alpha }{(𝒩+2k1)}K_+]|0,k$$
(34)
up to a normalization constant. Let $`G(𝒩)=1/(𝒩+2k)`$ in Eq.(26), Eq.(18) is recovered as we expected.
Actually we have
$$[G(𝒩)K_{},\frac{1}{G(𝒩1)(𝒩+2k1)}K_+]=1.$$
(35)
By this observation, Eq.(26) is naturally obtained.
## III Displaced number states and Laguerre polynomial states
As a generalization of the PCS, we define the displaced number state(DNS) for su(1,1) Lie algebra in analogous with the definition of the displaced number state in Fock space,
$$|\xi ,m,k_{DN}=S(\xi )|m,k=\underset{n=0}{\overset{\mathrm{}}{}}n,k|S(\xi )|m,k|n,k,\text{ }\xi =r\mathrm{exp}(i\theta ).$$
(36)
All the work left is to calculate the matrix elements $`S_{nm}^k(\xi )=n,k|S(\xi )|m,k.`$ Using the decomposed form of the displacement operator
$$S(\xi )=\mathrm{exp}(\alpha K_+)(1|\alpha |^2)^{K_0}\mathrm{exp}(\alpha ^{}K_{})\text{ },\alpha =\mathrm{exp}(i\theta )\mathrm{tanh}r$$
(37)
and the relation
$$\mathrm{exp}(\eta ^{}K_{})|m,k=\underset{q=0}{\overset{m}{}}\frac{(\eta ^{})^{mq}}{(mq)!}\sqrt{\frac{m!\mathrm{\Gamma }(2k+m)}{q!\mathrm{\Gamma }(2k+q)}}|q,k$$
(38)
we obtain the matrix elements as
$`S_{nm}^k(\xi )`$ $`=`$ $`(1|\alpha |^2)^k\alpha ^n(\alpha ^{})^m\sqrt{m!n!\mathrm{\Gamma }(2k+m)\mathrm{\Gamma }(2k+n)}`$ (40)
$`{\displaystyle \underset{q=0}{\overset{\mathrm{min}(m,n)}{}}}{\displaystyle \frac{(11/|\alpha |^2)^q}{q!(nq)!(mq)!\mathrm{\Gamma }(2k+q)}}.`$
Using the relations
$$(m)_q=(1)^q\frac{m!}{(mq)!},(n)^q=(1)^q\frac{n!}{(nq)!},(2k)_q=\frac{\mathrm{\Gamma }(2k+q)}{\mathrm{\Gamma }(2k)},$$
(41)
we can write the matrix elements in terms of hypergeometric function as
$`S_{nm}^k(\xi )`$ $`=`$ $`(1|\alpha |^2)^k\alpha ^n(\alpha ^{})^m`$ (43)
$`\sqrt{{\displaystyle \frac{\mathrm{\Gamma }(2k+m)\mathrm{\Gamma }(2k+n)}{\mathrm{\Gamma }(2k)\mathrm{\Gamma }(2k)m!n!}}}\text{ }_2F_1(m,n;2k;1{\displaystyle \frac{1}{|\alpha |^2}}).`$
Here the hypergeometric function
$${}_{2}{}^{}F_{1}^{}(\alpha ,\beta ;\gamma ;z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(\alpha )_n(\beta )_n}{n!(\gamma )_n}z^n,$$
(44)
and
$$(x)_n=x(x+1)\mathrm{}(x+n1),(x)_01.$$
(45)
The combination of Eqs.(29) and (34) gives the expansion of the DNS in terms of the basis state $`|n,k.`$ It is easily checked that Eq.(29) reduces to Eq.(4) when $`m=0.`$ The matrix elements abtained here are useful in the study of su(1,1) quantum states.
As one type of MUS for su(1,1) Lie algebra, the LPS is given by,
$$|\alpha ,k_{LP}=C_0S(\beta )L_M(\xi \frac{K_0k}{K_0+k1}K_+)|0,k.$$
(46)
Here $`\beta =r\mathrm{exp}(i\theta )`$ is determined by the equation $`\mathrm{exp}(2i\theta )\mathrm{tanh}^2r=\nu /\mu .`$ $`\xi =\mathrm{exp}(i\theta )\mathrm{tanh}(2r),`$ $`C_0`$ can be determined by normalization, and
$$L_M(x)=\underset{n=0}{\overset{M}{}}\frac{1}{n!}\left(\genfrac{}{}{0pt}{}{M}{Mn}\right)(1)^nx^n$$
(47)
is the Laguerre polynomial. Using Eq.(38), we obtain the expansion of the LPS as
$`|\alpha ,k_{LP}`$ $`=`$ $`C_0S(\beta ){\displaystyle \underset{m=0}{\overset{M}{}}}(\xi )^n{\displaystyle \frac{M!}{(Mm)!\sqrt{m!\mathrm{\Gamma }(2k+m)/\mathrm{\Gamma }(2k)}}}|m,k,`$ (48)
$`=`$ $`C_0{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left[{\displaystyle \underset{m=0}{\overset{M}{}}}(\xi )^n{\displaystyle \frac{M!S_{nm}^k(\beta )}{(Mm)!\sqrt{m!\mathrm{\Gamma }(2k+m)/\mathrm{\Gamma }(2k)}}}\right]|n,k.`$ (49)
The combination of Eq.(34) and (39) gives the complete expansion of the LPS.
## IV Some su(1,1) optical systems
In the previous two sections, we obtain the general results of several quantum states for su(1,1) Lie algebra. Now we want to investigate some interesting su(1,1) optical systems.
### A Density-dependent HP realization
The HP realization of the su(1,1) Lie algebra is
$$K_+=a^+\sqrt{N+2k},K_{}=\sqrt{N+2k}a,K_0=N+k.$$
(50)
where $`a^+,a,`$ and $`N=a^+a`$ are the creation, annihilation, and number operator of a single-mode electromagnetic field satisfying $`[a,a^+]=1.`$ On the Fock space $`|n=[a^{+n}/\sqrt{n!}]|0,`$ we have
$`K_+|n`$ $`=`$ $`\sqrt{(n+1)(2k+n)}|n+1,`$ (51)
$`K_{}|n`$ $`=`$ $`\sqrt{n(2k+n1)}|n1,`$ (52)
$`K_0|n`$ $`=`$ $`(n+k)|n.`$ (53)
In comparison with Eq.(2), we see that the HP realization gives rise to the discrete representation of su(1,1) Lie algebra on the usual Fock space. Therefore, by replacing the state $`|n,k`$ by $`|n,`$ we recover all the results in Sec.II and III.
By the replacement procedure described above, we obtain the PCS via HP realization as
$$|\alpha ,M_{NB}=(1|\alpha |^2)^{M/2}\underset{n=0}{\overset{\mathrm{}}{}}\left(\genfrac{}{}{0pt}{}{M+n1}{n}\right)^{1/2}\alpha ^n|n$$
(54)
This is just the well-known negative binomial state(NBS). Here $`M=2k.`$ Since the PCS admits displacement operator formalism, we naturally obtain the displacement operator formalism of the NBS from Eq.(4)
$$|\alpha ,M_{NB}=\mathrm{exp}[\eta a^+\sqrt{N+2k}\eta ^{}\sqrt{N+2k}a]|0.$$
(55)
The parameter $`\eta `$ is determined by the equation $`\eta /|\eta |\mathrm{tanh}|\eta |=\alpha .`$
From Eq.(16), the ladder operator formalism of the NBS is written as
$$\frac{1}{\sqrt{N+M}}a|\alpha ,M_{NB}=\alpha |\alpha ,M_{NB}$$
(56)
As seen from the above equation, we conclude that the NBS is a NLCS in Fock space as discussed in our previous paper. In addition, the su(1,1) displaced number states via HP realization are studied in detail by Fu and Wang.
It can be seen that some useful results of the NBS are conveniently extracted from the general results for su(1,1) Lie algebra .
### B Amplitude-squared realization
The amplitude-squared su(1,1) is given by
$$K_+=\frac{1}{2}a^{+2},K_{}=\frac{1}{2}a^2,K_0=\frac{1}{2}(N+\frac{1}{2}).$$
(57)
The representation on the usual Fock space is completely reducible and decomposes into a direct sum of the even Fock space ($`S_0`$) and odd Fock space ($`S_1`$),
$$S_j=\text{span}\{||n_j|2n+j|n=0,1,2,\mathrm{}\},\text{ }j=0,1.$$
(58)
Representations on $`S_j`$ can be written as
$`K_+||n_j`$ $`=`$ $`\sqrt{(n+1)(n+j+1/2)}||n+1_j,`$ (59)
$`K_{}||n_j`$ $`=`$ $`\sqrt{(n)(n+j1/2)}||n1_j,`$ (60)
$`K_0||n_j`$ $`=`$ $`(n+j/2+1/4)||n_j.`$ (61)
The Bargmann index $`k=1/4`$($`3/4`$) for even(odd) Fock space. From Eq.(4) we see that the PCSs in even/odd Fock space are squeezed vacuum state and squeezed first Fock state
$`|\xi _{SV}`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{\xi }{2}}a^{+2}{\displaystyle \frac{\xi ^{}}{2}}a^2)|0,`$ (62)
$`|\xi _{SF}`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{\xi }{2}}a^{+2}{\displaystyle \frac{\xi ^{}}{2}}a^2)|1,`$ (63)
respectively. The ladder operator formalisms of the squeezed vacuum state and squeezed first Fock state are easily obtained from Eq.(16)
$`{\displaystyle \frac{1}{N+1}}a^2|\xi _{SV}`$ $`=`$ $`\xi /|\xi |\mathrm{tanh}(|\xi |)|\xi _{SV,}`$ (64)
$`{\displaystyle \frac{1}{N+2}}a^2|\xi _{SF}`$ $`=`$ $`\xi /|\xi |\mathrm{tanh}(|\xi |)|\xi _{SF.}`$ (65)
We see that the the two states are the two-photon nonlinear coherent state $`|\alpha _{TP}`$ which is defined as
$$f(N)a^2|\alpha _{TP}=\alpha |\alpha _{TP}.$$
(66)
Here $`f(N)`$ is a real function of the operator $`N.`$
Now we consider the matrix elements $`S_{nm}^k(\xi )`$ (Eq.(34)) in the representation(Eq.(47)). Reminding that the Bargmann index $`k=1/4(3/4)`$ for even(odd) Fock space and substituting $`\alpha =\mathrm{exp}(i\theta )\mathrm{tanh}r`$ into the Eq.(34), we obtain the matrix elements in the representation as
$`S_{nm}^{1/4}(\xi )`$ $`=`$ $`{\displaystyle \frac{(1)^m}{m!n!}}\sqrt{{\displaystyle \frac{(2n)!(2m)!}{\mathrm{cosh}r}}}\mathrm{exp}[i(nm)\theta ](\mathrm{tanh}r/2)^{m+n}`$ (68)
$`{}_{2}{}^{}F_{1}^{}(m,n;1/2;1/\mathrm{sinh}^2r),`$
$`S_{nm}^{3/4}(\xi )`$ $`=`$ $`{\displaystyle \frac{(1)^m}{m!n!}}\sqrt{{\displaystyle \frac{(2n+1)!(2m+1)!}{\mathrm{cosh}^3r}}}\mathrm{exp}[i(nm)\theta ](\mathrm{tanh}r/2)^{m+n}`$ (70)
$`{}_{2}{}^{}F_{1}^{}(m,n;3/2;1/\mathrm{sinh}^2r).`$
As special cases of our general result(Eq.(34)), the above two equations with $`k=1/4(3/4)`$ have been obtained by Marian.
### C Two-mode realization
The two-mode photon operators
$$K_+=a^+b^+,K_{}=ab,K_0=\frac{1}{2}(N_1+N_2+1)$$
(71)
generate the su(1,1) Lie algebra. Here $`N_1=a^+a`$ and $`N_2=b^+b.`$ The Fock space $``$ of the two-mode states is decomposed into a direct sum of irreducible invariant subspaces $`_p^\pm `$
$``$ $`=`$ $`_0_1^\pm \mathrm{}_p^\pm \mathrm{},`$ (72)
$`_p^+`$ $``$ $`\text{span}\{||n_{+p}|n,n+p|n=0,1,2,\mathrm{}\},`$ (73)
$`_p^{}`$ $``$ $`\text{span}\{||n_p|n+p,n|n=0,1,2,\mathrm{}\}.`$ (74)
Representations on $`F_p^\pm `$ are isomorphic and take the form
$`K_+||n_{\pm p}`$ $`=`$ $`\sqrt{(n+1)(n+p+1)}||n+1_{\pm p},`$ (75)
$`K_{}||n_{\pm p}`$ $`=`$ $`\sqrt{n(n+p)}||n1_{\pm p},`$ (76)
$`K_0||n_{\pm p}`$ $`=`$ $`[n+(p+1)/2]||n_{\pm p}.`$ (77)
which are representation (2) with $`k=(p+1)/2.`$ Then by replacing $`|0,k`$ by $`||0_{\pm p}`$ and $`k`$ by $`(p+1)/2`$ in Eq.(4), we obtain the two-mode squeezed vacuum state
$$|\xi ,p_\pm =\mathrm{exp}(\xi a^+b^+\xi ^{}ab)||0_{\pm p}$$
(78)
From Eq.(16), ladder operator formalism of the two-mode squeezed vacuum state is
$$\frac{2}{(N_1+N_2)+p+2}ab|\xi ,p_\pm =\xi /|\xi |\mathrm{tanh}(|\xi |)|\xi ,p_\pm $$
(79)
We can define two-mode NLCS as
$$f(N_1,N_2)ab|\alpha _{TM}=\alpha |\alpha _{TM}.$$
(80)
Therefore, the two-mode squeezed vacuum state can be viewed as the two-mode NLCS. In addtion, the pair coherent state is a special case of two-mode NLCS with $`f(N_1,N_2)=1.`$
## V Conclusions
In this paper we have given the ladder operator formalism of a general quantum state for su(1,1) Lie algebra. The algebra involved in the general state is well-known GDO algebra. The ladder operator formalism of the PCS is found and it is a su(1,1) NLCS. The expansion and exponential form of the NLCS are given. The matrix elements of the su(1,1) squeezing operator is obtained in terms of hypergeometric functions. Using the matrix elements, expansions of the su(1,1) displaced number states and Laguerre polynomial states are obtained. As realizations of su(1,1) Lie algebra, some optical su(1,1) systems are considered. We obtain the ladder operator formalism of the negative binomial state, squeezed vacuum state, squeezed first Fock state, and two-mode squeezed vacuum state. We have generalized the notion of the NLCS in Fock space to the su(1,1) case. It is interesting to study further the su(1,1) NLCS in various quantum optical systems.
Acknowledgment: The author thanks for the discussions with Prof. H. C. Fu and the help of Prof. C. P. Sun, S. H. Pan and G. Z. Yang. The work is partially supported by the National Science Foundation of China with grant number:19875008. |
warning/0001/cond-mat0001181.html | ar5iv | text | # Ab Initio Theoretical Description of the Interrelation between Magnetocrystalline Anisotropy and Atomic Short-Range Order
## Abstract
The cubic lattice symmetry of ferromagnetic homogeneously disordered alloys is broken when a compositional modulation is imposed. This can have a profound influence on the magnetocrystalline anisotropy energy (MAE). We describe our ab initio theory of this effect and use the framework of concentration waves with the electronic structure described within the spin-polarised relativistic Korringa-Kohn-Rostoker coherent-potential approximation. We find that ordering produces a 2 order of magnitude increase in the MAE as well as altering the equilibrium direction of magnetisation. Using the same theoretical framework we also examine directional compositional order produced by magnetic annealing with an explicit study of permalloy.
Magnetocrystalline anisotropy (MCA) of ferromagnetic materials containing transition metals has become the subject of intensive theoretical and experimental research because of the technological implications for high-density magneto-optical storage media . Potential materials for these applications need to exhibit a substantial perpendicular magnetic anisotropy (PMA), which is evidently due to an intrinsic magnetic anisotropy in the crystal lattice strong enough to overcome the extrinsic macroscopic shape anisotropy favoring an in-plane magnetisation. Whereas in ultrathin films and multilayers the PMA is due to surface and interface effects respectively, in thick films of transition metal alloys it is an intrinsically bulk magnetic property which leads to PMA. Such systems are particularly interesting for magneto-optic recording because in addition to PMA they exhibit large Kerr effect signals as well as being chemically stable and easy to manufacture. There has been much effort directed towards an understanding of the mechanism of MCA from a first-principles electronic structure point of view to aid future magnetic material design but since MCA arises from spin-orbit coupling, which is essentially a relativistic effect this means that a fully relativistic electronic structure framework is desirable.
Empirically compositional structure is found to have a profound influence on both the magnitude of the magnetocrystalline anisotropy energy (MAE) as well as the equilibrium or easy magnetisation direction. Compositional order lowers the lattice symmetry of the homogeneously disordered alloy and this enhances its MAE. For example, the measured MAE ($`130\mu `$eV) of ordered CoPt alloy (CuAu or $`L1_0`$ type) is some two orders of magnitude larger than that of its disordered counterpart ($``$3 $`\mu `$eV), and with a different easy axis. In this letter we provide the first quantitative calculations of the effect of compositional order on MAE via a study of CoPt.
We also note that thick fcc-Co<sub>c</sub>Pt<sub>1-c</sub> and fcc-Co<sub>c</sub>Pd<sub>1-c</sub> films exhibit PMA of a size comparable with that of multilayers. Such a large PMA is quite unexpected for systems which should have effectively bulk cubic structure. Now whereas tensile strain and large negative magnetostriction coefficients may account for some of this in the Co<sub>c</sub>Pd<sub>1-c</sub> films, it cannot for the thick Co<sub>c</sub>Pt<sub>1-c</sub> ones . Rather recent experiments have confirmed that although Co<sub>c</sub>Pt<sub>1-c</sub> films seem to be nearly homogeneously disordered there are actually more Co-Co nearest neighbors in-plane and very few out-of-plane , i.e. there is presence of some atomic short-range order (ASRO). The existence of this compositional order produces internal interfaces analogous to Co/Pt multilayers and this has been suggested as a likely cause of the strong PMA in fcc-Co<sub>c</sub>Pt<sub>1-c</sub> films . This suggestion is strengthened by the fact that films grown at higher temperatures, in which enhanced bulk diffusion tends to destroy the in-plane compositional order, show no PMA. Very recently Kamp et al have shown from magnetic circular x-ray dichroism measurements that chemical ordering is also responsible for the enhanced MAE in thick Fe<sub>0.5</sub>Pd<sub>0.5</sub> films. These observations clearly underline a correlation between MCA and ASRO. In this letter we provide, for the first time, a detailed theoretical analysis based on first-principles calculations which shows how chemical order significantly influences MCA and compare the MAE of ordered and disordered CoPt alloys. Moreover, we use the same approach to model other hitherto unfabricated structures. This is pertinent now that it is possible to tailor compositionally modulated films to obtain better magneto-optic recording characteristics .
Upon ordering CoPt undergoes a modest tetragonal lattice distortion ($`c/a=0.98`$) which also lowers the symmetry. From magnetostriction data however we can estimate this to contribute only 15$`\%`$ of the MAE. The measured magnetostriction coefficient, $`\lambda `$, of pure fcc-Co ($`5\times 10^5`$) is of the same size as that of CoPt (-$`4\times 10^5`$) although with different sign. From ‘first-principles’ electronic structure calculations Wu et al have shown that the rate of change of the MAE with lattice strain is proportional to $`\lambda `$ and thus in fcc-Co and CoPt should be roughly of the same magnitude. Thus from Wu et al’s calculation of the MAE of Co for a range of $`c/a`$ ratios we estimate that a 2% tetragonalization will change the MAE only by about 20$`\mu `$eV. Consequently we conclude that it is the compositional order that is primarily responsible for the large MAE of CoPt.
In previous work, we presented a theory of MCA of disordered bulk cubic alloys within the framework of spin-polarised relativistic Korringa-Kohn-Rostoker coherent-potential approximation (SPR-KKR-CPA) . In this letter we set up a new framework to investigate the effects of compositional order, both short and long-ranged. Any compositionally modulated alloy can be specified by site-dependent concentrations $`\{c_i\}`$ which themselves can be written as a superposition of static concentration waves (CWs) , i.e.,
$`c_i=c+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐪}{}}\left[c_𝐪e^{i𝐪𝐑_i}+c_𝐪^{}e^{i𝐪𝐑_i}\right],`$with wave-vectors q and amplitudes $`c_𝐪`$. Usually only a few CWs are needed to describe a particular ordered structure. For example, the CuAu-like $`L1_0`$ structure (Fig. 1) is set up by a single CW with $`c_𝐪=\frac{1}{2}`$ and $`𝐪=(001)`$, and the -layered CuPt-like $`L1_1`$ structure is set up by a CW with $`c_𝐪=\frac{1}{2}`$ and $`𝐪=(\frac{1}{2}\frac{1}{2}\frac{1}{2})`$ (q is in units of $`\frac{2\pi }{a}`$, $`a`$ being the lattice parameter).
The grand-potential for the interacting electrons in an alloy with composition $`\{c_i\}`$ and magnetised along the direction $`𝐞_1`$ at a finite temperature $`T`$ is given by ,
$`\mathrm{\Omega }(\{c_i\};𝐞_1)=\nu Z`$ $``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\epsilon f(\epsilon ,\nu )N(\{c_i\},\epsilon ;𝐞_1)`$
$`+`$ $`\mathrm{\Omega }_{DC}(\{c_i\};𝐞_1),`$
where, $`\nu `$ is the chemical potential, $`Z`$ the total valence charge, $`f(\epsilon ,\nu )`$ the Fermi factor, $`N(\{c_i\},\epsilon ;𝐞_1)`$ the integrated electronic density of states, and $`\mathrm{\Omega }_{DC}(\{c_i\};𝐞_1)`$ the ‘double-counting’ correction. The MAE of the inhomogeneous alloy can be characterised by the difference
$`K(\{c_i\})=\mathrm{\Omega }(\{c_i\};𝐞_1)\mathrm{\Omega }(\{c_i\};𝐞_2).`$Assuming that $`\mathrm{\Omega }_{DC}(\{c_i\};𝐞)`$ is generally unaffected by the change in the magnetisation direction, we get,
$`K(\{c_i\})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\epsilon f(\epsilon ,\nu _1)\left[N(\{c_i\},\epsilon ;𝐞_1)N(\{c_i\},\epsilon ;𝐞_2)\right]`$
$`+`$ $`O(\nu _1\nu _2)^2.`$
Note that the correction due to the change in the chemical potential (from $`\nu _1`$ to $`\nu _2`$) with the magnetisation direction is of second order in $`(\nu _1\nu _2)`$, and can be shown to be very small compared to the first term . We now expand $`K(\{c_i\})`$ around $`K_{CPA}(c)`$, the MAE of the homogeneously disordered alloy $`A_cB_{1c}`$,
$`K(\{c_i\})`$ $`=`$ $`K_{CPA}(c)+{\displaystyle \underset{j}{}}{\displaystyle \frac{K(\{c_i\})}{c_j}}|_{\{c_i=c\}}\delta c_j`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k}{}}{\displaystyle \frac{^2K(\{c_i\})}{c_jc_k}}|_{\{c_i=c\}}\delta c_j\delta c_k+O(\delta c)^3.`$ (2)
The second term in Eq. (2) vanishes if the number of $`A`$ and $`B`$ atoms in the alloy is to be conserved. Now taking the Fourier transform, we get the MAE of the compositionally modulated alloy with wave-vector q,
$$K(𝐪)=K_{CPA}(c)+\frac{1}{2}|c_𝐪|^2\left[S^{(2)}(𝐪;𝐞_1)S^{(2)}(𝐪;𝐞_2)\right],$$
(3)
and so, for example, the MAE of the CuAu-type $`L1_0`$ ordered alloy is obtained by choosing $`𝐪=(001)`$ and $`c_𝐪=\frac{1}{2}`$. Here $`S^{(2)}(𝐪;𝐞)`$ are the Fourier transforms of the so-called direct correlation functions which determine the ASRO parameter $`\alpha (𝐪)`$ in the disordered phase, $`\alpha (𝐪)=\beta c(1c)/[1\beta c(1c)S^{(2)}(𝐪)]`$, ($`\beta =1/k_BT`$, $`k_B`$ being the Boltzmann constant). These have been calculated for many alloys, both non-magnetic and ferromagnetic, in which up to now relativistic effects were largely ignored and were compared with diffuse X-ray and neutron scattering data . $`S^{(2)}(𝐪)`$ is given by,
$`S^{(2)}(𝐪)`$ $`=`$ $`{\displaystyle \frac{Im}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\epsilon f(\epsilon ,\nu ){\displaystyle \underset{L_1L_2L_3L_4}{}}(X^AX^B)_{L_1L_2}`$
$`\times I_{L_2L_3;L_4L_1}(𝐪)\mathrm{\Lambda }_{L_3L_4}(𝐪),`$
where,
$`\mathrm{\Lambda }_{L_1L_2}(𝐪)=(X^AX^B)_{L_1L_2}`$
$`{\displaystyle \underset{L_3L_4L_5L_6}{}}X_{L_1L_5}^AI_{L_5L_3;L_4L_6}(𝐪)X_{L_6L_2}^B\mathrm{\Lambda }_{L_3L_4}(𝐪),`$
$`I_{L_5L_3;L_4L_6}(𝐪)={\displaystyle \frac{1}{V_{BZ}}}{\displaystyle }`$ $`d𝐤`$ $`\tau _{L_5L_3}(𝐤+𝐪)\tau _{L_4L_6}(𝐤)`$ (4)
$``$ $`\tau _{L_5L_3}^{00}\tau _{L_4L_6}^{00},`$ (5)
and $`X^{A(B)}=[(t_{A(B)}^1t_c^1)^1+\tau ^{00}]^1`$. Here, $`\tau ^{00}`$ is the site-diagonal path-operator matrix, $`t_{A(B)}`$ and $`t_c`$ are the t-matrices for electronic scattering from sites occupied by $`A(B)`$ atoms and the CPA effective potentials respectively, $`\tau (𝐤)=[t_c^1g(𝐤)]^1`$, and $`g(𝐤)`$ is the KKR structure constants matrix . The spinodal ordering temperature below which the alloy orders into a structure characterised by concentration wave-vector $`𝐪_{max}`$ is given by , $`T_c=c(1c)S^{(2)}(𝐪_{max};𝐞)/k_B`$, where $`𝐪_{max}`$ is the value at which $`S^{(2)}(𝐪;𝐞)`$ is maximal.
As with all MAE calculations, given the sizes of the energies involved, numerical computation of $`K(𝐪)`$ needs to be done very carefully. The energy integration is done using a complex contour described elsewhere . The Brillouin zone (BZ) integration is done using the adaptive grid method which has been found to be very efficient and accurate. In this method one can preset the level of accuracy of the integration by supplying an error parameter $`ϵ`$. The integration in Eq. (5) is done with $`ϵ=10^6`$ which means that $`S^{(2)}(𝐪;𝐞)`$ which are of the order of 0.1 eV are accurate up to 0.1 $`\mu `$eV. To achieve such level of accuracy, we had to sample a large number of $`𝐤`$-points in the BZ. Also, owing to the form of the integrand in Eq. (5) the integration has to be done using the full BZ. Typically, in our calculations, we needed around 35,000 $`𝐤`$-points for points on the energy contour 0.5 Ry above the real axis. When the energy was 0.0001 Ry above the real axis (and that is the closest point to the real axis, we need) we required as many as 3 million $`𝐤`$-points for the same level of accuracy. Furthermore, we have calculated $`S^{(2)}(𝐪;𝐞_1)`$ and $`S^{(2)}(𝐪;𝐞_2)`$ simultaneously ensuring that they are calculated on the same grid, and thus cancelling out the systematic errors if any. Therefore, we claim that the values of $`K(𝐪)`$ are accurate to within 0.1 $`\mu `$eV.
We summarise the results for fcc-Co<sub>0.5</sub>Pt<sub>0.5</sub> alloy in Tables I and II. We note that, $`S^{(2)}(𝐪)`$ is maximum for the $`L1_0`$ structure, implying that the alloy would order into this structure at 1360 K as it is cooled from high temperature in good agreement with experiment (ordering temperature of 1000 K ). From Table II we note that for $`𝐪=(001)`$ and $`𝐪=(100)`$ which correspond to CuAu-like $`L1_0`$ ordered structure, with Co and Pt layers stacked along the and directions respectively, the direction of spontaneous magnetisation is along the and directions respectively in excellent agreement with experiment. Also, the MAE (58.6 $`\mu `$eV) is comparable to the experimental value ($`130\mu `$eV ).
In order to probe the relationship between the compositional structure and MAE further, we also performed calculations for some hypothetical structures. These are also summarised in Table II. The case of $`𝐪=(\frac{1}{2}\frac{1}{2}\frac{1}{2})`$ corresponding to the CuPt-type $`L1_1`$ ordered structure with Co and Pt layers perpendicular to the direction of the crystal produces a spontaneous magnetisation along the direction of the crystal. This may be close to the structure of the thick oriented disordered CoPt alloy films exhibiting PMA which is attributed to the existence of internal interfaces , analogous to Co and Pt layers along the direction. Our result is clearly consistent with these observations. The point is that the magnetic anisotropy (152$`\mu `$eV) intrinsic to this structure is nearly 3 times larger than that of $`L1_0`$ structure. Indeed, we predict that a -oriented film will exhibit a markedly stronger PMA than that of a -oriented film.
The structure set up by CWs with $`𝐪=(10\frac{1}{2})`$ and $`(01\frac{1}{2})`$ is also a -oriented layered structure, but the layers are not alternately pure Co and Pt planes, rather they are layers of ordered Co and Pt (Fig. 1). Even in this case, we note that the magnetisation is perpendicular to the layered structure, and the magnitude of MAE is large. Similarly, for $`𝐪=(\frac{1}{2}01)`$ and $`(\frac{1}{2}10)`$ where the planes are stacked along the direction the magnetisation is also along the direction.
Evidently, the spontaneous magnetisation always tends to align itself perpendicular to any layering in the structure. In addition, the magnitude of MAE depends strongly on the symmetry of the system, i.e. it increases when the symmetry is lowered. The cubic symmetry of the homogeneously disordered alloy quenches the orbital magnetic moment. In the $`L1_0`$ ordered alloy the symmetry is lower and in the $`L1_1`$ layered structure it is lower still, thus increasing the MAE in each case. On detailed examination of the electronic structure of the disordered alloy around the Fermi energy we find that the large values of $`K(𝐪)`$ near the BZ boundary arise from van Hove singularities of the Bloch spectral function . The number, location, and occupation of these depend on the magnetisation direction and produce a large contribution to the difference in the convolution integrals (Eqs. (3) and (5)).
Our theory can also be used to produce the first quantitative description of the phenomenon of magnetic annealing. Here a soft magnetic material develops directional chemical order when annealed in a magnetic field . We demonstrate this effect for Ni<sub>0.75</sub>Fe<sub>0.25</sub> (permalloy) and Table III is a summary of the results. We calculate $`S^{(2)}(𝐪,𝐞)`$ for permalloy in an applied magnetic field of strength 600 Oe (same as used in the experiment ) orientated along $`𝐞=`$ , , and directions of the crystal. We find that when the magnetic field is along (column 2) $`S^{(2)}(𝐪)`$ is maximum for q=(001) confirming that ordering is favored along the direction of applied field. Similar is the case when the magnetic field is along . However, when the applied field is along (column 6) all the three orderings, namely, (100), (010), and (001) are favored equally. Thus, in this case, we will get a CuAu<sub>3</sub>-type $`L1_2`$ ordering. The calculated transition temperature 721 K is in good agreement with the experimental value of 820 K . Noting that the measured intensity in a scattering experiment is proportional to the ASRO parameter $`\alpha (𝐪)`$, we estimate that for an alloy cooled in a magnetic field along the direction the superlattice spot at q=(001) will be 20% more intense than that at q=(100) at a temperature 1 K above the transition temperature.
In conclusion, we have presented a fully relativistic electronic structure scheme to study the MCA of alloys and its dependence upon compositional structure. We applied this theory to fcc-Co<sub>0.5</sub>Pt<sub>0.5</sub> and found that when the system is cooled it tends to order into $`L1_0`$ layered-ordered structure and that the spontaneous magnetisation tends to align itself perpendicular to the layer stacking. These observations are in complete accord with experiment. We also found that if the layers are stacked along the direction then the MCA becomes larger which may be the case in -textured films. Within the same framework we have also been able to explain the appearance of directional order in Ni<sub>0.75</sub>Fe<sub>0.25</sub> when it is annealed in a magnetic field.
We thank B.L. Gyorffy for many helpful discussions. This research is supported by the Engineering and Physical Sciences Research Council (UK), National Science Foundation (USA), and the TMR Network on “Electronic structure calculation of materials properties and processes for industry and basic sciences”. |
warning/0001/hep-th0001054.html | ar5iv | text | # Contents
## 3.1 The Holographicly-Dual Pair
We start by reviewing the case with a hypersurface in flat space, considered in , and then introduce the orbifold.
A Hypersurface in Flat Space
A large class of singular manifolds $`X^{2n}`$ can be modeled by a hypersurface
$$W(z_1,\mathrm{},z_{n+1})=0$$
in $`^{n+1}`$, with $`W`$ a quasi-homogeneous polynomial, i.e., transforming as $`W\lambda W`$ under
$$z_a\lambda ^{r_a}z_a,\lambda ,$$
(3.1)
for some positive (rational) “weights” $`r_a`$. The authors of considered such a singularity which has the following additional properties:
* An isolated singularity (point-like; at $`z=0`$); the condition for this is that the polynomial $`W`$ is “transverse”
$$dW=0z=0.$$
(3.2)
* A singularity that appears at finite distance in the moduli space of (compact, smooth) CY manifolds; as shown in , the condition for this is $`r_\mathrm{\Omega }>0`$, where
$$r_\mathrm{\Omega }=\underset{a}{}r_a1.$$
(3.3)
It was proposed in that the dynamics at the singularity in the decoupling limit is described by type II string theory on
$$^{d1,1}\times _\varphi \times U(1)_Y\times LG_W,$$
(3.4)
where $`_\varphi \times U(1)_Y`$ is the $`N=2`$ Liouville theory and $`LG_W`$ is the $`N=2`$ Landau-Ginzburg (LG) model with the superpotential being the polynomial $`W`$. One way to understand this identification is as follows: to obtain better control on the model, one deforms $`X^{2n}`$ to a smooth hypersurface $`X_\mu ^{2n}`$ by changing $`W`$
$$WW_\mu =W+\mu V$$
(3.5)
(where $`V`$ is an appropriate quasi-homogeneous polynomial<sup>6</sup><sup>6</sup>6We assume quasi-homogeneity for simplicity. The arguments below generalize in a straight-forward way to a general polynomial, which is of the form $`V=_iV_i`$, where $`V_i`$ are quasi-homogeneous. Note that the $`z_0`$-dressing in eq. (3.6) will not be the same for all terms., with weight $`r_V`$). Next, this deformed hypersurface is embedded in a weighted projective space: $`^{n+1}`$ is embedded in $`𝕎_{r_\mathrm{\Omega },r_1,\mathrm{},r_{n+1}}^{n+1}`$, by introducing an additional coordinate $`z_0`$, with a negative weight $`r_\mathrm{\Omega }`$ (where $`r_\mathrm{\Omega }`$ is defined in eq. (3.3)); to embed $`X_\mu ^{2n}`$ in this space, $`W_\mu `$ is made quasi-homogeneous, by an appropriate $`z_0`$ “dressing”<sup>7</sup><sup>7</sup>7Note that, to deform the singularity, $`V`$ must dominate at $`z0`$. This means that $`r_V<1`$, so the power of $`z_0`$ in (3.6) is negative!:
$$W_\mu \widehat{W}_\mu =W+\mu \widehat{V},\widehat{V}=z_0^{(r_V1)/r_\mathrm{\Omega }}V.$$
(3.6)
Now one argues, following , that the SUSY NL$`\sigma `$M on $`X_\mu ^{2n}`$ is described, in the IR limit, by a LG model with the superpotential $`\widehat{W}_\mu `$ (eq. (3.6))<sup>8</sup><sup>8</sup>8More precisely, this is an orbifold of the above LG model, with a projection on integral $`U(1)`$ charges of the N=2 superconformal algebra. However, in the present context (of string theory on (3.4)), this orbifold is part of the GSO projection, so it can be ignored at this stage. <sup>9</sup><sup>9</sup>9In , this identification was obtained, for $`\mu =0`$, using an embedding in a gauged linear $`\sigma `$ model, following .. Finally, one identifies the $`z_0`$ CFT with $`N=2`$ Liouville theory<sup>10</sup><sup>10</sup>10Strictly speaking, this is Liouville theory when $`V=1`$, in which case, $`\delta \widehat{W}=\mu \widehat{V}=\mu e^{\widehat{\varphi }/Q}`$ is the Liouville interaction. This is the case considered in , and the $`z_0`$ CFT was identified there with the $`SL(2)/U(1)`$ coset SCFT (at level $`k=1/r_\mathrm{\Omega }`$; see for a review of the evidence for this identification). On the other hand, the analysis in suggested the identification with Liouville theory, as described above. This led to a proposal that these two theories are equivalent.. This theory has a single chiral superfield $`\widehat{\varphi }`$. Denoting the (scalar) bottom component of $`\widehat{\varphi }`$ by $`\varphi +iY`$, the dilaton $`\mathrm{\Phi }`$ is proportional to $`\varphi `$:
$$\mathrm{\Phi }=\frac{Q}{2}\varphi ,$$
(3.7)
where $`Q`$ is a positive parameter. Now, the identification is
$$z_0=e^{\frac{Q}{2}\widehat{\varphi }},Q=\sqrt{2r_\mathrm{\Omega }}.$$
(3.8)
Recall that $`r_\mathrm{\Omega }`$ was chosen to be positive (for the singularity to be at finite distance in the CY moduli space), and this is indeed required for the consistency of the above identification. Thus, combining these relations, one obtains the background (3.4).
After obtaining the above identification, one may want to remove the deformation $`\mu V`$, returning to the original, singular, hypersurface $`X^{2n}`$. However, one finds that the worldsheet formulation of the string theory in the background (3.4) becomes singular. Indeed, the string coupling $`g_s`$ is related to the dilaton: $`g_s=e^\mathrm{\Phi }`$, so eq. (3.7) implies that $`g_s`$ diverges at $`\varphi \mathrm{}`$. For $`\mu =0`$, the superpotential is independent of $`\varphi `$, so the region $`\varphi \mathrm{}`$ is accessible and the perturbative expansion is singular. This singularity is a reflection of the existence of massless solitons: D-branes wrapped on the vanishing cycles at the singularity . Therefore, to obtain a well-behaved perturbative description, it is necessary to make these solitons massive.
This is achieved by the double scaling limit (2.1), suggested in : the deformation discussed there was with $`V=1`$ in eq. (3.5) and the quantity held fixed was the ratio
$$m=\frac{\mu ^{r_\mathrm{\Omega }}}{g_s}$$
(3.9)
($`r=r_\mathrm{\Omega }`$ in eq. (2.1)). To understand the meaning of $`m`$, one notes that $`X_\mu ^{2n}`$ has holomorphic $`n`$-cycles that vanish for $`\mu =0`$. The volume $`V_C`$ of such a cycle $`C`$ is $`V_C=_C\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is the holomorphic $`n`$-form on $`X_\mu ^{2n}`$. From the explicit expression for $`\mathrm{\Omega }`$
$$\mathrm{\Omega }=\frac{dz_1\mathrm{}dz_{n+1}}{dW},$$
(3.10)
one finds that the weight of $`\mathrm{\Omega }`$ under (3.1) is $`r_\mathrm{\Omega }`$, defined in eq. (3.3) (and assumed to be positive). This implies that $`V_C\mu ^{r_\mathrm{\Omega }}`$ and, therefore, the $`d`$-dimensional mass/tension of a D-brane wrapped on $`C`$ is $`V_C/g_sm`$. Thus, in the double scaling limit with $`m0`$, there are no massless solitons, so one can expect that the worldsheet formulation will be under control and indeed it is: the strong coupling region $`\varphi \mathrm{}`$ is inaccessible, because of the deformed superpotential $`\widehat{W}`$ that diverges there. Using the relation $`g_s=e^\mathrm{\Phi }`$, one can show that the theory depends on $`g_s`$ and $`\mu `$ only through the ratio $`m`$ defined in eq. (3.9) and it can be studied perturbatively, $`1/m`$ being the string loop expansion parameter.
The Orbifold
Considering the above scenario, let $`\mathrm{\Gamma }`$ be a finite subgroup of $`SU(n+1)`$ which commutes with (3.1) (i.e., mixes only coordinates with the same weight) and leaves the polynomial $`W`$ invariant. The equation $`W=0`$ is, therefore, well defined in the orbifold $`^{n+1}/\mathrm{\Gamma }`$, and it defines there a hypersurface $`X^{2n}`$, with an isolated singularity at $`z=0`$. We are interested in the $`d`$-dimensional theory describing the decoupled dynamics near this singularity. Note that the singular point $`z=0`$ of $`X^{2n}`$ is at the orbifold singularity, which is what is needed for the decoupled theory to be affected by the orbifold. To obtain an anti-holographic formulation of this theory, we replace the $`LG_W`$ factor in the background (3.4) by the orbifold $`LG_W/\mathrm{\Gamma }`$. Note that $`\mathrm{\Gamma }`$ defines a symmetry group of $`LG_W`$ (acting trivially on the supercharges), so one can orbifold by it. We, therefore, propose (generalizing ) that the following theories are two formulations of the same $`d`$-dimensional theory:
* geometric formulation: string theory on
$$^{d1,1}\times X^{2n},$$
(3.11)
in the decoupling limit (as described in section 1.1);
* worldsheet formulation: string theory on
$$^{d1,1}\times _\varphi \times U(1)_Y\times LG_W/\mathrm{\Gamma }.$$
(3.12)
## 3.2 Aspects of the Holographic Relation
In this and the next subsections, properties of the decoupled $`d`$-dimensional theory are identified in the worldsheet formulation. This will be used, in the next section, to analyze specific four-dimensional theories. Some aspects of this identification depend only on the $`^{d1,1}\times _\varphi \times U(1)_Y`$ factor of the background (3.12), while the factor $`LG_W/\mathrm{\Gamma }`$ could be replaced by any 2D (2,2) SCFT. These aspects are, therefore, naturally the same here as in the unorbifolded case, discussed in . Other aspects depend also on properties of the $`LG_W/\mathrm{\Gamma }`$ factor and, accordingly, on the correct identification of the $`d`$ dimensional theory. As in , these aspects are compared to the information available from the geometric formulation and agreement is found, providing evidence for the duality. The results and their implications are summarized in section 4.4.
Space-Time Supersymmetry
We start with supersymmetry, showing that in both formulations there are $`2^{\frac{d}{2}+1}`$ preserved space-time supercharges.
In the worldsheet formulation, the orbifold action in $`LG_W`$ was chosen to commute with the supercharges and the transformation (3.1) so, in the IR CFT, it commutes with the full $`(2,2)`$ superconformal algebra (in which one of the $`U(1)`$ R-symmetries is identified with (3.1)). Therefore, the CFT (3.12) has $`(2,2)`$ supersymmetry and, as in any such situation, one can construct the space-time supercharges using the $`U(1)`$ currents of the superconformal algebra. This leads, as before orbifolding (see for more details), to $`2^{\frac{d}{2}+1}`$ space-time supercharges. The unimodularity of the elements of $`\mathrm{\Gamma }`$ (det($`g`$)=1, $`g\mathrm{\Gamma }`$) insures that, for an appropriate choice of the orbifold action<sup>11</sup><sup>11</sup>11See for more details (in the notation there, the choice is $`(1)^{K_g}=det(g)`$)., the space-time supercharges are $`\mathrm{\Gamma }`$-invariant.
In the geometric formulation, before orbifolding there were $`2^{\frac{d}{2}+1}`$ supercharges (as in any CY compactification), so it remains to check the effect of the orbifold. The supercharges are related to the holomorphic $`n`$-form $`\mathrm{\Omega }`$ on $`X^{2n}`$: indeed, $`\mathrm{\Omega }`$ can be written in terms of a covariantly constant spinor $`\eta `$ on $`X^{2n}`$ as $`\mathrm{\Omega }_{i_1\mathrm{}i_n}=\eta ^t\mathrm{\Gamma }_{i_1\mathrm{}i_n}\eta `$, so the orbifold actions on $`\mathrm{\Omega }`$ and on the supercharges are related. Using the explicit expression (3.10) for $`\mathrm{\Omega }`$, one finds that $`\mathrm{\Omega }`$ is $`\mathrm{\Gamma }`$-invariant: both numerator and denominator are invariant (the first invariance follows from the unimodularity of the elements of $`\mathrm{\Gamma }`$). Therefore, the orbifold can have at most a $`_2`$ action on the supercharges. This $`_2`$-ambiguity has the following meaning: the supercharges appear quadratically in the supersymmetry algebra, so with a given $`\mathrm{\Gamma }`$-action on the bosonic coordinates, its action on the supercharges always has a $`_2`$ ambiguity. The invariance of $`\mathrm{\Omega }`$ means that one can choose the $`\mathrm{\Gamma }`$-action on the supercharges to be trivial and, with this choice, the orbifold does not break supersymmetry<sup>12</sup><sup>12</sup>12This choice is the space-time analog of the choice made in the worldsheet formulation (see footnote 11)..
Deformations Preserving the Space-Time Supersymmetry
Next we discuss SUSY-preserving deformations.
Consider, for example, a deformation of the polynomial $`W`$: $`\delta W=\mu V`$ (where $`V`$ is a quasi homogeneous polynomial, with weight $`r_V`$; as in eq. (3.5)). In the geometric formulation, it induces a change in the manifold $`X^{2n}`$ and, therefore, also in the corresponding decoupled theory. More specifically, the coefficient (modulus) $`\mu `$ parametrizes a change in the complex structure of $`X^{2n}`$, a deformation that does not break supersymmetry. In the worldsheet formulation, a change in $`W`$ induces a change in the superpotential: $`\delta \widehat{W}=\mu \widehat{V}`$, as in eq. (3.6). This means that one adds to the worldsheet Lagrangian a top component of a chiral-chiral superfield (of the worldsheet $`N=2`$ SCFT). The $`z_0`$ dressing insures that this operator has dimension $`\mathrm{\Delta }=\stackrel{~}{\mathrm{\Delta }}=1`$. Such a change in the Lagrangian preserves the worldsheet superconformal symmetry and, therefore, defines a (truly) marginal perturbation, parametrized by the coefficient (coupling) $`\mu `$, which preserves the space-time supersymmetry.
This can be generalized for any chiral-chiral primary operator $`V`$ in the $`LG_W/\mathrm{\Gamma }`$ factor with equal left and right $`U(1)`$ charges $`f_V=\stackrel{~}{f}_V`$. $`V`$ can be “dressed” by the Liouville fields:
$$\widehat{V}=e^{\beta (\varphi +iY)}V.$$
(3.13)
This is a chiral-chiral primary of the worldsheet $`N=2`$ SCFT, which is the bottom component of a chiral-chiral superfield (we denote both by the same symbol $`\widehat{V}`$). The $`U(1)`$ charge of $`\widehat{V}`$ (both left and right) is $`f_V\beta Q`$ and, to obtain a marginal deformation, as described above, one chooses $`\beta `$ to set this charge to 1:
$$\beta =\frac{f_V1}{Q}.$$
(3.14)
This defines a map from the ring of chiral-chiral primary operators (the (c,c) ring) in the $`LG_W/\mathrm{\Gamma }`$ factor, to SUSY-preserving changes in the decoupled theory.
Before orbifolding, the (c,c) ring in $`LG_W`$ is spanned by quasi homogeneous polynomials $`V(z_a)`$, identified modulo $`_aW`$. Their $`U(1)`$ charges $`f_V,\stackrel{~}{f}_V`$ are equal and coincide with the weight $`r_V`$ of $`V`$ under (3.1). In the geometric formulation, the corresponding deformation was identified in with a deformation of the complex structure in $`X^{2n}`$, as described above. The orbifold projection truncates this set of deformations in the same way in both the worldsheet and geometric formulations, by restricting the polynomials $`V`$ to be $`\mathrm{\Gamma }`$-invariant. This can be viewed as further evidence for both the holographic relation and the geometric identification of these deformations. Additional (c,c) operators appear in twisted sectors of the orbifold. The corresponding deformations in the $`d`$-dimensional theory are related to its influence by the orbifold singularity. This will be seen in specific examples in the next section.
Couplings vs. Moduli
In the previous paragraph, possible deformations of the $`d`$-dimensional theory were discussed, as viewed in its two formulations. The parameter $`\mu `$ parametrizing such a deformation can be either
* a coupling: parametrizing a change in the theory (e.g., a coefficient of a term in the action); or
* a modulus: parametrizing a change in the vacuum (e.g., a vacuum expectation value (vev)).
In the holographic duality with $`AdS`$ vacua (for a review, see ), considered in the framework of semi-classical supergravity, the distinction between couplings and moduli is related to the behavior of the corresponding bulk fields near the boundary: non normalizable modes are related to couplings in the boundary theory while normalizable modes are related to vevs in the boundary (see also ). It is natural to expect the same distinction in the present holographic duality<sup>13</sup><sup>13</sup>13The relation to $`AdS_{d+1}`$ vacua can be made more concrete in the $`d=2`$ case, by using the correspondence, described in , between $`AdS_3`$ vacua and $`^{1,1}\times _\varphi `$ vacua.. Here, the “boundary” is at $`\varphi \mathrm{}`$ and, for a deformation defined by a vertex operator $`\widehat{V}`$, the deformation parameter $`\mu `$ is identified as a coupling iff the corresponding wave function is non-normalizable at $`\varphi \mathrm{}`$. For the vertex operators (3.13), this implies that for
$$\beta >\frac{Q}{2},$$
(3.15)
$`\mu `$ is a coupling and for $`\beta <\frac{Q}{2}`$, it is a modulus<sup>14</sup><sup>14</sup>14To obtain these conditions, one notes that the vertex operator $`\widehat{V}`$ is related to the wavefunction $`\mathrm{\Psi }`$ by $`\widehat{V}=g_s\mathrm{\Psi }`$, where $`g_s=e^\mathrm{\Phi }=e^{\frac{Q}{2}\varphi }`$.. The case $`\beta =\frac{Q}{2}`$ is more delicate and will be discussed below. The bound (3.15) appears also when the string theory on (3.12) is viewed as a 2D CFT coupled to quantum gravity. In that context, (3.15) is the condition for the operator (3.13) to exist as a local operator in the theory (see for more details).
Using eqs. (3.14),(3.8), one obtains that for
$$f_V>1r_\mathrm{\Omega },$$
(3.16)
$`\mu `$ is a coupling and for
$$f_V<1r_\mathrm{\Omega },$$
(3.17)
it is a modulus. This applies, in particular, to the deformations of the polynomial $`W`$: $`\delta W=\mu V`$. For these deformations, the distinction between couplings and moduli was investigated in the geometric formulation in (see also ). In this approach, a coupling is a 10D mode whose $`d`$-dimensional kinetic energy diverges and, consequently, its fluctuations are “frozen”. For $`r_V1r_\mathrm{\Omega }`$, this criterion led to the same conditions on $`r_V`$ as above, namely, the conditions (3.16),(3.17) with $`f_V=r_V`$. The deformation with $`r_V=1r_\mathrm{\Omega }`$ is found to be a coupling. It is interesting to note that also in this approach, $`\mu `$ was identified as a non-fluctuating coupling when the perturbation was supported (in some sense) far from the singularity (which is indeed identified here with $`\varphi \mathrm{}`$). The manifold considered in was a hypersurface in flat space. However, the aspects analyzed were convergence of certain integrals, and these aspects are not influenced by a finite group of identifications. Therefore, the results of apply also to the present case.
R-Symmetry
In the worldsheet formulation, the scalar $`Y`$ is free, and this leads to two $`U(1)`$ symmetries, with the following conserved charges
$$R=i\frac{2}{Q}Y,\stackrel{~}{R}=i\frac{2}{Q}\overline{}Y.$$
(3.18)
The space-time supercharges are charged under these symmetries, (see for details), so these are R-symmetries. Defining the linear combinations
$$R_\pm =R\pm \stackrel{~}{R},$$
(3.19)
all supercharges have charge $`|R_\pm |=1`$. For the deformations defined by $`\widehat{V}`$ in eq. (3.13), $`R=\stackrel{~}{R}`$, so
$$R_{}(\mu )=0,R_+(\mu )=R_+(\widehat{V})=2R(\widehat{V})=\frac{4\beta }{Q}=2\frac{1f_V}{r_\mathrm{\Omega }}$$
(3.20)
(using eqs. (3.13),(3.14),(3.8) and the fact that the worldsheet action is invariant). Combining eqs. (3.16),(3.17),(3.20), one obtains that $`\mu `$ is a coupling when $`R_+(\mu )<2`$ and a modulus when $`R_+(\mu )>2`$ .
In the geometric formulation, the transformation (3.1) with $`|\lambda |=1`$ is an isometry of $`X^{2n}`$ and, therefore, induces a $`U(1)`$ symmetry in the corresponding theory. To identify its action on supercharges, one considers again the holomorphic $`n`$-form $`\mathrm{\Omega }`$. From eq. (3.10) one finds that the weight of $`\mathrm{\Omega }`$ under (3.1) is $`r_\mathrm{\Omega }`$, defined in eq. (3.3), and the fact that it is non-zero means that the above $`U(1)`$ symmetry is an R-symmetry. With a normalization in which the supercharges have $`U(1)`$ charge $`\pm 1`$, the R-charge $`R_+^{}`$ is related to the weight $`r`$ by
$$R_+^{}=2\frac{r}{r_\mathrm{\Omega }}.$$
It was suggested in that $`R_+^{}`$ should be identified with $`R_+`$. This was checked by verifying that, for deformations of $`W`$: $`\delta W=\mu V`$ (where $`f_V=r_V`$), the geometric formulation indeed gives the same R-charge as obtained in eq. (3.20):
$$R_+^{}(\mu )=2\frac{r_\mu }{r_\mathrm{\Omega }}=2\frac{1r_V}{r_\mathrm{\Omega }}$$
(3.21)
(where the last equality follows from the fact that the polynomial $`W`$ has weight $`r_W=1`$).
## 3.3 Four-Dimensional Theories
We move now to a closer look at the case $`d=4`$. The number of supercharges is $`8`$, corresponding to $`N=2`$ supersymmetry in four dimensions.
Coupling-Moduli Pairing
The (c,c) ring of the $`LG_W/\mathrm{\Gamma }`$ CFT has a $`_2`$ “reflection” symmetry, relating operators with $`U(1)`$ charges $`f`$ and $`\widehat{c}f`$, where $`3\widehat{c}`$ is the central charge of the CFT:
$$\widehat{c}=\underset{a}{}(12r_a)=(n1)2r_\mathrm{\Omega }.$$
(3.22)
To identify this symmetry, one uses two bijective relations between the (c,c) ring and the (a,a) ring – the ring of antichiral-antichiral primary operators. One relation is charge conjugation and the other is obtained using the spectral flow of the $`N=2`$ superconformal algebra (the existence of this second relation is a consequence of the unimodularity of the elements of $`\mathrm{\Gamma }`$ ). Combining these two, one obtains a bijection in the (c,c) ring. As to the $`U(1)`$ charge, starting with a (c,c) operator with charge $`f`$, complex conjugation gives an (a,a) operator with charge $`f`$ and then the spectral flow, gives a (c,c) operator with charge $`\widehat{c}f`$.
This reflection symmetry is not specific to models related to four-dimensional singularities. However, its significance is enhanced for $`d=4`$, as we now explain. In this case, $`\widehat{c}=2(1r_\mathrm{\Omega })`$ (using eq. (3.22) with $`n=3`$). Comparing to the bounds (3.16),(3.17), one can see that, for operators with $`f\widehat{c}/2`$, the reflection symmetry induces a pairing between coupling deformations and moduli deformations. This pairing has a natural interpretation in the 4D theory (as was also observed independently in ). The deformations defined by (c,c) operators are related, in the 4D theory to scalars in vector superfields. When the deformation parameter $`\mu `$ is a coupling with $`R_+(\mu )<2`$, it corresponds to adding a top component $`𝒜_t`$ of a vector superfield $`𝒜`$ to the prepotential (with $`\mu `$ as a coefficient), while if it is a modulus with $`R_+(\mu )>2`$, the change is in the vev of a bottom component $`𝒜_b`$ of a vector superfield $`𝒜`$. Thus, each vector superfield $`𝒜`$ in the 4D theory defines two deformations - one coupling and one modulus – and it is natural to identify these pairs with the pairs seen in the worldsheet formulation.
Evidence for this identification is obtained by considering R-charges. The 4D $`N=2`$ supersymmetry algebra has a $`U(1)\times SU(2)`$ R symmetry group. In the previous subsections, a $`U(1)_+\times U(1)_{}`$ R-symmetry group was found (with charges (3.18). To identify the relation between these two groups, one notes that the scalars in vector superfields are neutral under the $`SU(2)`$ factor of the R-symmetry and charged under the $`U(1)`$ factor. The charges found in subsection 3.2 (eq. (3.20)) imply that $`R_+`$ should be identified with the $`U(1)`$ factor and $`R_{}`$ with a $`U(1)`$ subgroup of the $`SU(2)`$ factor. Now, considering a pair $`\mu _c,\mu _m`$ of deformation parameters, related by the reflection symmetry described above, one finds in both approaches that the sum of their R-charges is four, providing evidence for their correspondence to the same vector superfield $`𝒜`$. In the worldsheet formulation this follows from eq. (3.20), while in the 4D field theory this follows from<sup>15</sup><sup>15</sup>15The $`U(1)`$ factor of the R-symmetry is defined only up to a shift by a $`U(1)`$ symmetry that commutes with the supercharges. However, the relations (3.23) are invariant under such a shift, so they indeed provide evidence for the proposed interpretation of the pair of deformations.:
$$R_+(\mu _m)=R_+(𝒜_b),R_+(\mu _c)=R_+(𝒜_t)=[R_+(𝒜_b)4].$$
(3.23)
IR Conformal Dimensions and Unitarity
In the extreme IR limit, one obtains an $`N=2`$ superconformal field theory. The 4D $`N=2`$ superconformal algebra (SCA) includes a $`U(1)\times SU(2)`$ R-symmetry and implies a relation between the R-symmetry quantum numbers of an operator and its conformal dimension . In particular, the bottom component of a vector superfield is a chiral primary field of the SCA, a scalar of the $`SU(2)`$ R-symmetry and a Lorentz scalar. For such a field, the conformal dimension is $`D={\scriptscriptstyle \frac{1}{2}}R_+^{\prime \prime }`$, where $`R_+^{\prime \prime }`$ is the charge of the $`U(1)`$ R-symmetry. Identifying $`R_+^{\prime \prime }`$ with the charge $`R_+`$ (3.19)<sup>16</sup><sup>16</sup>16This identification will be discussed further at the end of this section., one can determine the conformal dimensions of the 4D vector superfield related to a deformation defined by a worldsheet (c,c) operator $`V`$: if $`\mu =\mu _c`$ is a coupling with $`R_+(\mu )<2`$,
$$D(𝒜_b)=\frac{1}{2}R_+(𝒜_b)=\frac{1}{2}[R_+(𝒜_t)+4]=\frac{1}{2}[R_+(\mu _c)+4]=\frac{f_V1}{r_\mathrm{\Omega }}+2$$
(3.24)
and if $`\mu =\mu _m`$ is a modulus with $`R_+(\mu )>2`$,
$$D(𝒜_b)=\frac{1}{2}R_+(𝒜_b)=\frac{1}{2}R_+(\mu _m)=\frac{1f_V}{r_\mathrm{\Omega }},$$
(3.25)
(where, in the last stage, we used (3.20)).
As observed in , the bound (3.16) for deformations that are couplings gives, when substituted in (3.24)), $`D(𝒜_b)>1`$, which is (almost; see below) the unitarity bound on the conformal dimension of $`𝒜_b`$. The same is true for deformations that are moduli: the bound is (3.17) and when it is substituted in (3.25)), it also gives<sup>17</sup><sup>17</sup>17This was observed independently also in . $`D(𝒜_b)>1`$.
Deformations With $`R_\mathbf{+}\mathbf{(}\mu \mathbf{)}\mathbf{=}\mathrm{𝟐}`$
So far, we only discussed the identification, in the 4D theory, of deformations with $`R_+(\mu )2`$. We turn now to those with $`R_+(\mu )=2`$. The precise unitarity bound is $`D(𝒜_b)1`$ so, for $`R_+(\mu )=2`$, it allows both possibilities – a coupling and a modulus. In either case, this leads to the identification of a vector superfield $`𝒜`$ with $`D(𝒜_b)=1`$. The SCA implies that such a superfield is free. The geometric formulation suggests that $`\mu `$ is a coupling (as described in the previous subsection) so its natural interpretation would be as that for $`R_+(\mu )<2`$: a coefficient of a term $`𝒜`$ in the superpotential. However, as observed in (using $`D(𝒜)=1`$ and the SCA), the corresponding contribution to the Lagrangian is a total derivative and, therefore, has no effect. Instead, we propose, following , that in this case, $`\mu `$ corresponds to a vev (of a bottom component) of a vector superfield $`𝒜`$ with the following properties:
* $`𝒜`$ is frozen in the IR limit, i.e., its kinetic energy diverges in this limit;
* $`𝒜`$ couples to a conserved current which is non-trivial in the IR SCFT, i.e., the corresponding symmetry, which becomes a global symmetry in the IR, acts non-trivially in this SCFT.
The simplest realization of such a situation is $`N=2`$ SQED: a single vector multiplet (leading to a $`U(1)`$ gauge theory) with a charged massless hypermultiplet. Because of the non-trivial charge, the gauge coupling vanishes in the IR limit, leading to a divergent kinetic energy.
According to the above proposal, $`\mu `$ is a coupling in the IR SCFT, in agreement with the analysis in the geometric formulation. However, it originates from a modulus. Moreover, it may become a modulus also in the IR, upon deformation: for example, in SQED, deforming the theory by giving a vev $`\mu `$ to (the bottom component of) the vector superfield leads to a massive hypermultiplet. In the deformed theory, the gauge coupling does not vanish in the IR and, consequently $`\mu `$ remains a vev of a fluctuating (although free) field. This mixed nature of $`\mu `$ is, in fact, suggested also by the $`\varphi `$ dressing in the worldsheet formulation, as described in the previous subsection. For $`\beta >\frac{Q}{2}`$, the wave function is exponentially supported at $`\varphi \mathrm{}`$ and vanishes at $`\varphi \mathrm{}`$, while for $`\beta <\frac{Q}{2}`$, the situation is reversed. The present case corresponds to $`\beta =\frac{Q}{2}`$ and it is special in having support at both regions.
The above identification also leads to the correct conformal dimension for $`\mu `$. This can be argued as follows : since a conserved charge is dimensionless, the corresponding conserved current has dimension three, so the vector field coupling to it has dimension one; the SCA now implies that the bottom component of the corresponding superfield has also dimension one. We end with some comments:
* The dimension of the parameter $`\mu `$ is, by definition, $`D(\mu )=D(𝒜_b)`$ if it is a modulus; and $`D(\mu )=4D(𝒜_t)`$ if it is a coupling. In both cases, this gives
$$D(\mu )=\frac{1}{2}R_+(\mu ),$$
(3.26)
therefore,
$`\mu `$ is a coupling $`D(\mu )1`$ .
* As explained above, a free vector superfield $`𝒜`$ (with $`D(𝒜)=1`$) leads only to a single deformation of the theory. In the worldsheet formulation this means that one should not expect to identify, in the $`LG_W/\mathrm{\Gamma }`$ CFT, a pairing between (c,c) operators with charge $`f=\widehat{c}`$. Indeed, one finds that the number of such operators is not always even.
* It was suggested in that a coupling $`\mu `$ with $`D(\mu )<1`$ can also be identified with a vev of a vector superfield $``$ that, in the IR limit, is frozen but dues not couple to a non-trivial conserved current. Such a superfield would couple to the SCFT through a term $`𝒜/\mathrm{\Lambda }^\delta `$ in the prepotential, where $`𝒜`$ is an interacting vector superfield in the SCFT and $`\mathrm{\Lambda }`$ is some scale in the underlying theory (note that $`\delta =D(𝒜)1>0`$, so this is an irrelevant interaction). The parameter $`\mu `$ is then identified as $`\mu =_b/\mathrm{\Lambda }^\delta `$. With this identification, all the deformations are identified with vevs of vector superfields: interacting for $`D(\mu )>1`$ and frozen for $`D(\mu )0`$.
Relevance
There is another property of the deformation that can be deduced from the dimension $`D`$ of the deformation parameter $`\mu `$: relevance. Marginal deformations correspond to $`D(\mu )=0`$, relevant deformations – to $`D(\mu )>0`$ and irrelevant – to $`D(\mu )<0`$. This distinction can also be identified in the worldsheet formulation<sup>18</sup><sup>18</sup>18In the geometric formulation, for deformations $`\delta W=\mu V`$, relevance means dominance of $`V`$ at $`z0`$: for a relevant deformation ($`D(\mu )>0`$), $`r_V<1`$ (see eq. (3.21)), so $`V`$ dominates over $`W`$ at $`z0`$ and makes a macroscopic change in the singularity; for a marginal deformation ($`D(\mu )=0`$), $`\mu `$ parametrizes a continuous change of the singularity; and for an irrelevant deformation, $`V`$ is negligible at $`z0`$ and has no effect on the singularity.: as described in the introduction, the coordinate $`\varphi `$ parametrizes the distance from the singularity. As in the holographic dualities with $`AdS`$ vacua, radial motion in the near-horizon geometry corresponds, in the decoupled theory, to RG flow (see also ): large distances ($`\varphi \mathrm{}`$) correspond to high energies (UV) and small distances ($`\varphi \mathrm{}`$) – to low energies (UV). Therefore, a vertex operator (3.13) with $`\beta <0`$ defines a perturbation that increases in the IR – a relevant perturbation. Similarly, $`\beta <0`$ corresponds to an irrelevant perturbation and $`\beta =0`$, to a marginal one. Since
$$D(\mu )=\frac{2\beta }{Q}$$
(see eqs. (3.26), (3.20)), these two approaches to identify relevance give the same result.
To summarize, we have made the identification $`R_+^{\prime \prime }=R_+`$ and it was used to reproduce, using the worldsheet formulation, two bounds in the decoupled IR theory: the unitarity bound and relevance. The supercharges indeed satisfy $`R_+^{\prime \prime }=R_+`$, so the non-trivial content of this identification is the charge of the vector superfields (which are the only superfields discussed in this work). Any one of the bounds could be used to deduce this equality. Then, the agreement in all other aspects, as described above, serves as additional evidence for the proposed holographic relation.
4. IR Fixed Points in 4D $`N\mathbf{=}\mathrm{𝟐}`$ SUSY Gauge Theories
In this section, we restrict attention to configurations relevant for the study of 4D $`N=2`$ supersymmetric gauge theories. Specifically, we consider SQCD: $`SU(N_c)`$ gauge theory with $`N_f`$ hypermultiplets (“quarks”) in the fundamental representation of the gauge group.
## 4.1 Realization of SQCD in String Theory
4D $`N=2`$ SQCD can be realized in string theory as follows (for a review, see ): one considers type IIA string theory with the following brane configuration: all branes are extended in the (0123) direction; there are two NS5 branes which are also extended in (45), and $`N_f`$ D6 branes which are also extended in (789); finally, there are D4 branes extended also in the (6) direction over finite intervals, ending on the other branes. When there are $`N_c`$ coinciding D4 branes extending between the two NS5 branes, the low-energy dynamics of the D4 branes is described by 4D $`N=2`$ SQCD, as defined above . The vector multiplets correspond to open strings between the D4 branes, and the hypermultiplets correspond to open strings between the D4 branes and the D6 branes. The dynamics of the NS5 and D6 branes is considered “frozen” in the 4D theory, since these branes have infinite extension in the “internal” directions. In particular, the $`U(k)`$ global (“flavor”) symmetry is the frozen gauge symmetry of the D6 branes. The relative location of these branes determines the parameters of the gauge theory. In particular, the (45) locations of the D6 branes determine the mass parameters of the hypermultiplets. In a given NS5-D6 configuration, the possible locations of the D4 branes (consistent with supersymmetry) parametrize the moduli space of the gauge theory. In particular, distributing the D4 branes in the (45) directions (along the NS5 branes), corresponds to the Coulomb branch, parametrized by vacuum expectation values for the scalars in the vector multiplets.
Lifting this configuration to M-theory , the D6 branes are identified as KK monopoles, corresponding to a non-trivial 4D transverse space $``$ (including the directions (456)) – the multi Taub-NUT space . One of the complex structures of $``$ is that of a hypersurface
$$zw=Q(x)$$
(4.1)
in $`^3`$, where $`Q(x)`$ is a polynomial of degree $`N_f`$ with the coefficients related to the mass parameters of the hypermultiplets. In configurations corresponding to the Coulomb branch of the gauge theory, the NS5 and D4 branes combine to a single M5 brane, wrapped on a 2D Riemann surface $`\mathrm{\Sigma }`$, embedded holomorphicly in $``$. This surface is identified with the Seiberg-Witten (SW) curve of the gauge theory <sup>19</sup><sup>19</sup>19This curve was determined, using field-theoretic considerations, in (for pure $`SU(N_c)`$ SYM) and (for SQCD).. In the representation (4.1) of $``$, $`\mathrm{\Sigma }`$ is the curve $`H=0`$, where<sup>20</sup><sup>20</sup>20Substituting $`w=Q(x)/z`$ and $`z=y+P(x)`$ in $`H=0`$, one obtains the familiar form
$$y^2=P(x)^2gQ(x).$$
$$H=z+gw2P(x).$$
(4.2)
Here $`g`$ is a function of the gauge coupling and $`P(x)`$ is a polynomial of degree $`N_c`$, with the coefficients being the moduli of the Coulomb branch.
Compactifying the (7) direction, one obtains back type IIA string theory, this time with $`N_f`$ KK monopoles – corresponding to the space $`^{3,1}\times \times `$ – and an NS5 brane extended in $`^{3,1}`$ and wrapped on $`\mathrm{\Sigma }`$. This is a situation of the type considered in the introduction.
Naively, the above relations seem to suggest that the decoupled dynamics on this NS5 brane is described by $`N=2`$ SQCD. This is not quite so, because the second description is valid in a range of parameters (of string theory) which is different from that in which the gauge theory was identified in the first configuration. However, there is evidence that some aspects of the low energy dynamics ,including those studied below, are not sensitive to the above changes and, therefore, are shared by the gauge theory and the dynamics of the NS5 brane. Moreover, there is another chain of dualities relating gauge theories (realized this time in heterotic string theory) and the dynamics of NS5 branes in the same sense as above (see for reviews).
## 4.2 Interacting IR Fixed Points in the Moduli Space
In most of the vacua in the Coulomb branch of $`N=2`$ SQCD, the massless fields are vector superfields corresponding to an Abelian gauge symmetry and, possibly, additional electrically-charged hypermultiplets. In these vacua, the dynamics is free in the IR. However, there are vacua with additional massless fields, for which the IR dynamics is non-trivial, defining an interacting superconformal field theory (SCFT). This is the case for SQCD with $`N_f=2N_c`$ massless quarks, at the origin of the Coulomb branch, where the additional massless fields can be identified, at weak coupling, as vector superfields, enhancing the gauge symmetry to a non-Abelian group. Other non-trivial SCFT’s are obtained at points in the Coulomb branch in which mutually non-local hypermultiplets become massless
All these vacua correspond to a SW curve $`\mathrm{\Sigma }`$ with an isolated singularity so, in the stringy realization, there is an NS5 brane wrapped on a surface $`\mathrm{\Sigma }`$ with an isolated singularity. As explained in section 1.1, the only part of the configuration that is relevant for the decoupled dynamics on the NS5 brane is the neighborhood of the singular point. This implies that the relevant KK monopoles are those that are at the singularity. With the singular point in $`\mathrm{\Sigma }`$ being at $`x=0`$, $`Q(x)`$ in (4.1) can be represented by $`Q(x)=x^k`$ ($`kN_f`$), corresponding to $`k`$ coinciding monopoles; in the gauge theory this corresponds to $`k`$ massless quarks. For the same reason, only the region near the center of the KK monopoles is relevant and there the geometry is that of an orbifold $`^2/_k`$. This orbifold can be parametrized by two flat coordinates $`z^{},w^{}`$, subject to the identification
$$z^{}\alpha z^{},w^{}w^{}/\alpha ,\alpha =e^{2\pi i/k},$$
(4.3)
and the coordinates $`z,w,x`$ in (4.1) are $`_k`$-invariant functions of $`z^{},w^{}`$:
$$z=z^k,w=w^k,x=z^{}w^{}.$$
(4.4)
To summarize, we can consider type IIA string theory on $`^{3,1}\times _{z^{}w^{}}^2/_k\times `$ with an NS5 brane on $`^4\times \mathrm{\Sigma }`$, where $`\mathrm{\Sigma }`$ is the 2D surface $`H(z,w,x)=0`$ in $`_{z^{}w^{}}^2/_k`$. This has, as a dual description (see footnote 5), type IIB string theory on $`^{3,1}\times X^6`$, where $`X^6`$ is a hypersurface $`W=0`$ in $`_{z^{}w^{}}^2/_k\times _{uv}^2`$, with the $`_k`$ action (4.3) and the polynomial
$$W=H(z,w,x)+uv.$$
(4.5)
The surface $`\mathrm{\Sigma }`$ ($`H=0`$) has an isolated singularity at the origin and, therefore, so does $`X^6`$ ($`W=0`$). For a quasi-homogeneous $`H`$, this looks like a configuration of the type considered in the previous section (with $`d=4`$, $`n=3`$ and $`\mathrm{\Gamma }=_k`$). Actually, there is a slight difference in the configurations but, as we shall argue below, it is irrelevant for the study of the 4D gauge theory. The difference is the following: the orbifold considered in the previous section has a non-singular worldsheet description and leads here to a $`U(1)^k`$ 6D gauge symmetry. In the present construction, one has in mind an orbifold with an enhanced $`U(k)`$ gauge symmetry, the extra massless states being related (in the type IIA picture) to D2 branes wrapped on the vanishing cycles. The difference in the backgrounds is that in the first one there is a non-vanishing $`B`$-field, leading to non-vanishing masses for the above wrapped D2 branes . However, in the present context, i.e., concentrating on the 4D dynamics, the 6D gauge symmetry on the orbifold is considered frozen anyway, therefore, it is reasonable to expect that the above difference in the 6D dynamics has no effect on the 4D theory. Therefore, one can identify the present configuration as one of those considered in the previous section.
Study of the Interacting SCFT’s
Using the above identification, one can use the anti-holographic description proposed in section 1.1 to obtain information about the interacting 4D SCFT’s. The general idea is to study deformations of the theory and through them, to identify operators in it and to obtain information about their dynamics. We concentrate, as in section 1.1, on SUSY-preserving deformations, defined in the worldsheet formulation by (c,c) operators of the $`LG_W/_k`$ CFT. As described in subsection 3.3, these are related to vector superfields in the 4D theory and this relation depends on the dimension $`D={\scriptscriptstyle \frac{1}{2}}R_+`$ of the corresponding deformation parameter $`\mu `$, as summarized below:
* $`D(\mu )>1`$:
$`\mu 𝒜_b`$, where $`𝒜_b`$ is the bottom component of a superfield $`𝒜`$, whose dimension is $`D(𝒜_b)=D(\mu )`$; the dimension $`D>1`$ indicates that this superfield is involved in a non-trivial interaction;
* $`D(\mu )=1`$
$`\mu _b`$, where $`_b`$ is the bottom component of a superfield $``$ which, in the IR, is frozen but couples to a non-trivial conserved current;
* $`D(\mu )<1`$:
the deformation is a change $`\mu 𝒜`$ in the prepotential of the 4D SCFT, where $`𝒜`$ is an interacting vector superfield of dimension $`D(𝒜)=2D(\mu )>1`$; $`\mu `$ is, possibly, related to a vev $`\mu /\mathrm{\Lambda }^\delta `$ of a free superfield $``$, which does not couple to a conserved current in the IR SCFT and its only interaction with this theory is through an irrelevant term $`𝒜/\mathrm{\Lambda }^\delta `$ in the superpotential (where $`\mathrm{\Lambda }`$ is a scale in the underlying 4D theory and $`\delta =1D(\mu )>0`$).
The deformations with $`D(\mu )1`$ are expected to appear in pairs ($`\mu _m,\mu _c`$), each pair corresponding to an interacting vector superfield $`𝒜`$, with the conformal dimensions related by
$$D(𝒜)=D(\mu _m)=2D(\mu _c).$$
(4.6)
Some of this analysis can be performed directly in the gauge theory. Indeed, for a deformation of the polynomial $`W`$, the dimension of the deformation parameter can be obtained from eq. (3.21):
$$D(\mu )=\frac{1}{2}R_+^{}(\mu )=\frac{r_\mu }{r_\mathrm{\Omega }}=\frac{1r_V}{r_\mathrm{\Omega }},$$
(4.7)
thus using only geometric information about $`X^6`$: the polynomial $`W`$ and the R-charge of the holomorphic 3-form $`\mathrm{\Omega }`$. This information can be translated to purely gauge-theoretical data: the deformations of $`W`$ (4.5) are the deformations of the SW curve and the SW differential can be expressed as an integral over the holomorphic 3-form $`\mathrm{\Omega }`$ , so they have the same R-charge. Therefore, one can use the SW theory to find dimensions of deformation parameters. This was indeed done in . However, the parametrization of the SW curve is not unique: for example, the polynomials $`P(x)`$ and $`Q(x)`$ could be parametrized either by their coefficients or by their roots, leading to a different set of conformal dimensions. The embedding in string theory and the explicit worldsheet formulation fixes this ambiguity (by identifying the deformation parameters with coefficients of terms in the worldsheet Lagrangian)<sup>21</sup><sup>21</sup>21As will be seen below, for the polynomials $`P(x),Q(x)`$, the correct parameters are the coefficients of $`P(x)`$ and the roots of $`Q(x)`$.. This demonstrates the usefulness of the holographic relation for the study of these IR SCFT’s.
We now analyze several families of SCFT’s, using the worldsheet formulation and compare the results to information from field theory. Note that, in all cases, the superfields $`u,v`$ in eq. (4.5) decouple (they are massive), so $`LG_W=LG_H`$.
## 4.3 Singular Points in SYM
For a singularity away from KK monopoles, one can set $`Q=1`$ in eq. (4.1). These are the singularities appearing in pure $`N=2`$ $`SU(N_c)`$ SYM theories<sup>22</sup><sup>22</sup>22Unlike the singularities discussed in the next subsection, these singularities appear only in the strong coupling region of the Coulomb moduli space and do not have a semi-classical interpretation..
In this case, $`g`$ in eq. (4.2) can be rescaled to 1 (by rescaling the coordinates $`z,w,x`$ and the parameters in $`P(x)`$). This is a reflection of the fact that the gauge coupling in the pure SYM theory is transmuted to a scale and is not a real parameter in the theory. The singularity appears for $`P1x^l`$ (with $`2lN_c`$). The singular point is at $`z=w=1`$, and expanding around it (with $`z=1+\stackrel{~}{z}`$), one obtains
$$H\stackrel{~}{z}^22(P1)\stackrel{~}{z}^22x^l.$$
(4.8)
In the worldsheet formulation, the $`LG_H`$ CFT is a minimal model. The (c,c) ring is spanned by
$$V=x^{lj},j=2,\mathrm{},l.$$
Applying eq. (3.3) to (4.5),(4.8), one obtains $`r_\mathrm{\Omega }=1/2+1/l`$, which gives , for the spectrum of conformal dimensions (using eq. (4.7)):
$$r_\mu =jr_x=\frac{j}{l}D(\mu )=\frac{2j}{l+2}.$$
This spectrum corresponds to $`[{\scriptscriptstyle \frac{1}{2}}(l1)]`$ interacting vector superfields (where $`[\mathrm{}]`$ denotes the integer part):
$$𝒰_j,D=\frac{2j}{l+2},j=\left[\frac{l}{2}+2\right],\mathrm{},l$$
(4.9)
and, for even $`l`$, an additional IR-free vector superfield (with $`D=1`$).
Choosing $`l=N_c`$ (i.e., the most singular point for a given $`N_c`$), one finds a nice agreement with the SW theory: this singularity occurs at a singe point in the moduli space, which is $`(l1)`$-dimensional, so one expects to find $`l1`$ different relevant deformations of the IR SCFT. Indeed, the present approach leads to $`l1`$ deformations and all of them are relevant: either vevs or relevant couplings ($`D(\mu )<0`$). The SW theory also provides the charges (under the unbroken $`U(1)^{l1}`$ gauge group) of the dyons that become massless at the singular point . In the present case one finds that (for an appropriate choice of duality frame) out of the $`l1`$ gauge fields, $`[{\scriptscriptstyle \frac{1}{2}}(l1)]`$ couple to mutually-non-local dyons and, for even $`l`$, there is one more gauge field with an electrically-charged hypermultiplet. This is in full agreement with the present results: the interacting superfields $`𝒰_j`$ are those coupled to mutually-non-local dyons and the free superfield with $`D=1`$ is the one coupled only electrically.
The case $`l=2`$ is special, since it corresponds to a free IR SCFT. Indeed, this singularity indicates in the gauge theory the appearance of a massless dyon. The low energy theory is SQED: a $`U(1)`$ vector multiplet with an electrically charged hypermultiplet<sup>23</sup><sup>23</sup>23This vector multiplet is related to the UV vector multiplet by a duality transformation .. There is a single deformation of this theory: a change of the vev of the above vector superfield, giving a mass to the dyon. Because of the electric charge, the $`U(1)`$ gauge coupling flows to zero in the IR, so the vector superfield is frozen in the IR and, correspondingly, the deformation parameter has $`D=1`$. This is exactly what was found in the worldsheet formulation.
## 4.4 Singular Points in SQCD
We turn to singularities with a non-trivial flavor symmetry ($`k2`$). The singular point is on the orbifold singularity $`z^{}=w^{}=0`$ and, neglecting sub-leading terms, one obtains
$$H=z+gw+x^l=z^k+gw^k2h(z^{}w^{})^l$$
(4.10)
(corresponding to $`P(x)hx^l`$). We now distinguish between three situations:
* $`k<2l`$
In this case the last term is sub-leading, so $`H`$ in eq. (4.10) becomes
$$H=z^k+gw^k.$$
(4.11)
Note that this is independent of $`l`$, therefore, so is the physics at this singularity. The $`LG_H`$ CFT is a product of two decoupled minimal models (they become coupled by the $`_k`$ orbifold action). In the untwisted sector, the (c,c) ring is spanned by $`_k`$-invariant polynomials in $`z^{},w^{}`$ (which is the same as arbitrary polynomials in $`z,w,x`$), with identifications modulo $`_z^{}H,_w^{}H`$. A possible choice for a basis is
$$V_j=x^j,j=0,\mathrm{}k2.$$
Applying eq. (3.3) to (4.5),(4.11), one obtains $`r_\mathrm{\Omega }=\frac{2}{k}=r_x`$, which gives , for the spectrum of conformal dimensions (using eq. (4.7)):
$$r_\mu =\left(\frac{k}{2}j\right)r_xD(\mu )=\frac{k}{2}j.$$
The twisted sector of a LG orbifold was analyzed in . Applying this analysis to the present case, one finds in each sector one state $`V_i^{}`$ ($`i=1,\mathrm{},2N1`$) with a $`U(1)`$ charge $`f_{V_i^{}}=1r_\mathrm{\Omega }`$, corresponding to $`D(\mu )=1`$. This leads to the identification of $`[{\scriptscriptstyle \frac{1}{2}}(k2)]`$ vector superfields (where $`[\mathrm{}]`$ denotes the integer part)
$$𝒰_D,\{\begin{array}{cc}D=2,3,\mathrm{},\frac{k}{2},\hfill & k\text{ even}\hfill \\ D=\frac{3}{2},\frac{5}{2},\mathrm{}\frac{k}{2},\hfill & k\text{ odd}\hfill \end{array}$$
(4.12)
and $`2[k/2]`$ vector superfields with $`D=1`$.
* $`k=2l`$
In this case, $`H`$ in eq. (4.10) is quasi homogeneous and, for $`gh`$, the surface $`H=0`$ has an isolated singularity (the case $`g=h`$ will be discussed below). As in the previous case, the $`LG_H`$ CFT is a product of two minimal models, the only difference being that they interact through the last term in $`H`$. This difference, however, has no effect on the chiral ring, which is, therefore, the same as for $`k<2l`$.
* $`k>2l`$
In this case, $`H`$ is not quasi-homogeneous. Naively, at least one of the first two terms is negligible (i.e., has a higher weight), however, neglecting it would lead to a non-isolated singularity (at $`z^{}=0`$ or $`w^{}=0`$). A closer look reveals that the first term is leading in the direction $`w^{}=0`$ and the second, in the direction $`z^{}=0`$, so the precise statement is that this kind of singularity cannot be described by a quasi-homogeneous polynomial and, therefore, does not have a worldsheet description of the form discussed in section 1.1. Note, however, that it can be realized indirectly, as a deformation of $`k2l`$ singularities
Comparing to Field Theory
Each of the above singularities appears in SQCD with $`lN_c`$, $`kN_f`$, where $`k`$ of the quarks have the same mass parameter. Moreover, each such singularity can be found in the semi-classical region of an asymptotically-free theory (with $`N_f<2N_c`$ and $`l<N_c`$), i.e., with mass parameters and vevs much larger then the strong coupling scale. There, using semi-classical considerations, it can be identified as the IR limit of $`SU(l)`$ gauge theory with $`k`$ massless quarks, at the origin of the moduli space.
Using this identification, one can demonstrate that the holographic duality considered here is a strong-weak coupling duality, in the sense that there is no situation in which both the gauge theory and the worldsheet CFT are weakly coupled. For $`k<2l`$, the $`LG_H`$ theory is a product of two decoupled minimal models, so the worldsheet CFT is solvable, but the gauge theory is asymptotically free and strongly coupled in the IR. For $`k>2l`$, the gauge theory is free in the IR but the $`LG_H`$ CFT is complicated. Finally, for $`k=2l`$ the situation depends on the interaction strength of the last term in $`H`$ (in eq. (4.10)). It is proportional to $`h/\sqrt{g}`$ (as can be seen by rescaling $`w^{}w^{}/g^{1/k}`$), so the minimal models are decoupled for $`h/\sqrt{g}0`$ while the gauge theory is free in the opposite limit (as will be described below).
As noted above, the singularity with $`k<2l`$ is independent of $`l`$. Moreover, for even $`k`$ it is the same as $`k=2l`$ with $`h0`$. This last relation can be understood also in field theory, considering SQCD with $`N_f=k<2l<2N_c`$ with even $`k`$, where the $`SU(N_c`$) gauge group is broken to $`SU(l)`$ at some scale $`M`$ and further to $`SU(l^{})`$ ($`l^{}={\scriptscriptstyle \frac{1}{2}}k`$) at a lower scale<sup>24</sup><sup>24</sup>24To be able to apply semi-classical consideration, as is done above, both scales should be sufficiently large, compared to the scale representing the $`SU(N_c)`$ gauge coupling. $`M^{}`$. This gives the $`(k,l^{})`$ singularity, with
$$h\left(M^{N_cl}M^{ll^{}}\right)^2.$$
Unbroken $`SU(l)`$ gauge group (corresponding to the $`(k,l)`$ singularity) is obtained in the limit $`M^{}0`$, which indeed implies $`h0`$.
Therefore, there are two families of singularities, each labeled by $`l`$, with $`k=2l`$ ($`l1`$) and $`k=2l1`$ ($`l2`$) respectively<sup>25</sup><sup>25</sup>25The singularity with $`k=2l=2`$ is the same as the $`l=2`$ singularity in SYM and, as there, it corresponds to SQED (see the end of the previous subsection).. Choosing $`N_c=l`$, the singularity appears at a single point in the moduli space, which is $`(l1)`$-dimensional, so, as in the SYM case, one expects $`l1`$ relevant deformations of the IR SCFT, and this is what is found. For even $`k`$, these are the vevs of the $`{\scriptscriptstyle \frac{1}{2}}(k2)=l1`$ vector superfields in (4.12), while for odd $`k`$, these are the $`[{\scriptscriptstyle \frac{1}{2}}(k2)]=l2`$ vevs and one relevant coupling (that with $`D(\mu )={\scriptscriptstyle \frac{1}{2}}`$). The distinction between couplings and vevs indicates that, for even $`k`$, all the $`l1`$ massless vector superfields are interacting, while for odd $`k`$, one becomes free in the IR.
In addition, there are $`k`$ mass parameters for the quarks in the gauge theory. As explained in subsection 4.1, these can be identified with vevs of the vector superfields in the Cartan subalgebra of the 6D $`U(k)`$ gauge theory on the coinciding KK monopoles and these vector superfields are frozen in the 4D theory because of the infinite extent of the KK monopoles in two additional (internal) directions. These vector superfields couple to conserved (flavor) currents, so they should correspond to deformations with $`D(\mu )=1`$. Strictly speaking, only the $`k1`$ non-diagonal flavor currents (those in $`SU(k)`$) have the above interpretation. The deformations corresponding to these currents are those found above in the twisted sectors.
The diagonal flavor current has a somewhat more complicated nature. If the gauge group was $`U(N_c)`$ instead of $`SU(N_c)`$, the diagonal flavor current would couple to the $`U(1)`$ factor. This is reflected also in the stringy embedding. In fact, the original brane configuration, that with D4 branes, realizes classically a $`U(N_c)`$ gauge theory, but it was argued in that in this stringy realization, the vector superfield corresponding to the diagonal $`U(1)`$ factor is frozen by quantum effects. In the configuration with NS5 branes and KK monopoles, the modulus of the $`U(1)`$ factor in 6D $`U(k)`$ gauge group corresponds to a collective motion of all the KK monopoles in the $`x`$ direction, while the corresponding modulus in the 4D $`U(N_c)`$ gauge group corresponds to such a motion of the NS5 brane. Since the quark masses depend on the relative displacement between the KK monopoles and the NS5 brane, the diagonal flavor current couples to a diagonal subgroup of the above two $`U(1)`$ factors. For odd $`k`$, we did not find a deformation corresponding to this current. This suggests that it decouples from the IR SCFT. For even $`k`$, the corresponding deformation is the $`D=1`$ deformation from the untwisted sector. Indeed, the deformation is $`\delta Wx^{l1}_xW`$, which is an $`x`$-translation, corresponding to a relative motion between the NS5 brane and the KK monopoles.
Additional deformations in the field theory are moduli of the Higgs branch, corresponding to vevs of quark scalar fields. We do not find these in the worldsheet description. Their absence can be understood as follows<sup>26</sup><sup>26</sup>26This explanation was suggested to the author by D. Kutasov.: at the singular point, to which the Higgs branch is connected, the quarks are massless. This corresponds, in the geometric formulation, to the singular hypersurface $`X^6`$, and in the worldsheet formulation, to a superpotential independent of the Liouville superfield $`\widehat{\varphi }`$. As explained in subsection 3.1, in this situation the worldsheet formulation suffers from a strong coupling singularity ($`g_s\mathrm{}`$ for $`\varphi \mathrm{}`$) and to avoid this singularity, one should consider the double scaling limit. In this limit, the mass $`m`$ of the quarks is kept non-zero, which means that one actually considers a point in the Coulomb branch near the singular one. By taking the limit $`m0`$, one can obtain information about the theory at the singularity, as was done above, but clearly this limit will miss aspects that appear discontinuously only at the singularity. This is why one does not see the Higgs moduli in this approach.
$`𝑺𝑼\mathbf{(}𝒍\mathbf{)}`$ with $`k\mathbf{=}\mathrm{𝟐}l`$ Massless Quarks
The $`SU(l)`$ gauge theory with $`k=2l`$ massless quarks is a conformally-invariant theory and the (complexified) gauge coupling $`\tau `$ is a modulus of the theory, which can be changed continuously. In particular, for $`\tau i\mathrm{}`$ one obtains a free field theory (at all scales). Therefore, unlike the other non trivial SCFT’s, the one with $`k=2l`$ is continuously related to a free theory, and the information implied by this relation can be compared to that obtained from the stringy embedding. This theory has an $`SL(2,)`$ duality, relating different values of $`\tau `$. The parameter $`g`$ in (4.2) is an $`SL(2,)`$-invariant function of $`\tau `$ (see for more details) and is, therefore, the quantity characterizing the strength of the interaction. Vanishing interaction corresponds to $`g0`$, while for non vanishing $`g`$, the theory is an interacting SCFT.
The origin of the moduli space corresponds to $`P(x)=x^l`$ in (4.2), so $`h=1`$ in (4.10). For $`g=1`$, $`W`$ takes the form
$$W=(z^Nw^N)^2+uv,$$
which has a (non-isolated) $`_2`$ singularity at $`z^N=w^N`$. The corresponding NS5 brane configuration includes two coinciding NS5 branes, which are the T-duals of the above $`_2`$ singularity. Thus, in this case, the decoupled dynamics includes a six-dimensional sector.
To obtain an isolated singularity we, therefore, assume $`g1`$. In fact, we would like to consider $`g1`$, where one can apply semi-classical considerations. Then, the vector superfield $`𝒜`$ in the adjoint representation defines the following gauge-invariant vector superfields
$$𝒰_j=\mathrm{tr}𝒜^j,j=2,\mathrm{},l$$
and their dimension is, classically (i.e., for $`g0`$), $`D(𝒰_j)=j`$. The gauge coupling $`\tau `$ is identified as the coefficient of a term tr$`𝒜^2`$ in the prepotential.
This information from the gauge theory agrees with the results obtained above from the stringy realization. The conformal dimensions found there were independent of the coupling, so should be the same as for $`g0`$ and, indeed, the correct spectrum is obtained (see eq. (4.12)). The deformations $`\delta H=\mu x^j`$ correspond, for $`jl1`$, to a vev $`\mu \mathrm{tr}𝒜^{lj}`$ and for $`jl`$, to a term $`\mu \mathrm{tr}𝒜^{2+jl}`$ in the prepotential. In particular, for $`j=l`$ (corresponding to a term $`\mu \mathrm{tr}𝒜^2`$), $`\mu `$ is a change in the gauge coupling<sup>27</sup><sup>27</sup>27Note that such a marginal coupling (with $`D(\mu )=0`$) does not appear in the other cases considered above, as expected from field theory.. This can also be seen directly in the polynomial $`H`$ (eq. (4.10) with $`h=1`$): the corresponding deformation is $`\delta H=\mu x^l`$ and by rescaling $`w^{}`$, this can be transformed to a change in $`g`$, which indeed represents the gauge coupling.
5. Summary
In this work, type II string theory in a background of the form
$$^{d1,1}\times _\varphi \times U(1)_Y\times 𝒞$$
(5.1)
was interpreted as a $`d`$-dimensional theory. For $`𝒞=LG_W`$ (a 2D $`N=2`$ Landau-Ginzburg SCFT), this $`d`$-dimensional theory was identified in as the description of the decoupled dynamics near an isolated singularity in type II string theory on
$$^{d1,1}\times X^{2n},$$
(5.2)
where $`X^{2n}`$ is a hypersurface $`W=0`$ in a flat space $`^{n+1}`$ (with $`d+2n=10`$). Here this identification was extended to $`𝒞=LG_W/\mathrm{\Gamma }`$ (a Landau-Ginzburg orbifold), in which case, $`X^{2n}`$ is a hypersurface in an orbifold $`^{n+1}/\mathrm{\Gamma }`$. Furthermore, for $`d=4`$, $`\mathrm{\Gamma }=_k`$ and special choices of the polynomial $`W`$, this four-dimensional theory was related to an interacting SCFT appearing in the moduli space of 4D $`N=2`$ SQCD: $`SU(N_c)`$ gauge theory with fundamental quarks. Properties of the $`d`$ dimensional theory were identified in the worldsheet formulation (of the string theory on (5.1)) and contrasted, when it was possible, with information from the geometric formulation (of string theory on (5.2); in section 1.1) and from the gauge theory (in section 3.3). Here we summarize this analysis.
We start with properties that are independent of the details of $`𝒞`$ and depend only on $`𝒞`$ being a 2D CFT with (2,2) supersymmetry and a spectral flow operator relating the (c,c) and (a,a) rings. For any such $`𝒞`$:
* The theory has $`2^{\frac{d}{2}+1}`$ supercharges;
this corresponding to $`N=2`$ supersymmetry in 4 dimensions.
* It has a $`U(1)_+\times U(1)_{}`$ R-symmetry, under which all the supercharges have charges $`|R_\pm |=1`$;
in $`d=4`$, where the R-symmetry group is $`U(1)\times SU(2)`$, $`R_+`$ was identified with the $`U(1)`$ factor and $`R_{}`$, with a $`U(1)`$ subgroup of the $`SU(2)`$ factor;
furthermore, in the 4D IR SCFT, $`R_+`$ was identified as the $`U(1)`$ R-charge appearing in the superconformal algebra (SCA)<sup>28</sup><sup>28</sup>28Analogous identifications can be made also for $`d4`$..
* Each operator in the (c,c) ring of $`𝒞`$ defines a deformation of the theory, parametrized by a continuous parameter $`\mu `$ (which is a coupling of a term in the worldsheet Lagrangian) with $`R_{}(\mu )=0`$, and this deformation preserves the supersymmetry of the theory;
in $`d=4`$, these deformations are related to scalars in vector multiplets (which are, indeed, invariant under the $`SU(2)`$ R-symmetry): $`\mu `$ is either a vev (modulus) of a bottom component $`𝒜_b`$ of a vector superfield $`𝒜`$; or a coefficient (coupling) of a term in the 4D Lagrangian, which is a top component $`𝒜_t`$ of such a superfield (corresponding to a term $`\mu 𝒜`$ in the prepotential);
this implies (using the SCA), that the conformal dimension of the deformation parameter is$`^{\text{28}}`$ (in both cases)
$$D(\mu )={\scriptscriptstyle \frac{1}{2}}R_+(\mu ).$$
(5.3)
* The value of $`R_+(\mu )`$ is related to the $`\varphi `$-dependence of the corresponding worldsheet vertex operator (see eqs. (3.13),(3.20)) and, consequently, provides information about the nature of the deformation: for $`R_+(\mu )>2`$, this is a modulus and for $`R_+(\mu )<2`$, this is a coupling; the coupling is relevant for $`R_+(\mu )>0`$, marginal for $`R_+(\mu )=0`$ and irrelevant for $`R_+(\mu )<0`$;
in the IR CFT, the same information is provided by the conformal dimension of $`\mu `$, where couplings are distinguished from moduli by using the representation theory of the SCA (including the “unitarity bound” $`D(𝒜_b)1`$ for a bottom components of a vector superfields $`𝒜`$);
using the identification (5.3), these two approaches can be compared and one finds identical distinctions.
* Deformations with $`R_+(\mu )2`$, were shown to appear in pairs; for $`d=4`$ this is a coupling-modulus pairing and this is identified with the pair of deformations that are defined by a given vector superfield.
Some of the evidence for the above identifications was obtained by considering specific cases, as described below. However, it is natural to expect that they have a more general range of validity.
We now turn to properties that are different for different $`𝒞`$ factors, corresponding to differences between $`d`$-dimensional theories (with $`2^{\frac{d}{2}+1}`$ supercharges). First, considering $`𝒞=LG_W/\mathrm{\Gamma }`$ in general (including the case $`\mathrm{\Gamma }=1`$ studied in ), the above identifications can be compared with the geometric formulation of the theory, with the following results:
* One is led to the same requirements on the elements of $`\mathrm{\Gamma }`$.
* One finds the same amount of supersymmetry: $`2^{\frac{d}{2}+1}`$ supercharges.
* The $`U(1)_+`$ R-symmetry is identified as a geometric isometry.
* The (c,c) deformations from the untwisted sector are naturally identified as complex structure deformations, induced by changes in the polynomial $`W`$; this identification was shown to be consistent with the multiplicity and the R-charges of the deformation parameters.
* For the above deformations, a distinction between parameters and moduli can be identified also in the geometric formulation and it was shown to agree with the distinction derived from the worldsheet formulation.
* Relevance of the above deformations translates, in the geometric description, to dominance of the change in $`W`$ at $`z0`$.
Finally, we considered specific four-dimensional SCFT’s, appearing as IR limits of SQCD, in singular points of its moduli space. They are labeled by $`(k,l)`$, where $`k`$ is the number of massless quarks at the singularity and $`l`$ is the degree of the singularity. We considered three families of these singularities, corresponding to $`k=0,2l1,2l`$, and chose $`l=N_c`$, (corresponding to the most singular point in the moduli space). For these theories, there is independent information, from field-theoretic considerations, and it was compared to the results form the worldsheet analysis:
* The holographic duality was shown to be a strong-weak coupling duality, in the sense that there is no situation in which both the gauge theory and the worldsheet CFT are weakly coupled.
* The effect of all the relevant and marginal (c,c) deformations (in both the untwisted and twisted sectors) was identified in the underlying gauge theory.
* The deformations expected in the gauge theory are $`N_c1`$ Coulomb moduli, $`k`$ mass parameters (couplings to the flavor currents) and, for $`k=2l`$, a gauge coupling<sup>29</sup><sup>29</sup>29The moduli of the Higgs branch are not expected, as explained in subsection 4.4.. The only deformation that was not identified is the coupling to the $`U(1)`$ flavor current for $`k=2l1`$. Its absence is an indication that this current decouples from the IR SCFT.
* In pure SYM ($`k=0`$), field-theoretic considerations predict the number of interacting superfields (with $`D>1`$), which are those coupled to mutually-non-local charges, and the number superfields (with $`D=1`$) coupled to (non-trivial) mutually local charges. This prediction agrees with the worldsheet results.
As to the other families of SCFT’s, the conformal dimensions (4.12) indicate that, for $`k=2l`$, all the $`N_c1`$ massless vector fields in the Coulomb branch are interacting and, for $`k=2l1`$, one of them is free.
* For $`k=2l`$, the SCFT is continuously connected to a free theory by changing the gauge coupling and the conformal dimensions were found to be independent of such a change. They should, therefore, be the same as in the free theory and indeed they were found to be so.
The above detailed agreement between results obtained using different approaches is a strong evidence for the duality proposed in and in the present work. Additional evidence is found in the study of the $`d=6`$ case, in . This duality can now be used to study the decoupled theory on the singularity, using the worldsheet formulation. In particular, one can calculate correlation functions of observables in the theory, as was done in . This will provide information about 4D $`N=2`$ SCFT. It is also of interest to extend this duality further, e.g., to other gauge groups and to theories with less supersymmetry. All this is left for future study.
Acknowledgment:
I wish to thank my collaborators in this subject, A. Giveon and D. Kutasov, for helpful discussions and exchange of ideas. I am especially grateful to D. Kutasov, for numerous suggestion and comments. I also benefitted from discussions with S. Kachru. This work is supported in part by DOE grant #DE-FG02-90ER40560 |
warning/0001/hep-ph0001189.html | ar5iv | text | # Toward the non-perturbative description of high energy processes E.V. Shuryak Department of Physics and Astronomy, State University of New York, Stony Brook, NY 11794-3800
## Abstract
General implications of existence of non-perturbative scales and hadronic sub-structure for high energy processes are discussed. We propose that the dependence of the cross section of $`\overline{q}q`$ dipoles on their size d should deviate from $`d^2`$ when d becomes comparable to substructure scale. Then we discuss Kharzeev-Levin pomeron model , based on ladder-type diagrams with $`scalar`$ resonances (scalar $`\pi \pi `$ or $`\sigma `$ and the scalar glueball $`G_0`$). This channel is truly unique, because instanton-induced attractive gg interaction leads to unusually small sizes and strong coupling constants of these states, supplemented by unusually large mass scale, $`M_04GeV`$, of the transition boundary to the perturbative regime. As pomeron is a small-size object by itself, these resonances may play a special role in its dynamics. Furthermore, we use more realistic description of the scalar gluonic spectral density without free parameters, and slightly modify the model to get correct chiral limit. We conclude that the non-perturbative part of the scalar contribution to the soft pomeron intercept is $`\mathrm{\Delta }=.05\pm 0.015`$, with comparable contributions from both $`\sigma `$ and $`G_0`$.
1.Significant progress of non-perturbative QCD has been mostly related with approaches based on its Euclidean formulations: lattice simulations, semi-classical theory based on instantons etc. During the last decade we have learned a lot about correlation functions and their spectral densities, hadronic wave functions and form-factors. Dramatically different features of different channels pointed out in was confirmed and studied in details : below we discuss one of the most striking cases, the gluonic $`J^{PC}=O^{++}`$ channel.
However, little of this progress has so far contributed toward understanding of high energy processes. True, it is difficult to translate many of those tools into Minkowski space. Elsewhere we will report some semi-classical calculations aiming to bridge this gap, while we start this Letter with more general discussion of some qualitative ideas related hadronic substructure to high energy scattering.
The most important lesson we learned is that the non-perturbative objects in the QCD vacuum (and inside hadrons) are not some shapeless soft fields, with typical momenta of the order of $`\mathrm{\Lambda }_{QCD}1fm^1`$, as it was assumed in 70’s. Instead we have semi-classical small-size instantons and very thin QCD strings. The instanton radius peaks around the $`\rho 1/3fm`$ (,for recent lattice data see and general review ). String energy (action) is concentrated in a radius of .2 fm (.4 fm) in transverse directions . Both are small compared to typical hadronic size, suggesting a $`substructure`$ inside hadrons.
A snapshot of parton distribution in a transverse plane inside the nucleon should look like indicated in Fig.(1), for different x regions. These parton clusters originate from “scars” in the vacuum, being perturbed by external objects – valence quarks and strings, and therefore they must have the same transverse dimensions. One expects that these images of constituent quarks, diquarks and strings should be best seen at some intermediate x, before hadrons become black disks at very small x (high energies). The non-perturbatively produced sea quarks supposed to be more concentrated inside the constituent quarks<sup>1</sup><sup>1</sup>1This concentration should be enhanced for the polarized part of the sea. Sea quarks are found to be polarized $`opposite`$ to valence quark (and the nucleon), as the instanton-based mechanism demands . , while the string is supposed to be gluonic. A diquark cluster is also believed to be an instanton effect . Finally, strong $`<\overline{q}q>`$ modification inside the nucleon should result in additional small density of sea quarks and gluons filling the whole disk (shown by light grey in Fig.1).
If one prefers to use the language of hadrons rather than fields, existence of two distinct components can be viewed as being due to two different scales for glueball and pion clouds, respectively. However, using hadronic description in transverse plane (or t-channel) is probably not very useful, because such complicated and coherent field configurations as instantons and strings can hardly be discussed well in this way.
These qualitative ideas were discussed in literature for long time. Constituent quarks as clusters were discussed e.g.in , and “scars” of two strings is behind Nachtmann-Dosch model of high energy hadron-hadron scattering. But they are still mostly ignored by high energy practitioners, who think about partons as being randomly distributed inside the hadronic disk.
How can we tell whether such substructure really exists experimentally? The simplest process we have is deep-inelastic scattering (DIS). In the target frame, it can be viewed as a scattering of a dipole-like $`\overline{q}q`$ objects with variable size $`d2/Q`$, where Q is the momentum transfer<sup>2</sup><sup>2</sup>2This estimate works better for longitudinally polarized virtual photons, while for transverse ones large d tail is more significant, see . . The small-d dipoles measure only the average gluonic field in the target<sup>3</sup><sup>3</sup>3For exact definition see ., but dipoles with d comparable to substructure scales $`\rho `$ indicated above should show a nontrivial behavior of their cross section $`\sigma (d)`$ on d. pQCD predicts that $`\sigma (d)/d^2`$ is constant at small d, but above some critical value $`d>d_c\rho `$ we expect this ratio to drop, because such large dipoles start to miss the “black spots” in the target.
Phenomenologically, when HERA data were translated into such dipole cross section similar behavior $`\sigma (d)`$ was indeed found, and at the right scale for d. In principle however, two different effects can be responsible for it. One, emphasized in especially for very small $`x`$ is $`saturation`$ of the cross section, when the unitarity bound (or “blackness”) is reached. Our idea suggest similar behavior of $`\sigma (d)`$, but even at larger $`x10^2`$ where the spots are still “grey” and the dipole cross section is not large enough to need shadowing corrections. In fact, the particular parameterization used in have exactly this feature for $`all`$ x. Further work is needed here to understand real magnitude of both effects, due to shadowing and substructure, respectively.
Another source of information about parton correlations in transverse plane is diffraction. Clearly an inhomogeneous distribution we advocate enhances it, as compared to the homogeneous disk: the blackness of spots in higher, and there are more edges at which diffraction may take place. Studies of how diffraction depends on $`Q^2`$ at HERA is therefore very important. One possible parameterization is hard-plus-soft pomerons , which places the boundary between soft and hard contributions at $`Q^2410GeV^2`$, or $`dd_c`$ we speak about.
More generally, like it happened for form-factors at similar $`Q^2`$ some time ago<sup>4</sup><sup>4</sup>4 Although $`Q^2`$ dependence of form-factors roughly follow perturbative power counting rules, their absolute magnitude significantly exceeds the pQCD predictions. Non-perturbative approaches, such as instanton-based calculation of the pion form-factor , are in quantitative agreement with data. , after new round of works we may be forced to reassess where exactly DIS is truly perturbative. The previous paradigm – the leading twist dominance down to Q as low as $`Q^20.5GeV^2`$ – seems now oversimplified.
Hadron-hadron collisions are of course much more complicated. Studies of inelastic diffraction in $`\pi p`$,pp, pA and collision of two nuclei definitely show very large O(1) fluctuations of the total nucleon cross section. It is clear that it is completely incompatible with the picture of a grey parton disk filled with multiple independent partons, with dozens of degrees of freedom involved. Only very few degrees of freedom may drive those fluctuations. One obvious suspect is the distance between constituent quarks, or the length of the string in Fig.1. Another is intermittent blackness of each constituent quark, presumably related to “twinkling” character of the field strength distribution in the instanton vacuum.
The pomeron itself is known to be a small-size object in transverse plane, as can be inferred from the fact that most of the t-dependence of pp scattering is explained by nucleon form-factors, and also from $`\alpha ^{}(2GeV)^2`$. It means when we see diffractive scattering of two nucleons, those are related with diffraction of their small parts<sup>5</sup><sup>5</sup>5Except at extremely very high energies, when the whole disk becomes black..
The history of pomeron goes back 40 years: it still works very well, new discoveries are being made (such as existence of a pomeron polarization vector ), but its microscopic dynamics still lacks theoretical understanding. pQCD promised to explain the “hard pomeron” , and although recent calculation of the next-to-leading correction may put it in doubt, hopefully there will be a way out <sup>6</sup><sup>6</sup>6Similar things happened before: let me just recall an optimistic example which may be unfamiliar to many high energy physicists. Free energy for high temperature QCD has negative $`O(\alpha _s)`$ correction of modest magnitude: but higher order ones show much larger corrections of different sign. However (unlike the pomeron intercept) one can calculate this free energy non-perturbatively on the lattice. The result is about 15% below free gas, close to the $`O(\alpha _s)`$ term. All high order corrections apparently canceled out! Furthermore, there are indications that some resummed or improved perturbation theory exists, in which those cancellations happen explicitly. . Non-perturbative approaches include incarnations of the old multi-peripheral model (e.g. ): but they neither provide clear cut explanations of why particular hadrons should be used, nor give convincing quantitative predictions. They also have no connection with Law-Nussinov 2-gluon exchange model, which seems to be a very natural starting point explaining constant (not growing with s) part of the cross section.
New model for growing cross section and soft pomeron have been recently proposed by Kharzeev and Levin (KL) . It includes (i) a ladder made of two t-channel longitudinal gluons, as in perturbative approach; (ii) while the s-channel “rungs” of the ladder being replaced by production of $`scalar`$ physical states. Their main motivation was to apply some known non-perturbative matrix elements, such as gluon-to- $`\pi \pi `$ transition near threshold. Furthermore, they inserted “realistic” scalar spectral density instead of perturbative one below some mass $`M_0`$, suggested a schematic model for it and estimated a resulting value for the pomeron intercept<sup>7</sup><sup>7</sup>7 $`\mathrm{\Delta }`$ enters the total hadronic cross section energy dependence as $`\sigma (s)s^\mathrm{\Delta }`$. $`\mathrm{\Delta }0.08`$, close to the phenomenological value . Their input was however very limited and therefore they have treated gluonic scalar spectral density in a very schematic way. It explains correctly qualitative features of the result (e.g. its dependence on the number of colors and flavors $`N_c,N_f`$), but one may question the accuracy of the resulting numbers. Below we (i) provide further motivation for the KL approach, and also (ii) improve on their crude schematic model in several respects. We explain why one should modify the expression (1) for the pomeron intercept, and get its meaningful chiral limit. We use other information about scalar spectral density, with different glueball parameters. We found that quark-related (pion) contribution to the pomeron intercept $`\mathrm{\Delta }`$ is reduced, and the glueball one is enhanced compared to KL numbers, making them comparable at the end.
2.To motivate the approach, let us start with the following question: How using hadrons in a ladder diagram can be consistent with the statements made above, that the pomeron exchanges take place between small parts of the colliding hadrons, such as constituent quarks or strings?
Well, it depends: different hadrons have different sizes! For example, multiple papers use chiral Lagrangians and pions propagating $`inside`$ the nucleon<sup>8</sup><sup>8</sup>8 Well known ultimate model of that kind was suggested by Skyrme: in it a nucleon is made out of pions entirely. As noted by Witten, large $`N_c`$ makes the nucleon static and pion field classical. Nevertheless, the $`R_N\mathrm{\Lambda }_\chi >>1`$ condition is still needed to justify the model, and because this parameter is not really large, Skyrme model cannot be very accurate at any $`N_c`$. . It follows from chiral Lagrangian that it can be used for momenta up to the so called chiral scale $`p<\mathrm{\Lambda }_\chi 1GeV`$. Note: it is the small pion size which matters here, not its mass. We will have similar situation for glueballs below.
With this in mind, we can move to the next question: Why is it reasonable to single out scalar $`O^{++}`$ gg channel? (Apart of the fact that we know few related coupling constants.) Because both its prominent resonances – the $`\overline{q}q`$ state $`\sigma `$ and the scalar glueball we call $`G_0`$ – are very small-size. $`\sigma `$ is the pion’s brother and its interaction is covered by the same chiral scale. Remarkably,$`G_0`$ is even smaller, with $`R_{G_0}.2fm`$ . It is the smallest hadron we know, setting a record of its kind. It should be possible to construct some effective $`G_0`$ Lagrangian, applicable below some momentum scale $`M_0`$. (For first attempts to build it see , for discussion of the magnitude of $`M_0`$ see below.) The qualitative reason why it is so compact is very strong instanton-induced attraction in the scalar channel due to small-size instantons. As shown in , in this case diluteness of the instanton vacuum $`(\rho /R)^4(1/3)^4`$ (where $`R1fm`$ is instanton mean separation) is compensated by classical enhancement factor $`(8\pi ^2/g^2(\rho ))^210^2`$.
Other channels do not have this feature. For example, another gg channel one may think of is tensor $`2^{++}`$. However instanton field do not have such component, and so they do not act in it. Consequently the tensor glueball has normal hadronic size, $`R_{2++}.8fm`$ according to , and it would be meaningless to consider its propagation inside a nucleon. As experience with quark vector channels (where the situation is similar) shows, in such cases we have multi-hadron spectral density dual to perturbative one down to low scale $`M_{2++}<<M_0`$.
KL have shown that the scalar channel contribute the following “non-perturbative part” to the pomeron intercept:
$`\mathrm{\Delta }={\displaystyle \frac{18\pi ^2}{b^2}}{\displaystyle \frac{dM^2}{M^6}(\rho _{phys}(M^2)\rho ^{pert}(M^2))}`$ (1)
where $`\rho _{phys}(M^2),\rho ^{pert}(M^2)`$ are physical and perturbative (gg cut) spectral densities respectively. The integrand in (1) is non-zero only for $`M<M_0`$ because at $`M>M_0`$ pQCD works and two spectral densities become identical.
3.Let us discuss properties of scalar spectral density $`\rho _{phys}(M^2)`$, first in gluodynamics and then in QCD.
We use the same normalization of the probing operator as KL
$$\theta _\mu ^\mu =\frac{\beta (g)}{2g}F^{a\alpha \beta }F_{\alpha \beta }^a\frac{bg^2}{32\pi ^2}F^{a\alpha \beta }F_{\alpha \beta }^a;$$
(2)
so that its perturbative spectral density is
$`\rho _{pert}(M)={\displaystyle \frac{1}{4096}}{\displaystyle \frac{b^2g^4(N_{c}^{}{}_{}{}^{2}1)M^4}{\pi ^6}}`$ (3)
The spectral density should obey the low energy theorem
$`{\displaystyle \frac{dM2}{M^2}[\rho _{\mathrm{phys}}(M^2)\rho _{\mathrm{pert}}(M^2)]}=40|\theta _\mu ^\mu (0)|0`$ (4)
As alsomphasized in , this relation has historically provided the first indication that there should be large non-perturbative scale in scalar gg channel. Later more direct and quantitative assessment of this effect was proposed , based on small-size instantons.
In gluodynamics $`\rho (M)`$ is dominated by the contribution of the scalar glueball $`G_0`$, the lightest (and therefore stable) particle of this theory. Although phenomenologically its assignment to the observed scalar resonances is confused by mixing with $`\overline{q}q`$ resonances and is still under debate, both multiple lattice works and the instanton model point toward $`M_{G_0}=1.51.7GeV`$. As mentioned already, even more important is its small size, which leads to a remarkably large coupling constant to gg current<sup>9</sup><sup>9</sup>9We remind the reader that the units in the gluodynamics is traditionally defined by setting the string tension to be the same as in QCD., which according to is
$`\lambda _0=<0|g^2G_{\mu \nu }^2|G_0>16.\pm 2GeV^3`$ (5)
Substituting glueball contribution to the spectral density $`\rho _{G_0}=(b/(32\pi ^2))^2\lambda _0^2\delta (M^2M_{G_0}^2)`$ one finds the following contribution to the pomeron intercept
$`\mathrm{\Delta }_{G_0}.03`$ (6)
This is not yet the complete non-perturbative part: one still has to subtract the “missing” perturbative contribution for $`M<M_0`$.
As in , we determine $`M_0`$ from a “duality” sum rule<sup>10</sup><sup>10</sup>10Rather than from the low energy theorem (4). With the simple sharp cutoff we use one cannot satisfy both, so we select duality because it closer to the integral we ultimately need. However the correlator calculated in is in exact agreement with this theorem.
$`{\displaystyle ^{M_0}}(\rho _{G_0}\rho _{pert}){\displaystyle \frac{dM^2}{M^4}}=0`$ (7)
which ensures that correlator at small distances is not changed by changing from $`\rho _{pert}`$ to $`\rho _{phys}`$. Solving it for $`M_0`$, one gets<sup>11</sup><sup>11</sup>11In KL paper significantly smaller number $`M_02.2GeV`$ was used for QCD. It is dual to only $`\sigma `$ meson contribution, without a glueball. $`M_04/,GeV`$. For a boundary between hadronic and partonic descriptions it is unexpectedly high scale indeed.
The resulting contribution of “missing perturbative states” below $`M_0`$ leads to negative contribution to $`\mathrm{\Delta }`$ of about -.01. In total, we got $`\mathrm{\Delta }_{gluodynamics}.02`$, about twice the value estimated by KL.
4.Now we return to the real world with light quarks, and consider the sigma (or $`\pi \pi `$) contribution. The major input of the KL paper is the $`\pi \pi `$ coupling at small M. It follows from the scale anomaly for chiral Lagrangian
$`\rho _{\pi \pi }^{M0}={\displaystyle \frac{3M^4}{32\pi ^2}}`$ (8)
which is larger than $`\rho _{pert}`$ because there is no $`g^2`$ . The KL schematic model assumed that $`\rho (M)=\rho _{\pi \pi }(M)`$ for $`all`$ $`M<M_0`$ (see the dashed line in Fig.(2)). But this assumption cannot be true for larger masses, because the pions are known to interact strongly in this channel, forming the famous scalar $`\sigma `$ resonance.
Since we know its parameters, low energy $`\pi \pi `$ contribution into the correlation function in question can be easily reconstructed. Consider the isovector vector ($`\rho `$-meson) channel, for which the spectral density is well known from $`e+e`$ collisions and $`\tau `$ lepton decay. In this case the $`\pi \pi `$ contribution at the threshold is trivial: the pion coupling is just the pion charge. Due to pion attractive interaction, the spectral density grows from threshold, till it reaches the peak - the $`\rho `$-meson. Its magnitude can therefore be fixed from the “vector dominance”, using the normalization to the M=0 point and known $`\rho `$ meson width.
Similar “sigma-dominance” should work even better, because this resonance is very wide $`m_\sigma \mathrm{\Gamma }_\sigma `$.
$`\rho _\sigma ={\displaystyle \frac{3\pi ^2M^4}{32}}{\displaystyle \frac{M_{\sigma }^{}{}_{}{}^{4}}{(M^2M_{\sigma }^{}{}_{}{}^{2})^2+M^2\mathrm{\Gamma }_\sigma ^2}}`$ (9)
Naive integration from the threshold ($`2m_\pi `$) gives large contribution to pomeron intercept $`\mathrm{\Delta }_\sigma ^{naive}.09`$, where we have used $`M_\sigma =.6GeV,\mathrm{\Gamma }_\sigma =.4GeV`$. (Coincidentally it is close to what KS got in their paper without sigma resonance, integrating till (their) $`M_0`$.)
However, this large contribution cannot be correct, because it rest heavily on smallness of $`m_\pi `$. In the chiral limit, when quarks and pions become massless, there is no threshold at $`2m_\pi `$, and the KL $`\mathrm{\Delta }`$ simply diverges.
To resolve this problem, we should know the distribution over the so called intrinsic parton transverse momenta<sup>12</sup><sup>12</sup>12Deduced e.g. from $`p_t`$ distribution of Drell-Yan pairs, after perturbative effects due to extra gluon emission at large dilepton mass M are subtracted.. If gluonic partons are indeed nothing else but expansion of classical field of an instanton, their transverse momenta are related to basic non-perturbative scale $`\rho `$, the mean size of QCD instantons. If so, one gets its right magnitude $`<p_t^2>(0.6GeV)^2`$ , and also predicts strong cut-offs, $`both`$ at high and low momenta. The former is because the instanton is a finite-size object, leading to a form-factor $`exp(p_t\rho )`$. The latter happens because the instanton field $`A_\mu ^a\eta _{\mu \nu }^ax_\nu `$ has a vortex-like shape with changing sign, so that its projection to constant (or long-wavelength) field vanishes<sup>13</sup><sup>13</sup>13 In fact, there are experimental indications from diffractive dissociation cross section for transverse virtual photon is dominated by large $`p_t`$, which probably implies that indeed $`xg(x,k_t)k_t^2`$ at small $`k_t`$, see details in .. Therefore, gg collisions with small invariant mass are suppressed.
Referring to quantitative calculation elsewhere , here we simply modify the KL expression of (1) with a logarithmic accuracy, introducing into the KL integral over transverse momentum $`k_t^2`$ of the exchanged gluons, $`𝑑k_t^2/(k_t^2+M^2)^2`$, a finite cutoff $`k_t^{min}`$. Now one of the factors $`1/M^2`$ in (1) is modified, and
$`\mathrm{\Delta }={\displaystyle \frac{18\pi ^2}{b^2}}{\displaystyle \frac{dM^2}{M^4}\frac{(\rho _{phys}(M^2)\rho ^{pert}(M^2))}{(M^2+(k_t^{min})^2)}}`$ (10)
The unphysical divergence in the chiral limit is no longer there. Although sensitivity to the cutoff value is formally logarithmic, it still has significant effect on the intercept. With this modification, and $`(k_t^{min})^2`$ for the $`pair`$ to be equal to single parton $`<p_t^2>`$ mentioned, we find the $`\sigma `$ (or two-pion) contribution to be
$`\mathrm{\Delta }_\sigma |_{realistic}.03`$ (11)
It is now comparable to the glueball contribution discussed above. This outcome is in fact more natural than the KS numbers, 0.08 and 0.01 for $`\sigma ,G_0`$ effects, because they are is $`O(N_f^2/N_c^2)`$ (for $`N_f=2`$) and O(1), respectively.
5.Summary and discussion. Finally, combining all contributions together we get our final estimate of the non-perturbative contribution to pomeron intercept, resulting from $`O^{++}`$ hadronic ladder with $`M<M_04GeV`$ to be
$`\mathrm{\Delta }=.05\pm 0.015`$ (12)
now with our guessed uncertainties.
How meaningful is this result? One still has to add to it the pQCD contribution, which is with the region $`M<M_0`$ in the scalar channel subtracted. With current uncertainties, we do not know its value, and some readers may be disappointed at this point. However, the progress is not zero. First of all, the contribution we discuss in non-perturbative $`O(\alpha _s^0)`$, enhanced by classical instanton effects $`O(\alpha _s^1)`$ compared to perturbative result. Second, it is sufficiently close to the phenomenological value $`\mathrm{\Delta }=0.08`$. Third, it can be experimentally tested by itself.
Let us briefly indicate how it can be done. The discussed model claims that the ladders made of scalar resonances explains most of the growing part of the cross section, namely $`\delta \sigma _NN\sigma _{NN}(s_0)0.05log(s/s_0)`$. It means quantitative predictions about multiplicities of these scalar resonances $`\sigma ,G_0`$, both in inelastic collision and their double diffractive production. Sigmas are wide and distorted by low $`p_t`$ cut we do not undertsnad well: so they are difficult to trace down. The glueball however shows up as relatively narrow resonance $`f_0(1500)`$, found in the double diffractive production. It has good experimental signature: large $`\eta \eta `$ and $`\eta \eta ^{}`$ decay modes. So its growing production with energy may even affect s-dependence of the $`\eta /\pi `$ ratio. But the hottest thing to understand is the azimuthal distribution of nucleons in double diffractive production: in fact the WA102 data (discussed in details in ) show surprisingly different distribution for “glueball-type” resonances (such as $`f_0(1500)`$ we discuss) from that for quark-based mesons.
Acknowledgements The author thanks D.Kharzeev,E.Levin, L.McLerran,M.Strikman and I.Zahed for helpful discussions. The work is partly supported by the US DOE grant No. DE-FG02-88ER40388. |
warning/0001/hep-ph0001039.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The solar and atmospheric neutrino problems have provided two milestones in the search for physics beyond the Standard Model, giving strong evidence for $`\nu _e`$ and $`\nu _\mu `$ conversions, respectively. Although flavour changing neutral current interactions (FCNC) and/or neutrino decays may play an important rôle in the interpretation of the data , we concentrate here on oscillations involving very small neutrino mass splittings, because they provide the simplest picture . Moreover, reconciling the LSND data and/or hot dark matter suggests the existence of a light sterile neutrino . Pending the confirmation of LSND results by future experiments, we choose however to use in this work just the solar and atmospheric data assuming the absence of sterile neutrinos, FCNC and/or of heavier neutrinos.
Within the simplest extensions of the Standard Model in which neutrino masses are introduced by an explicit breaking of lepton–number via the seesaw mechanism there are neutral-current-mediated neutrino decays $`\nu ^{}3\nu `$ , but these decays are extremely slow for the neutrino masses of interest for us here. If neutrino masses arise from the spontaneous violation of ungauged lepton number, the corresponding Goldstone boson brings in the possibility of new, potentially faster, 2-body neutrino decays ,
$$\nu ^{}\nu +J.$$
(1)
However it is well-known that in the simplest well–motivated majoron models the neutrino mass and coupling matrices are proportional, so that decays are highly suppressed *in vacuo*<sup>1</sup><sup>1</sup>1This suppression can be avoided in some models due to a judicious choice of the quantum numbers . In such case the effects discussed here, although important, are not essential. . As a result, neutrino decays become cosmologically and astrophysically irrelevant in these models. Moreover, for the small masses indicated by solar and atmospheric neutrino data, the annihilation channels $`\nu ^{}+\nu ^{}J+J`$ are also negligible.
The purpose of the present work is to discuss the possible impact of neutrino-majoron interactions on the supernova (SN) neutrino signal. Together with the early Universe SN are the only site where neutrinos are in thermal equilibrium and so abundant that neutrino-neutrino interactions become important. Therefore, SN physics has been one of the main tools to derive limits on neutrino parameters in majoron models .
The arguments used in the literature can be divided mainly into two classes: The first uses the fact that majoron interactions violate lepton number (like in $`\nu \overline{\nu }+J`$ or $`\nu +\nu \overline{\nu }+\overline{\nu }`$) and, thereby, can reduce the trapped electron lepton number fraction $`Y_{L_e}`$ in the SN core . Since in several SN simulations it was found that a minimum value of $`Y_{L_e}>0.38`$ is needed in order to have a successful bounce shock, the strength of the lepton number violating interactions can be correspondingly restricted. Unfortunately, the older works using this argument did not take into account the derivative nature of the coupling of a Goldstone boson calculating the scattering cross sections. Moreover, an approximation valid only in vacuum was used for the neutrino decay rates. In ref. the correct medium decay rates were used for the first time, giving the stringent limit $`g_{ee}<5\times 10^6`$.
The second argument used is a “classical” energy loss argument: the observed neutrino signal of SN 1987A should not be shortened too much by additional majoron emission . The most comprehensive work along this line is Ref. which also uses the correct derivative coupling of the majoron. The limits found there depend on the value of the $`(BL)`$ breaking scale. For example the range $`5\times 10^7<g<6\times 10^5`$ is excluded for $`v=20`$ GeV while $`1\times 10^5<g<7\times 10^5`$ is excluded for $`v=500`$ GeV. The tau-neutrino masses considered in Ref. are, compared to the values currently discussed, rather large, 100 eV $`<m_{\nu _\tau }<30`$ MeV. It is therefore of certain interest to extend their discussion to lower neutrino masses presently indicated by the solutions to the solar and atmospheric neutrino anomalies.
These astrophysical limits should be confronted with the available laboratory constraints. While there is a stringent limit from $`\beta \beta `$ experiments ,
$$\underset{i,j}{}g_{ij}U_{ei}U_{ej}<3\times 10^5,$$
(2)
the limits from pion and kaon decays are rather weak,
$$\underset{l=e,\mu ,\tau }{}g_{el}^2<3\times 10^5\mathrm{and}\underset{l=e,\mu ,\tau }{}g_{\mu l}^2<2.4\times 10^4.$$
(3)
Note also that individual couplings $`g_{ij}`$ could be larger than the limit from $`\beta \beta `$ experiments due to possible cancellations .
The main purpose of the present work is the study of the impact of neutrino-majoron interactions on the observable neutrino signal of a supernova. Since decays of the type $`\overline{\nu }_e\nu _l+J`$ reduce the $`\overline{\nu }_e`$ flux, a limit on the neutrino majoron coupling constants can be derived from the observed signal of SN 1987A. In so–doing one must include in the analysis the fact that massive neutrinos may oscillate on their way from the SN envelope to the detector. We do that and present the excluded regions for the three currently discussed solutions of the solar neutrino problem . Then we estimate the discovery potential of new experiments like Superkamiokande (SK) and the Sudbury Neutrino Observatory (SNO) in the case of a future galactic supernova. We find that these experiments could probe majoron neutrino coupling constants $`g`$ down to $`g>\mathrm{few}\times 10^5`$.
## 2 Matter effect on neutrino-majoron interactions
We consider the simplest class of models in which neutrinos acquire mass from the spontaneous violation of ungauged lepton number . In this case it is well-known that the massless Goldstone boson $`J`$ – the majoron – couples diagonally to the mass-eigenstate neutrinos $`\nu _i`$ to a very good approximation . In other words, after rotation from the weak basis $`\nu _\alpha `$ through $`\theta _0`$ angle(s) the original coupling matrix $`g_{\alpha \beta }`$ transforms into
$$g_{ij}\delta _{ij}g_i.$$
(4)
We denote by $`\nu _i^{(h_i)}`$ the 4-spinor describing the majorana neutrino with mass $`m_i`$ and helicity $`h=\pm 1`$.
In this section, we briefly review the effect of a thermal background on the neutrino majoron interactions. It was first demonstrated by Berezhiani and Vysotsky in Ref. that the effective mass induced by the interactions of neutrinos with background matter can break the proportionality between the mass matrix $`m_{ij}`$ and the coupling matrix $`g_{ij}`$ characteristic of the simplest majoron models *in vacuo* . A thermal background consists, except in the early Universe, only of particles of the first generation and, therefore distinguishes the electron flavour from the other flavours.
In Ref. , the Lagrangian describing neutrino majoron interaction in a thermal background was obtained in a relativistic approximation similar to the usual treatment of neutrino oscillations in matter. Later, the authors of Ref. solved this problem without using this approximation and confirmed the results of Ref. in the appropriate limiting cases. Here we are concerned with SN neutrinos, which have typical energies around $`1025`$ MeV. Therefore, we will take advantage of the simpler relativistic approximation and we will follow closely Ref. .
The Hamiltonian $`H_{\mathrm{tot}}H(x)`$ describing the evolution of neutrinos may be given as
$`H_{\mathrm{tot}}`$ $`=`$ $`H_0+H_{\mathrm{med}}+H_{\mathrm{int}}`$ (5)
$`=`$ $`H_0+{\displaystyle \underset{i,j}{}}{\displaystyle \underset{h_i,h_j}{}}\overline{\nu }_i^{(h_i)}V_{ij}^{h_i,h_j}\nu _j^{(h_j)}+g_{ij}\overline{\nu }_i^{(h_i)}\gamma _5\nu _j^{(h_j)}J,`$ (6)
where the free Hamiltonian $`H_0`$ describes the propagation in vacuo, $`H_{\mathrm{med}}`$ describes the effects of matter and $`H_{\mathrm{int}}`$ takes into account the presence of neutrino-majoron interactions which may lead to decays. Instead of using the eigenstates of the free Hamiltonian $`H_0`$ as basis for perturbation theory, we will use the eigenstates of $`H_0+H_{\mathrm{med}}`$. Therefore, as a first step the Dirac equation, $`(H_0+V)\nu _j^{(\pm )}=i_t\nu _j^{(\pm )}`$ has to be solved<sup>2</sup><sup>2</sup>2The majoron remains massless because the forward scattering amplitude of Goldstone bosons on matter vanishes.. In the case of ultra-relativistic neutrinos, it is well-known that a vector-like potential does not change the helicity of the neutrinos — its only effect is a rotation of the eigenstates of $`H_0+V`$ with respect to the mass basis . Therefore, the Dirac equation simplifies and reduces to the standard form known from neutrino oscillations,
$$i_t\nu _i^{(h)}=(H_{ij}^{\mathrm{rel}}+U_{i\alpha }V_{\alpha \beta }U_{\beta j}^{})\nu _j^{(h)}$$
(7)
where $`H_{ij}^{\mathrm{rel}}(p+m_i^2/(2p))\delta _{ij}`$ and $`V_{\alpha \beta }`$ is the potential matrix in the weak basis
$$V_{\alpha \beta }=\left(\begin{array}{ccc}V_C+V_N& 0& 0\\ 0& V_N& 0\\ 0& 0& V_N\end{array}\right).$$
(8)
The potentials induced by the charged and neutral currents are $`V_C=\sqrt{2}hG_Fn_B(Y_e+Y_{\nu _e})`$ and $`V_N=\sqrt{2}hG_Fn_B\left(\frac{1}{2}Y_N+Y_{\nu _e}\right)`$, where $`Y_i=(n_in_{\overline{i}})/n_B`$ and $`n_B`$ is the baryon density. Finally, $`U`$ is the mixing matrix relating mass and weak basis and defined through $`\nu _i=U_{i\alpha }\nu _\alpha `$. Diagonalizing $`H^{\mathrm{rel}}+UVU^{}`$ gives the medium states $`\stackrel{~}{\nu }_i^{(h)}=\stackrel{~}{U}_{ij}^{(h)}\nu _j^{(h)}`$.
In the case of a two-flavor neutrino system the mixing matrices, $`U`$ and $`\stackrel{~}{U}`$, can be parametrized by $`\theta _0`$ and $`\theta ^{(h)}`$ respectively. The diagonalization of the Hamiltonian leads us to the following expression for the effective mixing angle,
$$\mathrm{sin}^22\theta ^{(\pm )}(p)=\frac{\mathrm{sin}^22\theta _0}{\mathrm{sin}^22\theta _0+(\mathrm{cos}2\theta _0\xi ^{(\pm )})^2},$$
(9)
where
$$\xi ^{(h)}=h\frac{\mathrm{\Delta }_0}{2pG_F(Y_e+Y_{\nu _e})n_B},$$
(10)
and $`\mathrm{\Delta }_0=m_2^2m_1^2`$. The effective masses in the medium are
$$m_{1,2}^{(h)\mathrm{\hspace{0.33em}2}}=hp(2V_N+V_C)+\frac{1}{2}(m_1^2+m_2^2)\frac{1}{2}\mathrm{\Delta }^{(h)},$$
(11)
where the upper sign is for $`m_1`$, the lower for $`m_2`$ and
$$\mathrm{\Delta }^{(h)}=\sqrt{(\mathrm{\Delta }_0\mathrm{cos}2\theta _0+2hpV_C)^2+(\mathrm{\Delta }_0\mathrm{sin}2\theta _0)^2}.$$
(12)
For typical neutrino energies and densities near the neutrino–spheres, the parameter
$$\xi =6.53\times 10^8\frac{\mathrm{\Delta }_0}{10^3\mathrm{eV}^2}\frac{10\mathrm{M}\mathrm{e}\mathrm{V}}{p}\frac{10^{10}\mathrm{g}/\mathrm{cm}^3}{(Y_e+Y_{\nu _e})\rho }$$
(13)
is much smaller than one. This fact will give rise to the following simplification, $`\theta ^{(+)}(p)\theta _0`$ and $`\theta ^{()}(p)\pi /2\theta _0`$, which, as it will be shown later, will allow us to identify medium and weak interaction states.
In the three-flavor neutrino case the mixing matrix $`U`$ can be parametrized as $`U=U_{12}U_{13}U_{23}U_0`$, where the matrices $`U_{ij}=U_{ij}(\theta _{ij})`$ perform the rotation in the $`ij`$-plane by the angle $`\theta _{ij}`$ and $`U_0`$ includes possible CP-violation effects . In the following we will assume for simplicity CP conservation and $`\theta _{13}=0`$, the latter motivated both by detailed fits of the atmospheric neutrino anomaly, but also by the results of the Chooz experiment . This simplifies the mixing matrix to $`\nu _i=U_{i\alpha }\nu _\alpha =U_{12}U_{23}\nu _\alpha `$ and we can make the assignment, $`\theta _{12}=\theta _{}`$ and $`\theta _{23}=\theta _{\mathrm{atm}}`$. Notice that for light neutrinos near the neutrino spheres the condition $`|V_{\alpha \alpha }|m_i^2/(2p)`$ ($`\xi 1`$) holds and, since in the weak basis the potential is diagonal, the medium states can be identified with the weak ones up to a rotation in the $`\nu _\mu \nu _\tau `$ subspace. The expressions can be simplified by choosing this arbitrary rotation angle to coincide with $`\theta _{23}`$. Then, only one angle, $`\theta ^{(h)}`$, will be required to connect medium and mass eigenstates, and one will be able to *recover* the two-flavor neutrino case,
$$\stackrel{~}{\nu }_i^{(h)}=\stackrel{~}{U}_{ij}^{(h)}\nu _j^{(h)}=\stackrel{~}{U}_{ij}(\theta ^{(h)})\nu _j^{(h)}=\{\begin{array}{ccc}\stackrel{~}{\nu }_i^{(+)}\hfill & =& \stackrel{~}{U}_{ij}(\theta _{12})\nu _j^{(+)}\hfill \\ \stackrel{~}{\nu }_i^{()}\hfill & =& \stackrel{~}{U}_{ij}(\pi /2\theta _{12})\nu _j^{()}\hfill \end{array}.$$
(14)
Thus the medium states can be identified with weak interaction states according to Table 1 and the coupling matrix $`\stackrel{~}{g}_{ij}`$ in the medium basis can be approximated by the one in the weak basis, $`g_{\alpha \beta }`$.
Above, we have derived the transformation matrix $`U_{ij}(\theta ^{(\pm )})`$ between mass and medium eigenstates. The relation between the corresponding coupling matrices *in vacuo* $`g_{if}`$ and *in medium* $`\stackrel{~}{g}_{if}`$ is given by
$$\stackrel{~}{g}_{if}^{h_ih_f}=\stackrel{~}{U}_{in}(\theta ^{(h_i)})g_{nm}\stackrel{~}{U}_{mf}^T(\theta ^{(h_f)})=\left(\begin{array}{ccc}c_fc_ig_{11}+s_fs_ig_{22}& s_ic_fg_{22}c_is_fg_{11}& 0\\ c_is_fg_{22}s_ic_fg_{11}& s_fs_ig_{11}+c_fc_ig_{22}& 0\\ 0& 0& g_{33}\end{array}\right)$$
(15)
with $`c_{i,f}=\mathrm{cos}[\theta ^{(\pm )}(p_{i,f})]`$ and $`s_{i,f}=\mathrm{sin}[\theta ^{(\pm )}(p_{i,f})]`$. Notice that the coupling constant $`g_{33}`$ will only appear in decays not involving electron neutrinos or anti–neutrinos, so that it will not be important for us.
Within our relativistic approximation, only helicity–flipping neutrino decays $`\nu _i^\pm \nu _j^{}+J`$ occur. For these decays, the differential decay rate is
$$\frac{d\mathrm{\Gamma }}{dp_f}=\frac{\stackrel{~}{g}_{if}^2}{8\pi }\left(\frac{p_ip_f}{p_i^2}\right)\left(\frac{m_i^{(\pm )2}}{2p_i}\frac{m_f^{(\pm )2}}{2p_f}\right).$$
(16)
In the next section, we will apply this formula to supernova. Additional simplifications of the limit $`\xi 1`$ are in this case $`m_{i,f}^{(\pm )2}/(2p_{i,f})=V_{i,f}`$ and that the coupling constants do not depend on $`p_{i,f}`$. Hence, we can integrate Eq. (16) and obtain as total decay rate
$$\mathrm{\Gamma }=\frac{\stackrel{~}{g}_{if}^2}{16\pi }(V_iV_f)$$
(17)
and as average energy of the final neutrino $`E_f=E_i/3`$. Obviously, only those decays are possible for which $`V_iV_f>0`$.
Besides enhancing the rates for neutrino decay, the dispersive effects of the medium open also completely new decay channels of the majoron into neutrinos . The majoron decays $`J\nu _i^\pm +\nu _j^\pm `$ have the same matrix elements like the neutrino decays $`\nu _i^{}\nu _j^\pm +J`$ considered above. Their total decay rate is
$$\mathrm{\Gamma }=\frac{\stackrel{~}{g}_{ij}^2}{8\pi }(V_iV_j)S$$
(18)
and now those decays are possible for which $`V_i+V_j<0`$. The symmetry factor $`S`$ is $`S=1/2`$ if the two neutrinos are identical, and $`S=1`$ otherwise. Finally, we note that the neutrinos are emitted isotropically.
## 3 Supernova neutrinos and neutrino majoron decays
Here we first collect the available limits on majoron neutrino coupling constants from the observation of the neutrino signal of SN 1987A and comment on their validity. Then, we discuss the sensitivity of new experiments like SK or SNO to probe neutrino-majoron interactions in the case of a future galactic supernova.
### 3.1 Constraints from collapsing phase
Massive stars become inevitably unstable at the end of their life, when their iron core reaches the Chandrasekhar limit. The collapse of the iron core is only intercepted when nuclear density, $`\rho _03\times 10^{14}`$ g/cm<sup>3</sup>, is reached. At this point, the implosion is turned into an explosion: a shock wave forms at the edge of the core and moves outward. The strength of this bounce shock and its successful propagation is extremely sensitive to the trapped electron lepton fraction $`Y_{L_e}=Y_e+Y_{\nu _e}`$, attained by the core during its infall. A successful SN explosion occurs only if at least 90% of the initial $`Y_{L_e}`$ is still present , which translates to $`Y_L(t_{\mathrm{bounce}})>0.375`$ .
We now use this requirement in order to derive a limit on majoron decays. Such decays clearly change $`Y_{L_e}`$ either by two units ($`\nu _e\overline{\nu }_e+J`$) or by one ($`\nu _e\overline{\nu }_\mu +J`$). Since the allowed change in $`Y_L`$ is small, we can still use the profiles $`Y_e(t)`$ and $`Y_{\nu _e}(t)`$ from a “standard” SN simulation . One can easily check from Table 1 that only the first of the above decays takes place, since $`Y_e+3/2Y_{\nu _e}<1/2`$. Hence the deleptonization rate is governed by
$$\frac{dY_{L,\mathrm{decays}}}{dt}=2\mathrm{\Gamma }(\nu _e\overline{\nu }_e+J)Y_{\nu _e}.$$
(19)
We integrate Eq. (19) numerically from $`t_0`$, the time when neutrinos start to become trapped (at $`\rho (t_0)5\times 10^{11}`$ g/cm<sup>3</sup>), to the time of the bounce $`t_{\mathrm{bounce}}`$. Note that in the numerical simulation whose results we are using $`Y_L`$ increases from $`Y_L(t_0)0.37`$ up to $`Y_L(t_{\mathrm{bounce}})0.39`$. Requiring that $`Y_L(t_{\mathrm{bounce}})>0.375`$ , i.e. $`|\mathrm{\Delta }Y_{L,\mathrm{decays}}|<0.015`$, we obtain
$$g_{ee}=g_{11}\mathrm{cos}^2\theta _0+g_{22}\mathrm{sin}^2\theta _0<2\times 10^6.$$
(20)
Using the same argument, the limit $`g_{ee}<5\times 10^6`$ was derived in Ref. .
This limit relies on numerical modeling of SN explosions, in particular on the success of the explosion in specific models. However, it is probably fair to say that current supernova models have generally problems to produce successful explosions. Therefore, Eq. (20) cannot be viewed as a trustworthy limit as long we do not have a better understanding of SN dynamics.
### 3.2 Constraints from majoron luminosity
Numerical computations of the total amount of binding energy $`E_b`$ released in a supernova explosion yield, within a plausible range of progenitor star masses and somewhat depending on the equation of state used,
$$E_b=(1.54.5)\times 10^{53}\mathrm{erg}.$$
(21)
This range of values is confirmed by likelihood analysis of the observed $`\overline{\nu }_e`$ spectrum of SN 1987A under the hypothesis of small mixing of $`\overline{\nu }_e`$ with other neutrino flavours . Therefore, the parameter space of models which give rise to majoron luminosity large enough that the observed $`\overline{\nu }_e`$ signal is significantly shortened can be restricted. The most comprehensive analysis of SN cooling and majoron emission using this argument was given in Ref. .
Our following analysis has three main differences compared to Ref. : First, we are interested in relatively light neutrinos ($`m^2<\mathrm{eV}^22pV`$) as suggested by the simplest interpretation of data from solar and atmospheric neutrino experiments. Therefore, medium effects become important not only for decays like $`\nu \nu ^{}+J`$ but also for scattering processes like $`J+\nu J+\nu `$. Second, we do not include the process $`J+JJ+J`$ in the calculation of the majoron opacity. A discussion of why self-scattering processes do not contribute to the opacity was given in Ref. for the case of $`\nu +\nu \nu +\nu `$ scattering. Note that trapping due to $`J+JJ+J`$ scattering prevented Choi and Santamaria from excluding majoron–neutrino couplings for the case of light neutrino masses, $`m_{\nu _\tau }<100`$ eV. Third, we calculate the majoron luminosity in the trapping regime in a different way.
Let us discuss first the mean free path of majorons. The main source of opacity for majorons are the processes $`J+J\nu +\nu `$, $`J+\nu J+\nu `$, and $`J+\nu \nu `$ (cf. ). Taking into account the effective mass of the neutrinos, the corresponding mean free paths inside the SN core with radius $`r_010`$ km and density $`\rho \rho _{\mathrm{nuc}}`$ are given by,
$`l^1(J+J\nu +\nu )`$ $`=`$ $`1.9\times 10^{18}g^4\left({\displaystyle \frac{\mathrm{keV}}{m_{\mathrm{eff}}}}\right)^2\left({\displaystyle \frac{T}{25\mathrm{M}\mathrm{e}\mathrm{V}}}\right)^3\mathrm{cm}^1`$ (22)
$`l^1(J+\nu J+\nu )`$ $`=`$ $`1.5\times 10^{18}g^4\left({\displaystyle \frac{\mathrm{keV}}{m_{\mathrm{eff}}}}\right)^2\left({\displaystyle \frac{T}{25\mathrm{M}\mathrm{e}\mathrm{V}}}\right)^3\mathrm{cm}^1`$ (23)
$`l^1(\nu +J\nu )`$ $`=`$ $`1\times 10^4g^2\left({\displaystyle \frac{T}{25\mathrm{M}\mathrm{e}\mathrm{V}}}\right)\mathrm{cm}^1.`$ (24)
For the effective neutrino mass inside the SN core, $`m_{\mathrm{eff}}^22pV`$, we use $`m_{\mathrm{eff}}20`$ keV. Requiring that the mean free path of the majorons is $`>10`$ km, we obtain that they escape freely for $`g<3\times 10^6`$.
Next, we consider the luminosity for the case that majorons are not trapped. Then they are emitted by the whole core volume with luminosity
$`(\nu +\nu J+J)`$ $`=`$ $`5.5\times 10^{80}g^4\left({\displaystyle \frac{\mathrm{keV}}{m_{\mathrm{eff}}}}\right)^2\mathrm{erg}/\mathrm{s}`$ (25)
$`(\nu \nu +J)`$ $`=`$ $`4.8\times 10^{65}g^2\mathrm{erg}/\mathrm{s}.`$ (26)
We require now that the majoron luminosity during 10 s does not exceed the maximal theoretical value of the binding energy, $`_J<5\times 10^{52}`$erg/s, and obtain the limit $`g<3\times 10^7`$. In the trapping regime, the energy loss argument therefore excludes the band
$$3\times 10^7<g<3\times 10^6$$
(27)
for vacuum neutrino masses $`mm_{\mathrm{eff}}20`$ keV.
Let us now discuss the case that majorons are trapped, $`g>3\times 10^6`$. We recall first the standard treatment as presented, e.g., in Ref. . The main assumption is that $`_J`$ can be approximated by blackbody surface emission
$$_J=\frac{\pi ^3}{30}R_J^2T^4$$
(28)
of a thermal majoron–sphere which radius $`R_J`$ is defined to be at optical depth 2/3,
$$_{R_J}^{\mathrm{}}𝑑rl^1(r)=\frac{2}{3}.$$
(29)
Assuming furthermore a temperature profile $`T(r)`$ outside the supernova core, one can determine for which range of coupling constants $`_J`$ exceeds a certain critical value.
There are two shortcomings in this argumentation. First, the profile $`T(r)r^2`$ used for $`r>r_0`$ in represents the temperature of ordinary matter (nucleons, $`e^\pm `$, photons). However, majorons couple mainly to neutrinos and would at best be in thermal equilibrium with neutrinos, if at all, but not with nucleons. Second, the use of Eqs. (22-24) which were derived for isotropic distributions inside the core is not justified for the highly non–isotropic distributions outside the neutrino– or majoron–spheres. Using for the density $`n_i(r)`$ of a particle species $`i`$ with sphere radius $`R_i`$ the expression
$$n_i(r)=\frac{_i}{E_i}\frac{1}{4\pi r^2},$$
(30)
and introducing the additional factor $`1\mathrm{cos}\theta (R_i/r)^2`$ which represents the averaging over the angle $`\theta `$ between the momenta of the radially outgoing test majoron and the particle $`i`$ in the mean free path one finds that $`n_i(r)1/r^4`$. This renders support to our claim that there is a sharp drop in the majoron opacity in crossing the neutrino spheres. Note also that it was recently stressed in Ref. that a naive application of the Stefan-Boltzmann law can be dangerous for the calculation of neutrino–sphere radii.
We prefer therefore not to rely on the argumentation presented above. Instead, we calculate the *volume* rather than the *surface* luminosity, but taking into account that only majorons emitted from the shell $`[r_0l:r_0]`$ can escape . We find that the majoron luminosity drops below $`_{\nu ,\mathrm{tot}}E_b/10`$ s for $`g>2\times 10^5`$. Going even to larger values of $`g`$, majorons and neutrinos are becoming so strongly coupled that it is reasonable to assume that the SN core emits roughly the same luminosity in form of majorons as in one neutrino species, ie. $`_J_{\nu ,\mathrm{tot}}/6`$. Therefore, majoron emission in this regime is not constrained at all by the energy loss argument. Combining the both limits obtained, majoron couplings one concludes that the range
$$3\times 10^7<g<2\times 10^5$$
(31)
is excluded for vacuum neutrino masses $`mm_{\mathrm{eff}}20`$ keV.
A final remark is in order. We have been using rather loosely only one coupling constant $`g`$ in this section. More precisely, we mean by $`g`$ the element of the coupling matrix $`g_{\alpha \beta }`$ in the weak basis with the largest absolute value.
### 3.3 Constraints from neutrino spectra
During the Kelvin-Helmholtz cooling phase, $`t110`$ s after core bounce, the protoneutron star slowly contracts and cools by neutrino emission. The neutrino luminosities are governed by the energy loss of the core and are, therefore, approximately equal for each type of (anti-) neutrinos. Since the opacity of, e.g., $`\overline{\nu }_{\mu ,\tau }`$ is smaller than of $`\overline{\nu }_e`$, due to their smaller cross section, their energy-exchanging reactions already freeze out in the denser part of the protoneutron star. Hence, one expects the spectral temperatures of $`\overline{\nu }_e`$ to be smaller than the one of $`\nu _h=\{\nu _{\mu ,\tau },\overline{\nu }_{\mu ,\tau }\}`$. Typically, the average energies $`E_i`$ found in simulations are
$$E_{\nu _e}11\mathrm{MeV},E_{\overline{\nu }_e}16\mathrm{MeV},E_{\nu _h}25\mathrm{MeV}.$$
(32)
It is convenient to define three energy spheres $`R_{E,\nu _e}`$, $`R_{E,\overline{\nu }_e}`$ and $`R_{E,\nu _h}`$, outside of which only energy-conserving reactions contribute to the neutrino opacity<sup>3</sup><sup>3</sup>3Neutrino-matter interactions are energy–dependent and, therefore, the concept of a neutrino–sphere is evidently an over-simplification.. These reactions, mainly neutral current neutrino-nucleon and neutrino-nucleus scattering, do not change the neutrino spectra, although the neutrinos still undergo several scattering as they diffuse outward. Eventually, the neutrinos reach the transport sphere at $`R_{t,\nu _i}`$, where also the energy-conserving reactions freeze out, and escape freely. This sphere, i.e. the surface of last scattering, is what most authors mean with “the” neutrino–sphere.
For the numerical evaluation of the decay rates, we use profiles for $`\rho (r)`$ and $`Y_e(r)`$ from Wilson’s SN model at the time $`t=6`$ s after core bounce. For all radii of interest, only anti-neutrinos can decay into neutrinos. Correspondingly, the allowed majoron decay channels can be determined from Table 1. For the average position of the various neutrino–spheres, we use the implicit definition
$$\rho (R_{E,\nu _h})=2\times 10^{13}\mathrm{g}/\mathrm{cm}^3,\rho (R_{E,\overline{\nu }_e})=2\times 10^{12}\mathrm{g}/\mathrm{cm}^3,\rho (R_{E,\nu _e})=2\times 10^{11}\mathrm{g}/\mathrm{cm}^3,$$
(33)
and
$$\rho (R_{T,\nu _h})=\rho (R_{T,\overline{\nu }_e})=\rho (R_{T,\nu _e})=2\times 10^{11}\mathrm{g}/\mathrm{cm}^3.$$
(34)
In calculating the effect of majoron–emitting neutrino decays on the emitted $`\nu `$ spectra we must distinguish three regions. If a $`\overline{\nu }_i`$ is produced inside its own energy sphere, $`r<R_{E,\overline{\nu }_i}`$, it is still in thermal and chemical equilibrium with the stellar medium. Therefore, the production (or the decay) of a $`\overline{\nu }_i`$ within $`r<R_{E,\overline{\nu }_i}`$ does not influence the emitted $`\overline{\nu }_i`$ spectrum. In contrast, the spectrum of $`\overline{\nu }_i`$ neutrinos is fixed for $`r>R_{E,\overline{\nu }_i}`$ and both its decay or its production changes its spectra. The only difference between the two regions $`R_{E,\overline{\nu }_i}<r<R_{t,\overline{\nu }_i}`$ and $`r>R_{t,\overline{\nu }_i}`$ is the different “effective” velocity $`v`$ of the decaying neutrino. In the latter, the neutrino can escape freely ($`v=1`$), while in the first it diffuses outward with $`v\lambda /(R_{E,\overline{\nu }_i}R_{t,\overline{\nu }_i})`$. Here, $`\lambda `$ is the average mean free path of the neutrino $`\overline{\nu }_i`$. Hence, we can compute the survival probability $`N(\overline{\nu }_i)`$ of a $`\overline{\nu }_i`$ neutrino emitted from its energy sphere as
$$N(\overline{\nu }_i)=\mathrm{exp}\left\{_{R_{E,\overline{\nu }_i}}^{R_{t,\overline{\nu }_i}}\frac{dr^{}}{v}\mathrm{\Gamma }_{\overline{\nu }_i}(r^{})_{R_{t,\overline{\nu }_i}}^{\mathrm{}}𝑑r^{}\mathrm{\Gamma }_{\overline{\nu }_i}(r^{})\right\},$$
(35)
where
$$\mathrm{\Gamma }_{\overline{\nu }_i}=\underset{l=e,\mu ,\tau }{}\mathrm{\Gamma }(\overline{\nu }_i\nu _l+J)$$
(36)
is its total decay rate. Denoting with $`p_{ij}[r_1,r_2]`$ the probability that the decay $`\overline{\nu }_i\nu _j+J`$ happens in between $`r_1`$ and $`r_2`$, it follows that
$$p_{ij}[r_1,r_2]=_{r_1}^{r_2}\frac{dr^{}}{v}\mathrm{\Gamma }(\overline{\nu }_i\nu _j+J;r^{})N(r^{}).$$
(37)
and $`_jp_{ij}[r_1,r_2]+N=1`$, where N denotes the neutrino survival probability.
Finally, we have to consider the effect of majoron decays $`J\nu _i+\nu _j`$ on the neutrino spectra. If $`g<3\times 10^6`$, majorons escape freely. Thus they are emitted by the core with luminosity given by Eq. (25). For their mean energy we assume $`E_J=3T50`$ MeV, so that the number of emitted majorons per unit time is $`N_J=_J/E_J`$. In the intermediate regime, $`3\times 10^6<g<2\times 10^5`$, we use the same average energy for the majorons but use the reduced luminosity due to emission from the shell $`[r_0l:r_0]`$. Finally, for $`g>2\times 10^5`$, we identify simply the majoron–sphere $`R_J`$ with the energy sphere of the neutrino species which has the largest coupling to the majoron. In the following, we will use $`R_J=R_{E,\nu _h}`$ and we will also assume that its average energy is similar to the one of $`\nu _h`$, $`E_JE_{\nu _h}25\mathrm{MeV}`$.
The decay probability $`P=1N`$ of a majoron is given by
$$N(J)=\mathrm{exp}\left\{_{R_J}^{\mathrm{}}𝑑r^{}\mathrm{\Gamma }_J(r^{})\right\}$$
(38)
Together with the majoron luminosity $`_J`$ this allows us to calculate the spectra of the produced neutrinos.
#### 3.3.1 Neutrino signal from SN1987A
The existing observational data on the supernova SN 1987A were already used in Ref. , in order to constrain the allowed permutation between the $`\overline{\nu }_e`$ and all other types of neutrinos as a result of neutrino oscillations. At 99% (95%) confidence level, it was found that no more than 35% (23%) of the $`\overline{\nu }_e`$ flux could be converted into $`\overline{\nu }_{\mu ,\tau }`$ or $`\nu _{\mu ,\tau }`$. This limit can be applied to any process which diminishes the flux of $`\overline{\nu }_e`$ and is only weakly energy dependent, such as the case of medium-induced neutrino decays. Thus one can apply their results in order to restrict the neutrino decay models under consideration here.
Since $`\nu _\mu `$ and $`\nu _\tau `$ feel the same potential, the decay rates $`\mathrm{\Gamma }(\overline{\nu }_e\nu _\mu +J)`$ and $`\mathrm{\Gamma }(\overline{\nu }_e\nu _\tau +J)`$ differ only due to their different coupling constants. We define therefore $`g_{eh}^2=g_{e\mu }^2+g_{e\tau }^2`$ and present the region excluded by the SN 1987A signal in Fig. 1 in the $`g_{ee}g_{eh}`$ plane. The region excluded at 95% confidence level corresponds approximately to the one given by the condition
$$g_{ee}^2+g_{e\mu }^2+g_{e\tau }^2>1\times 10^7.$$
(39)
Up to now, we have neglected that massive neutrinos do not simply decay but may also oscillate on their way. In the high-density region where the decays occur this approximation is legitimate because the medium states coincide essentially with the weak flavour eigenstates.
However, the SN 1987A neutrinos must also propagate first through the SN envelope, then through vacuum before they cross the Earth on the way to the detectors. Let us consider now what will be the impact of neutrino oscillations for the three popular solutions the solar neutrino problem, namely small-angle (SMA) MSW, large-angle (LMA) MSW and the just-so or vacuum oscillations<sup>4</sup><sup>4</sup>4We follow closely the discussion given in Refs. .. These solutions are characterized by particular values of $`\mathrm{\Delta }_0`$ and $`\mathrm{sin}^22\theta _0`$, which define different regimes in the dynamics of neutrino propagation through the supernova.
Let us first analyze the two MSW solutions. For the usual neutrino mass hierarchy, anti-neutrinos will not encounter an MSW resonance on their way through the SN envelope. Furthermore, they propagate adiabatically, because in both cases $`\mathrm{\Delta }_010^5`$ eV<sup>2</sup>. Therefore, neutrinos that had been created as $`\stackrel{~}{\nu }_1^+(\rho )\overline{\nu }_e`$ leave the star as $`\stackrel{~}{\nu }_1^+(\rho =0)=\nu _1^+`$, i.e. as a definite mass eigenstate. Consequently, no neutrino oscillations take place on the way from the SN to the Earth and the probability of $`\overline{\nu }_e\overline{\nu }_\mu `$ transitions (permutation factor) equals $`P_{\mathrm{osc}}=|\overline{\nu }_\mu |\stackrel{~}{\nu }_1^+|^2=\mathrm{sin}^2\theta _0`$, as long as matter effects in the Earth can be neglected. Since decays $`\overline{\nu }_\alpha \nu _\beta `$ and oscillations $`\overline{\nu }_\alpha \overline{\nu }_\beta `$ are decoupled, the total survival probability can be written as
$$N=N_{\mathrm{decay}}N_{\mathrm{osc}}=N_{\mathrm{decay}}(1P_{\mathrm{osc}}).$$
(40)
where $`P_{\mathrm{osc}}`$ is the neutrino conversion probability due to oscillations.
In the SMA-MSW case, $`\mathrm{sin}^22\theta _07\times 10^3`$, oscillations can be neglected both in the SN envelope and in the Earth, so that $`P_{\mathrm{osc}}0`$. Thus, one can use directly the results obtained above, shown in Fig. 1.
In contrast, the matter effect inside the Earth has to be taken into account in the LMA-MSW case, for which $`\mathrm{sin}^22\theta _00.6`$ . The permutation probabilities due to oscillations is given by,
$$P_{\mathrm{osc}}=\mathrm{sin}^22\theta _0\mathrm{sin}2\theta \mathrm{sin}(2\theta _02\theta )\mathrm{sin}^2(\pi d/l_{\mathrm{osc}}).$$
(41)
The distance $`d`$ traveled inside the Earth by the neutrinos and the average density $`\rho `$ are different for Kamiokande and IMB detectors. For Kamiokande we have $`d=3900`$ km and $`\rho =3.4`$ g/cm<sup>3</sup> while for IMB $`d=8400`$ km and $`\rho =4.6`$ g/cm<sup>3</sup>. As an approximation we will use therefore the average value $`P_{\mathrm{osc}}=(12P_{\mathrm{osc},\mathrm{Kam}}+8P_{\mathrm{osc},\mathrm{IBM}})/20`$ according to the number of events detected in each detector. In Fig. 2, we show the excluded region for $`\mathrm{sin}^22\theta _0=0.6`$, $`\mathrm{\Delta }_0=10^5`$ eV<sup>2</sup>, and in addition the experimental limit Eq. (3). The coupling $`g_{22}`$ is fixed by
$$g_{22}=g_{11}\sqrt{1+\frac{\mathrm{\Delta }_0}{m_1^2}}.$$
(42)
Therefore, for $`m_1/\mathrm{\Delta }_00`$, $`g_{22}`$ becomes larger for constant $`g_{11}`$. Hence, the limit on $`g_{11}`$ becomes stronger. Note also that we can not fix $`g_{33}`$ by the solar neutrino data. Therefore, we have set conservatively $`g_{33}=0`$. As can be seen from the Fig. 2, for the range of masses involved, the values $`g_{11}>10^4`$ are excluded at $`95\%`$ C.L.
The remaining case to consider is the ’just-so’ solution, in which $`\mathrm{\Delta }_010^{10}`$ eV<sup>2</sup>. For such small values of the mass splitting, the neutrino propagation out of the SN is non-adiabatic. Therefore, neutrinos leave the SN envelope as flavor eigenstates which then oscillate on their way to the Earth. Taking into account that in this case the Earth effect is unimportant, their averaged permutation probability due to oscillation is simply
$$P_{\mathrm{osc}}=\frac{1}{2}\mathrm{sin}^22\theta _0.$$
(43)
Let us consider now the typical value $`\mathrm{sin}^22\theta _00.9`$ . Then, according to the analysis of Ref. , this case is already disfavored assuming only oscillation. In any case for completeness we have plotted in Fig. 3 the regions corresponding to $`P=0.55,0.7`$ and the experimental limit Eq. (3).
To summarize, the limits we have obtained in this subsection from the observed $`\overline{\nu }_e`$ signal of SN 1987A are more than one order of magnitude stronger than the limit from pion decay. They require as theoretical input some knowledge about the spectral shape of the emitted neutrino fluences (e.g. $`E_{\nu _i}`$) and are therefore to a certain extent model dependent. However this model dependence is much weaker than that implicit in the limit from the collapsing phase (20).
#### 3.3.2 Constraints from a future galactic supernova
The neutrino signal from SN 1987A observed by Kamiokande and IBM confirmed the general astrophysical picture of a supernova explosion. However the small number of observed neutrinos prevents a sensitive probe of neutrino properties or of the equation of state of super–dense matter. Meanwhile, several new experiments like Super–Kamiokande or the Sudbury Neutrino Observatory (SNO) have been constructed and it seems therefore interesting to estimate their sensitivity to new physics, such as the majoron–emitting neutrino decays, through their ability to determine the spectra of SN neutrinos.
All detectors in the near future can detect SN only in our own Galaxy and in the Large and Small Magellanic Cloud, with however a much smaller rate than expected for the Milky Way. Therefore, we use in the following 10 kpc as the reference distance to the hypothetical Supernova. Then, according to Ref. , the number of $`\overline{\nu }_e+pn+e^+`$ events expected in SK is 5310, compared to 220 events in all other channels. While a light-water detector like SK acts like a $`\overline{\nu }_e`$ filter, a heavy-water detector has a broad response to all types of neutrinos: from the total of 781 events expected at SNO, 200 events are neutral-current reactions of $`\nu _h`$ (=$`\nu _\mu ,\overline{\nu }_\mu ,\nu _\tau `$ or $`\overline{\nu }_\tau `$) with deuterium. In addition, there are 82 charged-current reactions of $`\nu _e`$ with deuterium. The combined spectral information of these reactions should contain the major part of the information extractable from all data. In the following, we will consider therefore only the reaction $`\overline{\nu }_e+pn+e^+`$ for SK, and $`\nu _h+d\nu _h+n+p`$ and $`\nu _e+dp+p+e^{}`$ for SNO.
The predicted number of events has a rather strong dependence on how the cooling of the SN is modeled. In Ref. , e.g., the expected number of $`\nu _h+d\nu _h+n+p`$ reactions is 485 in SNO. An important question is therefore how well different astrophysical and/or particle physics processes (e.g. neutrino oscillations versus decays) can be disentangled using observational data. We will comment shortly on this question at the end of this section.
The neutrino spectrum $`n_\nu (E_\nu )`$ emitted by the SN can be be inferred in a rather indirect way from the experimentally measured energies $`E_i`$ of the observed events. Let us consider first the reaction $`\overline{\nu }_e+pn+e^+`$ and $`\nu _e+dp+p+e^{}`$. Experimentally, the neutrino energy $`E_\nu `$ is reconstructed in both reactions from the number of photo-multipliers that have been triggered by Cherenkov photons emitted by the $`e^\pm `$. The probability that a certain energy $`E_i`$ is ascribed to an event $`i`$ where the $`e^\pm `$ has the energy $`E_e`$ depends in principle on the geometry and the details of the detector. For our purposes it is enough to take into account the Poissonian nature of the emission and detection of the Cherenkov photons. In this case the probability distribution function that the energy $`E_i`$ is ascribed to an $`e^\pm `$ with energy $`E_e`$ is a Gaussian distribution,
$$G(E_i,E_e)=\frac{1}{\sqrt{2\pi }\sigma (E_e)}\mathrm{exp}\left(\frac{(E_eE_i)^2}{2\sigma ^2(E_e)}\right)$$
(44)
with $`\sigma (E_e)=\sqrt{E_eE_\sigma }`$, where $`E_\sigma `$ is the energy resolution. Following Refs. , we use $`E_\sigma =0.22`$ MeV for SK and $`E_\sigma =0.20`$ MeV for SNO. Apart from its resolution $`E_\sigma `$, the other main feature of a detector is its efficiency $`ϵ(E)`$. For both experiments, we use the approximation $`ϵ(E)=ϵ_0\theta (EE_{\mathrm{cut}})`$ with $`ϵ_0=1`$ and a conservative value for $`E_{\mathrm{cut}}=7`$ MeV for SK and $`E_{\mathrm{cut}}=5`$ MeV for SNO. Combining the Gaussian kernel $`G(E_i,E_e)`$ with the detector efficiency $`ϵ(E)`$, the relation between the $`e^\pm `$ energy $`E_e`$ and the measured energy $`E_i`$ is
$$n(E_i)=ϵ_0_{E_{cut}}^{\mathrm{}}𝑑EG(E_i,E_e)n_e(E_e).$$
(45)
In contrast to the two reactions considered above, the neutral-current reaction $`\nu _h+d\nu _h+n+p`$ used to detect $`\nu _h`$ in SNO does not allow to measure the energy $`E_\nu `$. Nevertheless, it is possible to reconstruct partially the $`\nu _h`$ energy spectrum comparing the observed signal in different detector materials with different energy responses (cf., e.g., Fig. 3 of Ref. ).
As next ingredient, we need the relation between the time–integrated neutrino spectra $`n_i(E)`$ and the spectra of the secondary $`e^\pm `$’s. For the reaction $`\overline{\nu }_e+pn+e^+`$, the positron spectrum is given by
$$n_e(E_e)=\frac{N_p}{4\pi D^2}\sigma _{\overline{\nu }_e+pn+e^+}(E+\mu )n_{\overline{\nu }_e}(E+\mu )$$
(46)
where $`N_p`$ is the number of target protons in SK and $`D`$ the distance to the supernova. For the cross section, we use
$$\sigma _{\overline{\nu }_e+pn+e^+}=\sigma _0E^2\left(1\frac{\mu }{E}\right)\left(1\frac{2\mu }{E}+\frac{\mu ^2+m_e^2}{E^2}\right)^{1/2}$$
(47)
with $`\sigma _0=2.3\times 10^{44}`$ cm<sup>2</sup>, where $`\mu =1.227`$ MeV is the proton-neutron mass difference and $`m_e`$ the electron mass. In the case of the two neutrino-deuterium reactions, we use the cross-sections tabulated in Ref. and replace $`N_p`$ by the number of deuterium nuclei.
Numerical simulations of the neutrino transport show that the instantaneous neutrino spectra can be described as Fermi-Dirac distributions with an effective degeneracy parameters $`\eta _i`$ . The instantaneous neutrino spectra found are pinched, i.e. their low- and high-energy parts are suppressed relative to a Maxwell-Boltzmann distribution. Janka and Hillebrandt found that during the cooling process the effect of pinching becomes less important. Moreover, the pinching of the instantaneous neutrino spectra is compensated by the superposition of different spectra with decreasing temperatures.
In the following, we assume for the time-integrated energy spectra of the different neutrino types
$$n_i(E)=N_i\frac{E^2}{e^{E/T_i\eta _i}+1}$$
(48)
with $`E_{\nu _e}=11`$ MeV, $`E_{\overline{\nu }_e}=16`$ MeV and $`E_{\nu _h}=25`$ MeV. For the degeneracy parameter, we use the values of the lower end of the range found by Janka and Hillebrandt, namely $`\eta =3`$ for $`\nu _e`$, $`\eta =2`$ for $`\overline{\nu }_e`$ and $`\eta =0`$ for $`\nu _h`$. Then the relation between the effective temperature $`T_i`$ is given by $`3T0.751E_{\nu _e}`$ for $`\nu _e`$, $`3T0.832E_{\overline{\nu }_e}`$ for $`\overline{\nu }_e`$ and $`3T=E_{\nu _h}`$ for $`\nu _h`$.
We have performed a simulation of the expected neutrino signals in SK and SNO in the case of a galactic Supernova. Here have assumed $`g_{11}g_{22}g_{33}`$, as suggested by a scheme with all three neutrino masses are quasi-degenerate, in the eV range . Moreover, here we have chosen the solar and atmospheric mixing angles to be maximal . The resulting average neutrino signals disregarding neutrino oscillations are shown for several values of $`g_{11}`$ in Fig. 4 and 5. As one could expect from the discussion in Section 3.3.1, for the chosen values of $`g`$ neutrino decays do not influence drastically the $`\overline{\nu }_e`$ signal (Fig. 4). In contrast the $`\nu _h`$ signal shown in Fig. 5 shows two main features. First, in the presence of majoron interactions, the spectra show a surplus of low-energy ($`E<18`$ MeV) and a deficit of high-energy ($`E>18`$ MeV) $`\nu _h`$ compared to the reference Standard Model spectrum with $`g=0`$. One sees that the decays $`\overline{\nu }_h\nu _l+J`$ reduce the $`\nu _h`$ fluence by at most a factor 2. Since the decays $`\overline{\nu }_h\nu _l+J`$ and the majoron decays $`J\nu +\nu ^{}`$ both produce also low-energy $`\nu _{\tau ,\mu }`$, the $`\nu _h`$ signal shows a surplus of events at low energies.
We have not shown the $`\nu _e`$ spectra because they are basically unchanged with respect to the reference spectrum with $`g=0`$. Since the $`\nu _e`$ energy sphere at $`R_{E,\nu _e}`$ is located at a much smaller density than those of $`\overline{\nu }_e`$ and $`\nu _h`$, most ($`95\%`$) $`\overline{\nu }_e`$ and $`\nu _h`$ decay inside $`R_{E,\nu _e}`$ where the $`\nu _e`$ are still in chemical equilibrium. Therefore, its spectrum is hardly affected by $`\overline{\nu }_e`$ and $`\nu _h`$ decays.
Finally, we want to comment on the dependence of our results on the assumed astrophysical parameters. First note that, an overall suppression of all three neutrino signals could be explained more naturally by an astrophysical reason than by neutrino decays. Similarly if a suppression with respect to the expectations of a given SN model were observed only in the $`\nu _h`$ signal, this would indicate an astrophysical explanation, since the $`\nu _h`$ temperature has the largest uncertainty. It is therefore of importance that part of the spectral information of the $`\nu _h`$ signal can be recovered by comparing the signal in different detector materials. The signature for majoron neutrino decays would be two-fold. First, a reduction of the observed $`\nu _h`$ fluence compared to the one expected from the observed $`\nu _h`$ temperature. Second, a non-thermal $`\nu _h`$ spectrum with a surplus of low-energy and a deficit of high-energy $`\nu _h`$.
A more quantitative study of the restrictions on the majoron–neutrino couplings attainable at SK and SNO would require a detailed likelihood analysis. However, from Fig. 5 one expects that the sensitivity of these experiments will be restricted to $`g<5\times 10^5`$. Therefore, we conclude that although these experiments could narrow down the allowed window for majoron–neutrino interactions, but not close it completely.
## 4 Conclusions
We have reconsidered the influence of majoron neutrino decays on the neutrino signal of supernovae in the full range of allowed neutrino masses, in the light of recent Super–Kamiokande data on solar and atmospheric neutrinos. In the high–density supernova medium the effects of Majoron–emitting neutrino decays are important even if they are suppressed *in vacuo* by small neutrino masses and/or off-diagonal couplings. In contrast to previous works, we have considered scattering and decay processes, taking into account medium effects for both kinds of processes. The dispersive effects of the dense SN core are particularly important since the currently favoured interpretation of the solar and atmospheric neutrino data points towards light neutrinos.
The observation of SN 1987A excludes two parts of the possible range of neutrino majoron coupling constants. In the range $`3\times 10^7<g<2\times 10^5`$, where $`g`$ is the largest element of the coupling matrix $`g_{\alpha \beta }`$, the supernova looses to much energy into majorons, thereby shortening the neutrino signal too much. For larger couplings, the fluence of escaping $`\overline{\nu }_e`$ is reduced due to decays $`\overline{\nu }_e\nu _l+J`$. Depending on which is the correct solution to the solar neutrino problem, different values of $`g_{\mathrm{eff}}`$ can be excluded: $`g_{ee}^2+g_{eh}^2>1\times 10^7`$ for SMA-MSW (Fig. 1) and $`g_{11}>1\times 10^4`$ for LMA-MSW (Fig. 2). In the case of vacuum oscillations, the predicted number of $`\overline{\nu }_e`$ events is already disfavored (even in the absence of decays) by the observed number of events and by the current theoretical understanding of supernova explosions. The corresponding sensitivities are displayed in Fig. 3.
Finally, we have discussed the potential of Superkamiokande and the Sudbury Neutrino Observatory to detect majoron neutrino interactions in the case of a future galactic supernova. We have found that although these experiments could narrow down the allowed window for majoron–neutrino interactions, but not close it completely, reaching a sensitivity at the few $`\times 10^5`$ level, as seen from Fig. 5.
## Acknowledgements
We are grateful to J.F. Beacom, Z. Berezhiani, J.A. Pons and A. Rossi for helpful comments. We would like to thank especially H.-T. Janka for correspondence and for sending his supernova simulation data. This work was supported by DGICYT grant PB95-1077 and by the EEC under the TMR contract ERBFMRX-CT96-0090. MK has been supported by a Marie-Curie grant and RT by a fellowship from Generalitat Valenciana. |
warning/0001/gr-qc0001081.html | ar5iv | text | # Exact Solution for the Exterior Field of a Rotating Neutron Star
## I Introduction
Observations of binary pulsars indicate that the individual neutron stars (NS) in such systems have masses very close to the Chandrasekhar limit of 1.4 $`M_{}`$ of white dwarfs. Theoretically, the issue of the maximum mass of a NS hings strongly on the equation of state (EOS) and the particle interactions at the high density of the center. Models of NS with strong kaon condensation or even a quark nucleus (‘strange star’ ) can have at most 1.5 – 2 $`M_{}`$, which would leave range for lower mass black holes. However, even modest differential rotation may easily increase the maximum mass 2 $`M_{}`$ of a nonrotating NS to above 3 $`M_{}`$ for a nascient NS in a transient phase of a supernova. Moreover, a mass-quadrupole moment $`Q`$ is also important for achieving correspondence with numerical results .
In order to model analytically the exterior field of a NS in the framework of general relativity (GR), one needs an exact asymptotically flat solution of the Einstein-Maxwell equations (electrovac spacetimes) possessing at least four arbitrary physical parameters which are the mass $`M`$, angular momentum $`J`$, magnetic dipole $`\mu `$ and mass-quadrupole moment $`Q`$. The simplest solution, besides, can be envisaged as axially symmetric, the magnetic field sharing the symmetry of the mass and angular momentum distributions; an additional reflection symmetry with respect to the equatorial plane which is expected of the self-gravitating objects (see, e.g., and references therein) must be also imposed.
Our paper aims at presenting an exact solution which does satisfy the above requirements and, what is most important, is a mathematically very simple solution admitting a representation exclusively in terms of the rational functions of spheroidal coordinates (previous effort in this direction only led either to the solution which had no reflection symmetry or to solutions which did not permit the rational function representation, and consequently could not be written in a concise form).
## II Four–parameter exact solution
The reported solution has been constructed with the aid of Sibgatullin’s method according to which the complex potentials $``$ and $`\mathrm{\Phi }`$ satisfying Ernst’s equations are defined, for specified axis data $`e(z):=(z,\rho =0)`$ and $`f(z):=\mathrm{\Phi }(z,\rho =0)`$, by the integrals
$$(z,\rho )=\frac{1}{\pi }_1^1\frac{e(\xi )\mu (\sigma )d\sigma }{\sqrt{1\sigma ^2}},\mathrm{\Phi }(z,\rho )=\frac{1}{\pi }_1^1\frac{f(\xi )\mu (\sigma )d\sigma }{\sqrt{1\sigma ^2}}.$$
(1)
The unknown function $`\mu (\sigma )`$ is to be found from the singular integral equation
$$_1^1\frac{\mu (\sigma )[e(\xi )+\stackrel{~}{e}(\eta )+2f(\xi )\stackrel{~}{f}(\eta )]d\sigma }{(\sigma \tau )\sqrt{1\sigma ^2}}=0$$
(2)
with the normalizing condition
$$_1^1\frac{\mu (\sigma )d\sigma }{\sqrt{1\sigma ^2}}=\pi ,$$
(3)
where $`\xi =z+\mathrm{i}\rho \sigma `$, $`\eta =z+\mathrm{i}\rho \tau `$, $`\rho `$ and $`z`$ being the Weyl-Papapetrou cylindrical coordinates and $`\sigma ,\tau [1,1]`$; $`\stackrel{~}{e}(\eta ):=\overline{e(\overline{\eta })}`$, $`\stackrel{~}{f}(\eta ):=\overline{f(\overline{\eta })}`$, and the overbar stands for complex conjugation.
In what follows, the axis data $`e(z)`$ and $`f(z)`$ are chosen in the form
$`e(z)`$ $`=`$ $`{\displaystyle \frac{(zM\mathrm{i}a)(z+\mathrm{i}b)+d\delta ab}{(z+M\mathrm{i}a)(z+\mathrm{i}b)+d\delta ab}},`$ (4)
$`f(z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}\mu }{(z+M\mathrm{i}a)(z+\mathrm{i}b)+d\delta ab}},`$ (5)
$`\delta `$ $`:=`$ $`{\displaystyle \frac{\mu ^2M^2b^2}{M^2(ab)^2}},d:={\displaystyle \frac{1}{4}}[M^2(ab)^2],`$ (6)
such that the algebraic equation
$$e(z)+\overline{e}(z)+2f(z)\overline{f}(z)=0$$
(7)
will have a pair of distinct roots of multiplicity two. This is a key point for having the rational form of the final expressions for $`(\rho ,z)`$ and $`\mathrm{\Phi }(\rho ,z)`$ after performing the Riemann-Hilbert procedure of the analytic continuation of the functions $`e(z)`$, $`f(z)`$ into the complex plane $`(\rho ,z)`$. The resulting expressions for the potentials $`(\rho ,z)`$ and $`\mathrm{\Phi }(\rho ,z)`$ obtained from Eqs. (1)-(7) are of the following polynomial form:<sup>*</sup><sup>*</sup>*All the formulas of this paper have been checked with the aid of the Mathematica 3.0 computer program .
$``$ $`=`$ $`(A2MB)/(A+2MB),\mathrm{\Phi }=2\mathrm{i}\mu C/(A+2MB),`$ (8)
$`A`$ $`=`$ $`4[(k^2x^2\delta y^2)^2d^2\mathrm{i}k^3xy(ab)(x^21)]`$ (10)
$`(1y^2)[(ab)(d\delta )M^2b][(ab)(y^2+1)+4\mathrm{i}kxy],`$
$`B`$ $`=`$ $`kx\{2k^2(x^21)+[b(ab)+2\delta ](1y^2)\}`$ (12)
$`+\mathrm{i}y\{2k^2b(x^21)[k^2(ab)M^2b2a\delta ](1y^2)\},`$
$`C`$ $`=`$ $`2k^2y(x^21)+[2\delta y\mathrm{i}kx(ab)](1y^2),`$ (13)
where we have introduced the generalized spheroidal coordinates
$`x`$ $`=`$ $`{\displaystyle \frac{1}{2k}}(r_++r_{})\mathrm{and}y={\displaystyle \frac{1}{2k}}(r_+r_{}),`$ (14)
$`r_\pm `$ $`:=`$ $`\sqrt{\rho ^2+(z\pm k)^2},k:=\sqrt{d+\delta }.`$ (15)
The four arbitrary real parameters entering the solution are the total mass $`M`$, total angular momentum per unit mass $`a:=J/M`$, magnetic dipole moment $`\mu `$ and mass-quadrupole moment
$$Q=\frac{M}{4[M^2(ab)^2]}\left[M^4+2M^2(a^2+b^2)(3a+b)(ab)^34\mu ^2\right]$$
(16)
of the source.
The corresponding complete metric is given by the axisymmetric line elementThroughout the paper, natural units are used in which the gravitational constant and the velocity of light are equal to unity.
$$\mathrm{d}s^2=f(\mathrm{d}t\omega \mathrm{d}\phi )^2+k^2f^1\left[e^{2\gamma }(x^2y^2)\left(\frac{\mathrm{d}x^2}{x^21}+\frac{\mathrm{d}y^2}{1y^2}\right)+(x^21)(1y^2)\mathrm{d}\phi ^2\right],$$
(17)
in which the metric coefficients $`f`$, $`\gamma `$ and $`\omega `$ are the following rational functions of the coordinates $`x`$ and $`y`$ (see, e.g., for detais of Sibgatullin’s method):
$`f`$ $`=`$ $`E/D,\mathrm{e}^{2\gamma }=E/16k^8(x^2y^2)^4,\omega =(y^21)L/E,`$ (18)
$`E`$ $`=`$ $`\{4[k^2(x^21)+\delta (1y^2)]^2+(ab)[(ab)(d\delta )M^2b](1y^2)^2\}^2`$ (20)
$`16k^2(x^21)(1y^2)\{(ab)[k^2(x^2y^2)+2\delta y^2]+M^2by^2\}^2,`$
$`D`$ $`=`$ $`\{4(k^2x^2\delta y^2)^2+2kMx[2k^2(x^21)+(2\delta +abb^2)(1y^2)]+(ab)[(ab)`$ (23)
$`\times (d\delta )M^2b](y^41)4d^2\}^2+4y^2\{2k^2(x^21)[kx(ab)Mb]`$
$`2Mb\delta (1y^2)+[(ab)(k^22\delta )M^2b](2kx+M)(1y^2)\}^2,`$
$`L`$ $`=`$ $`8k^2(x^21)\{(ab)[k^2(x^2y^2)+2\delta y^2]+M^2by^2\}`$ (28)
$`\times \{kMx[(2kx+M)^2a^2+b^22y^2(2\delta +abb^2)]2y^2(4\delta dM^2b^2)\}`$
$`\{4[k^2(x^21)+\delta (1y^2)]^2+(ab)[(ab)(d\delta )M^2b](1y^2)^2\}`$
$`\times ((1y^2)\{2M(2kx+M)[(ab)(d\delta )b(M^2+2\delta )]4M^2b\delta `$
$`+(ab)(4\delta dM^2b^2)\}8k^2Mb(kx+M)(x^21)).`$
### Special cases
Eqs. (6) and (10) admit several well-known classical limits:
1. In the absence of the magnetic field and vanishing arbitrary quadrupole deformation, i.e $`\mu =b=0`$, only the Ernst potential $``$ survives which is readily recognizable as that of the Tomimatsu–Sato $`\delta =2`$ solution with the mass quadrupole $`Q=\frac{1}{4}(M^3+3J^2/M)`$.
2. The stationary pure vacuum limit with a non-vanishing quadrupole parameter $`b`$ is a particular three-parameter specialization of the Kinnersley–Chitre solution .
3. The magnetostatic limit ($`a=b=0`$) is represented by Bonnor’s solution for a massive magnetic dipole. For this solution, the quadrupole moment is $`Q=\mu ^2/M\frac{1}{4}M^3`$.
4. Reduction to the Kerr metric with total mass $`M`$ and total angular momentum per unit mass $`a:=J/M`$ is achieved by setting $`\mu =0`$ and then formally choosing $`b^2=a^2M^2`$. The values of $`M`$ and $`a`$ remain independent, in particular, $`a^2<M^2`$ can be imposed since in this special limit the complex continuation $`b\mathrm{i}b`$ is admitted. It should be stressed, however, that there exist general arguments for the interior of the Kerr metric consisting of a perfect fluid according to which a large rotational flattening of the body necessarily implies a large absolute value of the quadrupole moment.
5. It is remarkable that the hyperextreme part of our solution corresponding to pure imaginary values of $`k`$ belongs to the Chen–Guo–Ernst family of hyperextreme spacetimes . This branch might represent exterior fields of relativistic disks. Their importance for astrophysics was shown by Bardeen and Wagoner , cf. also . In the absence of the electromagnetic field an exact global solution for an infinitesimally thin disk of dust has been constructed by Neugebauer and Meinel , cf. .
Since neutron stars are known to be ‘slowly’ rotating astrophysical objects (even at the Kepler frequency $`\omega _\mathrm{K}=\sqrt{GM/R^3}0.5`$ ms of a millisecond pulsar would the equator rotate only with a speed of $`v_\mathrm{K}c/4`$), therefore, it is the subextreme part of the metric (17) which should be used for their description. At the same time, we still need to know the location of singularities in our solution to support the physical interpretation we are attributing to it. In Fig. 1 we have plotted in coordinates $`\rho `$ and $`z`$ the typical shapes of the infinite redshift surface which one has for the real-valued $`k`$, the dots indicating the position of singularities. The two point singularities on the symmetry axis (the poles $`x=1`$, $`y=\pm 1`$) belong to the stationary limit surface, while the ring singularity in each case lies at the equatorial plane between the symmetry axis and ergosphere; no singularity outside the infinite red shift surface arises, and hence the metric (17) is describing the exterior fields of compact objects such as neutron stars.
We shall conclude the presentation of our solution by writing out Kinnersley’s potential $`𝒦`$ the real part of which gives the magnetic component $`A_\varphi `$ of the electromagnetic four potential:
$`𝒦`$ $`:=`$ $`A_\varphi +\mathrm{i}A_t^{}=\mu (1y^2)K/(A+2MB),`$ (29)
$`K`$ $`=`$ $`2k^2(x^21)[2kx+3M+\mathrm{i}y(ab)](ab)[2MaMb(1y^2)4\mathrm{i}\delta y]`$ (31)
$`+2(2kx+M)[\delta (1y^2)+M(M\mathrm{i}by)]+4M\delta ,`$
where $`A_t^{}`$ is the electric component of a vector potential associated with the dual electromagnetic field tensor. The knowledge of $`A_\varphi `$ is a necessary basis for the investigation of plasma-dynamical effects around neutron stars. In Fig. 2 the magnetic lines of force are plotted in cylindrical coordinates for two characteristic cases.
## III Matching to numerical models of neutron stars
Our exact axisymmetric solution needs to be matched to interior solutions of neutron stars, in order to be realistic. Since the junction conditions on the surface of the NS depend very much on the material details such as equation of state, conductivity etc., we will restrict ourselves here only to the identification of asymptotically conserved quantities. The identification of mass $`M`$ and angular momentum $`J`$ of our solution agrees already with the standard parameters of asymptotically flat spacetimes in GR and in the astrophysics of NSs.
### A Magnetic field
Normally, the NS’ magnetic field $`\stackrel{}{B}`$, predicted already in 1964 by Hoyle, Narlikar, and Wheeler and reaching high values below the upper limit of $`\stackrel{}{B}3\times 10^9`$ Tesla, is ignored in numerical studies of rapidly rotating NS. More recently, however, axisymmetric solutions of the Einstein–Maxwell equations have been studied using a pseudo-spectral method involving Chebyshev-Legendre polynomials in terms of maximal slicing quasi-isotropic coordinates. Since the magnetic axis is aligned with the rotation axis and only poloidal fields are permitted, this numerical work is particularly suited for a comparison of the electrovac spacetime outside the NS with our exact solution:
For a star close to a sphere and small polar fields $`\stackrel{}{B}10^6`$ Tesla, these numerical results are within an error of $`10^3`$ in agreement with Ferraro’s solution
$$A_\varphi =4\pi \rho _0f_0\frac{R_{}^5}{15r}\mathrm{sin}^2\theta ,r>R_{}$$
(32)
where $`R_{}`$ is the radius of the star, $`\rho _0`$ is the mass density, and $`f_0`$ is the constant value taken by the electric current function .
The equation above enables us to see how the parameter $`\mu `$ of our solution may depend on the parameters of the corresponding interior metric for small values of $`Q`$. Indeed, introducing the Boyer-Lindquist-like coordinates $`R`$ and $`\theta `$ via the formulas
$$kx=RM,y=\mathrm{cos}\theta ,$$
(33)
we easily find from (11), in the limit $`R\mathrm{}`$, that
$$A_\varphi =\frac{\mu \mathrm{sin}^2\theta }{R}+O\left(\frac{1}{R^2}\right)=\frac{\mu \mathrm{sin}^2\theta }{r}+O\left(\frac{1}{r^2}\right),$$
(34)
since $`R`$ has a representation $`R=r(1+O(r^1))`$ in terms of the isotropic coordinate $`r`$ . From (12) and (14) the desired relation of $`\mu `$ to the parameters determining the interior of a NS follows immediately.
### B Quadrupole moment
An operational way of defining the quadrupole moment of an axially symmetric body in GR has been developed by Ryan . For NS one finds quite generally that the numerical simulations are rather well accounted by the simple quadratic relation
$$Q=c(M,\mathrm{EOS})\frac{J^2}{M},$$
(35)
where the constant $`c=c(M,\mathrm{EOS})`$ depends only on the mass $`M`$ and the equation of state (EOS) for the interior of the NS. For NSs of 1.4 $`M_{}`$ the range of this constant is $`c=2`$ to 7.4. It is intriguing that this simple relation quadratic in $`J`$ holds also for fast rotating NS.
If we reparametrize the arbitrary quadrupole parameter $`b`$ of our exact solution rather by the dimensionless parameter $`\mathrm{\Delta }`$ via $`b=\pm \sqrt{a^2+2aM\mathrm{\Delta }M^2}`$ we obtain from (16) for the quadrupole moment ($`\mu =0`$):
$$Q=\left(1+\frac{1}{2}\mathrm{\Delta }^2\frac{a\mathrm{\Delta }^2(a\mathrm{\Delta }^2M\mathrm{\Delta }\pm \sqrt{a^2+2aM\mathrm{\Delta }M^2})}{2(M^2a^2+a^2\mathrm{\Delta }^22aM\mathrm{\Delta })}\right)\frac{J^2}{M}.$$
(36)
Our reparametrization is such that for $`\mathrm{\Delta }=0`$ we recover the mass-quadrupole moment $`Q=J^2/M`$ of the Kerr metric. Thus it is worth pointing out that the mass-quadrupole parameter in our solution is also intimately related with the angular momentum dipole and octupole moments, and this means that the deformations of the source are mainly due to rotation. In a particular case, for instance, when $`M=1.4M_{}`$ and $`a=0.625M_{}`$, the values of $`\mathrm{\Delta }`$ covering the NS range $`2<c<7.4`$ are given by the interval $`0.986<\mathrm{\Delta }<2.068`$.
In comparison with (35), the constant $`c`$ for our solution depends not only on the mass $`M`$, but also on $`\mathrm{\Delta }>(M^2a^2)/2aM`$ which can be adjusted to different EOS. The additional piece depending on angular momentum per unit mass, however, arises only in higher order of $`\mathrm{\Delta }`$. Thus the quadrupole moment of our exact solution accounts rather well to the simple quadratic law (35) of NSs. Moreover, in the ‘extreme’ limit $`aM`$ we find
$$Q=\left(1\frac{\mathrm{\Delta }(\mathrm{\Delta }\pm \sqrt{2\mathrm{\Delta }})}{2(\mathrm{\Delta }2)}\right)M^3,$$
(37)
whence it can be seen that $`Q`$ can assume arbitrary values for any given value of $`M`$, unlike in the case of the extreme Kerr metric for which $`Q=M^3`$. In this particular limit, the ‘NS interval’ corresponds to $`0.931<\mathrm{\Delta }<1.584`$.
Futher study is needed to exhibit the astrophysical significance of our solution in more detail. In view of a simple analytical form of the new electrovacuum metric and clear physical interpretation of the characteristic parameters it possesses, it is anticipated that the solution will prove itself suitable for the use in concrete astrophysical applications involving neutron stars.
## Acknowledgments
It is our pleasure to thank Jerzy Plebański and Alfredo Macías for stimulating discussions and support. We are grateful to Olga Manko for her help in establishing the factor structure of the metric coefficients (10). This work was partially supported by Conacyt, grant No. 28339–E, as well as the joint German–Mexican project DLR–Conacyt E130–1148 and MXI 010/98 OTH. One of us (J.D.S-G.) acknowledges financial support from Colciencias of Colombia. |
warning/0001/hep-ph0001142.html | ar5iv | text | # Bloch-Nordsieck violating electroweak corrections to inclusive TeV scale hard processes
## Abstract
We point out that, since the colliders initial states ( $`e^+e^{},pp,p\overline{p}`$, … ) carry a definite nonabelian flavor, electroweak radiative corrections to inclusive hard cross sections at the TeV scale are affected by peculiar Bloch-Nordsieck violating double logs. We recall the setup of soft cancellation theorems, and we analyze the magnitude of the noncancelling terms in the example of electron - positron annihilation into hadrons.
Interest in logarithmically enhanced electroweak corrections at NLC energies has recently arisen , after the observation - made by two of us \- that double and single logarithms of “soft” Sudakov type are present and sizeable in fixed angle fermion-antifermion annihilation processes at the TeV scale. Such logarithms occur because at energies $`\sqrt{s}`$ much larger than the EW scale $`M_WM_ZM`$, the latter acts as a cutoff for the collinear and infrared (IR) divergences that would be present in the vanishing $`M`$ limit.
The study of Sudakov form factors in the fixed angle, high energy, regime where the expansion parameter is $`\alpha _W(\mathrm{log}\frac{s}{M^2})^2`$ is actually a challenging problem, because one has to investigate the electroweak theory in a transition region between broken and unbroken theory, in which two mass parameters, the effective photon mass and the weak bosons mass $`M_WM_ZM`$ may both be important. Attempts to generalize by various approximation methods the one loop results to higher orders have been made, but with somewhat controversial results<sup>*</sup><sup>*</sup>*For instance, Refs. and agree on the observation that photon and weak boson contributions have to be considered together to restore the gauge symmetry at very large energies, but differ in the evaluation of the symmetry breaking contributions.. Further study is then needed to fully clarify this point.
The fact remains, however, that the NLC regime is one in which the electroweak theory acquires, because of the enlarged phase space, a full-fledged non abelian structure, therefore giving rise to novel physical features with respect to the LEP regime. In this note we wish to point out the fact that even inclusive TeV scale cross sections are affected by Sudakov double logs, due to a lack of cancellation of virtual and real emission electroweak contributions.
In principle, this lack of cancellation is well known, and is due to the violation of the Bloch-Nordsieck (BN) theorem in nonabelian theories. It was initially pointed out for QCD , where it would imply genuine IR divergences at partonic level, but was found eventually to have no physical consequences because of the color averaging of the initial partonic states, forced by the coupling to colorless hadrons. For instance, non factorized IR singular terms do occur in Drell-Yan processes , but at next-to-leading level only, where they are possibly suppressed by a Sudakov form factor , and eventually turn out to be higher twist .
On the other hand, in the electroweak case the initial state is fixed and carries a non abelian charge. Therefore, no averaging is possible, so that the BN violating terms - though finite because of the symmetry breaking scale $`M`$\- are unavoidable for any inclusive observable and need to be carefully computed.
In order to understand this problem, let us start recalling the structure of soft interactions accompanying a hard SM process, of type
$$\{\alpha _1^Ip_1^I,\alpha _2^Ip_2^I\}\{\alpha _1^Fp_1^F,\alpha _2^Fp_2^F,\mathrm{},\alpha _n^Fp_n^F\}$$
(1)
where $`\alpha ,p`$ denote flavor/color and momentum indices of the initial and final states, that we collectively denote by $`\{\alpha _Ip_I\}`$ and $`\{\alpha _Fp_F\}`$. The S matrix for such a process can be written as an operator in the soft Hilbert space $`_S`$, that collects the states which are almost degenerate with the hard ones, in the form
$$S=𝒰_{\alpha _F\beta _F}^F(a_s,a_s^{})S_{\beta _F\beta _I}^H(p_F,p_I)𝒰_{\alpha _I\beta _I}^I(a_s,a_s^{})$$
(2)
where $`𝒰^F`$ and $`𝒰^I`$ are operator functionals of the soft emission operators $`a_s,a_s^{}`$.
Eq. (2) is supposed to be of general validity , because it rests essentially on the separation of long-time interactions (the initial and final ones described by the $`𝒰`$’s), and the short-time hard interaction, described by $`S^H`$. The real problem is to find the form of the $`𝒰`$’s , which is well known in QED , has been widely investigated in QCD , and is under debate in the electroweak case . Their only general property is unitarity in the soft Hilbert space $`_S`$, i.e.
$$𝒰_{\alpha \beta }𝒰_{\beta \alpha ^{}}^{}=𝒰_{\alpha \beta }^{}𝒰_{\beta \alpha ^{}}=\delta _{\alpha \alpha ^{}}$$
(3)
The key cancellation theorem satisfied by eq. (2) is due to Lee, Nauenberg and Kinoshita , and states that soft singularities cancel upon summation over initial and final soft states which are degenerate with the hard ones:
$$\underset{i\mathrm{\Delta }(p_I)}{\overset{f\mathrm{\Delta }(p_F)}{}}|f|S|i|^2=\mathrm{Tr}__S(𝒰^IS^H𝒰^F𝒰^FS^H𝒰^I)=\mathrm{Tr}_{\alpha _I}(S^H(p_F,p_I)S^H(p_F,p_I))$$
(4)
where $`\mathrm{\Delta }(p_I,p_F)`$ denote the sets of such soft states, and we have used the unitarity property (3).
Although general, the KLN theorem is hardly of direct use, because it involves the sum over the initial degenerate set, which is not available experimentally. In the QED case, however, there is only an abelian charge index, so that $`𝒰_I`$ commutes with $`S^HS^H`$, and cancels out by sum over the final degenerate set only. This is the BN theorem: observables which are inclusive over soft final states are infrared safe.
If the theory is non abelian, like QCD or the electroweak one under consideration, the BN theorem is generally violated, because the initial state interaction is not canceled, i.e., by working in color space,
$$\underset{f\mathrm{\Delta }(p_F)}{}|f|S|i|^2=_S0|𝒰_{\alpha _I\beta _I^{}}^I(S^HS^H)_{\beta _I^{}\beta _I}𝒰_{\beta _I\alpha _I}^I|0_S=(S^HS^H)_{\alpha _I\alpha _I}+\mathrm{\Delta }\sigma _{\alpha _I}$$
(5)
where the $`\alpha _I`$ indices are not summed over and $`\mathrm{\Delta }\sigma _{\alpha _I}`$ is, in general, nonvanishing and IR singular. Fortunately, in QCD the BN cancellation is essentially recovered because of two features: (i) the need of initial color averaging, because hadrons are colorless, and (ii) the commutativity of the leading order coherent state operators ($`𝒰^l`$) for any given color indices :
$$𝒰^l=𝒰^l(a_sa_s^{}),[𝒰_{\alpha \beta }^l,𝒰_{\alpha ^{}\beta ^{}}^l]=0$$
(6)
We obtain therefore
$$\underset{color}{}𝒰_{\alpha _I\beta _I^{}}^l(S^HS^H)_{\beta _I^{}\beta _I}𝒰_{\beta _I\alpha _I}^l=\underset{color}{}(S^HS^H)_{\beta _I^{}\beta _I}𝒰_{\beta _I\alpha _I}^l𝒰_{\alpha _I\beta _I^{}}^l=Tr_{color}S^HS^H$$
(7)
thus recovering an infrared safe result (for subleading features, see Refs. ).
In the electroweak case, in which $`M`$ provides the physical infrared cutoff, there is no way out, because the initial state is prepared with a fixed non abelian charge. Therefore eq. (5) applies, and double log corrections $`\alpha _W\mathrm{log}^2\frac{s}{M^2}`$ must affect any observable associated with a hard process, even the ones which are inclusive over final soft bosons. This fact is surprising, because one would have expected such observables to depend only on energy and on running couplings, while the double logs represent an explicit $`M`$ (infrared cutoff) dependence not yet found before.
In order to compute the uncanceled double logs, a few preliminary remarks are in order. First, we assume the underlying process to be hard, involving a scale much larger than $`M`$. Therefore the lowest order soft contributions to $`\mathrm{\Delta }\sigma `$ in eq. (5) can be simply described by the external (initial) line insertions of the eikonal current
$$J_a^\mu =\frac{p_1^\mu }{p_1k}t_1^a+\frac{p_2^\mu }{p_2k}t_2^a$$
(8)
as depicted in Fig. 1. Secondly, we start from the first non trivial order, where the effect is present and easily understandable. Then, the calculation is obviously gauge invariant, because in the hard (Born) cross section the symmetry is restored, and weak isospin is conserved. For instance, the weak isospin charge resulting from the squared insertion current (8) in the Feynman gauge and from Fig. 1 is
$$(𝒕_1𝒕_1^{})(𝒕_2𝒕_2^{})=(𝒕_1𝒕_1^{})^2=2𝒕_1𝒕_1^{}𝒕_1^2𝒕_1^2$$
(9)
The last expression, obtained using isospin conservation, is identical to the axial gauge result. From this form of the charge it is clear that the $`Z_0`$ and $`\gamma `$ contributions cancel out between real and virtual terms, and only the $`W`$ contribution remains, which is coupled to left handed fermions only. By adding the obvious eikonal radiation factor, we finally obtain the following formulas for the corrections to the Born cross sections $`\sigma _{e^+e^{}}`$ for the hard process defined by eq. (1):
$`\mathrm{\Delta }\sigma _{e^+e^{}}^{RR}`$ $`=`$ $`0`$ (10)
$`\mathrm{\Delta }\sigma _{e^+e^{}}^{LL}`$ $`=`$ $`\mathrm{\Delta }\sigma _{e^{}\overline{\nu }}^{LL}=𝒜_W(s)(\sigma _{e^+\nu }^{LL}\sigma _{e^+e^{}}^{LL}),𝒜_W(s)={\displaystyle \frac{\alpha _W}{4\pi }}\mathrm{log}^2{\displaystyle \frac{s}{M^2}}`$ (11)
where $`s`$ is the c.m. energy squared, $`\alpha _W=g^2/(4\pi )`$, $`L,R`$ refer to the initial fermions chiralities and where use has been made of the isospin conservation constraints $`\sigma _{e^+e^{}}^{LL}=\sigma _{\nu \overline{\nu }}^{LL}`$, $`\sigma _{e^+\nu }=\sigma _{e^{}\overline{\nu }}`$.
Eq.(11) provides a rather general result because the $`W`$ coupling is universal and because only the initial state needs to be specified, so that it applies to various kinds of hard processes of the type of eq. (1). It is also clear that the double log cancellation is recovered by summing over flavors ($`\mathrm{\Delta }\sigma _{e^+e^{}}+\mathrm{\Delta }\sigma _{\nu e^+}=0`$). The actual magnitude of the effect is however dependent on the hard process, because the cross section difference between initial flavors appears in the r.h.s. of eq.(11). To be definite, consider the example of $`e^+e^{}q\overline{q}+X`$, i.e. the total cross section associated to two jets with large transverse energyWe are not concerned here with experimental subtleties, coming from contamination with four jet events coming from two boson production, or with the competing boson fusion mechanism, which starts at higher perturbative order.. The Born amplitude in the hard symmetric limit is $`(g^2𝒕𝑻+g^2yY)/s`$ where $`T`$ (and $`Y`$) denote the quark weak isospin (hypercharge). By squaring and summing over final flavors, we obtain the cross section difference ($`N_f`$ is the number of familiesThe special case of the top quark, requiring a heavy mass cutoff, was considered in, but leads to no important differences at the double log level we are working. and $`N_c`$ is the number of colors):
$$\sigma _{e^+\nu }^{LL}\sigma _{e^+e^{}}^{LL}=\frac{N_fN_c}{12\pi s}(\frac{g^4}{8}\frac{g^4}{4}Y^2)(Y^2=2Y_L^2+Y_R^2+Y_R^{}^2=2\frac{1}{36}+\frac{4}{9}+\frac{1}{9})$$
(12)
which occurs in eq.(11), and yields
$$\frac{\mathrm{\Delta }\sigma _{e^+e^{}}^{LL}}{\sigma _{e^+e^{}}^{LL}}0.8𝒜_W(s)$$
(13)
In this case the non canceling terms are positive for initial $`e^+e^{}`$ beams and no particular suppression is noticed. Given the size of this effect, it is advisable to compute higher orders as well. We think that, due to the inclusive nature of the measurement, the QED scale cannot play, in this case, an important role. It is then tempting to compute higher orders on the basis of the leading form (6) of the coherent state operator with gauge group $`SU(2)\times U(1)`$ and only one cutoff, that is $`M`$. With this assumption, a straightforward calculation shows that $`\gamma `$ and $`Z`$ contributions cancel out in the general case as well, and the W contribution exponentiates in eq.(11) in the form
$$𝒜_W(s)\frac{1}{2}(1e^{2𝒜_W(s)})$$
(14)
This means that, when the energy increases, the noncancellation in eq. (14) becomes maximal and the initial state Sudakov effects equalize eventually the electron and neutrino beam cross sections. The universal exponent appearing in eq. (14) is that of the Sudakov form factor in the adjoint representation, whose relevance was noticed for QCD by Mueller and by Catani et al , and is proved for the present case in .
So far, we have concentrated on lepton colliders, because EW double logs are more directly relevant in such a case. It is amusing to note that the non canceling terms affect hadron colliders as well, because hadrons carry EW charges too. For instance, at parton level, a formula like (11) holds with an initial quark doublet also, while EW effects wash out in the quark-gluon and, of course, in the gluon-gluon cases. We think therefore that such EW double logs are to be seriously considered for LHC too.
To sum up, TeV scale accelerators open up a regime in which we really see non abelian charges at work: even inclusive cross sections have large electroweak corrections. For instance, in our example of $`e^+e^{}`$ into hadrons, while QCD corrections are O($`\frac{\alpha _s}{\pi }`$), the electroweak ones are O($`\frac{\alpha _W}{4\pi }\mathrm{log}^2\frac{s}{M^2}`$), which increases with energy and is already 7% at the TeV threshold. |
warning/0001/cond-mat0001311.html | ar5iv | text | # Off-Equilibrium Effective Temperature in Monatomic Lennard-Jones Glass
\[
## Abstract
The off-equilibrium dynamics of a monatomic Lennard-Jones glass is investigated after sudden isothermal density jumps (crunch) from well equilibrated liquid configurations towards the glassy state. The generalized fluctuation-dissipation relation has been studied and the temperature dependence of the violation factor $`m`$ is found in agreement with the one step replica symmetry breaking scenario, i. e. at low temperature $`m(T)`$ is found proportional to $`T`$ up to an off-equilibrium effective temperature $`T_{eff}`$, where $`m(T_{eff})`$=1. We report $`T_{eff}`$ as a function of the density and compare it with the glass transition temperatures $`T_g`$ as determined by equilibrium calculations.
\]
In the last few years the off-equilibrium dynamics of glassy systems has been the object of intensive analytical and computational studies, due to the very rich phenomenology exhibited and the possibility to obtain detailed information on the phase-space properties of the glassy system itself. One of the most important features of the off-equilibrium regime is the generalization of the fluctuation-dissipation relation . At equilibrium the Fluctuation-Dissipation Theorem (FDT) states that the susceptibility, $`\chi _{AB}(t)`$, that describes the response of the variable $`A(t)`$ to an external field conjugate to the variable $`B(t)`$ is proportional, through the factor $`\beta `$=$`(K_BT)^1`$, to the derivative of the cross correlation function of the variables $`A(t)`$ and $`B(t)`$ themselves, $`C_{AB}(t)`$. In the off-equilibrium case both susceptibilities and correlation functions are two-times quantities and the FDT no longer holds. However, as suggested by Cugliandolo and Kurchan through analytical studies of soluble spin glass models , the FDT can be generalized: the ratio between $`\chi _{AB}(t_1,t_2)`$ and $`\beta \dot{C}_{AB}(t_1,t_2)`$ defines a function, $`X(t_1,t_2)`$, that depends on the times $`t_1,t_2`$ only through the actual value of the correlation function $`C_{AB}(t_1,t_2)`$, $`X(t_1,t_2)`$$``$$`X(C_{AB}(t_1,t_2))`$. The behavior of $`X(C_{AB})`$ gives information on the equilibrium order parameter in the low temperature phase , in agreement with numerical experiments on spin glasses .
The conjecture of the similarity between structural glasses and some spin glass model has been the basis to extend these analytical results on spin glasses to structural glasses. According to the latter conjecture one expects, in structural glasses, that $`X(C_{AB})`$ is a two values function, with $`X(C_{AB})`$=1 at short times and $`X(C_{AB})`$=$`m<1`$ in the non trivial long time region (the so called $`\mathrm{𝑎𝑔𝑖𝑛𝑔}`$ region). The constant $`m`$ is expected to be proportional to the temperature of the system up to an effective temperature $`T_{eff}`$, above which $`m`$=1 and, therefore, the FDT is recovered. In the numerical tests of this scenario performed so far , the system is brought off-equilibrium at constant density by sudden temperature jumps (quench), and the effective temperature $`T_{eff}`$ is found to coincide with the glass transition temperature $`T_g`$. This is explained assuming that during the quench the system ”captures” the properties of the 3$`N`$-dimensional potential energy surface at that temperature where the relaxation time becomes comparable with the simulation time ($`T_g`$) and, at lower $`T`$, it is no longer able to equilibrate. In this way the actual aging dynamics coincides with those at $`T_g`$ even if the “true” temperature is $`T`$$`<`$$`T_g`$. This is a reasonable scenario, but it is still unclear what happens if the system is brought out of equilibrium without crossing the thermodynamic state $`(T_g,\rho )`$.
In this Letter we numerically investigate the off-equilibrium fluctuation-dissipation relation in a simple monoatomic glass, whose particles interacts via a Lennard-Jones (LJ) potential modified in such a way to avoid the crystalization occurring in standard LJ. At variance with the usual constant density temperature jumps method (quench), we use density jumps at constant temperature (crunch) in order to produce off-equilibrium configurations. As expected, at low temperature, we find violation of the equilibrium FDT in agreement with previous computational results in other model glasses . We also find a linear temperature dependence of the off-equilibrium parameter $`m`$ in the aging region, thus confirming the conjecture of one step replica symmetry breaking scenario for real glasses . More important, the effective temperature $`T_{eff}`$, below which $`m`$$`<`$1, is found very close to the glass transition temperature $`T_g`$ at the density reached after the jump. While in the usual quench from high $`T`$ at constant $`\rho `$ it seems natural to visualize the process of freezing at $`T_g`$ during the quench, in our crunch method the identification of the two temperature $`T_{eff}`$ and $`T_g`$ evidences how the effective temperature depends only on the density of the final state and does not depends neither on the initial state nor on the path in $`(T,\rho )`$ plane. This observation gives a strong support to the idea that the aging process can be considered as a pathway in the 3$`N`$-d energy landscape from the off-equilibrium situation towards the glassy minima.
The investigated system is a Modified Lennard-Jones (MLJ) model glass, with $`N`$=256 particles. The potential energy is $`V`$=$`V_{LJ}`$+$`\delta V`$, where $`V_{LJ}`$ is the usual 6-12 Lennard-Jones interaction (in the following all the dimensional quantities are expressed in reduced units) and $`\delta V`$ is a many-body term that inhibits crystalization:
$$\delta V=\alpha \mathrm{\Sigma }_\stackrel{}{q}\theta (S(\stackrel{}{q})S_0)\left[S(\stackrel{}{q})S_0\right]^2.$$
(1)
Here $`S(\stackrel{}{q})`$ is the static structure factor and the parameters in $`\delta V`$ ($`\alpha ^{}`$=0.8 and $`S_0`$=10) have been tuned in order to avoid crystalization without introducing a perturbation on the ”true” LJ dynamics . The sum is made over all $`\stackrel{}{q}`$ with $`q_{max}`$$``$$`\mathrm{\Delta }`$$`<`$$`|\stackrel{}{q}|`$$`<`$$`q_{max}`$+$`\mathrm{\Delta }`$, where $`q_{max}^{}`$=7.12$`(\rho ^{})^{1/3}`$ and $`\mathrm{\Delta }^{}`$=0.34. With this potential we can investigate all the phase space without falling down in spatially ordered configurations and with negligible corrections to the equation of state of monatomic LJ.
First we study equilibrium properties of the system, determining the liquid-glass transition line in the ($`T,\rho `$) plane. The system is equilibrated at high temperature at different densities, from $`\rho ^{}`$=0.87 to $`\rho ^{}`$=1.21, then it is slowly cooled ($`dT^{}/dt^{}`$=4.2$``$10<sup>-4</sup>) measuring the potential energy as a function of $`T`$.
In Fig. 1 we report an example of the potential energy vs $`T^{}`$ at density $`\rho ^{}`$=0.99. It is evident the existence of a temperature $`T_g^{}`$ that marks the transition between two smooth regimes: at high $`T`$, in the liquid, as conjecture by Rosenfeld and Tarazona the potential energy follows a $`T^{3/5}`$ low while at low $`T`$, in the glass, the linear behaviour expected for harmonic glasses is recovered. In this way we are able to reconstruct the liquid-glass transition line $`T_g^{}(\rho ^{})`$, as reported in Fig. 2. Below $`\rho ^{}`$=0.87 a spinodal decomposition takes place, evidenced by the appearance of “bubbles” in the sample.
We focus now on the fluctuation-dissipation relation. This relation connects the correlation function to the response to an external field. Let be $`A(t)`$ and $`B(t)`$ two generic microscopic quantities, and $`\mathrm{\Delta }`$=$`\lambda h(t)B(t)`$ a perturbation added to the Hamiltonian, with $`h(t)`$=0 for $`t`$$`<`$$`t_w`$ and $`\lambda `$ an adimensional control parameter. If we define the correlation function at zero perturbation, $`C_{AB}(t,t_w)`$, as:
$$C_{AB}(t,t_w)=A(t)B(t_w)_{_{\lambda =0}},$$
(2)
and the corresponding response function $`\chi _{AB}`$:
$$\chi _{AB}(t,t_w)=\underset{\lambda 0}{lim}\frac{1}{\lambda }\frac{\delta A(t)__\lambda }{\delta h(t_w)},$$
(3)
(hereafter the subscript $`\lambda `$ indicates that the average is performed in presence of the perturbing term in the Hamiltonian) then the equilibrium fluctuation-dissipation relation takes the form:
$$\chi _{AB}(t,t_w)=\beta C_{AB}(t,t_w)/t_w.$$
(4)
We notice that even if at equilibrium the presence of $`t_w`$ is redundant, as $`\chi _{AB}(t,t_w)`$ and $`C_{AB}(t,t_w)`$ depend only on $`tt_w`$, we choose to use a two times formalism in order to easily generalize the above formulas to off-equilibrium case.
Introducing the integrated response function $`R_{AB}(t,t_w)`$,
$$R_{AB}(t,t_w)=_{t_w}^tdt^{}\chi _{AB}(t,t^{}),$$
(5)
we obtain from eq. (4) the fluctuation-dissipation relation in the integrated form:
$$R_{AB}(t,t_w)=\beta \left[C_{AB}(t,t_w)C_{AB}(t,t)\right].$$
(6)
In the present work we use the quantities $`A(t)`$ and $`B(t)`$:
$`A(t)`$ $`=`$ $`\sqrt{2}(qN)^1\mathrm{\Sigma }_i\mathrm{sin}(\stackrel{}{q}_i\stackrel{}{r}_i(t)+\varphi _i),`$ (7)
$`B(t)`$ $`=`$ $`\sqrt{2}q^1\mathrm{\Sigma }_i\mathrm{sin}(\stackrel{}{q}_i\stackrel{}{r}_i(t)+\varphi _i),`$ (8)
where $`\varphi _i`$ is a uniform random phase, $`\stackrel{}{q}_i=q\stackrel{}{s}_i`$, $`\stackrel{}{s}_i`$ a vector whose Cartesian components are random variables $`\pm 1`$, and $`q`$=$`2\pi n/L`$, with $`n`$ an integer and $`L`$ the sample size. Moreover we choose $`h(t)=h_o\theta (tt_w)`$ ($`h_o^{}`$=1 hereafter). From Eq. (2), after an average over the random phases $`\varphi _i`$, one gets:
$$C_{AB}(t,t_w)=(q^2N)^1\mathrm{\Sigma }_i\mathrm{cos}(\stackrel{}{q}_i\delta \stackrel{}{r}_i(t,t_w))_{_{\lambda =0}},$$
(9)
where $`\delta \stackrel{}{r}_i(t,t_w)`$=$`\stackrel{}{r}_i(t)`$-$`\stackrel{}{r}_i(t_w)`$, and, according to Eqs. (3) and (5),
$`R_{AB}(t,t_w)`$ $`=`$ $`\sqrt{2}(q\lambda N)^1\mathrm{\Sigma }_i[\mathrm{sin}(\stackrel{}{q}_i\stackrel{}{r}_i(t)+\varphi _i)__\lambda +`$ (10)
$``$ $`\mathrm{sin}(\stackrel{}{q}_i\stackrel{}{r}_i(t)+\varphi _i)_{_{\lambda =0}}].`$ (11)
(Fortunately in the limit $`N\mathrm{}`$ the r.h.s of the equations becomes self-averaging). A more suitable and easy to compute form is obtained in the limit $`\stackrel{}{q}_i\delta \stackrel{}{r}_i1`$. In this limit, and averaging over $`\stackrel{}{s}_i`$, Eq. (9) becomes:
$$C_{AB}(t,t_w)q^2+\mathrm{\Delta }(t,t_w)/2,$$
(12)
having defined the mean square displacement:
$$\mathrm{\Delta }(t,t_w)=N^1\mathrm{\Sigma }_i|\stackrel{}{r}_i(t)\stackrel{}{r}_i(t_w)|^2_{_{\lambda =0}}.$$
(13)
Similarly, Eq. (11) becomes:
$$R_{AB}(t,t_w)(\lambda N)^1\mathrm{\Sigma }_i[\stackrel{}{f}_i\stackrel{}{r}_i(t)__\lambda \stackrel{}{f}_i\stackrel{}{r}_i(t)_{_{\lambda =0}}]$$
(14)
where we have introduced the force $`\lambda \stackrel{}{f}_i`$ acting on the particle $`i`$ due to the perturbation term in the Hamiltonian:
$$\stackrel{}{f}_i=\sqrt{2}\stackrel{}{s}_i\mathrm{cos}(\stackrel{}{q}_i\stackrel{}{r}_i(t_w)+\varphi _i),$$
(15)
this is a random variable due to the randomness of $`\varphi _i`$ and $`\stackrel{}{s}_i`$. Therefore, from Eqs. (12,14), Eq. (6) becomes
$$2R(t,t_w)=\beta \mathrm{\Delta }(t,t_w).$$
(16)
In the off-equilibrium case one can always define a violation factor $`X(t,t_w)`$ in such a way to rewrite the Eq. (4) as follows:
$$\chi _{AB}(t,t_w)=\beta X_{AB}(t,t_w)C_{AB}(t,t_w)/t_w.$$
(17)
In some recent papers it has been conjectured that, in the aging region, $`X_{AB}(t,t_w)`$ depends on its arguments only through the correlation function $`C_{AB}(t,t_w)`$. This allows to generalize the fluctuation dissipation relation (6) as:
$$R(C_{AB})=\beta _{C_{AB}(t,t_w)}^{C_{AB}(t,t)}dC_{AB}X(C_{AB}),$$
(18)
or
$$dR(C_{AB})/dC_{AB}=\beta X(C_{AB}).$$
(19)
In our specific case we obtain:
$$2dR(\mathrm{\Delta })/d\mathrm{\Delta }=\beta X(\mathrm{\Delta }),$$
(20)
which generalizes Eq. (16). In this work we have studied the generalized FDT, in the form expressed by Eq. (20). For sake of semplicity rather than perturbing the system with the random force $`\stackrel{}{f}_i`$ in Eq. (15), we have used a random force $`\stackrel{}{f}_{}^{}{}_{i}{}^{}`$ with the same variance of $`\stackrel{}{f}_i`$ ($`f^2=1`$) and Cartesian components $`\pm 1`$ with equal probability. Equation (16) has been numerically tested at equilibrium (see the upper part of Fig. 3).
The off-equilibrium states are obtained by crunches from an initial equilibrated configuration at density $`\rho _0^{}`$=0.95 at a given temperature $`T`$ to a final state at the same $`T`$ and higher density ($`\rho _1^{}`$=1.14 and $`\rho _2^{}`$=1.24). We set the time of the crunch as $`t`$=0. For each investigated $`T`$ and $`\rho `$ we perform an isothermal molecular dynamics simulations with $`\lambda `$=0 and calculate $`\mathrm{\Delta }(t,t_w)`$ as a function of $`t`$ and $`t_w`$. Starting from the same $`t`$=0 configurations, after a waiting time $`t_w`$ we switch on the external field ($`\lambda `$$``$0) and measure the response $`R(t,t_w)`$. The strength of the external field $`\lambda ^{}`$=0.27 was chosen after extensive tests of the linear dependence of the response on the perturbation. All the measured quantities are averaged over $`50`$ initial configurations and $`10`$ random extractions of the variables $`\stackrel{}{f^{}}_i`$. The $`t_w`$ values investigated are $`t_w^{}`$=1, 5, 10. As an example, in Fig. 3 we show the quantities $`\mathrm{\Delta }`$ and $`2R/\beta `$ after a crunch from $`\rho _0^{}`$=0.95 to $`\rho _2^{}`$=1.24 for three different temperatures, $`T^{}`$=1.76, 0.96, 0.48, and for $`t_w^{}`$=5. In the left hand part the two quantities are plotted versus $`\mathrm{log}(t/t_w)`$ while in the right side $`2R/\beta `$ is plotted as a function of $`\mathrm{\Delta }`$. All the data are compatible with the assumption that violation factor $`X`$ depends on $`t`$ and $`t_w`$ only through $`\mathrm{\Delta }(t,t_w)`$, as verified comparing the results obtained for different values of the waiting time. Moreover the measured values $`X(\mathrm{\Delta })`$ are well represented by a piecewise constant behavior ($`X`$=1 for short times and $`X`$$`<`$1 at longer times) suggesting a one step replica symmetry breaking scenario for structural glasses as conjectured in and already found in in the cases of quenching systems of soft spheres and LJ binary mixture.
In Fig. 4 we show the $`T`$ dependence of the violation factor $`m(T)`$ in the aging region, for the two final densities analyzed. It is evident a linear behavior up to a temperature $`T_{eff}`$, as found at low temperature in all know cases in which one step replica symmetry holds. Fitting the data with the expression $`m(T)`$=$`aT\theta (T_{eff}T)`$+$`\theta (TT_{eff})`$ we extract the corresponding $`T_{eff}`$ values, which turn out to correspond with the previously measured equilibrium values of $`T_g`$ as evidenced Fig. 1. The coincidence between $`T_{eff}`$ at a given final density and $`T_g`$ at the same equilibrium density evidences supports the hypothesis that, in the aging region, the point representing the system in the 3$`N`$-d energy landscape moves from the region pertaining to ”high temperatures”, where it has been brought by the sudden density change, towards those configurations where the characteristic time for structural rearrangement ”diverges” .
In conclusion we have, for the first time, studied the off equilibrium dynamics of a simple monatomic glassy system: a LJ potential where the crystalization has been inhibited adding a small many-body term. The system has been brought off-equilibrium via isothermal density jumps and the generalized fluctuation-dissipation relation has been studied. The found piecewise constant behavior of the the violation factor and its linear temperature dependence agrees with one step replica symmetry scenario conjectured for structural glasses. For a given density, the emergent off-equilibrium effective temperature $`T_{eff}`$ seems to be very close the glass transition temperature $`T_g`$, determined by equilibrium molecular dynamics measures. Then, even if brought out of equilibrium by crunches at constant $`T`$ at final density $`\rho `$ without crossing the $`T_g(\rho )`$ transition point, the system captures the properties of energy surface at $`T_{eff}T_g(\rho )`$. The off-equilibrium effective temperature $`T_{eff}`$ is then independent on the particular path in $`(T,\rho )`$ plane crossing the transition line, and depends only on the final reached density. |
warning/0001/astro-ph0001463.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Since 1992, when the first Extra-solar Planetary System (EPS) was discovered<sup>?</sup>, a number of such systems, composed of one or more planets orbiting around a main sequence solar-type star has been observed<sup>?</sup>, and it is conceivable that many others will be found in the near future. Some of the discovered EPS’s are well characterized, since the mass and radius of the main star, the mass of the observed planets and their orbital parameters can be deduced from observations. With this information it is possible to attempt a first estimate of the characteristics of the gravitational radiation emitted by these systems, that, being at a distance of a few tens of parsecs, are very close to us. We shall select a set of EPS’s for which all needed information is available, and compute the gravitational signal emitted by each couple star-planet due to the orbital motion, using the quadrupole formalism. The estimated frequency and amplitude of the waves impinging on Earth will be compared with the emission of the binary pulsar PSR 1913+16<sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>.
We shall then consider a further mechanism of gravitational emission which may play an interesting role in these systems, i.e. the resonant excitation of the modes of oscillation of the star induced by tidal interaction with the orbiting planet. This problem has been studied in the literature for solar type stars with a planet, and for neutron star-neutron star close binary systems by using different approaches, based essentially on the deformation induced by the dynamical tides raised by the companion<sup>?</sup><sup>,</sup><sup>?</sup><sup>,</sup><sup>?</sup>. The emission of gravitational waves associated with the resonant excitation of the modes of a star has been studied by Kojima<sup>?</sup>. He considered a small mass in circular orbit around a compact star, and integrated the equations describing the perturbation of the star excited by the small mass, in general relativity. He showed that a sharp resonance occurs if the frequency of the fundamental mode of the star is twice the orbital frequency, and that the characteristic wave amplitude emitted at that frequency can be up to 100 times larger than that evaluated by the quadrupole formula. In this view, it is interesting to check whether the conditions of resonant excitation can be fulfilled in planetary systems for some modes of the central star, and in particular if this is the case for any of the EPS’s listed in table 1. This will be done in section 4. In addition, we shall discuss the possibility for a planet to get sufficiently close to a star to excite the lowest order $`g`$-modes, and possibly the fundamental one, without being disrupted by tidal forces. We will show that for solar type stars the resonant excitation of the low-order $`g`$-modes may, in principle, be possible.
## 2 Main characteristics of the extrasolar planetary systems
Among the EPS’s discovered up to now, we have selected those for which the parameters of the central star and of the planets, which we need to estimate the gravitational emission, have been determined with sufficient accuracy. These parameters are tabulated in table 1 and 2. Here and in the following, data will be given with the corresponding errors, when available in the literature. In the first column of table 1 we list the selected EPS’s with the corresponding bibliography, and in column 2 the spectral class of the central stars. Most of them belong to a class similar to that of the Sun, which is a G2 V star.
In column 3, 4 and 5 we tabulate, respectively, the mass of the central star, $`M_{},`$ its distance from Earth, $`D`$, \- taken from the Extrasolar Planets Catalog<sup>?</sup>\- and the quantity $`\sqrt{\frac{GM_{}}{R_{}^3}}`$, which will be used in section 4. It should be mentioned that the presence of planets has been discovered mainly by using accelerometric techniques, that is to say, by measuring the variations of the radial velocity of the star caused by its gravitational interaction with the planet. These techniques do not allow to determine the mass of the planet, $`M_p`$, but only the product $`M_p\mathrm{sin}i`$, where $`i`$ is the angle between the line of sight and the normal to the orbital plane. In column 6 we list the values $`M_p\mathrm{sin}i`$ for each planet, normalized to the mass of the central star. From the data of table 1 we see that the closest system is at $`4.70`$ pc, the farthest at $`45`$ pc, whereas the mass of the central star ranges within $`[0.3,1.37]M_{}.`$
In table 2 we tabulate the orbital parameters of the planets orbiting around the stars listed in table 1, i.e. the orbital period $`P`$, the semimajor axis $`a`$, and the eccentricity $`e`$. From the data in column 2 we see that in some cases the period is very short, indicating that the planet gets very close to the central star. For instance, it can be as short as 1.79 days for the planet orbiting around HD 283750.
## 3 Gravitational wave emission due to the orbital motion
We shall first evaluate the amount of gravitational radiation that each couple star-planet emits because of the orbital motion. Since we are not interested in the detailed form of the signal, the quadrupole formalism is sufficient to derive the required information. It is known that two pointlike masses revolving around their common center of mass emit gravitational energy because of their time varying quadrupole moment. To describe their motion it is convenient to choose a coordinate system with the origin located at the center of mass of the system, the orbital plane coincident with the $`x`$-$`y`$ plane and the $`x`$-axis oriented along the relative position vector, $`𝐗𝐱_{}𝐱_p`$, when the planet is at the periastron. The vector $`𝐗`$ describes, in general, an ellipse with semi-major axis $`a`$ and eccentricity $`e`$, and its evolution is described by the parametric equations:
$`\rho `$ $`=`$ $`a(1e\mathrm{cos}u)`$ (1)
$`\alpha `$ $`=`$ $`2\mathrm{arctan}\left[\left({\displaystyle \frac{1+e}{1e}}\right)^{1/2}\mathrm{tan}{\displaystyle \frac{u}{2}}\right].`$ (2)
Here $`\rho =\left|𝐗\right|`$, $`\alpha `$ is the angle between $`𝐗`$ and the $`x`$axis, and $`u`$ is the eccentric anomaly, related to time by the Kepler equation:
$$\omega _kt=ue\mathrm{sin}u,$$
(3)
where $`\omega _k`$ is the keplerian orbital frequency
$$\omega _k=\frac{2\pi }{P}=\left(\frac{GM}{a^3}\right)^{1/2},$$
and $`M=M_{}+M_p`$ is the total mass of the system. From the reduced quadrupole moment $`Q_{kl}`$, given by:
$$Q_{kl}=\mu \left(X^kX^l\frac{1}{3}\delta _l^k\left|𝐗\right|^2\right)$$
(4)
where $`\mu =M_{}M_p/M`$ is the reduced mass, it is straightforward to compute the amplitudes of the metric perturbation, projected on a sphere at distance $`r`$ from the source:
$$h_{ij}(t,r,\theta ,\varphi )=\frac{2G}{c^4r}\left(P_{ik}P_{jl}\frac{1}{2}P_{ij}P_{kl}\right)\ddot{Q}_{kl}(\tau )|_{\tau =tr/c}$$
(5)
where $`P_{ij}=\delta _j^in_in_j`$ is the projector onto the plane orthogonal to the radial unit vector $`𝐧`$. The explicit expressions of $`h_{\theta \theta }`$ and $`h_{\theta \varphi }`$ evaluated by this procedure are
$`h_{\theta \theta }(u,r,\theta ,\varphi )`$ $`=`$ $`{\displaystyle \frac{2G^2M\mu }{c^4ra(1e\mathrm{cos}u)^3}}\times `$
$`\times `$ $`\{(e^21)(1+\mathrm{cos}^2\theta )(12\mathrm{cos}^2\varphi )`$
$`+`$ $`\mathrm{cos}^2u(e\mathrm{cos}u2)[(1e^2)(\mathrm{cos}^2\varphi +\mathrm{cos}^2\theta \mathrm{cos}^2\varphi \mathrm{cos}^2\theta )`$
$`+`$ $`\mathrm{cos}^2\varphi +\mathrm{cos}^2\theta \mathrm{cos}^2\varphi 1]+e\mathrm{cos}u(\mathrm{cos}^2\varphi +\mathrm{cos}^2\theta \mathrm{cos}^2\varphi 1)`$
$`+`$ $`2\sqrt{1e^2}\mathrm{sin}\varphi \mathrm{cos}\varphi (1+\mathrm{cos}^2\theta )\mathrm{sin}u(e\mathrm{cos}^2u2\mathrm{cos}u+e)\},`$
$`h_{\theta \varphi }(u,r,\theta ,\varphi )`$ $`=`$ $`{\displaystyle \frac{4G^2M\mu \mathrm{cos}\theta }{c^4ra(1e\mathrm{cos}u)^3}}\times `$
$`\times `$ $`\{\mathrm{sin}\varphi \mathrm{cos}\varphi [(e^22)\mathrm{cos}^2u(e\mathrm{cos}u2)e\mathrm{cos}u+2(e^21)]`$
$`+`$ $`(2\mathrm{cos}^2\varphi 1)\sqrt{1e^2}\mathrm{sin}u(e\mathrm{cos}^2u2\mathrm{cos}u+e)\}.`$
If the orbit is circular the radiation is emitted at twice the orbital frequency $`(\omega ^{GW}=2\omega _k).`$ By Fourier-transforming the wave amplitudes, and by averaging over the solid angle, it is easy to show that if the orbit is eccentric, waves will be emitted at frequencies multiple of $`\omega _k,`$ and the number of equally spaced spectral lines will increase with the eccentricity.
The average energy flux $`F_n`$ relative to the $`n`$th harmonic is given by:
$$F_n=\frac{c^3(n\omega _k)^2}{8\pi G}\left[\stackrel{~}{h}_{\theta \theta }^{(n)}^2+\stackrel{~}{h}_{\theta \varphi }^{(n)}^2\right]$$
(8)
where $`\stackrel{~}{h}_{\theta \theta }^{(n)}^2`$ and $`\stackrel{~}{h}_{\theta \varphi }^{(n)}^2`$, are the square of the $`n`$th Fourier component of the two independent polarizations, averaged over the solid angle
$$\stackrel{~}{h}_{\theta \theta }^{(n)}^2=\frac{1}{4\pi }𝑑\mathrm{\Omega }\left|h_{\theta \theta }(n\omega _k,r,\theta ,\varphi )\right|^2,\stackrel{~}{h}_{\theta \varphi }^{(n)}^2=\frac{1}{4\pi }𝑑\mathrm{\Omega }\left|h_{\theta \varphi }(n\omega _k,r,\theta ,\varphi )\right|^2.$$
An estimate of the characteristic amplitude of the gravitational waves emitted by a planetary system can now be given by using the well known formula<sup>?</sup>
$$h_c(n\omega _k,r)=\sqrt{\frac{2}{3}}\left[\stackrel{~}{h}_{\theta \theta }^{(n)}^2+\stackrel{~}{h}_{\theta \varphi }^{(n)}^2\right]^{1/2},$$
(9)
where the factor $`\sqrt{2/3}`$ takes into account the average over orientation. The aforementioned procedure, which is equivalent to that introduced by Peters and Mathews<sup>?</sup>, has been applied to compute the characteristic wave amplitude $`h_c`$ impinging on Earth, emitted by our set of EPS’s and, for comparison, by the binary system PSR 1913+16. Some results are shown in figure 1, where $`h_c`$ is plotted as a function of the harmonic index $`n`$. It should be reminded that PSR 1913+16 is composed of two very compact stars with masses $`m_1=1.4411M_{}`$ and $`m_2=1.3874M_{},`$ revolving around their center of mass with an eccentric orbit ($`e=0.617139`$), semimajor axis $`a=1.949010^{12}`$ cm, and keplerian frequency $`\nu _k=3.58310^5Hz.`$ The binary system is at a distance $`D=5`$ kpc from Earth. In the upper panel of figure 1, we show two systems in which the planet moves in a nearly circular orbit around the central star, HD 283750 ($`e=0.02`$) and $`\tau `$ Boo ($`e=0.018`$).
As expected, the emission is concentrated at twice the keplerian frequency, and the amplitude is comparable to the maximum wave amplitude reached by PSR 1913+16, which is shown for comparison at the bottom of the figure. However, the frequency is about ten times lower than that of the binary pulsar. In the lower panel we show the characteristic emission of two EPS’s with high eccentricity, $`16`$ Cyg B ($`e=0.634`$) and HD $`89707`$ ($`e=0.93`$). In this case the gravitational emission is spread over a larger set of frequencies, all multiple of $`\nu _k,`$ as for PSR 1913+16.
In table 3 we tabulate the keplerian frequency, the maximum characteristic amplitude on Earth, and the corresponding frequency, for each EPS. In the last row the same data are listed for PSR 1913+16. For most systems the maximum emission frequency appears to be extremely low, in general smaller than $`10^6Hz`$ except the case of HD 283750, for which $`\nu _{max}=1.310^5Hz.`$
## 4 Resonant excitation of the modes of a star by the orbiting planet
We shall now consider another mechanism through which gravitational waves can be emitted by a system composed of a star and a planet. Since the masses of planets are much smaller than those of stars, in what follows we shall treat the planet as a pointlike mass, whereas the central star will be assumed to have a structure, and to possess a set of eigenmodes of oscillation, the quasi-normal modes, that are associated to the emission of gravitational radiation. We want to establish whether these modes can be excited by a planet.
As mentioned in the introduction, perturbative calculations in the framework of general relativity made by Kojima<sup>?</sup> have shown that, when a small mass moves in a circular orbit with orbital frequency $`\omega _k`$ around a compact star, the fundamental mode is excited if its frequency is $`\omega _f=2\omega _k.`$ In this case, a sharp resonance occurs, and the characteristic wave amplitude can be considerably larger than that evaluated by the quadrupole formula. It may be reminded that the quadrupole radiation by a planet in circular orbit around a star is also emitted at twice the orbital frequency, thus the condition of resonant excitation can also be formulated in terms of the quadrupole emission frequencies.
Similar calculations have never been done for solar type stars; however, Kojima’s results suggest that the resonant excitation of a mode may enhance the gravitational emission also for non compact stars. For this reason, it is interesting to check whether the conditions of resonant excitation can be fulfilled in the planetary systems listed in table 1. We shall first verify whether the maximum quadrupole emission frequency of the planets of our set of EPS’s, $`\nu _{max}`$ (table 3, column 4) is close enough to any of the frequencies of the modes of the central star.
The oscillation frequencies of a star can be computed if we make an assumption on its internal structure, i.e. on the equation of state prevailing in the interior. We shall consider, as an example, a very simple, polytropic model of star, with polytropic index $`n=2.`$ The oscillation frequencies of newtonian polytropic stars are known to scale with the mean density of the star<sup>?</sup>. In table 4 we tabulate the dimensionless eigenfrequencies of the modes, $`\nu _{mode}/(GM_{}/R_{}^3)^{1/2},`$ for the chosen value of $`n`$. In order to explicitely compute the frequency of a given mode for a given star, the entries of table 4 have to be multiplied by the entries of table 1, column 5. For instance, for the star HD 89707 the frequency of the $`f`$-mode is given by $`\nu =0.28\times 6.110^4=1.710^4Hz.`$ This number has to be compared whith $`\nu _{max}=2.210^6Hz,`$ given in table 3 for the same star, which is much smaller. This means that the planet cannot excite the $`f`$-mode of the central star.
By repeating the same calculation for all systems and all modes, we find that it is unlikely that the stellar modes are excited in the EPS’s we consider. This is because the angular velocities reached by the planets are too low, and consequently the frequencies of the radiation they emit are lower than those of the $`f`$-mode or of the lowest-order $`g`$-modes of the star. Higher order $`g`$-modes could be excited, but the efficiency in producing gravitational radiation by this process would be too low.
However, it should be reminded that according to recent theories on the evolution of planetary systems<sup>?</sup>, there could exist planets moving on orbits even closer to the central star than those observed until now. In view of this possibility, it is interesting to investigate whether it is possible for a planet to approach a star at such a short distance that its angular velocity is high enough to excite the $`f`$-mode or the lowest-order $`g`$-modes, without being disrupted by tidal forces. In addition, we shall impose that the star does not accrete matter onto the planet. These conditions are equivalent to impose that neither the planet nor the star overflow their Roche lobe. We would like to stress that this limit takes into account only the gravitational interaction between the planet and the star. There may exist other processes that would prevent the planet to reach the innermost orbit allowed by the Roche-lobe analysis. However, as far as we know, the present knowledge on the formation of planets and on their possible migration toward the central star, still does not allow to firmly establish what is the minimum distance from a star at which a planet can safely sit.
We shall assume, for simplicity, that the planet flies on a circular orbit with keplerian angular velocity $`\omega _k,`$ and consequently emits radiation at the frequency $`\omega ^{GW}=2\omega _k.`$ Let us indicate the dimensionless frequency tabulated in table 4 multiplied by $`2\pi `$ as $`k_{mode}:`$ for instance, $`k_{g_1}=0.750.`$ The condition that a mode of the star is excited by the resonant interaction with the planet, $`\omega _{mode}=2\omega _k,`$ therefore becomes
$$k_{mode}\left[\frac{GM_{}}{R_{}^3}\right]^{1/2}=2\left[\frac{G\left(M_p+M_{}\right)}{a^3}\right]^{1/2},$$
(10)
which can be written as
$$a=\left[\frac{4}{k_{mode}^2}\left(1+\frac{M_p}{M_{}}\right)\right]^{1/3}R_{}.$$
(11)
This equation gives the value the separation star-planet must have for a given mode to be excited.
The further condition that the planet lies inside its Roche lobe, $`R_p<R_{RL},`$ is equivalent to the following constraint on its density:
$$\rho _p>\rho _{RL},$$
where $`\rho _{RL}={\displaystyle \frac{M_p}{\frac{4}{3}\pi R_{RL}^3}}`$ is the critical density. If we now introduce the dimensionless quantity $`\overline{R}_{RL}`$ as given by
$$R_{RL}=a\overline{R}_{RL},$$
(12)
i.e. we set to 1 the radius of the orbit, by the use of eq. (11) we find
$$\rho _{RL}=\frac{M_p}{\frac{4}{3}\pi R_{}^3\left[{\displaystyle \frac{4}{k_{mode}^2}}\left(1+{\displaystyle \frac{M_p}{M_{}}}\right)\right]\overline{R}_{RL}^3},$$
(13)
which can be rewritten as
$$\frac{\rho _{RL}}{\rho _{}}=k_{mode}^2\frac{M_p/M_{}}{4\left(1+M_p/M_{}\right)\overline{R}_{RL}^3}$$
(14)
where $`\rho _{}`$ is the mean density of the central star. Thus, a planet can excite a mode of the star corresponding to an assigned $`k_{mode},`$ without overflowing its Roche lobe, if the ratio between its mean density and that of the central star exceeds the critical ratio (14).
We have computed the dimensionless radius of the Roche lobe $`\overline{R}_{RL}`$ and the corresponding critical ratio (14), for assigned values of the ratio $`M_p/M_{},`$ and the results are shown in table 5, which has to be read as follows.
Suppose that the ratio between the mass of a planet and that of the central star is $`M_p/M_{}=10^3.`$ The frequency of the quadrupole radiation emitted by the planet will coincide with that of the first $`g`$-mode of the star, and the planet will not be disrupted, only if its density is higher that $`1.49\rho _{}.`$ We have also checked whether the star overflows its Roche lobe and accretes matter onto the planet, and excluded from table 5 the corresponding cases. It should be reminded that the ratio between the mean density of the planets of the solar system and that of the Sun is 3.9 for Mercury and the Earth, 3.7 for Venus, 2.8 for Mars, 0.9 for Jupiter, etc. A comparison of these values with the data of table 5 suggests that in principle, there can exist EPS’s in which the first $`g`$-modes could be excited by a resonant process. For instance, a planet like the Earth could approach a polytropic star ($`n=2`$) with the mass of the sun at a distance close enough to excite the mode $`g_1`$ without being disrupted by the tidal interaction, whereas a planet like Jupiter could only be at a distance good to excite the second $`g`$-mode.
By the Roche-lobe analysis we can also deduce another interesting information. Suppose that a planetary system made of a star with the mass of the Sun and a planet in circular orbit, is located at a fiducial distance $`D=10pc.`$ We do not make any assumption on the internal structure of the star, but assign the values of the mass and of the mean density of the planet. In particular we consider four planets, with mass and density equal to that of Mercury, of the Earth, of Jupiter and one with a mass equal to 13 times the mass of Jupiter and similar mean density<sup>?</sup>. We want to answer the following questions:
* what is the minimum radius of the orbit?
* what is the corresponding quadrupole emission frequency $`\nu ^{GW}=2\nu _k\mathrm{?}`$
* what is the corresponding characteristic amplitude on Earth?
The answer is in table 6, where the required data are given for the four planets.
From table 6 we see that planets like Jupiter or bigger could emit quadrupole radiation at a frequency in the bandwidth of space interferometers, and with an amplitude which could be even ten times bigger than that emitted by the binary pulsar PSR 1913+16. In addition, if the quadrupole radiation resonates with a mode of the star, the amount of emitted energy could be even larger.
## 5 Concluding Remarks
In this paper we have studied some of the EPS’s discovered up to now, with the aim of characterizing their gravitational emission.
As far as the quadrupole emission is concerned, we have shown that among those systems there is one, HD 283750, that emits a signal whith a maximum amplitude on Earth even higher than that of the binary pulsar PSR 1913+16, but at a frequency which is six times smaller. Moreover, since the orbit of the planet is nearly circular, the radiation is almost entirely emitted at twice the keplerian frequency, whereas for systems with high eccentricity, as the binary pulsar, the radiation is emitted also at higher multiples of $`\omega _k.`$
A further mechanism of gravitational emission is the resonant excitation of a mode of the star due to the tidal interaction with a planet. For compact stars, this mechanism has proved to be much more efficient than the quadrupole emission, thus it is interesting to check whether the conditions needed for the resonant excitation of the $`f`$\- and $`g`$-modes can be fulfilled in any of the observed EPS’s. Although there exist planets that have a very short orbital period (up to 1.79 days), and orbital frequencies close to $`10^5Hz,`$ these values are too low to induce a resonant excitation of the $`f`$-mode or the lowest $`g`$-modes of a solar type star.
However, since many new EPS’s will certainly be discovered in the near future, we wanted to understand whether more favourable conditions for the emission of gravitational waves may occur in these systems. In particular, since a higher emission frequency would be desirabile both for the excitation of the lowest order $`g`$-modes of the central star, and for a possible detection by space interferometers, we have investigated how close can a planet move on a circular orbit around a star without being tidally disrupted, and without accreting matter from the star.
It should be stressed that the limits we establish by the Roche-lobe analysis do not take into account other processes which may destabilize the orbit of the planets, whose study is beyond the scopes of this paper. Having this caveat in mind, we have established that, in principle, there could exist systems in which the excitability conditions of the lowest order $`g`$-modes could be fulfilled. In this case radiation would be emitted at frequencies of the order of $`10^4Hz.`$ The amplitude of the emitted gravitational signals depends on the ratio $`M_p/M_{},`$ and for a planet like Jupiter or bigger, located at a distance of $`10pc`$, it would range between $`10^{23}10^{22}.`$ The emitted radiation could be even larger, if the system is in a condition of resonant excitation of a mode of the star, and we plan to investigate this problem in detail in a subsequent paper.
Acknowledgements
We would like to thank K. Kokkotas and L. Stella for their valuable suggestions.
References |
warning/0001/hep-th0001151.html | ar5iv | text | # 1 Introduction.
## 1 Introduction.
What is usually defined as a 1+1 dimensional Integrable System is a classical or quantum field theory with the property to have an infinite number of local integrals of motion in involution (LIMI), among which the hamiltonian (energy) operator . This kind of symmetry does not allow the determination of the most intriguing and interesting features of a system because of its abelian character. Instead, the presence of an infinite dimensional non-abelian algebra could complete the abelian algebra giving rise to the possibility of building its representations, i.e. the spectrum (of the energy) and the spectrum of fields. We may name this non-commuting algebra a spectrum generating algebra. In different models and in a mysterious way the presence of this spectrum generating symmetry is very often connected to the abelian one. This is the case of the simplest integrable quantum theories – the two Dimensional Conformal Field Theories (2D-CFT’s) – their common crucial property being covariance under the infinite dimensional Virasoro symmetry, a true spectrum generating symmetry. Indeed, the Verma modules (highest weight representations) of this algebra classify all the local fields in 2D-CFT’s and turn out to be reducible because of the occurrence of vectors of null hermitian product with all other vectors, the so called null-vectors. The factorization by the modules generated over the null-vectors leads to a number of very interesting algebraic-geometrical properties such as fusion algebras, differential equations for correlation functions, etc. .
Unfortunately this beautiful picture collapses when one pushes the system away from criticality by perturbing the original CFT with some relevant local field $`\mathrm{\Phi }`$:
$$S=S_{CFT}+\mu d^2z\mathrm{\Phi }(z,\overline{z}),$$
(1.1)
and from the infinite dimensional Virasoro symmetry only the Poincaré subalgebra survives the perturbation. Consequently, one of the most important open problems in two Dimensional Quantum Field Theories (2D-IQFT’s) is the construction of the spectrum of the local fileds and consequently the computation of their correlation functions. Actually, the CFT possesses a bigger $`𝒲`$-like symmetry and in particular it is invariant under an infinite dimensional abelian subalgebra of the latter . With suitable deformations, this abelian subalgebra survives the perturbation (1.1), resulting in the so-called LIMI. As said before, this symmetry does not carry sufficient information, and in particular one cannot build the spectrum of an integrable theory of type (1.1) by means of LIMI alone.
In the literature, there are several attempts to find spectrum generating algebras at least at conformal point. For instance in , some progress was made in arranging the spectrum generated through the action of objects called spinons in representations of deformed algebras. The kind of structure built by spinons eliminates automatically all null vectors and consequently is unclear how to get (equations for) correlation functions. Besides, the extension of this method in the scaling limit out of criticality is up to now unknown. More recenlty, it has been conjectured in that one could add to LIMI $`I_{2m+1}`$ additional non-commuting charges $`J_{2m}`$ in such a way that the resulting algebra (actually it is not clear from if these objects close an algebra) would be sufficient to create all the states of a particular class of perturbed theory (Restricted Sine-Gordon Theory) of type (1.1). Therein it was also discovered that a sort of null-vector condition appears in the above procedure leading to certain equations for the form factors. However, what remains unclear in are the field theory expressions of null vectors, the general procedure for finding them and above all the symmetry structures lying behind their arising. Besides, heavy use is made of the very specific form of the form factors of the RSG model, and it is not clear how to extend this procedure to other integrable theories of the form (1.1). What is promising in this work is the link the authors created between (quantum) Form Factors and classical solutions of the equation of motion.
With this connection in mind, we present here a general framework to investigate symmetries and related charges in 2D integrable classical field theories. The plan of the paper is as follows.
In Section 2 we present the central idea, i.e. basing the whole construction on a generalization of the so-called Dressing Symmetry Transformations connecting these to the usual way of finding integrable systems . In fact, our basic objects will be the transfer matrix $`T(x;\lambda )`$, which generates the dressing, and the resolvent $`Z(x;\lambda )`$, the dressed generator of the underlying symmetry. Although it is clear from the construction that our method is applicable to any generalized KdV hierarchy , we will be concerned with the semiclassical limit of minimal CFT’s , namely the $`A_1^{(1)}`$-KdV and the $`A_2^{(2)}`$-KdV systems . Besides, we will show how to obtain in a geometrical way from these conformal hierarchies the non conformal Sine-Gordon Model (SGM), i.e. the semiclassical limit of Perturbed Conformal Field Theories (PCFT’s) (1.1). This kind of geometrical point of view on the SGM (Toda field theories, in general), is very useful in derivation of new symmetries and more general Toda field theories are linked to generalized KdV equations.
In Section 3 we will present an alternative approach to the description of the spectrum of the local fields in the classical limit of the 2D Integrable Field Theories. We will build a systematic and geometric method for deriving constraints or classical null vectors without the use of Virasoro algebra.
In Section 4 we will propose a further generalization of the aformentioned dressing transformations. The central idea is that we may dress not only the generators of the underlying Kac-Moody algebra but also differential operators in the spectral parameter, $`\lambda ^m_\lambda ^n`$, forming a $`w_{\mathrm{}}`$ algebra. The corresponding vector fields close a $`w_{\mathrm{}}`$ algebra as well with a Virasoro subalgebra (realized for $`n=1`$) made up of quasi-local and non-local transformations. The regular non-local ones are expressed in terms of vertex operators and complete the quasi-local asymptotic ones to the full Virasoro algebra. All these vector fields do not commute with the KdV hierarchy flows, but have a sort of spectrum generating action on them. Besides, it is very intriguing that only the positive index ones are exact symmetries of the SGM (apparently the negative index ones do not matter particularly in this theory). The apparition of a Virasoro symmetry, with its rich and well known structure, is particularly useful and interesting.
In Section 5 we will deal with the $`A_2^{(2)}`$-KdV showing, as an example of generalization, the building of the spectrum of local fields through the same geometrical lines as before.
In Section 6 we will summarize our results giving some hint on the meaning and quantisation of these symmetries in critical and off–critical theories.
## 2 The usual $`A_1^{(1)}`$–mKdV.
### 2.1 Introductory remarks on integrability.
As already observed in , the classical limit ($`c\mathrm{}`$) of CFT’s is described by the second Hamiltonian structure of the (usual) KdV which is built through the centerless Kac-Moody algebra $`A_1^{(1)}`$ in the Drinfeld-Sokolov scheme . In the classification a generalized modified-KdV (mKdV) hierarchy is attached to each affine Kac-Moody loop algebra $`𝒢`$. Various Miura transformations relate it to the generalized KdV hierarchies, each one classified by the choice of a node $`c_m`$ of the Dynkin diagram of $`𝒢`$. However, nodes symmetrical under automorphisms of the Dynkin diagram lead to the same hierarchy. The classical Poisson structure of this hierarchy is a classical $`w(\stackrel{~}{𝒢})`$-algebra, where $`\stackrel{~}{𝒢}`$ is the finite dimensional Lie algebra obtained by deleting the $`c_m`$ node. In the simplest case the usual mKdV equation
$$_tv=\frac{3}{2}v^2v^{}\frac{1}{4}v^{\prime \prime \prime }$$
(2.1)
describes the temporal flow for the spatial derivative
$$v=\varphi ^{}$$
(2.2)
of a Darboux field $`\varphi `$ defined on a spatial interval $`x[0,L]`$. Assuming quasi-periodic boundary conditions on $`\varphi (x)`$, it verifies by definition the Poisson bracket
$$\{\varphi (x),\varphi (y)\}=\frac{1}{2}s(xy),$$
(2.3)
with the quasi-periodic extension of the sign-function $`s(x)`$ defined as
$$s(x)=2n+1,nL<x<(n+1)L.$$
(2.4)
As all the generalized mKdV, the simplest one (2.1) can be re-written as a null curvature condition
$$[_tA_t,_xA_x]=0$$
(2.5)
for connections belonging to the $`A_1^{(1)}`$ loop algebra
$`A_x`$ $`=`$ $`vh+(e_0+e_1),`$
$`A_t`$ $`=`$ $`\lambda ^2(e_0+e_1vh){\displaystyle \frac{1}{2}}[(v^2v^{})e_0+(v^2+v^{})e_1]{\displaystyle \frac{1}{2}}({\displaystyle \frac{v^{\prime \prime }}{2}}v^3)h`$ (2.6)
where the generators $`e_0,e_1,h`$ are chosen in the canonical gradation of the $`A_1^{(1)}`$ loop algebra
$$e_0=\lambda E,e_1=\lambda F,h=H,$$
(2.7)
with $`E,F,H`$ generators of $`A_1`$ Lie algebra:
$$[H,E]=2E,[H,F]=2F,[E,F]=H.$$
(2.8)
For reasons of simplicity we will deal with the fundamental representation
$$e_0=\left(\begin{array}{cc}0& \lambda \\ 0& 0\end{array}\right),e_1=\left(\begin{array}{cc}0& 0\\ \lambda & 0\end{array}\right),h=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(2.9)
The KdV variable $`u(x,t)`$ is related to the mKdV variable $`\varphi (x)`$ by the Miura transformation
$$u(x)=\varphi ^{}(x)^2\varphi ^{\prime \prime }(x),$$
(2.10)
which is the classical counterpart of the quantum Feigin-Fuchs transformation . A remarkable geometrical interest is obviously attached to the transfer matrix which performs the parallel transport along the $`x`$-axis, i.e. the solution of the boundary value problem
$`_xM(x;\lambda )`$ $`=`$ $`A_x(x;\lambda )M(x;\lambda )`$
$`M(0;\lambda )`$ $`=`$ $`\text{1}.`$ (2.11)
The formal solution of the previous equation
$$M(x,\lambda )=𝒫e^{_0^xA_x(y,\lambda )𝑑y}$$
(2.12)
can be expressed by means of a series of non-negative powers of $`\lambda `$ with an infinite convergence radius and non-local coefficients. The main property of the solution (2.12) is that it allows one to calculate the equal time Poisson brackets between the entries of the monodromy matrix
$$M(\lambda )=M(L;\lambda )=𝒫e^{_0^LA_x(y,\lambda )𝑑y},$$
(2.13)
provided those among the entries of the connection $`A_x`$ are known . The result of this calculation is that the Poisson brackets of the entries of the monodromy matrix are fixed by the so called classical $`r`$-matrix
$$\{M(\lambda )\stackrel{}{,}M(\mu )\}=[r(\lambda \mu ^1),M(\lambda )M(\mu )].$$
(2.14)
In our particular case the $`r`$-matrix is the trigonometric one (calculated, possibly, in the fundamental representation of $`sl(2)`$):
$$r(\lambda )=\frac{\lambda +\lambda ^1}{\lambda \lambda ^1}\frac{HH}{2}+\frac{2}{\lambda \lambda ^1}(EF+FE).$$
(2.15)
By carrying through the trace on both members of the Poisson brackets (2.14) we are allowed to conclude immediately that
$$\tau (\lambda )=trM(\lambda )$$
(2.16)
Poisson-commute with itself for different values of the spectral parameter
$$\{\tau (\lambda ),\tau (\mu )\}=0.$$
(2.17)
In other words, $`\tau (\lambda )`$ is the generating function of the conserved charges in involution, i.e. it guarantees the integrability of the model à la Liouville. Now, it is possible to expand the generating function $`\tau (\lambda )`$ in two different independent ways in order to obtain two different sets of conserved charges in involution. One is the regular expansion in non-negative powers of $`\lambda `$, the other is the asymptotic expansion in negative powers. In the first case the coefficients in the Taylor series are non local charges of the theory, instead in the second one the coefficients in the asymptotic series are the LIMI. Likewise the transfer matrix $`T(x,\lambda )`$ can be expanded in the two ways just mentioned giving rise to different algebraic and geometric structures, as we will see in the following. The regular expansion is typically employed in the derivation of Poisson-Lie structures for Dressing Symmetries . Instead, the second type of expansion plays a crucial rôle in obtaining the flows of the integrable hierarchy and the local integrals of motion that generate symplectically these flows . In this section we will see how the aforementioned approaches are actually reconducible to a single geometrical procedure which, moreover, produces two different kinds of symmetries.
### 2.2 Regular expansion of the transfer matrix and the Dressing Transformations.
The coefficients of the regular expansion of $`\tau `$ (2.16) are non-local integrals of motion in involution (NLIMI). However, these latter may be included in a larger non-abelian algebra of conserved charges, i.e. commuting with local hamiltonians of the mKdV (2.1), but not all between themselves. Actually to get those in a suitable form, we use a slightly different procedure than the usual one , considering a solution of the associated linear problem (2.11) with a different initial condition. Explicitly, we select the following solution of the first equation in (2.12) which contains the fundamental primary field $`e^\varphi `$ :
$$T_{reg}(x;\lambda )=e^{H\varphi (x)}𝒫\mathrm{exp}\left(\lambda _0^x𝑑y(e^{2\varphi (y)}E+e^{2\varphi (y)}F)\right)$$
(2.18)
or, equivalently, defining $`K(x)=e^{2\varphi (x)}E+e^{2\varphi (x)}F`$,
$$T_{reg}(x;\lambda )=e^{H\varphi (x)}\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^k_{xx_1x_2\mathrm{}x_k0}K(x_1)K(x_2)\mathrm{}K(x_k)𝑑x_1𝑑x_2\mathrm{}𝑑x_k.$$
(2.19)
Now we apply the usual dressing techniques using the previous expression $`T_{reg}(x;\lambda )`$ and stressing the point of contact with the derivation of an integrable hierarchy given in . If we define a resolvent $`Z^X(x,\lambda )`$ for the Lax operator $`=_xA_x`$ (2.6) as a solution of the equation
$$[,Z^X(x;\lambda )]=0,$$
(2.20)
it turns out that we may build several solutions by mean of a dressing reformulation of the first equation of (2.11)
$$T_xT^1=_xA_x.$$
(2.21)
In the specific case of the regular expansion (2.19), if $`X=H,E,F`$ is one of the generators (2.8) we get a regular resolvent by dressing
$$Z^X(x,\lambda )=(T_{reg}XT_{reg}^1)(x,\lambda )=\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^kZ_k^X.$$
(2.22)
The definition (2.20) of the resolvent is the key property for the construction of a symmetry algebra, since, once the gauge connection
$$\mathrm{\Theta }_n^X(x;\lambda )=(\lambda ^nZ^X(x;\lambda ))_{}=\underset{k=0}{\overset{n1}{}}\lambda ^{kn}Z_k^X$$
(2.23)
is constructed, the commutator $`[,\mathrm{\Theta }_n^X(x;\lambda )]`$ is in $`\lambda `$ of the same degree of $`[,(\lambda ^nZ^X(x;\lambda ))_+]`$ and hence of degree zero. Therefore, to get a self-consistent gauge transformation
$$\mathrm{\Delta }_n^X=[\mathrm{\Theta }_n^X(x;\lambda ),],$$
(2.24)
we have to require only that the r.h.s. in (2.24) is proportional to $`H`$. This depends, for X fixed, on whether $`n`$ is even or odd. Indeed, a recursive relation between the terms $`Z_n^X`$ in (2.22) follows straightforward from the definition (2.20)
$`_xZ_0^X`$ $`=`$ $`\varphi ^{}[H,Z_0^X]`$
$`_xZ_n^X`$ $`=`$ $`\varphi ^{}[H,Z_n^X]+[E+F,Z_{n1}^X].`$ (2.25)
and allows us to find the modes $`Z_n^X`$ once the different initial conditions are established by inserting the first term of the expansion (2.19) into (2.22)
$$Z_0^H=H,Z_0^E=e^{2\varphi }E,Z_0^F=e^{2\varphi }F.$$
(2.26)
The previous two relations yield the various terms of the expansion of the resolvent in the form
$`Z_{2m}^H(x)=a_{2m}^H(x)H,Z_{2m+1}^H(x)=b_{2m+1}^H(x)E+c_{2m+1}^H(x)F`$
$`Z_{2r}^E(x)=b_{2r}^E(x)E+c_{2r}^E(x)F,Z_{2r+1}^E(x)=a_{2r+1}^E(x)H`$
$`Z_{2p}^F(x)=b_{2p}^F(x)E+c_{2p}^F(x)F,Z_{2p+1}^F(x)=a_{2p+1}^F(x)H`$ (2.27)
where $`a_n^X`$,$`b_n^X`$,$`c_n^X`$ are non-local integral expressions, containing exponentials of the field $`\varphi `$. In addition, the variation (2.24) may be explicitly calculated as
$$\mathrm{\Delta }_n^XA_x=[Z_{n1}^X,E+F]$$
(2.28)
and hence it is clear that $`Z_{n1}^X`$ cannot contain any term proportional to $`H`$. The conclusions about the parity of $`n`$ of (2.28) are that
* in the $`Z^H`$ case, $`n`$ in (2.24) must be even,
* in the $`Z^E`$ and $`Z^F`$ case, $`n`$ must conversely be odd.
From (2.25) it is possible to work out simple recursion relations for the coefficients $`a_n^X`$,$`b_n^X`$,$`c_n^X`$ in (2.27)
$`a_n^X=c_{n1}^Xb_{n1}^X,`$
$`b_{n+1}^X2\varphi ^{}b_{n+1}^X+2a_n^X=0,`$
$`c_{n+1}^X+2\varphi ^{}c_{n+1}^X2a_n^X=0`$ (2.29)
where $`n`$ is even for $`X=H`$ and odd for $`X=E,F`$. In fact, another way to obtain this expansion could be to substitute directly the regular expansion (2.19) in (2.22), but the recursive relations (2.29) (with the initial conditions (2.26)) will provide a contact with the use of the recursive operator in the theory of integrable hierarchies. Now, the action (2.28) of the symmetry generators on the bosonic field is given in terms of $`a_n^X(x)`$ by making use of (2.29)
$$\mathrm{\Delta }_n^X\varphi ^{}=_xa_n^X(x)$$
(2.30)
in which $`n`$ is even for $`X=H`$ and odd for $`X=E,F`$. Instead, using the previous equations of motions and the relations (2.29) it is simple to show that the action on the classical stress-energy tensor $`u(x)`$ (2.10) is given in terms of $`b_{n1}^X(x)`$
$$\mathrm{\Delta }_n^Xu=2_xb_{n1}^X(x),$$
(2.31)
in which $`n`$ is even for $`X=H`$ and odd for $`X=E,F`$. Now, we are interested in finding a recursive relation giving $`a_n^X`$ in terms of $`a_{n+2}^X`$ and vice versa. If we indicate with $`_x^1=_0^x𝑑y`$, the system of equations (2.29) yields easily
$`a_n^X`$ $`=`$ $`_aa_{n+2}^X,_a=v_x^1v_x+{\displaystyle \frac{1}{4}}_x^2,`$
$`a_{n+2}^X`$ $`=`$ $`_a^1a_n^X,_a^1=2(_x^1e^{2\varphi }_x^1e^{2\varphi }+_x^1e^{2\varphi }_x^1e^{2\varphi }),`$ (2.32)
in which $`n`$ is even for $`X=H`$ and odd for $`X=E,F`$. Similarly, for $`b_{n1}^X`$ and $`b_{n+1}^X`$
$`b_{n1}^X`$ $`=`$ $`_bb_{n+1}^X,_b={\displaystyle \frac{1}{2}}(u+_x^1u_x+{\displaystyle \frac{1}{2}}_x^2),`$
$`b_{n+1}^X`$ $`=`$ $`_b^1b_{n1}^X,_b^1=2_x^1{\displaystyle \frac{1}{2}}_a^1_x,`$ (2.33)
in which $`n`$ is even for $`X=H`$ and odd for $`X=E,F`$. The linear differential operators $`_a`$, $`_b`$ are called recursive operators and they generate the integrable flows of an hierarchy (next Section). We have proven here that the proper dressing transformation (2.30),(2.31) can be thought of as generated by the inverse power of the recursive operators, i.e., in a compact notation,
$$\mathrm{\Delta }_n^X\varphi ^{}=_x_a^{\frac{n\nu }{2}}a_\nu ^X,\mathrm{\Delta }_n^Xu=2_x_b^{\frac{n1\nu ^{}}{2}}b_\nu ^{}^X$$
(2.34)
where $`\nu (H)=0`$, $`\nu (E)=\nu (F)=1`$ and $`\nu ^{}(H)=1`$, $`\nu ^{}(E)=\nu ^{}(F)=0`$.
Now we develop a general scheme to find the algebra of the infinitesimal dressing transformations (2.24) and we will use the same procedure to find the commutation relations for the whole symmetry algebra we will discuss in the next sections. The procedure is based on three steps.
###### Lemma 2.1
The equations of motion of the resolvents (2.22) under the flows (2.24) have the form :
$$\mathrm{\Delta }_n^XZ^Y=[\mathrm{\Theta }_n^X,Z^Y]\lambda ^nZ^{[X,Y]}.$$
(2.35)
Proof. As first step we prove that
$$\stackrel{~}{Z}=\mathrm{\Delta }_n^XZ^Y[\mathrm{\Theta }_n^X,Z^Y]$$
(2.36)
is a resolvent (of $``$):
$`0`$ $`=`$ $`\mathrm{\Delta }_n^X[,Z^Y]=[[\mathrm{\Theta }_n^X,],Z^Y]+[,\mathrm{\Delta }_n^XZ^Y]=[[Z^Y,\mathrm{\Theta }_n^X],][[,Z^Y],\mathrm{\Theta }_n^X]+`$ (2.37)
$`+`$ $`[,\mathrm{\Delta }_n^XZ^Y]=[,\mathrm{\Delta }_n^XZ^Y[\mathrm{\Theta }_n^X,Z^Y]].`$
From definition (2.36) it has the form
$$\stackrel{~}{Z}=\underset{l=0}{\overset{\mathrm{}}{}}\lambda ^l(\delta _n^XZ_l^Y\underset{m=0}{\overset{n}{}}[Z_m^X,Z_{l+nm}^X])\underset{l=n}{\overset{1}{}}\lambda ^lZ_{l+n}^{[X,Y]}.$$
(2.38)
Now, we must distinguish two cases. In the first case $`XY`$ and the first negative powers are
$$\stackrel{~}{Z}=\lambda ^rZ_0^{[X,Y]}\mathrm{}.$$
(2.39)
But, given the first term of a resolvent, it is completely determined by the recursive relations (2.25). In the second case $`X=Y`$ and the last term in the previous equation vanishes. Therefore $`\stackrel{~}{Z}`$ is expressed by the series
$$\stackrel{~}{Z}=\underset{l=0}{\overset{\mathrm{}}{}}\lambda ^l(\delta _n^XZ_l^Y\underset{m=0}{\overset{n}{}}[Z_m^X,Z_{l+nm}^X])$$
(2.40)
and the first non-zero term must be a linear combination of $`Z_0^H`$,$`Z_0^E`$,$`Z_0^F`$ (see the first of equations (2.25)). It is clear that this is not possible and consequently $`\stackrel{~}{Z}=0`$, q.e.m..
###### Lemma 2.2
The equations of motion of the connections (2.23) under the flows (2.24) have the form :
$$\mathrm{\Delta }_n^X\mathrm{\Theta }_s^Y\mathrm{\Delta }_s^Y\mathrm{\Theta }_n^X=[\mathrm{\Theta }_n^X,\mathrm{\Theta }_s^Y]\mathrm{\Theta }_{n+s}^{[X,Y]}.$$
(2.41)
Proof. By using the definition of $`\mathrm{\Theta }_n`$ (2.23) and the previous Lemma 2.35, we obtain
$`\mathrm{\Delta }_n^X\mathrm{\Theta }_s^Y`$ $``$ $`\mathrm{\Delta }_s^Y\mathrm{\Theta }_n^X=`$ (2.42)
$`=`$ $`(\lambda ^s[\mathrm{\Theta }_n^X,Z^Y])_{}(\lambda ^{ns}Z^{[X,Y]})_{}(\lambda ^n[\mathrm{\Theta }_s^Y,Z^X])_{}+(\lambda ^{ns}Z^{[Y,X]})_{}=`$
$`=`$ $`([\lambda ^sZ^Y,(\lambda ^nZ^X)_{}])_{}([(\lambda ^sZ^Y)_{},\lambda ^nZ^X])_{}2\mathrm{\Theta }_{n+s}^{[X,Y]}=`$
$`=`$ $`([\lambda ^sZ^Y,(\lambda ^nZ^X)_+\lambda ^nZ^X])_{}([(\lambda ^sZ^Y)_{},\lambda ^nZ^X])_{}2\mathrm{\Theta }_{n+s}^{[X,Y]}=`$
$`=`$ $`([(\lambda ^sZ^Y)_++(\lambda ^sZ^Y)_{},(\lambda ^nZ^X)_+])_{}([(\lambda ^sZ^Y),\lambda ^nZ^X])_{}\mathrm{\Theta }_{n+s}^{[X,Y]}`$
$`=`$ $`([(\lambda ^sZ^Y)_{},(\lambda ^nZ^X)_+])_{}([(\lambda ^sZ^Y),\lambda ^nZ^X])_{}\mathrm{\Theta }_{n+s}^{[X,Y]},`$
from which the claim follows very simply, q.e.m..
###### Theorem 2.1
The algebra of the vector fields (2.24) form a representation of (twisted) Borel subalgebra $`A_1𝐂`$ (of the loop algebra $`A_1^{(1)}`$):
$$[\mathrm{\Delta }_n^X,\mathrm{\Delta }_s^Y]=\mathrm{\Delta }_{n+s}^{[X,Y]};X,Y=H,E,F.$$
(2.43)
Proof. We have to evaluate the action of the commutator in the l.h.s. on $``$ by using the equation of motion of $``$ (2.24), the previous Lemma (2.41) and the Jacobi identity:
$`[\mathrm{\Delta }_n^X,\mathrm{\Delta }_s^Y]`$ $`=`$ $`\mathrm{\Delta }_n^X[\mathrm{\Theta }_s^Y,]\mathrm{\Delta }_s^Y[\mathrm{\Theta }_n^X,]=`$ (2.44)
$`=`$ $`[\mathrm{\Delta }_n^X\mathrm{\Theta }_s^Y\mathrm{\Delta }_s^Y\mathrm{\Theta }_n^X,]+[\mathrm{\Theta }_s^Y,[\mathrm{\Theta }_n^X,]][\mathrm{\Theta }_n^X,[\mathrm{\Theta }_s^Y,]]=`$
$`=`$ $`[\mathrm{\Theta }_{n+s}^{[X,Y]},],`$
which is exactly the claim, q.e.m..
To get from the previous Theorem 2.43 the usual form of the algebra, it is enough to undertake the replacement $`\mathrm{\Delta }\mathrm{\Delta }`$ and untwist. This kind of transformations are historically called dressing transformations . In consideration of the fact that all our symmetries will be obtained by dressing, we will call them proper dressing transformations (or flows). In the case of mKdV, these flows are non local except the first ones which have the form of a Liouville model equation of motion:
$$\mathrm{\Delta }_1^E\varphi ^{}(x)=e^{2\varphi (x)},\mathrm{\Delta }_1^F\varphi ^{}(x)=e^{2\varphi (x)}.$$
(2.45)
In particular, the Theorem 2.43 means that these infinitesimal variations (2.45) generate by successive commutations all the proper dressing flows. In addition, from them it is simple to get the Sine-Gordon equation in light-cone coordinates $`x_\pm `$ for the boson
$$\varphi \frac{i}{2}\varphi ,$$
(2.46)
if we define
$$x_{}=x,\frac{}{x_+}=\frac{1}{2i}(\mathrm{\Delta }_1^E+\mathrm{\Delta }_1^F).$$
(2.47)
Indeed, it comes from (2.45) that
$$_+_{}\varphi =\mathrm{sin}\varphi $$
(2.48)
where it has been defined $`_\pm =\frac{}{x_\pm }`$.
The currents, originating from this symmetry algebra, can easily be found by applying the transformations (2.30) to both members of the continuity equation (2.1)
$$J_{t,n}^X=_xa_n^X(x),J_{x,n}^X=\mathrm{\Delta }_n^X(\frac{1}{2}v^3\frac{1}{4}v^{\prime \prime }).$$
(2.49)
To the $`J_{t,n}^X`$ correspond the non-local charges
$$Q_n^X=_0^LJ_{t,n}^X=a_n^X(L),$$
(2.50)
which are not necessarily conserved (depending on the boundary conditions), due to non-locality.
It is possible to verify by explicit calculations or from the Poisson brackets (2.14) that the charges themselves close a (twisted) Borel subalgebra $`A_1𝐂`$ (of the loop algebra $`A_1^{(1)}`$).
It is interesting to note that the action by which these charges generate the transformations (2.30) is not always symplectic, but only in the case of the variations $`\mathrm{\Delta }_1^E,\mathrm{\Delta }_1^F`$. For instance the following Poisson brackets
$`\mathrm{\Delta }_1^Ev=\{Q_1^E,v\},\mathrm{\Delta }_1^Fv=\{Q_1^F,v\}`$
$`\mathrm{\Delta }_2^Hv=\{Q_2^H,v\}+Q_1^E\{Q_1^F,v\}Q_1^E\{Q_1^F,v\}`$ (2.51)
denote how in the first case the action is symplectic, while in the second it is of Poisson-Lie type .
As a matter of fact, we will compute in the next Section how the transformations (2.24) act on $`_tA_t`$ (and the other higher time Lax operators of the hierarchy), finding that they do as a gauge transformations.
### 2.3 The Integrable Hierarchy and the Asymptotic Dressing.
It is however well-known that besides the regular expansion of the transfer matrix an asymptotic expansion exists for the latter . Since this will play an essential role in our construction, we shall review a few important points in the procedure to obtain the asymptotic expansion. The main idea is to apply a gauge transformation $`S(x)`$ on the Lax operator $``$ in such a way that its new connection $`D(x;\lambda )`$ will be diagonal :
$$(_xA_x(x))S(x)=S(x)(_x+D(x)).$$
(2.52)
Because of the previous equation $`T(x;\lambda )`$ takes the form
$$T(x;\lambda )=KG(x;\lambda )e^{_0^x𝑑yD(y)}$$
(2.53)
where we put $`S=KG`$ with
$$K=\frac{\sqrt{2}}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right),$$
(2.54)
while $`G`$ verifies the following equation
$$_xG+\stackrel{~}{A}_xG=GD,\stackrel{~}{A_x}=K^1A_xK.$$
(2.55)
It is clear now that the previous equation can be solved by finding the asymptotic expansion for $`D(x;\lambda )`$
$$D(x;\lambda )=\left(\begin{array}{cc}d_+& 0\\ 0& d_{}\end{array}\right)=\underset{i=1}{\overset{\mathrm{}}{}}\lambda ^id_i(x)H^i,$$
(2.56)
and expressing the asymptotic expansion of $`G(x;\lambda )`$ in terms of off-diagonal matrices
$$G(x;\lambda )=\left(\begin{array}{cc}1& g_+\\ g_{}& 1\end{array}\right)=H+\underset{j=1}{\overset{\mathrm{}}{}}\lambda ^jG_j(x),$$
(2.57)
where the matrices $`G_j(x)`$ are off-diagonal with entries
$$(G_j(x))_{12}=g_j(x),(G_j(x))_{21}=(1)^{j+1}g_j(x).$$
(2.58)
In addition, the off-diagonal part, $`g_j(x)`$, can be separated obeying an equation of Riccati type. The latter is solved by a recurrence formula for the $`g_j(x)`$
$$g_1=\frac{v}{2},g_{j+1}=\frac{1}{2}(g_j^{}+v\underset{k=1}{\overset{j1}{}}g_{ij}g_j).$$
(2.59)
In addition, it is simple to see that the diagonal part $`d_j(x)`$, $`j>0`$, is related to $`g_j(x)`$ by
$$d_j=(1)^{j+1}vg_j,$$
(2.60)
and is given by $`d_1=1`$, $`d_0=0`$. Note that the $`d_{2n}(x)`$ are exactly the charge densities (of the mKdV equation) resulting from the asymptotic expansion of
$$\tau (\lambda )=trM(\lambda ),M(\lambda )=T(L,\lambda )G^1(0,\lambda )K^1,$$
(2.61)
if we impose quasi-periodic boundary conditions on $`\varphi `$.
On the other hand, it is likewise known that the construction of the mKdV flows goes through the definition of the asymptotic expansion of a resolvent
$$Z^H(x,\lambda )=\underset{k=0}{\overset{\mathrm{}}{}}\lambda ^kZ_k^H,Z_0^H=H.$$
(2.62)
defined through the following property
$$[,Z^H(x;\lambda )]=0.$$
(2.63)
The previous equation may be translated into a recursive system of differential equations for the entries of $`Z^H(x;\lambda )`$ and the solution turns out to have the form
$$Z_{2k}^H(x)=b_{2k}(x)E+c_{2k}(x)F,Z_{2k+1}^H(x)=a_{2k+1}(x)H,$$
(2.64)
where
$`a_{2k+1}=\varphi ^{}b_{2k}{\displaystyle \frac{1}{2}}b_{2k}^{}`$
$`a_{2k+1}=\varphi ^{}c_{2k}+{\displaystyle \frac{1}{2}}c_{2k}^{}`$
$`a_{2k1}^{}=c_{2k}b_{2k}.`$ (2.65)
In a way similar to what has been done for proper dressing transformation, we build through $`Z^H`$ the hierarchy of commuting mKdV flows defining the gauge connections
$$\theta _{2k+1}^H(x;\lambda )=(\lambda ^{2k+1}Z^H(x;\lambda ))_+=\underset{j=0}{\overset{2k+1}{}}\lambda ^{2k+1j}Z_j^H(x),kN$$
(2.66)
and their induced transformation
$$\delta _{2k+1}^HA_x=[\theta _{2k+1}^H(x;\lambda ),].$$
(2.67)
The form of $`Z_{2k+1}^H`$ given by equation (2.64) imposes the self-consistency requirement $`[\theta _n^H(x;\lambda ),]H`$ satisfied only for odd subscript $`n=2k+1`$.
The action (2.67) of the mkdV-flows on the bosonic field can be re-cast in terms of $`a_{2k+1}`$ by using the recursive system (2.65)
$$\delta _{2k+1}\varphi ^{}=_xa_{2k+1}.$$
(2.68)
Instead, using the previous equations of motion and the relations (2.65) it is simple to show that the action on the classical stress-energy tensor $`u(x)`$ (2.10) is given in terms of $`b_{2k}(x)`$
$$\delta _{2k+1}u=2_xb_{2k+2}.$$
(2.69)
These relations have a form similar to that of proper dressing symmetries and consequently we are again interested in separating the recursive system (2.65) into one single recursion relation for $`a_{2k+1}`$. It is simple to show that the desired equation involves exactly the same recursion operator $`_a`$ of equation (2.32):
$$a_{2k+1}(x)=_aa_{2k1}(x),_a=v_x^1v_x+\frac{1}{4}_x^2.$$
(2.70)
This equation determines uniquely $`a_{2k+1}(x)`$, once the initial value of $`a_1(x)`$ has been given. For a similar reason we obtain a recursive differential equation for $`b_{2k+2}(x)`$
$$b_{2k+2}^{}(x)=\frac{1}{2}u^{}b_{2k}+ub_{2k}^{}+\frac{1}{4}b_{2k}^{\prime \prime \prime }.$$
(2.71)
This equation determines uniquely $`b_{2k+2}(x)`$, once the initial value of $`b_0`$ has been given. Indeed it implies
$$b_{2k+2}=_bb_{2k},$$
(2.72)
where $`_b`$ is the same as in equation (2.33). The arbitrariness in the initial condition for $`a_{2k+1}`$ and $`b_{2k}`$ will be fixed in the following using the geometrical interpretation of the resolvent (equation (2.74)). The recursive operators $`_a`$, $`_b`$ generate the integrable flows of the hierarchy as implied by (2.67,2.68):
$$\delta _{2k+1}\varphi ^{}=_x_a^k\varphi ^{},\delta _{2k+1}u=2_x_b^k1.$$
(2.73)
Now, like in the previous Section, it is interesting to interpret this solution $`Z^H`$ to equation (2.63) as generated by dressing through the asymptotic expansion of $`T(x,\lambda )`$ (2.53),(2.56),(2.60)
$$Z^H(x,\lambda )=(THT^1)(x,\lambda ).$$
(2.74)
The previous similarity transformation fixes the initial conditions
$$a_1(x)=\varphi ^{},b_0=1$$
(2.75)
throughout which all the other $`a_{2k+1}`$ and $`b_{2k}`$ can be determined via (2.72) and (2.70). As a consequence of this fact the $`b_{2k}`$ are the densities of the LIMI. Indeed, the differential relation (2.71) (or equivalently (2.72)) coincides (up to a normalization factor of $`u`$) with that satisfied by the expansion modes of the diagonal of the resolvent of the Sturm-Liuoville operator $`_xu`$ . The initial condition $`b_0=1`$ makes the $`b_{2k}`$ proportional to the aforementioned modes .
The observation (2.74) makes evident the same geometrical origin of integrable hierarchies and of their proper dressing symmetries. In addition it will allow us in the sequel to build a more general kind of symmetries and find out their algebra. Indeed, the first generalization of (2.74) consists in the construction of the flows deriving from the resolvents
$$Z^E(x,\lambda )=(TET^1)(x,\lambda ),Z^F(x,\lambda )=(TFT^1)(x,\lambda ).$$
(2.76)
Unlike the previous case, these resolvents possess an expansion in all the powers of $`\lambda `$
$$Z^E(x,\lambda )=\underset{i=\mathrm{}}{\overset{+\mathrm{}}{}}\lambda ^iZ_i^E,Z^F(x,\lambda )=\underset{j=\mathrm{}}{\overset{+\mathrm{}}{}}\lambda ^jZ_j^F.$$
(2.77)
In terms of the data (2.56),(2.57) of the asymptotic transfer matrix, they take the form
$$Z^E(x;\lambda )=\frac{1}{2(1+g_+g_{})}e^{2I(x)}\left(\begin{array}{cc}g_{}^21& (g_{}+1)^2\\ (g_{}1)^2& 1g_{}^2\end{array}\right)$$
(2.78)
and
$$Z^F(x;\lambda )=\frac{1}{2(1+g_+g_{})}e^{2I(x)}\left(\begin{array}{cc}g_+^21& (g_+1)^2\\ (g_++1)^2& 1g_+^2\end{array}\right),$$
(2.79)
after defining the function
$$I(x)=\frac{1}{2}_0^x(d_{}(y)d_+(y))𝑑y=\underset{k=1}{\overset{\mathrm{}}{}}\lambda ^{2k1}_0^xd_{2k+1}(y)𝑑y,$$
(2.80)
which generates the LIMI once calculated in $`x=L`$. Now, it is easy one to convince himself that the entries of the resolvents (2.78),(2.79) admit an expansion in all the (positive and negative) powers of $`\lambda `$, i.e. that the modes $`Z_i^E`$ and $`Z_j^F`$ of the series (2.77) are made up of linear combinations of all the three Lie algebra generators $`E`$,$`F`$,$`H`$. This implies the impossibility to satisfy the self-consistency condition $`\delta H`$. Nevertheless , we can go over this difficulty by defining two other resolvents, combinations of the previous ones
$$Z^+(x,\lambda )=(T(E+F)T^1)(x,\lambda ),Z^{}(x,\lambda )=(T(EF)T^1)(x,\lambda ).$$
(2.81)
By using the expressions (2.78),(2.79), we obtain the following formulæ for the entries of $`Z^\pm `$
$`(Z^+)_{11}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_{}^2g_+^22)\mathrm{cosh}2I+(g_+^2g_{}^2)\mathrm{sinh}2I],`$
$`(Z^+)_{12}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_{}^2g_+^2+2g_{}+2g_+)\mathrm{cosh}2I`$
$``$ $`(g_{}^2+g_+^22g_{}2g_++2)\mathrm{sinh}2I],`$
$`(Z^+)_{21}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_+^2g_{}^2+2g_{}+2g_+)\mathrm{cosh}2I+`$
$`+`$ $`(g_{}^2+g_+^22g_{}+2g_++2)\mathrm{sinh}2I],`$
$`(Z^+)_{22}`$ $`=`$ $`(Z^+)_{11};`$ (2.82)
and
$`(Z^{})_{11}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_{}^2g_+^2)\mathrm{cosh}2I(g_+^2+g_{}^22)\mathrm{sinh}2I],`$
$`(Z^{})_{12}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_{}^2+g_+^2+2g_{}2g_++2)\mathrm{cosh}2I`$
$``$ $`(g_+^2g_{}^22g_{}2g_+)\mathrm{sinh}2I],`$
$`(Z^{})_{21}`$ $`=`$ $`{\displaystyle \frac{1}{2(1+g_+g_{})}}[(g_+^2+g_{}^2+2g_+2g_++2)\mathrm{cosh}2I+`$
$`+`$ $`(g_{}^2g_+^22g_{}2g_+)\mathrm{sinh}2I],`$
$`(Z^{})_{22}`$ $`=`$ $`(Z^{})_{11}.`$ (2.83)
If we assume to denote by $`(e)`$ and $`(o)`$ series with only even and odd powers of $`\lambda `$ respectively, we have
$`g_+g_{}=(e)`$ $`,`$ $`g_+^2+g_{}^2=(e),g_+^2g_{}^2=(o),`$
$`g_++g_{}=(o)`$ $`,`$ $`g_+g_{}=(e).`$ (2.84)
Consequently, the parity of the entries of $`Z^\pm `$ is given by
$$Z^+=\left(\begin{array}{cc}(e)& (o)\\ (o)& (e)\end{array}\right),Z^{}=\left(\begin{array}{cc}(o)& (e)\\ (e)& (o)\end{array}\right),$$
(2.85)
or equivalently by
$`Z^+={\displaystyle \underset{i=\mathrm{}}{\overset{+\mathrm{}}{}}}\lambda ^{2i1}(b_{2i+1}^+E+c_{2i+1}^+F)+{\displaystyle \underset{i=\mathrm{}}{\overset{+\mathrm{}}{}}}\lambda ^{2i}a_{2i}^+H`$
$`Z^{}={\displaystyle \underset{j=\mathrm{}}{\overset{+\mathrm{}}{}}}\lambda ^{2j}(b_{2i}^{}E+c_{2i}^{}F)+{\displaystyle \underset{j=\mathrm{}}{\overset{+\mathrm{}}{}}}\lambda ^{2j1}a_{2i+1}^{}H.`$ (2.86)
It follows that self-consistency requirement may now be satisfied and we are allowed to define two new series of dressing transformations through the connections
$`\theta _{2i}^+(x;\lambda )=(\lambda ^{2i}Z^+(x;\lambda ))_+={\displaystyle \underset{l=\mathrm{}}{\overset{2i}{}}}\lambda ^{2il}Z_l^+(x),iZ`$
$`\theta _{2j+1}^{}(x;\lambda )=(\lambda ^{2j+1}Z^{}(x;\lambda ))_+={\displaystyle \underset{l=\mathrm{}}{\overset{2j+1}{}}}\lambda ^{2j+1l}Z_l^{}(x),jZ`$ (2.87)
which are no more finite sums. Finally, the following gauge transformations of the Lax operator $``$
$$\delta _{2i}^+A_x=[\theta _{2i}^+(x;\lambda ),],\delta _{2j+1}^{}A_x=[\theta _{2j+1}^{}(x;\lambda ),]$$
(2.88)
yield this compact form for the additional mKdV flows
$$\delta _{2i}\varphi ^{}=_xa_{2i}^+,\delta _{2j+1}\varphi ^{}=_xa_{2j+1}^{}.$$
(2.89)
These flows are complicated series in $`x`$ with quasi-local coefficients, so that it would be very difficult to find their commutation relations by direct computation.
Therefore, to find the algebra of these additional dressing transformations (2.88), we use the previous procedure based on three steps.
###### Lemma 2.3
The equations of motion of the resolvents (2.81),(2.74) under the flows (2.88) have the form :
$$\delta _n^XZ^Y=[\theta _n^X(x;\lambda ),Z^Y]\lambda ^nZ^{[X,Y]},$$
(2.90)
where now $`X,Y=H,E+F,EF`$.
Proof. We omit the specific proof because it can be carried out along the lines of the analogous Lemma (2.35).
###### Lemma 2.4
The equations of motion of the connections (2.66),(2.87) under the flows (2.88) have the form :
$$\delta _n^X\theta _s\delta _s^Y\theta _n=[\theta _n^X,\theta _s^Y]\theta _{n+s}^{[X,Y]}.$$
(2.91)
Proof. The proof is analogous to that of Lemma (2.41).
###### Theorem 2.2
The algebra of the asymptotic dressing vector fields (2.67),(2.88) is:
$`[\delta _{2k+1}^H,\delta _{2i}^+]`$ $`=`$ $`2\delta _{2k+2i+1}^{},kN,iZ,`$
$`[\delta _{2k+1}^H,\delta _{2j+1}^{}]`$ $`=`$ $`2\delta _{2k+2j+2}^+,jZ,`$
$`[\delta _{2k+1}^H,\delta _{2l+1}^H]`$ $`=`$ $`0,lN,`$
$`[\delta _{2i}^+,\delta _{2j+1}^{}]`$ $`=`$ $`2\delta _{2i+2j+1}^H,`$ (2.92)
where in the last relation we have defined $`\delta _{2k+1}^H=0`$ if $`k<0`$.
Proof. As in the proof of Theorem (2.43), the action on $``$ of the commutators in the l.h.s. can be calculated by using the equation of motion of $``$ (2.88),(2.67), the previous Lemma (2.91) and the Jacobi identity, q.e.m..
The previous Theorem 2.2 proves that the KdV flows form an hierarchy ( they commute with each other). Besides, they are local and the Lemma 2.91 ensures that each Lax connection transforms in a gauge way under a generic flow. This is why we may attach a time $`t_k`$ to each flow $`\delta _{2k+1}^H`$ and think to each flow as a true symmetry of all the others. Instead, the additional asymptotic flows $`\delta _{2i}^+`$ and $`\delta _{2j+1}^{}`$ do not commute with the hierarchy flows, but close an algebra in which they are (in some sense) spectrum generating symmetries.
It is easy to prove that the proper dressing transformation are true symmetries of the hierarchy as well.
###### Lemma 2.5
The transformation of the resolvent (2.74) under the regular flows (2.24) and the evolution with the times $`t_k`$ of the regular resolvents (2.22) have the same form of the hierarchy flow of $``$:
$$\delta _n^XZ^H=[\mathrm{\Theta }_n^X,Z^H],\delta _{2k+1}^HZ^X=\theta _{2k+1}^H,Z^X],$$
(2.93)
where for the regular resolvents we have $`X=H,F,F`$.
###### Lemma 2.6
The mKdV flows (2.67) act as gauge transformations on the connections (2.23) of the proper dressing flows:
$$\delta _{2k+1}^H\mathrm{\Theta }_n\mathrm{\Delta }_n^X\theta _{2k+1}=[\theta _{2k+1}^H,\mathrm{\Theta }_n^X].$$
(2.94)
###### Theorem 2.3
The proper dressing vector fields (2.24) commute with the mKdV flows (2.67):
$`[\delta _{2k+1}^H,\mathrm{\Delta }_n^X]`$ $`=`$ $`0.`$ (2.95)
In particular, the previous Theorem 2.3 implies that the light cone evolution $`_+`$ commutes with all the KdV flows, i.e. a different way to say that the KdV hierarchy is a symmetry of the light-cone SG. In particular, the symmetry generator $`\delta _{2k+1}^H`$ maps, at infinitesimal level, solution of SG into solution.
In consideration of the fact that these theories are classical limits of CFT’s and PCFT’s, let us concentrate our attention on the phase spaces of mKdV and KdV systems, i.e. those objects which at the quantum level constitute the spectrum of fields.
## 3 The spectrum of fields in the $`A_1^{(1)}`$ framework.
In the mKdV theory a local field is a polynomial in $`v=\varphi ^{}`$ and its derivatives and the space spanned by these polynomials is the space (Verma module) of the descendant of the identity. Instead, the action (simple product in the classical theory) of these polynomials on a primary field $`e^{m\varphi }`$ generates the space (Verma module) of the descendants of this primary field. In our approach to the spectrum of this classical limit of CFT’s, we propose here to treat the gauge fields, i.e. the entries of $`Z^H`$ as fundamental fields. Let us start by considering the composite fields $`a_{2n+1}`$, $`b_{2n}`$, and $`c_{2n}`$ of (2.64). In this Section we will suppress the index H. The differential equations (2.65) tell us immediately that not all of them are independent, we may for example express the $`c_{2n}`$ in terms of the basic fields $`b_{2n}`$ and $`a_{2n+1}`$. We use now Lemma 2.3
$$\delta _{2k+1}Z=[\theta _{2k+1},Z]$$
(3.1)
which allows us to establish the action of each $`\delta _{2k+1}`$ on these fields
$`\delta _{2k+1}a_{2n+1}={\displaystyle \underset{i=0}{\overset{n}{}}}(a_{2n+2k2i+1}^{}b_{2i}a_{2i1}^{}b_{2m+2k2i+2}),`$
$`\delta _{2k+1}b_{2n}=2{\displaystyle \underset{i=0}{\overset{n1}{}}}(a_{2n+2k+1}b_{2n+2k2i}a_{2n+2k2i+1}b_{2i})`$ (3.2)
Therefore, according to our conjecture the linear generators of the mKdV identity Verma module $`𝒱_\text{1}^{mKdV}`$ are made up of the repeated actions of the $`\delta _{2k+1}`$ on the polynomials $`𝒫(b_2,b_4,\mathrm{},b_{2N},a_1,a_3,\mathrm{},a_{2P+1})\}`$ in the $`b_{2n}`$ and the $`a_{2k+1}`$
$$𝒱_\text{1}^{mKdV}=\{linearcombinationsof\delta _{2k_1+1}\delta _{2k_2+1}\mathrm{}\delta _{2k_M+1}𝒫\},$$
(3.3)
with a natural gradation provided by the subscripts. Actually, the Verma module $`𝒱_\text{1}^{mKdV}`$ exhibits several null vectors, i.e. polynomials in the $`b_{2n}`$ and the $`a_{2k+1}`$ which are zero. This is due to the very simple constraint on $`Z^H`$
$$(Z^H)^2=\text{1}$$
(3.4)
originating from the dressing relation with the transfer matrix $`T`$ (2.74). The constraints (3.4) may be rewritten through the modes of $`a_{2k+1}`$ and $`b_{2k}`$
$$𝒞_{2n}=\underset{i=0}{\overset{n}{}}b_{2n2i}(b_{2i}+a_{2i1}^{})+\underset{i=0}{\overset{n1}{}}a_{2n2i1}a_{2i+1}=0,$$
(3.5)
and produce null-vectors under the application of mKdV flows $`\delta _{2k+1}`$. These latter generate linearly the graded vector space (Verma module) of all null vectors
$$𝒩=\{linearcombinationsof\delta _{2k_1+1}\delta _{2k_2+1}\delta _{2k_3+1}\mathrm{}\delta _{2k_Q+1}𝒞_{2n}\}.$$
(3.6)
In conclusion our conjecture is that the (conformal) family of the identity $`[\text{1}]^{mKdV}`$ of the mKdV hierarchy is obtained as a factor space:
$$[\text{1}]^{mKdV}=𝒱_\text{1}^{mKdV}/𝒩.$$
(3.7)
On the other hand, in order to deduce the form of the Verma module $`𝒱_\text{1}^{KdV}`$ of the identity for the KdV hierarchy we have to make three observations:
1. the recursive formula (2.72) proves that $`b_{2n}`$ are polinomials of the KdV field $`u(x)`$ and its derivatives, whereas the $`a_{2k+1}`$ do not enjoy this property;
2. the variation of $`b_{2n}`$ in (3.2) can be written accidentally in terms of the $`b_{2k}`$ alone, using the relationships (2.65) between $`a_{2k+1}`$ and $`b_{2k}`$
$$\delta _{2k+1}^Hb_{2n}=\underset{i=0}{\overset{n1}{}}(b_{2n+2k2i}^{}b_{2i}b_{2i}^{}b_{2m+2k2i});$$
(3.8)
3. also the null vector space $`𝒩`$ can be spanned by the $`b_{2k}`$ alone
$$𝒞_{2n}=b_{2n}+\underset{i=1}{\overset{n}{}}[b_{2n2i}b_{2i}2b_2b_{2n2i}b_{2i2}\frac{1}{2}b_{2n2i}b_{2j2}^{}+\frac{1}{4}b_{2n2j}^{}b_{2j2}^{}]=0,$$
(3.9)
using the relationships (2.65) between $`a_{2k+1}`$ and $`b_{2k}`$.
Therefore, we conjecture that the Verma module $`𝒱_\text{1}^{KdV}`$ shall be linearly generated by elements given by repeated actions of the $`\delta _{2k+1}`$ on the polynomials $`𝒫(b_2b_4\mathrm{}b_{2N})`$ in the $`b_{2n}`$
$$𝒱_\text{1}^{KdV}=\{linearcombinationsof\delta _{2k_1+1}\delta _{2k_2+1}\mathrm{}\delta _{2k_M+1}𝒫\}.$$
(3.10)
It turns out to be a sort of reduction of the Verma module $`𝒱_\text{1}^{mKdV}`$ (3.3) of the mKdV hierarchy. As for the mKdV case the (conformal) family of the identity $`[\text{1}]^{KdV}`$ of the KdV hierarchy is obtained as a factor space of $`𝒱_\text{1}^{KdV}`$ over $`𝒩`$ given by (3.6) and (3.9):
$$[\text{1}]^{KdV}=𝒱_\text{1}^{KdV}/𝒩.$$
(3.11)
Therefore we are led to the same scenario that arises also in the classical limit of the construction . Nevertheless, in our approach the generation of null-vectors is automatic and geometrical (see equations (3.8) and (3.9)). In addition , our approach is applicable to any other integrable system, based on a Lax pair formulation. We will illustrate this fact below by using the example of the $`A_2^{(2)}`$-mKdV system. Other local fields of the mKdV system are the primary fields, i.e. the exponential $`e^{m\varphi },m=0,1,2,3,\mathrm{}`$ of the bosonic field. Indeed, for $`m=0`$ we obtain just the identity 1, the fundamental primary field $`e^\varphi `$ (m=1) appears in the regular expansion (2.19) of the transfer matrix $`T(x;\lambda )`$ and the other primary fields $`e^{m\varphi },m>1`$ are the ingredients of the regular expansion of the power $`T^m(x;\lambda )`$. The previous construction of the identity operator family suggests the following form for the Verma module $`𝒱_m^{mKdV}`$ of the primary $`e^{m\varphi },m=0,1,2,3,\mathrm{}`$:
$$𝒱_m^{mKdV}=\{linearcombinationsof\delta _{2k_1+1}\delta _{2k_2+1}\mathrm{}\delta _{2k_M+1}[𝒫e^{m\varphi }]\},$$
(3.12)
where $`𝒫(b_2,b_4,\mathrm{},b_{2N},a_1,a_3,\mathrm{},a_{2P+1})`$ are polinomials in the $`b_{2n}`$ and the $`a_{2k+1}`$. As for the identity family, we have to subtract all the null-vectors (3.2) and (3.9). Besides, in this case, we have to take into account the null-vectors coming from the equations of motion of the power $`T_{reg}^m(x;\lambda )`$ of the regular expansion
$$\delta _{2k+1}T_{reg}^m=\underset{j=1}{\overset{m}{}}T_{reg}^j\theta _{2k+1}T_{reg}^{mj}.$$
(3.13)
By successive applications of (3.2), (3.9) and (3.13) we obtain the whole null-vector set $`𝒩_𝐦^{KdV}`$. In conclusion the spectrum is again the factor space
$$[𝐦]=𝒱_m^{mKdV}/𝒩_𝐦^{KdV}.$$
(3.14)
Similarly, the construction of the identity operator family suggests the following form for the Verma module $`𝒱_m^{KdV}`$ of the primary $`e^{m\varphi },m=0,1,2,3,\mathrm{}`$:
$$𝒱_m^{KdV}=\{linearcombinationsof\delta _{2k_1+1}\delta _{2k_2+1}\mathrm{}\delta _{2k_M+1}[𝒫(b_2,b_4,\mathrm{},b_{2N})e^{m\varphi }]\},$$
(3.15)
with $`𝒫(b_2,b_4,\mathrm{},b_{2N})`$ polinomials in $`b_{2n}`$. Again, by successive applications of (3.8), (3.9) and (3.13) we obtain the whole null-vector linear space $`𝒩_𝐦^{KdV}`$. In conclusion, the spectrum is again a factor space
$$[𝐦]=𝒱_m^{KdV}/𝒩_𝐦^{KdV}.$$
(3.16)
Of course, we have checked all our conjectures up to high gradation of the null vectors. Nevertheless, we did not manage to generate the spectrum of fields only through the asymptotic symmetry of Theorem 2.2. For dimensional arguments it is plausible to make the substitution
$`a_{2k+1}`$ $`\delta _{2k+1}^+,kZ`$ (3.17)
$`b_{2k}`$ $`\delta _{2k}^{},`$
but now the null vector meaning and origin should be completely different.
## 4 A non local Virasoro symmetry by dressing.
At this point we are in a position to construct in a natural way more general kinds of dressing-like symmetries. It is well known that the vector fields $`l_m=\lambda ^{m+1}_\lambda `$ on the circumpherence realize the centerless Virasoro algebra
$$[l_m,l_n]=(mn)l_{m+n}.$$
(4.1)
A very natural dressing is represented by the resolvents
$`Z_m^V=T_{reg}l_mT_{reg}^1,m<0`$
$`Z_m^V=T_{asy}l_mT_{asy}^1,m0`$ (4.2)
where we have to use the different regular and asymptotic tranfer matrices, $`T_{reg}`$ and $`T_{asy}`$. Of course, they satisfy the usual definition of resolvent
$$[,Z_m^V(x;\lambda )]=0$$
(4.3)
and, as in the previous cases, they have two different kinds of expansions
$`(Z_1^V)_{reg}=T_{reg}l_1T_{reg}^1={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\lambda ^nZ_{n+1}^{reg}_\lambda `$
$`(Z_1^V)_{asy}=T_{asy}l_1T_{asy}^1={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\lambda ^nZ_{n1}^{asy}_\lambda `$ (4.4)
and consquently the mode expansion of the more general Virasoro resolvent (4.2). In the same way, (4.3) authorizes us to define gauge connections
$`\theta _m^V=(Z_m^V)_{}={\displaystyle \underset{n=0}{\overset{m2}{}}}\lambda ^{n+1+m}Z_{n+1}^{reg}\lambda ^{m+1}_\lambda ,m<0`$
$`\theta _m^V=(Z_m^V)_+={\displaystyle \underset{n=0}{\overset{m+1}{}}}\lambda ^{m+1n}Z_{n1}^{asy}\lambda ^{m+1}_\lambda ,m0`$ (4.5)
and the relative gauge transformations
$$\delta _m^VA_x=[\theta _m^V(x;\lambda ),].$$
(4.6)
Finally, we have to verify the consistency of this gauge transformation requiring $`\delta _m^VA_x=H\delta _m^V\varphi ^{}`$ for positive and negative $`m`$. It is very easy to see that this requirement imposes $`m`$ to be even. Indeed, from (4.3) or (4.2) it is simple to derive the form of the generic term of the expansions (4.4)
$`Z_{2n1}^{reg}=b_{2n1}^VE+c_{2n1}^VF`$ , $`Z_{2n}^{reg}=a_{2n}^VH,n>0`$
$`Z_{2n3}^{asy}=\beta _{2n3}^VE+\gamma _{2n3}^VF`$ , $`Z_{2n2}^{asy}=\alpha _{2n2}^VH,n>0.`$ (4.7)
In addition, we can easily find recursive relations for the regular coefficients
$$b_{2k+1}^V(x)=_b^1b_{2k1}^V(x),a_{2n+2}^V(x)=_a^1a_{2n}^V(x),$$
(4.8)
and for the asymptotic coefficients
$$\beta _{2n1}^V(x)=_b\beta _{2n3}^V(x),\alpha _{2n}^V(x)=_a\alpha _{2n2}^V(x),$$
(4.9)
where the recursive operators $`_a`$,$`_b`$ are given in (2.32),(2.33) and (4.2) fixes the initial conditions:
$`b_1^V`$ $`=e^{2\varphi }_x^1e^{2\varphi },a_2={\displaystyle _0^x}dx_2{\displaystyle _0^{x_2}}dx_12\mathrm{cosh}[\varphi (x_1)\varphi (x_2)],`$ (4.10)
$`\beta _1^V`$ $`=x,\alpha _0^V=x\varphi ^{}.`$
In conclusion the Virasoro mKdV flows are given by:
$`\delta _{2m}^V\varphi ^{}=_xa_{2m}^V,m<0`$
$`\delta _{2m}^V\varphi ^{}=_x\alpha _{2m}^V,m0,`$ (4.11)
and explicit examples of the first flows are
$`\delta _2^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}e^{2\varphi (x)}{\displaystyle _0^x}𝑑ye^{2\varphi (y)}=e^{2\varphi (x)}B_1e^{2\varphi (x)}C_1`$
$`\delta _4^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}(3B_3(x)A_2(x)B_1(x))e^{2\varphi (x)}(3C_3(x)D_2(x)C_1(x))`$
$`\delta _6^V\varphi ^{}`$ $`=`$ $`e^{2\varphi (x)}(5B_5(x)3A_4(x)B_1(x)+A_2(x)B_3(x))`$
$``$ $`e^{2\varphi (x)}(5C_5(x)3D_4(x)C_1(x)+D_2(x)C_3(x))`$
$`\delta _0^V\varphi ^{}`$ $`=`$ $`\varphi ^{}+x\varphi ^{\prime \prime }`$
$`\delta _2^V\varphi ^{}`$ $`=`$ $`2xa_3^{}+6g_32g_1^3+2g_1^{}{\displaystyle _0^x}d_1,`$ (4.12)
where we have used a convenient notation for the entries of the regular expansion
$$T_{reg}(x,\lambda )=\left(\begin{array}{cc}A& B\\ C& D\end{array}\right),$$
(4.13)
with $`A=e^\varphi (1+_1^{\mathrm{}}\lambda ^{2n}A_{2n})`$, $`B=e^\varphi _0^{\mathrm{}}\lambda ^{2n+1}B_{2n+1}`$, $`C[\varphi ]=B[\varphi ]`$ and $`D[\varphi ]=A[\varphi ]`$. We stress that negative subscript variations have a form very similar to that of the regular dressing flows ((2.24) with $`X=H`$) $`\mathrm{\Delta }_{2r}^H`$. Nevertheless, in spite of the commutativity $`[\mathrm{\Delta }_{2r}^H,\mathrm{\Delta }_{2s}^H]=0`$ we will see that they obey instead Virasoro commutation relations. From the actions (4.12) the transformations of the classical primary fields $`e^\varphi `$ follow. For example:
$`\delta _2^Ve^\varphi `$ $`=`$ $`(D_2A_2)e^\varphi `$
$`\delta _4^Ve^\varphi `$ $`=`$ $`[(3D_4C_3B_1)(3A_4B_3C_1)]e^\varphi ,`$
$`\delta _0^Ve^\varphi `$ $`=`$ $`(x_x+\mathrm{\Delta })e^\varphi `$
$`\delta _2^Ve^\varphi `$ $`=`$ $`(2xa_3+2g_2+2g_1{\displaystyle _0^x}d_1)e^\varphi .`$ (4.14)
It is understood of course that these fields are primary with respect to the usual space-time Virasoro symmetry. The actions of the variations (4.11) on the generator of this symmetry can be easily calculated as usual by means of Miura transformation (2.10) and the recursive relations implied by (4.3)
$`\delta _{2m}^Vu=2_xb_{2m1}^V,m<0`$
$`\delta _{2m}^Vu=2_x\beta _{2m+1}^V,m0.`$ (4.15)
For example
$$\delta _2^Vu=x(u^{\prime \prime \prime }\frac{3}{2}uu^{})+u^{\prime \prime }2u^2\frac{1}{2}u^{}_0^xu.$$
(4.16)
We now apply our usual procedure in three parts to find transformation equations and symmetry algebra in this case. We will omit the proofs in consideration of the fact that they are very similar to the previous ones.
###### Lemma 4.1
The equations of motion of the resolvents (4.2) under the flows (4.6) have the form :
$$\delta _{2n}^VZ_{2m}^V=[\theta _{2n}^V,Z_{2m}^V](2n2m)Z_{2n+2m}^V,m,nZ.$$
(4.17)
###### Lemma 4.2
The transformations of the connections (4.5) under the flows (4.6) have the form :
$$\delta _{2n}^V\theta _{2m}^V\delta _{2m}^V\theta _{2n}^V=[\theta _{2n}^V,\theta _{2m}^V](2n2m)\theta _{2n+2m}^V,m,nZ.$$
(4.18)
###### Theorem 4.1
The algebra of the vector fields on $``$ (4.6) forms a representation of the centerless Virasoro algebra:
$$[\delta _{2m}^V,\delta _{2n}^V]=(2m2n)\delta _{2m+2n}^V,m,nZ,$$
(4.19)
after the redefinition $`\delta ^V\delta ^V`$.
The Theorem 4.19 gives us a very non-trivial information because of the different character of the asymptotic and regular Virasoro vector fields. Indeed, the asymptotic ones are quasi-local (they can be made local after differentiating a certain number of times), the regular ones instead are essentially non-local being expressed in terms of vertex operators. In addition, it is easy to compute the most simple relations $`[\delta _0,\delta _{2n}]=2n\delta _{2n}`$, $`nZ`$, which means that $`\delta _0`$ counts the dimension or level. We want to stress once more that this Virasoro symmetry is different from the space-time one and is essentially non-local. The additional symmetries coming from the regular dressing are very important for applications. They complete the asymptotic ones forming an entire Virasoro algebra and provide a possibility of a central extension in the (generalized) KdV hierarchy, which is the classical limit of CFT’s . However, this central term may appear only in the algebra of the hamiltonians of the above transformations, as it is for the case of CFT’ s as well.
With the aim of understanding the classical and quantum structure of integrable systems, we present here the complete algebra of symmetries. The Virasoro flows commute neither with the mkdV hierarchy (2.67) nor with the (proper) regular dressing flows (2.24). In fact one can show, following the lines of the three steps procedure, these statements.
###### Lemma 4.3
The equations of motion of the resolvents (4.2) under the mKdV flows (2.67) and of the mKdV resolvent (2.74) under the Virasoro flows (4.6) have the same form of the variation of $``$:
$`\delta _{2k+1}^HZ_{2m}^V`$ $`=`$ $`[\theta _{2k+1}^H,Z_{2m}^V],kN,mZ`$
$`\delta _{2m}^VZ_{2k+1}^H`$ $`=`$ $`[\theta _{2m}^V,Z_{2k+1}^H].`$ (4.20)
###### Lemma 4.4
The mixed transformations of the connections (4.5) under the mKdV flows (2.67) are not of gauge type:
$$\delta _{2k+1}^H\theta _{2m}^V\delta _{2m}^V\theta _{2k+1}=[\theta _{2k+1}^H,\theta _{2m}^V](2k+1)\theta _{2k+2m+1}^H.$$
(4.21)
###### Theorem 4.2
The algebra of the hierarchy flows and of the Virasoro flows is not abelian:
$$[\delta _{2k+1}^H,\delta _{2m}^V]=(2k+1)\delta _{2m+2k+1}^H,$$
(4.22)
where we have put $`\delta _{2k+1}^H=0`$ if $`k<0`$ in the r.h.s..
The content of the previous theorem is that Virasoro symmetry shifts along the KdV hierarchy.
Likewise, it may be proven that
$$[\mathrm{\Delta }_n^X,\delta _{2m}^V]=n\mathrm{\Delta }_{n2m}^X.$$
(4.23)
after putting the proper dressing flows (2.24) with negative $`n2m`$: $`\mathrm{\Delta }_{n2m}^X=0`$ in the r.h.s.. As a very important consequence of this fact, the light-cone Sine-Gordon flow $`_+`$ (2.47),(2.48) commutes with all the positive Virasoro modes, i.e. we have obtained a half Virasoro algebra as exact symmetry of SGM. We note again that this infinitesimal transformation is quasi-local in the boson $`\varphi `$.
One remark is necessary at this stage. It is quite interesting to have a Virasoro algebra not commuting (spectrum generating) with the KdV flows, but one may transform the Virasoro flows into true symmetries commuting with the mKdV hierarchy by adding a term containing all the times $`t_{2k+1}`$
$$\delta _{2m}^V\delta _{2m}^V\underset{k=1}{\overset{\mathrm{}}{}}(2k+1)t_{2k+1}\delta _{2m+2k+1}^H,$$
(4.24)
From the view point of CFT it is very difficult to give a physical meaning to these times, but from the restriction of the action of the positive part of the Virasoro algebra on $`u`$ (4.15) we can check that the previous formula yields the half Virasoro algebra described in , by using the pseudodifferential operator method. Actually, it plays an important role in the study of the matrix models where it leads to the so called Virasoro constraints: $`L_m\tau =0,m>0`$. Here $`\tau `$ is the $`\tau `$-function of the hierarchy and is connected to the partition function of the matrix model. Moreover, it seems that it should play an important role also in the context of the Matrix String Theory , which is now intensively studied. Note also that these Virasoro constraints are the conditions for the highest weight state and, because we also have $`L_0\tau \tau `$ , the $`\tau `$-function is a primary state for the Virasoro algebra. But, we uncovered the negative modes of the Virasoro algebra, which build the highest weight representation over the $`\tau `$-function. We are analyzing this intruiging scenario even in off-critical theories like Sine-Gordon .
Let us also note that actually the symmetry of mKdV is much larger. Indeed, the differential operators $`l_{2m,2n}=\lambda ^{2m+1}_\lambda ^{2n+1}`$ close a (twisted) $`w_{\mathrm{}}`$ which is isomorphic to its dressed version
$$\delta _{2m,2n}A_x=[\theta _{2m,2n}(x;\lambda ),],$$
(4.25)
where we have defined the connections and the resolvents
$`\theta _{2m,2n}=(Z_{2m,2n}^V)_{},Z_{2m,2n}=T_{reg}l_{2m,2n}T_{reg}^1,m<0`$
$`\theta _{2m,2n}=(Z_{2m,2n}^V)_+,Z_{2m,2n}=T_{asy}l_{2m,2n}T_{asy}^1,m0.`$ (4.26)
In particular from the commutations of the diagonal differential operators $`ł_{2n,2n}=\lambda ^{2n+1}_\lambda ^{2n+1}`$
$$[ł_{2n,2n},ł_{2m,2m}]=0$$
(4.27)
we deduce the existence of a quasi-local hierarchy of the diagonal flows
$$[\delta _{2n,2n},\delta _{2m,2m}]=0,n,m>0.$$
(4.28)
As far as we know, this observation is new and we suggest that the diagonal flows are connected to the higher Calogero-Sutherland hamiltonian flows in their collective field theory description. This could give a geometrical explicit explanation of the misterious connection between Calogero-Sutherland systems and KdV hierarchy .
## 5 Generalization: the $`A_2^{(2)}`$-KdV.
Let us show that our approach is easily applicable to other integrable systems. Here we consider the case of the $`A_2^{(2)}`$-KdV equation. The reason is that it can be considered as a different classical limit of the CFT’s . Consider the matrix representation of the $`A_2^{(2)}`$-KdV equation:
$$_t=[,A_t]$$
(5.1)
where
$$=_xA_x,A_x=\varphi ^{}h+(e_0+e_1),$$
(5.2)
and
$$e_0=\left(\begin{array}{ccc}0& 0& \lambda \\ 0& 0& 0\\ 0& 0& 0\end{array}\right),e_1=\left(\begin{array}{ccc}0& 0& 0\\ \lambda & 0& 0\\ 0& \lambda & 0\end{array}\right),h=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right)$$
(5.3)
are the generators of the Borel subalgebra of $`A_2^{(2)}`$ and $`A_t`$ is a certain connection that can be found for example in . Again, the two $`A_2^{(2)}`$-KdV equations are given by Miura transformations; the one of our interest is again $`u(x)=\varphi ^{}(x)^2\varphi ^{\prime \prime }(x)`$. As shown before, central role is played by the transfer matrix, wich is a solution of the associated linear problem $`(_xA_x(x;\lambda ))T(x;\lambda )=0`$. The formal solution is in this case
$$T_{reg}(x,\lambda )=e^{h\varphi (x)}𝒫\mathrm{exp}\left(_0^x𝑑y(e^{2\varphi (y)}e_0+e^{\varphi (y)}e_1)\right).$$
(5.4)
The equation (5.4) defines $`T`$ as an entire function of $`\lambda `$ with an essential singularity at $`\lambda =\mathrm{}`$. The corresponding proper dressing symmetries may be worked out in a way similar to the $`A_1^{(1)}`$ case, but we are mainly interested in the spectrum of local fields. As for the $`A_1^{(1)}`$ case, the asymptotic expansion is easily written by following the general procedure . The result is:
$$T_{asy}(x;\lambda )=\left(\begin{array}{ccc}1& h_1^+& h_2^+\\ 0& 1+h_3^0& h_1^0\\ 0& h_2^{}& 1+h_3^{}\end{array}\right)\mathrm{exp}\left(_0^x\underset{i=0}{\overset{\mathrm{}}{}}f_i\mathrm{\Lambda }^i\right).$$
(5.5)
where $`h_i^\pm ,h_i^0`$ are certain polinomials in $`\varphi ^{}`$, $`f_{6k},f_{6k+2}`$ are the densities of the local conserved charges of the $`A_2^{(2)}`$-KdV and $`\mathrm{\Lambda }=e_0+e_1`$. Complying with our approach let us introduce the asymptotic resolvents
$`Z_1=T_{asy}\mathrm{\Lambda }T_{asy}^1;`$
$`Z_2=T_{asy}\mathrm{\Lambda }^2T_{asy}^1`$ (5.6)
satisfying as before the equations (2.20) $`[_xA_x,Z_i(x;\lambda )]=0`$, $`i=1,2`$ with $`A_x`$ now given by (5.2) and (5.3). These have the form:
$`Z_1=\left(\begin{array}{ccc}\lambda ^1a_1^{(1)}+\lambda ^4a_4^{(1)}+\mathrm{}& \lambda ^2b_2^{(1)}+\lambda ^5b_5^{(1)}+\mathrm{}& 1+\lambda ^6b_6^{(1)}+\mathrm{}\\ 1+\lambda ^3c_3^{(1)}+\mathrm{}& 2\lambda ^4a_4^{(1)}+\mathrm{}& \lambda ^2b_2^{(1)}\lambda ^5b_5^{(1)}+\mathrm{}\\ \lambda ^2c_2^{(1)}+\mathrm{}& 1\lambda ^3c_3^{(1)}+\lambda ^6c_6^{(1)}+\mathrm{}& \lambda ^1a_1^{(1)}+\lambda ^4a_4^{(1)}+\mathrm{}\end{array}\right),`$ (5.10)
and
$`Z_2=\left(\begin{array}{ccc}\lambda ^2a_2^{(2)}+\lambda ^5a_5^{(2)}+\mathrm{}& 1+\lambda ^3b_3^{(2)}+\lambda ^6b_6^{(2)}+\mathrm{}& \lambda ^4b_4^{(2)}+\mathrm{}\\ \lambda ^1c_1^{(2)}+\lambda ^4c_4^{(2)}\mathrm{}& 2\lambda ^2a_2^{(2)}+\mathrm{}& 1\lambda ^3b_3^{(2)}\lambda ^6b_6^{(2)}+\mathrm{}\\ 1+\lambda ^6c_6^{(2)}+\mathrm{}& \lambda ^1c_1^{(2)}+\lambda ^4c_4^{(2)}+\mathrm{}& \lambda ^2a_2^{(2)}\lambda ^5a_5^{(2)}+\mathrm{}\end{array}\right).`$ (5.14)
For example, some expressions for the fiels in the entries of $`Z_i,i=1,2`$, as a function of $`v=\varphi ^{}`$ and its derivative, are:
$`a_1^{(1)}`$ $`=`$ $`v;`$
$`b_2^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}v^2+{\displaystyle \frac{1}{3}}v^{},c_2^{(1)}={\displaystyle \frac{1}{3}}v^2{\displaystyle \frac{2}{3}}v^{};`$
$`c_3^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}v^3+{\displaystyle \frac{1}{3}}vv^{}{\displaystyle \frac{1}{3}}v^{\prime \prime },b_3^{(2)}={\displaystyle \frac{2}{3}}vv^{}+{\displaystyle \frac{1}{3}}v^{\prime \prime };`$
$`b_4^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{9}}v^4+{\displaystyle \frac{1}{9}}v^2{\displaystyle \frac{2}{9}}vv^{\prime \prime }+{\displaystyle \frac{2}{9}}v^{}v^2{\displaystyle \frac{1}{9}}v^{\prime \prime \prime },etc..`$ (5.15)
The equation (5.1) is invariant under a gauge transformation of the form (2.67) . This latter will be a true symmetry provided the variation is proportional to $`h`$ : $`\delta A_x=h\delta \varphi ^{}`$. We construct the appropriate gauge parameters by means of the resolvents (5.6) in a way similar to what we did in the $`A_1^{(1)}`$-mKdV case:
$$\theta _{6k+1}(x;\lambda )=(\lambda ^{6k+1}Z_1(x;\lambda ))_+,\theta _{6k1}(x;\lambda )=(\lambda ^{6k1}Z_2(x;\lambda ))_+$$
(5.16)
which results in the following transformations for the $`A_2^{(2)}`$-mKdV field
$$\delta _{6k+1}\varphi ^{}=_xa_{6k+1}^{(1)},\delta _{6k1}\varphi ^{}=_xa_{6k1}^{(2)}.$$
(5.17)
One can easily recognize in (5.17) the infinite tower of the commuting $`A_2^{(2)}`$-mKdV flows.
Now, in accordance with our geometrical conjecture, we would like to treat the entries of the transfer matrix $`T`$ and of the resolvents $`Z_i,i=1,2`$ as independent fields and to build the spectrum of the local fields of $`A_2^{(2)}`$-KdV by means of them alone. As in the $`A_1^{(1)}`$ case, it turns out that not all of them are independent. If the defining relations of the resolvents are used, it is easy to see, that the entries of the lower triangle of both $`Z_i`$ can be expressed in terms of the rest. Therefore, taking also into account the gauge symmetry of the system, one is led to the following proposal about the construction of the Verma module of the identity:
$$𝒱_\mathrm{𝟎}^{mKdV}=\{l.c.o.\delta _{6k_1+1}\mathrm{}\delta _{6k_M+1}\delta _{6l_11}\mathrm{}\delta _{6l_N1}𝒫(b_i^{(1)},b_j^{(2)},a_k^{(1)},a_l^{(2)})\},$$
(5.18)
where $`l.c.o.`$ means linear combinations of. Again, null-vectors appear in the r.h.s. of the (5.18) due to the constraints:
$$Z_1^2=Z_2,Z_1Z_2=\text{1}$$
(5.19)
and the equations of motion
$$\delta _{6k\pm 1}Z_i=[\theta _{6k\pm 1},Z_i],i=1,2.$$
(5.20)
One can further realize that just as in the $`A_1^{(1)}`$ case, there is a subalgebra consisting of the upper triangular entries of $`Z_i,i=1,2`$, closed under the action of the gauge transformations $`\delta _{6k\pm 1}`$. The constraints (5.19) and (5.20) are consistent with such reduction giving a closed subalgebra of null vectors. The first non-trivial examples are :
$`level3`$ $`:`$ $`b_3^{(2)}_xb_2^{(1)}=0;`$
$`level4`$ $`:`$ $`b_4^{(2)}(b_2^{(1)})^2+{\displaystyle \frac{2}{3}}_xb_3^{(2)}=0;`$
$`level6`$ $`:`$ $`b_6^{(1)}2b_6^{(2)}+2b_2^{(1)}b_4^{(2)}+(b_3^{(2)})^2=0;`$ (5.21)
$`2`$ $`b_6^{(1)}b_6^{(2)}+2_xb_5^{(1)}+{\displaystyle \frac{1}{2}}_x^2b_4^{(2)}+b_2^{(1)}b_4^{(2)}(b_2^{(1)})^3+(b_3^{(2)})^2=0.`$
Therefore, in order to obtain the true spectrum of the family of the identity (i.e. the $`A_2^{(2)}`$-KdV spectrum), one has to factor out, from the linearly generated Verma module
$$𝒱_\mathrm{𝟎}^{KdV}=\{l.c.o.\delta _{6k_1+1}\mathrm{}\delta _{6k_M+1}\delta _{6l_11}\mathrm{}\delta _{6l_N1}𝒫(b_i^{(1)},b_j^{(2)})\},$$
(5.22)
the Verma module of null-vectors $`𝒩_\mathrm{𝟎}^{KdV}`$, i.e.
$$[\mathrm{𝟎}]=𝒱_0^{KdV}/𝒩_\mathrm{𝟎}^{KdV}.$$
(5.23)
Let us turn to the classical limit of the primary fields $`e^{m\varphi },m=0,1,2,3,\mathrm{}`$ . By virtue of the above reasoning we conjecture for their Verma modules the expression
$$𝒱_\mathrm{𝟎}^{mKdV}=\{l.c.o.\delta _{6k_1+1}\mathrm{}\delta _{6k_M+1}\delta _{6l_11}\mathrm{}\delta _{6l_N1}[𝒫(b_i^{(1)},b_j^{(2)})e^{m\varphi }]\}.$$
(5.24)
Again, we have to add to the null-vectors coming from (5.19) and (5.20) the new ones coming from (repeated) application of the equations of motion of the power $`T^m(x;\lambda )`$
$$\delta _{6k\pm 1}T^m=\underset{j=1}{\overset{m}{}}T^j\theta _{6k\pm 1}T^{mj}$$
(5.25)
obtaining the whole set of null-vectors $`𝒩_𝐦^{KdV}`$. The first non-trivial examples of these additional null-vectors are given by:
$`level`$ $`2`$ $`:(_x^2+3b_2^{(1)})e^\varphi =0`$
$`level`$ $`3`$ $`:(_x^36b_3^{(2)}+12_xb_2^{(1)})e^{2\varphi }=0`$
$`level`$ $`4`$ $`:(_x^4+{\displaystyle \frac{135}{2}}(b_2^{(1)})^2+30_x^2b_2^{(1)}{\displaystyle \frac{27}{2}}b_4^{(2)}30_xb_3^{(2)})e^{3\varphi }=0`$ (5.26)
where the operator $`_x`$ acts on all the fields to its right. As a result, the (conformal) family of the primary field $`e^{m\varphi },m=0,1,2,3,\mathrm{}`$ is conjectured to be in this case
$$[𝐦]=𝒱_m^{KdV}/𝒩_𝐦^{KdV}.$$
(5.27)
## 6 Conclusions and perspectives: off-critical theories and quantisation.
We have derived different types of symmetries of classical integrable systems within a general framework. The unifying geometrical idea bases on dressing simple objects by means of the transfer matrix $`T`$ with the aim to get the resolvent $`Z`$. In particular, the regular transfer matrix gives rise to a Poisson–Lie symmetry of non-commuting flows. The construction of the spectrum employing this symmetry is an interesting open problem, clearly connected with the spinon basis defined in , since the current density $`e^\varphi `$ is exactly the spinon of after quantisation. Moreover, it is of great interest the new way to look at the Sine-Gordon light-cone evolution as generated by the sum of the first two regular vector fields. It is also important to note that a similar half Virasoro symmetry can be obtained just underchanging the rôle of $`x_{}`$ and $`x_+`$. We leave for a work in progress the very important question on the whole algebra obtained from the union of both half Virasoro algebras .
On the contrary, the asymptotic formula for the transfer matrix provides the integrable hierarchy of (generalized) KdV flows and, in the case of dressing of non-cartan generators, two new infinite series of flows closing an algebraic structure with the integrable ones. A new costruction of the spectrum of classical Virasoro algebra has been given in terms of these ingredients.
Even in view of quantisation, it is worth extending our constructions out of criticality in another way. We may define the anti-chiral transfer matrix
$$\overline{T}(\overline{x},\lambda )=𝒫\mathrm{exp}\left(\lambda _x^0𝑑y(e^{2\overline{\varphi }(y)}E+e^{2\overline{\varphi }(y)}F)\right)e^{H\overline{\varphi }(\overline{x})}$$
(6.1)
which solves the linear problem
$$_{\overline{x}}\overline{T}(\overline{x};\lambda )=\overline{T}(\overline{x};\lambda )\overline{A}_{\overline{x}}(\overline{x};\lambda )$$
(6.2)
where $`\overline{A}_x(x;\lambda )`$ is obtained from $`A_x`$ by substitution $`\varphi \overline{\varphi }`$. In accordance with , we suggest for the off-critical transfer matrix
$$𝐓(x,\overline{x};\lambda ;\mu )=\overline{T}(\overline{x};\mu /\lambda )T(x;\lambda ),$$
(6.3)
and correspondingly for $`𝐙(x,\overline{x};\mu ;\lambda )`$.
Following the basic work , one quantizes the corresponding mKdV system by replacing the Kac-Moody algebra with the corresponding quantum group and the mKdV field $`\varphi `$ with the Feigin–Fuchs–Dotsenko–Fateev free field . As explained in the importance of considering also the $`A_2^{(2)}`$-mKdV hierarchies is due to the fact that the quantisation of this second semiclassical system exhausts the integrability directions of theories of type (1.1) starting from Minimal Models of . For different kinds of CFT it is sufficient to considere hierarchies attached to different Kac-Moody algebras (in the Drinfeld-Sokolov scheme ).
In conclusion, we have presented a generalization of the dressing symmetry construction leading to a non-local Virasoro symmetry of the mKdV hierarchy and SGM. We stress that it has nothing to do with the space-time Virasoro, generated at the classical level by the moments of the classical stress tensor $`x^nu(x)𝑑x`$. It is obtained instead by dressing the differential operator $`\lambda ^{n+1}_\lambda `$. In view of the relation between the spectral parameter and the on-shell rapidity $`\lambda =e^\theta `$, it is generated probably by diffeomorphisms in the momentum space and in this sense is dual to the space-time Virasoro symmetry. Although we presented a construction only in the case of mKdV, it can be easily extended for the generalized KdV theories as well. Of particular interest is the $`A_2^{(2)}`$ hierarchy, connected with the $`\varphi _{1,2}`$ perturbation of CFT models .
Furthermore, such a symmetry appears also in the study of Calogero-Sutherland model whose connection with the matrix models and CFT is well known. Moreover, it is known that in the q-deformed case it becomes a deformed Virasoro algebra . It is natural to suppose that in the same way our construction is deformed off-critically.
We suggest that this Virasoro symmetry could be of great importance for the study of 2D-IQFT. First of all it should provide a new set of conserved charges closing a non-abelian algebra, thus carrying nessesarily more information about the theory. Furthermore one might quantise these charges at conformal and off-critical level. We have reasons to belive that the perturbed version (i.e. SGM) should be closely related to the aforementioned DVA.
Recently, Babelon, Bernard and Smirnov constructed certain null-vectors off-criticality in the context of the form-factor approach. They showed that there is a deep connection, at the classical level, between their construction and the finite zone solutions of KdV and the Witham theory of averaged KdV. On the other hand the Virasoro algebra presented above has a natural action on the finite zone solutions, changing the complex structures of the corresponding hyperelliptic Riemann surfaces, and on the basic objects of the Witham hierarchy . This suggests we may found the quantum action of our symmetries (in particular the Virasoro one) in SGM using the form factor formalism developed in .
Acknowledgments \- We are indebted to E. Corrigan, P. Dorey, G. Mussardo, I. Sachs and F. Smirnov for discussions and interest in this work. D.F. thanks the I.N.F.N.–S.I.S.S.A., the Mathematical Sciences Departement in Durham and the EC Commission (TMR Contract ERBFMRXCT960012) for financial support. M.S. acknowledges S.I.S.S.A. for the warm hospitality over part of this work. |
warning/0001/gr-qc0001042.html | ar5iv | text | # Can a Particle’s Velocity Exceed the Speed of Light in Empty Space?
## I Introduction
As well known, the special relativity theory (SR) is the theory of inertial systems and for such systems the answer to the question posed in the title of this paper is negative. But in the nature ideal inertial systems do not exist. It allows to raise the following problem: is it possible to develop a theory of systems close to inertial (”almost” inertial systems) that would include, as a special case, special relativity but, at the same time, would allow for a motion of particles with any velocity? Obviously, in order to invent such a theory it is necessary to refuse from the rigorous validity of any of the SR postulates: the homogenity of space and time, the invariance of speed of light and the Galilee invariance principle.
The present paper suggests an example of such a theory, based on the concepts of time and space with fractional dimensions (FD) developed in the theory of multifractal time and space . We begin with the first of the mentioned SR principles. In an inhomogeneous space and time, if the inhomogenities are small enough, any motion will be close to that in homogeneous space (”almost” inertial), but the velocity of light can alter slightly, being thus ”almost” constant. Mere assumption that the values of fractional dimensions of space and time are close to integer leads then to conclusion that the motions of particles with any velocities become possible. Other main assumptions made are the light velocity invariance and invariance with respect to modified Lorentz transformations. The results our theory gives for the velocities less than the velocity of light $`c`$ almost coincide with SR, but it does not contain singularities at $`v=c`$. Stress, that this theory is not a generalization of SR theory, because any such generalization in the domain of SR validity (inertial systems) is absurd. Our theory describes relative movements only in the ”almost” inertial systems, and thus does not contradict to SR.
## II Multifractal time
Following , we will consider both time and space as the only material fields existing in the world and generating all other physical fields. Assume that every of them consists of a continuous, but not differentiable bounded set of small elements (elementary intervals, further treated as ”points”). Consider the set of small time elements $`S_t`$. Let time be defined on multifractal subsets of such elements, defined on certain measure carrier $`^n`$. Each element of these subsets (or ”points”) is characterized by the local fractional (fractal) dimension (FD) $`d_t(𝐫(t),t)`$ and for different elements FD are different. In this case the classical mathematical calculus or fractional (say, Riemann - Liouville) calculus can not be applied to describe a small changes of a continuous function of physical values $`f(t)`$, defined on time subsets $`S_t`$, because the fractional exponent depends on the coordinates and time. Therefore, we have to introduce integral functionals (both left-sided and right-sided) which are suitable to describe the dynamics of functions defined on multifractal sets (see ). Actually, this functionals are simple and natural generalization the Riemann-Liouville fractional derivatives and integrals:
$$D_{+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(1)
$$D_{,t}^df(t)=(1)^n\left(\frac{d}{dt}\right)^n_t^b\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}$$
(2)
where $`\mathrm{\Gamma }(x)`$ is Euler’s gamma function, and $`a`$ and $`b`$ are some constants from $`[0,\mathrm{})`$. In these definitions, as usually, $`n=\{d\}+1`$ , where $`\{d\}`$ is the integer part of $`d`$ if $`d0`$ (i.e. $`n1d<n`$) and $`n=0`$ for $`d<0`$. If $`d=const`$, the generalized fractional derivatives (GFD) (1)-(2) coincide with the Riemann - Liouville fractional derivatives ($`d0`$) or fractional integrals ($`d<0`$). When $`d=n+\epsilon (t),\epsilon (t)0`$, GFD can be represented by means of integer derivatives and integrals. For $`n=1`$, that is, $`d=1+\epsilon `$, $`\left|\epsilon \right|<<1`$ it is possible to obtain:
$$D_{+,t}^{1+\epsilon }f(t)\frac{}{t}f(t)+a\frac{}{t}\left[\epsilon (r(t),t)f(t)\right]$$
(3)
where $`a`$ is constant and defined by the choice of the rules of regularization of integrals (1)-(2) (for more detailed see ). The selection of the rule of regularization that gives a real additives for usual derivative in (3) yeilds $`a=0.5`$ for $`d<1`$ and $`a=1.077`$ for $`d>1`$ . The functions under integral sign in (1)-(2) we consider as the generalized functions defined on the set of the finite functions . The notions of GFD, similar to (1)-(2), can also be defined for the space variables $`𝐫`$.
The definitions of GFD (1)-(2) are formal until the connections between fractal dimensions of time $`d_t(𝐫(t),t)`$ and certain characteristics of physical fields (say, potentials $`\mathrm{\Phi }_i(𝐫(t),t),i=1,2,..)`$ or densities of Lagrangians $`L_i`$) are determined. Following , we define this connection by the relation
$$d_t(𝐫(t),t)=1+\underset{i}{}\beta _iL_i(\mathrm{\Phi }_i(𝐫(t),t))$$
(4)
where $`L_i`$ are densities of energy of physical fields, $`\beta _i`$ are dimensional constants with physical dimension of $`[L_i]^1`$ (it is worth to choose $`\beta _i^{}`$ in the form $`\beta _i^{}=a^1\beta _i`$ for the sake of independence from regularization constant). The definition of time as the system of subsets and definition (4) put the value of fractional (fractal) dimensionality $`d_t(r(t),t)`$ into accordance with every time instant $`t`$. The latter depends both on time $`t`$ and coordinates $`𝐫`$. If $`d_t=1`$ (absence of physical fields) the set of time has topological dimensionality equal to unity. The multifractal model of time allows, as will be shown below, to consider the divergence of energy of masses moving with speed of light in the SR theory, as the result of the requirement of rigorous validity, rather than approximate fulfillment, of the laws pointed out in the beginning of this paper in the presence of physical fields.
## III The principle of the velocity of light invariance
Because of the inhomogeneity of time in our multifractal model, the speed of light, just as in the general relativity theory, depends on potentials of physical fields that define the fractal dimensionality of time $`d_t(𝐫(t),t)`$ (see (4)). If fractal dimensionality $`d_t(𝐫(t),t)`$ is close enough to unity ($`d_t(r(t),t)=1+\epsilon ,\left|\epsilon \right|<<1`$), the difference of the speed of light in moving (with velocity $`v`$ along the $`x`$ axis) and fixed frame of reference will be small. In the systems that move with respect to each other with almost constant velocity (stationary velocities do not exist in the mathematical theory based on definitions of GFD (1) - (2)) the speed of light can not be taken as a fundamental constant. In the multifractal time theory the principle of the speed of light invariance can be considered only as approximate. But if $`\epsilon `$ is small, it allows to consider a nonlinear coordinates transformations from the fixed frame to the moving (replacing the transformations of Galilee in inhomogeneous time and space), as close to linear (weakly nonlinear) transformations and, thus, makes it possible to preserve the conservation laws, and all the invariants of the Minkowski space, as the approximate laws. Then the way of reasoning and argumentation accepted in SR theory (see for example, ) can also remain valid. Designating the coordinates in the moving and fixed frames of reference through $`x^{}`$ and $`x`$, accordingly, we write down
$`x^{}`$ $`=`$ $`\alpha (t,x)[xv(x,t)t(x(t),t]`$ (5)
$`x`$ $`=`$ $`\alpha ^{}(t,x)[x^{}+v^{}(x^{}(t^{}),t^{}),t^{}(x^{}(t^{}),t^{})`$ (6)
In (5) $`\alpha \alpha ^{}`$ and the velocities $`v^{}`$ and $`v`$ (as well as $`t`$ and $`t^{}`$) are not equal (it follows from the inhomogenity of multifractal time). Place clocks in origins of both the frames of reference and let the light signal be emitted in the moment, when the origins of the fixed and moving frames coincide in space and time at the instant $`t^{}=t=0`$ and in points $`x^{}=x=0`$. The propagation of light in moving and fixed frames of reference is then determined by equations
$$x^{}=c^{}t^{}x=ct$$
(7)
These characterize the propagation of light in both of the frames of reference at every moment. Due to the time inhomogenity $`c^{}c`$, but since $`|\epsilon <<1|`$ the difference between velocities of light in the two frames of reference will be small. For this case we can neglect the distinction between $`\alpha ^{}`$ and $`\alpha `$ and, for different frames of reference write the expressions for velocities of light, using (3) to define velocity (denote $`f(t)=x,dx/dt=c_0`$). Thus we obtain
$$c=D_{+,t}^{1+\epsilon }x=c_0(1\epsilon )\frac{d\epsilon }{dt}x$$
(8)
$$c^{}=D_{+,t}^{1+\epsilon ^{}}x^{}=c_0(1\epsilon )+\frac{d\epsilon }{dt}x^{}$$
(9)
$$c_1=c_0(1\epsilon )\frac{d\epsilon }{dt}x^{}$$
(10)
$$c_1^{}=c_0(1\epsilon )+\frac{d\epsilon }{dt}x$$
(11)
The equalities (10) and (11) appear in our model of multifractal time as the result of the fact, that in this model all the frames of reference are absolute frames of reference (because of material character of the time field) and the speed of light depends on the state of frames: if the frame of reference is a moving or a fixed one, if the object under consideration in this frame moves or not. This dependence disappears only when $`\epsilon =0`$. Before substitution the relations (5) in the equalities (8) - (11) (with $`\alpha ^{}\alpha `$) it is necessary to find out how $`d\epsilon /dt`$ depends on $`\alpha `$. Using for this purpose Eq.(4) we obtain:
$$\frac{d\epsilon }{dt}=\frac{d\epsilon }{d𝐫}𝐯\underset{i}{}\beta _i(𝐅_i𝐯+\frac{L_i}{t})$$
(12)
where $`𝐅_i=dL_i/d𝐫`$. Since the forces for moving frames of reference are proportional to $`\alpha `$ we get (for the case when there is no explicit dependence of $`L_i`$ on time)
$$\frac{d\epsilon }{dt}\underset{i}{}\beta _i𝐅_{0i}𝐯\alpha $$
(13)
where $`F_{0i}`$ are the corresponding forces at zero velocity. Multiplying (8) - (11) on the corresponding times $`t,t^{},t_1,t_1^{}`$ yields the following expressions
$$c^{}t^{}=c_0t\left[1+\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(14)
$$ct=c_0t^{}\left[1+\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1+\frac{v}{c})\right]$$
(15)
$$c_1^{}t_1^{}=c_0t_1\left[1\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1\frac{v}{c})\right]$$
(16)
$$c_1t_1=c_0t_1^{}\left[1\frac{v_i\beta _iF_{0i}}{c_0}\alpha ^2ct(1+\frac{v}{c})\right]$$
(17)
Since in our model the motion and frames of reference are absolute, the times $`t_1`$ and $`t_1^{}`$ correspond to the cases, when the moving and fixed frames of reference exchange their roles - the moving one becomes fixed and vice versa. These times coincide only when $`\epsilon =0`$. The times in square brackets, as well as the velocities, are taken to equal, because the terms containing them are small as compared to unity. The principle of invariance of the velocity of light for transition between the moving and fixed frames of reference in multifractal time model is approximate (though quite natural, because the frames of reference are absolute frames of reference). Taking into account (5), the relations (14) - (17) take the form
$$c^{}t^{}=c\alpha t(1\frac{v}{c}),c_1^{}t_1^{}=c\alpha t_1(1\frac{v}{c})$$
(18)
$$ct=c\alpha t^{}(1+\frac{v}{c}),c_1t_1=c\alpha t^{}(1+\frac{v}{c})$$
(19)
Once again we note, that the four equations for $`c_1^{}t_1^{}`$ and $`c_1t_1`$, instead of the two equations in special relativity, appear as the consequence of the absolute character of the motion and frames of reference in the model of multifractal time. In the right-hand side of (18) - (19) the dependence of velocity of light on fractal dimensions of time is not taken into account (just as in the equations (14) - (17)). Actually, this dependence leads to pretty unwieldy expressions. But if we retain only the terms that depend on $`\beta =\sqrt{|1v^2/c^2|}`$ or $`a_0`$ and neglect unessential terms containing the products $`\beta \alpha _0`$, utilizing (14) - (17) after the multiplication of the four equalities (18) - (19), we receive the following equation for $`\alpha `$ (it satisfies to all four equations):
$$4a_0^4\beta ^4\alpha ^84a_0^2\alpha ^4+1=\beta ^4\alpha ^4+4a_0^4\beta ^4\alpha ^8$$
(20)
where
$$\beta =\sqrt{\left|1\frac{v^2}{c^2}\right|}$$
(21)
$$a_0=\underset{i}{}\beta _iF_{0i}\frac{v}{c}ct$$
(22)
From (20) follows
$$\alpha _1\beta ^^1=\frac{1}{\sqrt[4]{\beta ^4+4a_0^2}}$$
(23)
The solutions $`\alpha _{2,3,4}`$ are given by $`\alpha _2=\alpha _1,\alpha _{3,4}=\pm i\alpha `$. Applicability of above obtained results is restricted by requirement $`|\epsilon |1`$
## IV Lorentz transformations and transformations of length <br>and time in multifractal time model
The Lorentz transformations, as well as transformations of coordinate frames of reference, in the multifractal model of time are nonlinear due to the dependence of the fractional dimensions of time $`d_t(𝐫,t)`$ on coordinates and time. Since the nonlinear corrections to Lorentz transformation rules are very small for $`\epsilon 1`$, we shall take into account only the corrections that eliminate the singularity at the velocity $`v=c`$. It yields in the replacement of the factor $`\beta ^1`$ in Lorentz transformations by the modified factor $`\alpha =1/\beta ^{}`$ given by (23). The Lorentz transformation rules (for the motion along the $`x`$ axis) take the form
$$x^{}=\frac{1}{\beta ^{}}(xvt),t^{}=\frac{1}{\beta ^{}}(tx\frac{v}{c^2})$$
(24)
In the equations (23) and (24) the velocities $`v`$ and $`c`$ weakly depend on $`x`$ and $`t`$ and their contribution to the singular terms is small. Hence, we can neglect this dependence. The transformations from fixed system to moving system are almost orthogonal (for $`\epsilon 1`$ ), and the squares of almost four-dimensional vectors of Minkowski space vary under the coordinates transformations very slightly (i.e. they are almost invariant). Then it is possible to neglect the correction terms of order about $`O(\epsilon ,\dot{\epsilon })`$, which, for not equal to infinity variables, are very small too. From (23) - (24) the possibility of arbitrary velocity motion of bodies with nonzero rest mass follows. With the corrections of order $`O(\epsilon ,\dot{\epsilon })`$ in nonsingular terms being neglected, the momentum and energy of a body with a nonzero rest mass in the frame of reference moving along the $`x`$ axis ($`E_0=m_0c^2)`$ equal to
$$p=\frac{1}{\beta ^{}}m_0v=\frac{m_0v}{\sqrt[4]{\beta ^4+4a_0^2}},E=E_0\sqrt{\frac{v^2c^2}{\sqrt{\beta ^4+4a_0^2}}+1}$$
(25)
The energy of such a body reaches its maximal value at $`v=c`$ and is equal then $`E_{v=c}E_0/\sqrt{2\alpha _0}`$. When $`v\mathrm{}`$ the energy is finite an tends to $`E_0\sqrt{2}`$. For $`vc`$ the total energy of a body is represented by the expression
$$E\frac{E_0}{\sqrt[4]{\beta ^4+4a_0^2}}=mc^2,m=\frac{m_0}{\beta ^{}}$$
(26)
For $`vc`$, total energy, defined by (25), is given by
$`E`$ $`=`$ $`mc^2,\beta ^2={\displaystyle \frac{v^2}{c^2}}1`$ (27)
$`m`$ $`=`$ $`{\displaystyle \frac{m_0}{\beta ^{}}}\sqrt{2v^2/c^2+4a_0^2}`$ (28)
If we are to take into account only the gravitational field of Earth (here, as in , gravitational field is a real field) and neglect the influence of all the other fields), the parameter $`a_0(t)`$ can be estimated to be $`a_0=r_0r^3x_Ect`$, where $`r_0`$ is the gravitational radius of Earth, $`r`$ is the distance from the Earth’s surface to its center ($`\epsilon =0.5\beta _g\mathrm{\Phi }_g,\beta _g=2c^2,x_Er,v=c`$). For energy maximum we get $`E_{max}E_010^3t^{0.5}sec^{0.5}`$.
Shortening of lengths and time intervals in moving frames of reference in the model of multifractal time also have several peculiarities. Let $`l`$ and $`t`$ be the length and time interval in a fixed frame of reference. In a moving frame
$$l^{}=\beta ^{}l,t^{}=\beta ^{}t$$
(29)
Thus, there exist the maximal shortening of length when the body’s velocity equals to the speed of light. With the further increasing of velocity (if it is possible to fulfill some requirements for a motion in this region with constant velocity without radiating), the length of a body begins to grow and at infinitely large velocity is also infinite. The slowing-down of time, from the point of view of the observer in the fixed frame (maximal shortening equals to $`t^{}=t\sqrt{2a_0}`$) is replaced, with the further increase of velocity beyond the speed of light, by acceleration of time passing ($`t0`$ when $`v\mathrm{}`$).
The rule for velocities transformation retains its form, but $`\beta `$ is replaced by $`\beta ^{}`$
$$u_x=\frac{u_x^{}+v}{1+\frac{u_x^{}v}{c^2}},u_y=\frac{u_y^{}\beta ^{}}{1+\frac{u_y^{}v}{c^2}},u_z=\frac{u_z^{}\beta ^{}}{1+\frac{u_z^{}v}{c^2}}$$
(30)
Since there is no law that prohibits velocities greater than that of light, the velocities in (30) can also exceed the speed of light. The electrodynamics of moving media in the model of multifractal time can be obtained, in most cases, by the substitution $`\beta \beta ^{}`$.
## V Conclusions
To conclude, the theory of relative motions in almost inertial systems based on the multifractal time theory is invented. This theory describes open systems (for statistical theory of open systems see in ) and in this theory motion with any velocity is possible. The theory coincides with special relativity after transition to inertial systems (if we neglect the fractional dimensions of time) or almost coincides (the differences are negligible) for velocities $`v<c`$. Movement of bodies with velocities that exceed the speed of light is accompanied by a number of physical effect’s which can be found experimentally (these effects will be considered in the separate paper in more detail). |
warning/0001/astro-ph0001423.html | ar5iv | text | # BeppoSAX and Chandra Observations of SAX J0103.2–7209 = 2E 0101.5–7225: a new Persistent 345 s X–ray Pulsar in the SMC
## 1. Introduction
During 1997–1998 the number of X–ray pulsars found in the Small Magellanic Cloud (SMC) rapidly increased from three (SMC X–1, RX J0053.8–7226, and 2E 0050.1–7247; Lucke et al 1976; Hughes 1994; Israel et al. 1997) up 14 (for a review see Yokogawa et al. 1998) thanks to sensitive observations of the large area detectors on board the RXTE and ASCA satellites. The majority were found to be associated with massive Be spectral type stars showing intense H$`\alpha `$ emission lines. Only SMC X–1, which is associated with a super–giant B0 spectral type star in a 3.9 d orbital period binary system, is a persistent (although moderately variable) X–ray pulsar. For all the remaining X–ray pulsars pronounced variability (a factor $``$ 50) or, more often, transient behaviour has been definitively proven.
The source 2E 0101.5–7225 was detected at a nearly constant flux level in all the Einstein, ROSAT and ASCA pointings which surveyed the relevant region of the SMC (see Hughes & Smith 1994), but pulsations were not found due to poor statistics. Based on the accurate position obtained with the ROSAT HRI, these authors found that 2E 0101.5–7225 is very likely associated with a Be spectral–type star (R.A. = 01<sup>h</sup>03<sup>m</sup>13$`\stackrel{\mathrm{s}}{\mathrm{.}}`$86, Dec. = –72°09′14$`\stackrel{}{\mathrm{.}}`$1; equinox J2000). 2E 0101.5–7225 is located near the optical limb of the supernova remnant SNR 0101–72.4. Hughes & Smith (1994) present several arguments that make the Be/X–ray binary – SNR association unlikely.
In this Letter we report the discovery of 345 s pulsations from the source SAX J0103.2–7209 during a BeppoSAX observation of the SMC. The comparison with the data of past X–ray missions allows us to conclude that SAX J0103.2–7209 and 2E 0101.5–7225 are the same object, a persistent source with moderate variability (within a factor of 5–10). We also report the results of the timing analysis of a recent public Chandra observation and discuss optical observations carried out at ESO.
## 2. Observations and Data Analysis
### 2.1. BeppoSAX observation
The SMC field including the position of the 2E 0101.5–7225 was observed by the Narrow Field Instruments (NFIs) on board the BeppoSAX satellite (Boella et al. 1997a) on 1998 July 26–27 (effective exposure time of 40320 s). We used data from the Medium Energy (MECS; Boella et al. 1997b) and Low Energy (LECS; Parmar et al. 1997) instruments. A bright X–ray source ($``$ 3.7 $`\times `$ 10<sup>-2</sup> ct s<sup>-1</sup>, 1–10 keV) was detected on–axis in the MECS, at R.A. = 01<sup>h</sup>03<sup>m</sup>13<sup>s</sup>, Dec. = –72°09′16″ (J2000; 90% confidence uncertainty radius of 30″ ). The MECS event list and spectrum were extracted from a circular region of 4′ radius (corresponding to an encircled energy of $``$ 90%) around the X–ray position.
A 4′ extraction radius ($``$ 85% encircled energy) was also used for the LECS in order to minimize the contamination from the soft X–ray emission of the bright nearby source 2E 0102.3–7217. The local background was measured in a region of the MECS and LECS images far from any detected field sources.
The arrival times of the $``$ 1500 photons were corrected to the barycenter of the solar system and background subtracted light curves were accumulated in 0.5 s bins. A single power spectrum was calculated over the entire time span covered by the observation in order to maximize the sensitivity to coherent pulsations. Significant power spectrum peaks were searched for using the algorithm described in Israel & Stella (1996). A highly significant peak ($``$ 9.4 $`\sigma `$ based on the fundamental only; see upper left panel of Fig. 1) was found at a frequency of 0.002889 Hz, corresponding to a period of 345.2 s. An accurate determination of the period was obtained by fitting the phases of the modulation over 4 different intervals of $``$ 16000 s each. The scatter of the phase residuals was consistent with a strictly periodic modulation at the best period of 345.2 $`\pm `$ 0.3 s (90% confidence). A comparison of the 1–4 keV (panel S in Fig. 1) and 4–10 keV (panel H) folded light curves provides marginal evidence for an energy dependent pulse profile: nearly sinusoidal at low energies, while double–horned above 4 keV. However the pulsed fraction (semi–amplitude of modulation divided by the mean source count rate) of $``$ 45% is constant over the two energy intervals considered.
A spectral analysis was performed in the 0.7–6.5 and 1.6–10 keV energy ranges for the LECS and MECS, respectively. The spectra were rebinned to have at least 20 counts in each energy channel, such that $`\chi ^2`$ fitting techniques could be used. No single component model was found to fit the data well (a power–law gave a $`\chi ^2/dof`$ = 41/15; where $`dof`$ is the degree of freedom). Among double component models, an absorbed power–law plus a black body gave the best fit ($`\chi ^2/dof`$ = 15/15) for a photon index of 1.0 $`\pm _{0.1}^{0.2}`$, N<sub>H</sub>$``$ 3.8$`\pm _{3.8}^{7.5}`$$`\times `$ 10<sup>21</sup> cm<sup>-2</sup>, and a black body temperature of 0.11 $`\pm `$ 0.03 keV (uncertainties refer to 1$`\sigma `$). The observed flux in the 2–10 keV energy band was 2.7 $`\times `$ 10<sup>-12</sup> erg s<sup>-1</sup> cm<sup>-2</sup> (the soft component accounts for $``$ 15% of the total) corresponding to an unabsorbed 2–10 keV X–ray luminosity of 1.2 $`\times `$ 10<sup>36</sup> erg s<sup>-1</sup> assuming a distance of 62 kpc (Laney & Stobie 1994).
### 2.2. Archival X–ray Observations
After the discovery of 345 s pulsations in 2E 0101.5–7225 (Israel et al. 1998), Yokogawa & Koyama (1998) found a signal at a period of 348.9 s during a 1996 May ASCA observation (absorbed luminosity of 5 $`\times `$ 10<sup>35</sup> erg s<sup>-1</sup>; 2–10 keV) implying a period derivative (with respect to BeppoSAX) of –1.7 s yr<sup>-1</sup>. The source was also detected by ASCA on 1993 May 12 and 1997 November 14 at an absorbed luminosity level of 6 and 5 $`\times `$ 10<sup>35</sup> erg s<sup>-1</sup>, respectively. A re–analysis of the latter ASCA datasets allowed us to infer an upper limit (3$`\sigma `$ confidence) in the 60–80% range on the pulsed fraction for a period of $``$ 345 s.
To better constrain the absorption we used archival data from the ROSAT satellite (PSPC; 0.1–2 keV band) which observed the field of 2E 0101.5–7225 between 1991 October 8 and 1992 April 28 (sequence 600195; effective exposure time 26630 s). Assuming that the source showed the same flux and spectrum during the two observations, we fitted the ROSAT and BeppoSAX data together. We found that the two component model derived above again gave the best fit; a $`\chi ^2/dof`$ = 23/25 was obtained for a photon index of 1.1 $`\pm `$ 0.1 with N<sub>H</sub> of 3.9$`\pm _{2.4}^{4.0}`$$`\times `$ 10<sup>21</sup> cm<sup>-2</sup> and a black body temperature of 0.12$`\pm _{0.03}^{0.04}`$ keV corresponding to an equivalent black body radius R<sub>bb</sub> in the 2–60 km range (see Fig. 2; 1$`\sigma `$ uncertainties). The corresponding power spectrum shows no significant (3$`\sigma `$ confidence) peak in the 339–359 s period interval with a $``$ 60% upper limit on the pulsed fraction.
### 2.3. Chandra observation
The field including 2E 0101.5–7225 was observed by the NASA Advanced X–ray Astrophysics Facility satellite Chandra on 1999 August 23 with the Imaging Spectrometer (ACIS; Garmire et al. 1992, Bautz et al. 1998 and references therein) in the high–resolution imaging mode for an effective exposure time of 19551 s. The source was detected at an off–axis angle of about 9′ in the S4 (front illuminated) CCD of the ACIS–S detector with a count rate of $``$ 9 $`\times `$ 10<sup>-2</sup> ct s<sup>-1</sup> (in the 0.5–10 keV energy range; see Fig. 3 left panel). The source shows two emission peaks separated by $``$ 5″. However the size of the emission region is comparable with the 90% encircled energy region of a source at 9′ off–axis angle.
We extracted a 3.24 s binned light curve of the source from a region within 8″ from the center of the emission: R.A. = 01<sup>h</sup>03<sup>m</sup>14$`\stackrel{\mathrm{s}}{\mathrm{.}}`$06, Dec. = –72°09′15$`\stackrel{}{\mathrm{.}}`$25 (equinox J2000). The statistical uncertainty radius is only $``$ 0.5″ but at this stage the uncertainty in the absolute positioning at such an off–axis angle might be as large as 90″. However the detection of the pulsations clearly associates the Chandra source with that of BeppoSAX. We also note that the Chandra coordinates are consistent with the ROSAT HRI uncertainty circle and differ by less than 2″ from those of the proposed optical counterpart (see below) making the association very likely. According to the Chandra source naming convention we designated it as CXO J010314.1–720915.
A single power spectrum was calculated over the entire observation. The search was performed over a period interval around that detected by BeppoSAX and assuming a maximum $`|`$$`|`$ of $``$3 s yr<sup>-1</sup> which translates into a search over only two Fourier frequencies (see Fig. 1 lower panels). A peak was detected at a significance level of 7.5$`\sigma `$. A refined period was determined by means of the phase fitting technique; this gave a value of 343.5 $`\pm `$ 0.5 s (see Table 1). We note that the Chandra data were not corrected to the barycenter of the solar system. However, for a relatively long period pulsar in the direction of the SMC, the effect of the spacecraft and of the Earth motion, would cause a maximum correction of a factor of 10 smaller than the statistical uncertainty of the period given above. The pulse profile is sinusoidal (over the whole Chandra energy band) and the pulsed fraction is $``$ 45% (see Fig. 1). This result implies that the pulsar is continuing to spin–up at a constant rate of –1.7 s yr<sup>-1</sup> since 1996 (see Fig. 1; right panel).
In order to further address the issue of the double peaked Chandra image (see Fig. 3), we extracted two separate light curves for each of the two peaks. We found that the signal at 343.5 s was present with a similar pulsed fraction in both light curves; this result is consistent with the hypothesis of a point–like source. An independent confirmation was also obtained through a raytracing simulation, encompassing a monochromatic source (1.49 keV) and models for the mirror assembly and ACIS detector. The results indicate that expected PSF at an off–axis angle of 9″ is artificially elongated in the direction connecting the two peaks detected in the Chandra image around the position of 2E 0101.5–7225; the two peaks are probably artefacts resulting from the small number of X–rays in the image. We conclude that the source is consistent with being point–like.
### 2.4. Optical Observations
Optical images (H$`\alpha `$ and H$`\alpha `$ red–continuum; 200 s each) of the BeppoSAX error circle of SAX J0103.2–7209 were obtained on 1998 July 24 at the Danish 1.5 m telescope with the Danish Faint Object Spectrometer Camera (DFOSC) at La Silla (Chile) in order to search for emission–line stars. The data were reduced using standard ESO–MIDAS procedures for bias subtraction, flat–field correction, aperture photometry and one dimensional stellar and sky spectra extraction. Within the BeppoSAX position circle we found only one H$`\alpha `$ active object: this is the O9–B1 III–Ve m<sub>V</sub>=14.8 star originally suggested as the optical counterpart of 2E 0101.5–7225 (Hughes & Smith 1994). A 10 Ȧ resolution 3800–8500 Ȧ spectrum (2$`\stackrel{}{\mathrm{.}}`$0 slit) of this star was obtained with the same instrument on 1998 October 20. Strong H$`\alpha `$ and H$`\beta `$ emission–lines were detected, with equivalent width –20 $`\pm `$ 2 Ȧ, and –1.8 $`\pm `$ 0.2 Ȧ, respectively. These results are in good agreement with those obtained by Hughes & Smith (1994) during observations carried out on December 1992.
## 3. Discussion
Although 2E 0101.5–7225 shares several characteristics with other X–ray pulsators in the Magellanic Clouds (MCs) and in the Galactic plane (i.e., the spin period and pulsed fraction, the spectral shape at high energy and the presence of a soft thermal component, the association with a high mass companion, etc.), it has also some peculiar differences which make it unusual.
The absence of long–term variability greater than a factor of $``$5–10 over an interval of $``$ 20 years is strongly suggestive of a persistent X–ray pulsar, the second in the SMC since the discovery of pulsations from SMC X–1. Moreover, the association of the companion with a Be spectral–type star makes 2E 0101.5–7225 the first example of a persistent (main sequence) Be/X–ray binary system in the SMC. 2E 0101.5–7225 is also a relatively low–luminosity (10<sup>35</sup>–10<sup>36</sup> erg s<sup>-1</sup>) X–ray pulsar. Finally, the inferred period derivative of –1.7 s yr<sup>-1</sup> corresponds to a secular spin–up time–scale of $``$ 200 yr which is the shortest of any known X–ray pulsar in a High Mass X–ray binary.
By using the P and Ṗ measurements and an X–ray luminosity in the 0.5–2 $`\times `$ 10<sup>36</sup> erg s<sup>-1</sup> range, we infer a magnetic field of 4–12 $`\times `$ 10<sup>12</sup> Gauss and, correspondingly, a maximum magnetospheric radius of $`r_\mathrm{m}`$ $``$ 4 $`\times `$ 10<sup>9</sup> cm (see Lamb et al. 1973). Since the derived corotation radius $`r_{\mathrm{co}}`$ is $``$ 8 $`\times `$ 10<sup>9</sup> cm, the $`r_\mathrm{m}`$ is considerably smaller than $`r_{\mathrm{co}}`$, and accretion on the surface of the neutron star proceeds unaffected by the magnetosphere’s centrifugal drag, as long as the luminosity of the source is larger than $``$ 2 $`\times `$ 10<sup>35</sup> erg s<sup>-1</sup>. Such a luminosity is not unusual for a bright X–ray pulsar in a nearly circular orbit around a giant OB companion star where the intense stellar wind continuously supplies matter to the compact object. Assuming a main sequence star, an orbital period of $``$ 300 d is expected for 2E 0101.5–7225 based on the pulse period – orbital period correlation of Be star/X–ray pulsar binaries (Corbet 1984). A timing analysis of I band photometric measurements from the Optical Gravitational Lensing Experiment (OGLE) revealed no evidence of any coherent signal in the 1–50 d period interval (see Coe & Orosz 1999) suggesting that 2E 0101.5–7225 is a long period system.
The characteristics of 2E 0101.5–7225 have their closest analogy to those of the persistent low–luminosity (10<sup>35</sup>–10<sup>36</sup> erg s<sup>-1</sup>; 2–10 keV band) X–ray pulsar RX J0146.9+6121, a 1455 s (spin–up timescale of $``$ 250 yr) spinning neutron star in a binary system with a B1Ve companion star (see Reig et al. 1997; Haberl et al. 1998; Mereghetti et al. 2000). Two other recently discovered long period Be/X–ray pulsars in the SMC, namely AX J0051–733 (P<sub>s</sub>=323 s; Imanishi et al. 1999) and AX J0058–7203 (P<sub>s</sub>=280 s; Tsujimoto 1999), although highly variable (a factor of 10–100), have a maximum luminosity level close to that of 2E 0101.5–7225.
In conclusion, we discovered a new X–ray pulsar with a period of $``$ 345 s in the SMC which is persistent, has a relatively low–luminosity, is rapidly spinning–up and is associated with a O9–B1 III–Ve star. The relatively high value inferred for the magnetic field and period derivative point to a young object. All these findings make 2E 0101.5–7225 an unusual X–ray pulsar which deserves more detailed investigations with instrumentation ranging from the X–ray to the IR band.
This work was partially supported through ASI and CNAA grants. We thank the Chandra Data Archive (CDA) of the Chandra X–Ray Observatory Science Center (CXC) at CfA for a prompt release of the data. We thank R. Ragazzoni for obtaining the optical spectrum of the counterpart of 2E 0101.5–7225. We thank H. Tananbaum, R.J. Edgar and L. Burderi for their helpful comments. This work has been supported in part by NASA contract NAS8-39073. |
warning/0001/nlin0001001.html | ar5iv | text | # Contents
### Introduction
Hypergeometric functions play an important role both, in physics and in mathematics . Many special functions and polynomials (such as $`q`$-Askey-Wilson polynomials, $`q`$-Jacobi polynomials, $`q`$-Gegenbauer polynomials , $`q`$-Racah polynomials , $`q`$-Hahn polynomials , expressions for Clebsch-Gordan coefficients) are just certain hypergeometric functions evaluated at special values of parameters. In physics hypergeometric functions and their $`q`$-deformed counterparts sometimes play the role of wave functions and correlation functions for quantum integrable systems. In the present paper we shall construct hypergeometric functions as tau-functions $`\tau `$ of the Kadomtsev-Petviashvili (KP) hierarchy of equations. It is interesting that the KP equation
$$4_{t_1}_{t_3}u=_{t_1}^4u+3_{t_2}^2u+3_{t_1}^2u^2\left(u=2_{t_1}^2\mathrm{log}\tau \right),$$
(0.0.1)
which originally served in plasma physics now plays a very important role both, in physics (see ; see review in for modern applications) and in mathematics. The peculiarity of Kadomtsev-Petviashvili equation appeared in the paper where L-A pair of KP equation was presented, and mainly in the paper of V.E.Zakharov and A.B.Shabat in 1974 where this equation was integrated by the dressing method. Actually it was the paper where so-called hierarchy of higher KP equations appeared. Another very important equation is the two-dimensional Toda lattice (TL) integrated first in . In the present paper we use these equations to construct hypergeometric functions which depend on many variables, these variables are KP and Toda lattice higher times. Here we shall use the general approach to integrable hierarchies of Kyoto school , see also . Especially a set of papers about Toda lattice is important for us.
About the structure of the paper. To learn the main result, which is the fermionic representation of the hypergeometric functions, one needs read only Sections 1 and Subsections 2.1 and 2.2 - to learn the notations, and then read Subsection 3.1, and Examples 3-6. The linear equations (constraints) for the tau-function, which generalize familiar Gauss equation for the well-known Gauss hypergeometric function see in Subsection 3.7. For determinant representation and integral representation see Subsections 3.8 and 3.9. For hypergeometric function as group two-cocycle see the Remark in the end of Appendix “Gauss factorization problem etc.”. All other material is just setting the topic into the theory of integrable systems.
Few words about notations. The symbols $``$ and $`\overline{}`$ do not denote the complex conjugation. Symbol does not denote the derivative. Bold $`𝐧`$ stands for partitions. Bold $`𝐭`$ and $`𝐭^{}`$ stand for collections of KP and TL higher time variables. Bold $`𝐱_{\left(N\right)}`$ and $`𝐲_{\left(N\right)}`$ stand for Miwa variables.
### 1 Milne’s hypergeometric series
#### 1.1 Ordinary hypergeometric functions
First let us remember that generalized hypergeometric function of one variable $`x`$ is defined as
$${}_{p}{}^{}F_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;x)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\left(a_1\right)_n\mathrm{}\left(a_p\right)_n}{\left(b_1\right)_n\mathrm{}\left(b_s\right)_n}\frac{x^n}{n!}.$$
(1.1.1)
Here $`\left(a\right)_n`$ is Pochhammer’s symbol:
$$\left(a\right)_n=\frac{\mathrm{\Gamma }\left(a+n\right)}{\mathrm{\Gamma }\left(a\right)}=a\left(a+1\right)\mathrm{}\left(a+n1\right).$$
(1.1.2)
Given number $`q`$, $`\left|q\right|<1`$, the so-called basic hypergeometric series of one variable is defined as
$${}_{p}{}^{}\mathrm{\Phi }_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;q,x)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{(q^{a_1};q)_n\mathrm{}(q^{a_p};q)_n}{(q^{b_1};q)_n\mathrm{}(q^{b_s};q)_n}\frac{x^n}{(q;q)_n}.$$
(1.1.3)
Here $`(q^a,q)_n`$ is $`q`$-deformed Pochhammer’s symbol:
$$(q^a;q)_0=1,(q^a;q)_n=\left(1q^a\right)\left(1q^{a+1}\right)\mathrm{}\left(1q^{a+n1}\right).$$
(1.1.4)
Both series converge for all $`x`$ in case $`p<s+1`$. In case $`p=s+1`$ they converge for $`\left|x\right|<1`$. We refer these well-known hypergeometric functions as ordinary hypergeometric functions.
#### 1.2 The multiple basic hypergeometric series related to Schur polynomials
There are several well-known different multivariable generalizations of hypergeometric series of one variable . Let $`\left|q\right|<1`$ and let $`𝐱_{\left(N\right)}=(x_1,\mathrm{},x_N)`$ be indeterminates. Let $`s_𝐧(x_1,x_2,\mathrm{},x_N)`$ be the Schur polynomial corresponding to a partition $`𝐧`$ . $`s_𝐧(x_1,x_2,\mathrm{},x_N)`$ is a symmetric function of variables $`x_k`$. The multiple basic hypergeometric series related to Schur polynomials were introduced by S.Milne as
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;q,𝐱_{\left(N\right)})={\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{(q^{a_1};q)_𝐧\mathrm{}(q^{a_p};q)_𝐧}{(q^{b_1};q)_𝐧\mathrm{}(q^{b_s};q)_𝐧}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}s_𝐧\left(𝐱_{\left(N\right)}\right),`$ (1.2.1)
where the sum is over all different partitions $`𝐧=(n_1,n_2,\mathrm{},n_r)`$, where $`n_1n_2\mathrm{}n_r`$, $`r\left|𝐧\right|`$, $`\left|𝐧\right|=n_1+\mathrm{}+n_r`$ and whose length $`l\left(𝐧\right)=rN`$. Schur polynomial $`s_𝐧\left(𝐱_{\left(N\right)}\right)`$, with $`Nl\left(𝐧\right)`$, is a symmetric function of variables $`𝐱_{\left(N\right)}`$ and defined as follows :
$$s_𝐧\left(𝐱_{\left(N\right)}\right)=\frac{a_{𝐧+\delta }}{a_\delta },a_𝐧=det\left(x_i^{n_j}\right)_{1i,jN},\delta =(N1,N2,\mathrm{},1,0).$$
(1.2.2)
Coefficient $`(q^c;q)_𝐧`$ associated with partition $`𝐧`$ is expressed in terms of the $`q`$-deformed Pochhammer’s Symbols $`(q^c;q)_n`$ (1.1.4):
$$(q^c;q)_𝐧=(q^c;q)_{n_1}(q^{c1};q)_{n_2}\mathrm{}(q^{cl+1};q)_{n_l}.$$
(1.2.3)
The multiple $`q^{n\left(𝐧\right)}`$ defined on the partition $`𝐧`$:
$$q^{n\left(𝐧\right)}=q^{_{i=1}^N\left(i1\right)n_i},$$
(1.2.4)
and $`q`$-deformed ’hook polynomial’ $`H_𝐧\left(q\right)`$ is
$`H_𝐧\left(q\right)={\displaystyle \underset{(i,j)𝐧}{}}\left(1q^{h_{ij}}\right),h_{ij}=\left(n_i+n_j^{}ij+1\right),`$ (1.2.5)
where $`𝐧^{}`$ is the conjugated partition (for the definition see ). For $`N=1`$ we get (1.1.3).
Another generalization of hypergeometric series is so-called hypergeometric function of matrix argument $`𝐗`$ with indices $`𝐚`$ and $`𝐛`$ :
$`{}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1,\mathrm{},a_p}{b_1,\mathrm{},b_s}}|𝐗\right)={\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\left(a_1\right)_𝐧\mathrm{}\left(a_p\right)_𝐧}{\left(b_1\right)_𝐧\mathrm{}\left(b_s\right)_𝐧}}{\displaystyle \frac{Z_𝐧\left(𝐗\right)}{\left|𝐧\right|!}}.`$ (1.2.6)
Here $`𝐗`$ is a Hermitian $`N\times N`$ matrix, and $`Z_𝐧\left(𝐗\right)`$ is zonal spherical polynomial for the symmetric space $`GL(N,C)/U\left(N\right)`$, see . Let us note that in the limit $`q1`$ series (1.2.1) coincides with (1.2.6), see .
#### 1.3 Hypergeometric series of double set of arguments
The formula
$${}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left(\genfrac{}{}{0pt}{}{a_1,\mathrm{},a_p}{b_1,\mathrm{},b_s}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=$$
$$\underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}\frac{(q^{a_1};q)_𝐧\mathrm{}(q^{a_p};q)_𝐧}{(q^{b_1};q)_𝐧\mathrm{}(q^{b_s};q)_𝐧}\frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}\frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{s_𝐧(1,q,q^2,\mathrm{},q^{N1})}$$
(1.3.1)
defines the multiple basic hypergeometric function of two sets of variables which was also introduced by S.Milne, see , .
Another generalization of hypergeometric series is so-called hypergeometric function of matrix arguments $`𝐗,𝐘`$ with indices $`𝐚`$ and $`𝐛`$:
$`{}_{p}{}^{}_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;𝐗,𝐘)={\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\left(a_1\right)_𝐧\mathrm{}\left(a_p\right)_𝐧}{\left(b_1\right)_𝐧\mathrm{}\left(b_s\right)_𝐧}}{\displaystyle \frac{Z_𝐧\left(𝐗\right)Z_𝐧\left(𝐘\right)}{\left|𝐧\right|!Z_𝐧\left(𝐈_n\right)}}.`$ (1.3.2)
Here $`𝐗,𝐘`$ are Hermitian $`N\times N`$ matrices and $`Z_𝐧\left(𝐗\right),Z_𝐧\left(𝐘\right)`$ are zonal spherical polynomials for the symmetric spaces $`GL(N,C)/U\left(N\right)`$, $`GL(N,R)/SO\left(N\right)`$ and $`GL(N,H)/Sp\left(N\right)`$ see . In our paper we shall consider only the first case; different hypergeometric functions related to zonal polynomials for symmetric spaces $`GL(N,C)/U\left(N\right)`$, $`GL(N,R)/SO\left(N\right)`$ and $`GL(N,H)/Sp\left(N\right)`$ will not be considered.
There are also hypergeometric functions related to Jack polynomials $`C_𝐧^{\left(d\right)}`$ :
$`{}_{p}{}^{}_{s}^{}{}_{}{}^{\left(d\right)}(a_1,\mathrm{},a_p;b_1,\mathrm{}b_s;𝐱_{\left(N\right)},𝐲_{\left(N\right)})=`$
$`{\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\left(a_1\right)_𝐧^{\left(d\right)}\mathrm{}\left(a_p\right)_𝐧^{\left(d\right)}}{\left(b_1\right)_𝐧^{\left(d\right)}\mathrm{}\left(b_s\right)_𝐧^{\left(d\right)}}}{\displaystyle \frac{C_𝐧^{\left(d\right)}\left(𝐱_{\left(N\right)}\right)C_𝐧^{\left(d\right)}\left(𝐲_{\left(N\right)}\right)}{\left|𝐧\right|!C_𝐧^{\left(d\right)}\left(1^n\right)}},`$ (1.3.3)
where
$$\left(a\right)_𝐧^{\left(d\right)}=\underset{i=1}{\overset{l\left(𝐧\right)}{}}\left(a\frac{d}{2}\left(i1\right)\right)_{n_i}.$$
(1.3.4)
Here $`\left(c\right)_k=c\left(c+1\right)\mathrm{}\left(c+k1\right)`$. It is known that for the special value $`d=2`$ the last expression (1.3) coincides with (1.2.1), and coincides with (1.3.2) as $`\left|q\right|1`$. These last cases we shall consider below.
### 2 A brief introduction to the fermionic description of the KP and TL hierarchies
#### 2.1 Fermionic operators and Fock space
We have fermionic fields:
$$\psi \left(z\right)=\underset{k}{}\psi _kz^k,\psi ^{}\left(z\right)=\underset{k}{}\psi _k^{}z^{k1}dz,$$
(2.1.1)
where fermionic operators satisfy the canonical anti-commutation relations:
$$[\psi _m,\psi _n]_+=[\psi _m^{},\psi _n^{}]_+=0;[\psi _m,\psi _n^{}]_+=\delta _{mn}.$$
(2.1.2)
Let us introduce left and right vacuums by the properties:
$`\psi _m|0=0\left(m<0\right),\psi _m^{}|0=0\left(m0\right),`$ (2.1.3)
$`0|\psi _m=0\left(m0\right),0|\psi _m^{}=0\left(m<0\right).`$ (2.1.4)
The vacuum expectation value is defined by relations:
$`0\left|1\right|0=1,0\left|\psi _m\psi _m^{}\right|0=1m<0,0\left|\psi _m^{}\psi _m\right|0=1m0,`$ (2.1.5)
$$0\left|\psi _m\psi _n\right|0=0\left|\psi _m^{}\psi _n^{}\right|0=0,0\left|\psi _m\psi _n^{}\right|0=0mn.$$
(2.1.6)
Let us notice that relations (2.1.2)-(2.1.6) are invariant under the transformation
$$\psi _ne^{T_n}\psi _n,\psi _n^{}e^{T_n}\psi _n^{}\left(T_nC\right).$$
(2.1.7)
Consider infinite matrices $`\left(a_{ij}\right)_{i,jZ}`$ satisfying the condition: there exists an $`N`$ such that $`a_{ij}=0`$ for $`\left|ij\right|>N`$. Let us take the set of linear combinations of quadratic elements $`a_{ij}:\psi _i\psi _j^{}:`$, where $`::`$ means the normal ordering $`:\psi _i\psi _j^{}:=\psi _i\psi _j^{}0\left|\psi _i\psi _j^{}\right|0`$. These elements together with $`1`$ span an infinite dimensional Lie algebra $`\widehat{gl}\left(\mathrm{}\right)`$:
$$[a_{ij}:\psi _i\psi _j^{}:,b_{ij}:\psi _i\psi _j^{}:]=c_{ij}:\psi _i\psi _j^{}:+c_0,$$
(2.1.8)
$$c_{ij}=\underset{k}{}a_{ik}b_{kj}\underset{k}{}b_{ik}a_{kj},$$
(2.1.9)
$$c_0=\underset{i<0,j0}{}a_{ij}b_{ji}\underset{i0,j<0}{}a_{ij}b_{ji}.$$
(2.1.10)
Now we define the operator $`g`$ which is an element of the group corresponding to the Lie algebra $`\widehat{gl}\left(\mathrm{}\right)`$:
$$g\psi _ng^1=\underset{m}{}\psi _ma_{mn},g^1\psi _n^{}g=\underset{m}{}a_{nm}\psi _m^{}.$$
(2.1.11)
#### 2.2 The KP and Toda tau functions
First let us define the vacuum vectors labeled by the integer $`M`$:
$`M|=0|\mathrm{\Psi }_M^{},|M=\mathrm{\Psi }_M|0,`$ (2.2.1)
$`\mathrm{\Psi }_M=\psi _{M1}\mathrm{}\psi _1\psi _0M>0,\mathrm{\Psi }_M=\psi _M^{}\mathrm{}\psi _2^{}\psi _1^{}M<0,`$
$`\mathrm{\Psi }_M^{}=\psi _0^{}\psi _1^{}\mathrm{}\psi _{M1}^{}M>0,\mathrm{\Psi }_M^{}=\psi _1\psi _2\mathrm{}\psi _MM<0.`$ (2.2.2)
The tau-function of the KP equation and the tau-function of the two-dimensional Toda lattice (TL) sometimes are defined as
$$\tau _{KP}(M,𝐭)=M\left|e^{H\left(𝐭\right)}g\right|M,$$
(2.2.3)
$$\tau _{TL}(M,𝐭,𝐭^{})=M\left|e^{H\left(𝐭\right)}ge^{H^{}\left(𝐭^{}\right)}\right|M.$$
(2.2.4)
According to the integer $`M`$ in (2.2.4) plays the role of discrete Toda lattice variable.
The times $`𝐭=(t_1,t_2,\mathrm{})`$ and $`𝐭^{}=(t_1^{},t_2^{},\mathrm{})`$ are called higher Toda lattice times (the first set $`𝐭`$ is in the same time the set of higher KP times. The first times of this set $`t_1,t_2,t_3`$ are independent variables for KP equation (0.0.1), which is the first nontrivial equation in the KP hierarchy). $`H\left(𝐭\right)`$ and $`H^{}\left(𝐭^{}\right)`$ belong to the following $`\widehat{gl}\left(\mathrm{}\right)`$ Cartan subalgebras:
$$H\left(𝐭\right)=\underset{n=1}{\overset{+\mathrm{}}{}}t_nH_n,H^{}\left(𝐭^{}\right)=\underset{n=1}{\overset{+\mathrm{}}{}}t_n^{}H_n,H_n=\frac{1}{2\pi i}:z^n\psi \left(z\right)\psi ^{}\left(z\right):.$$
(2.2.5)
For the Hamiltonians we have Heisenberg algebra commutation relations:
$$[H_n,H_m]=n\delta _{m+n,0}.$$
(2.2.6)
The action of $`e^{H\left(𝐭\right)}`$ on the fermions:
$`e^{H\left(𝐭\right)}\psi _ie^{H\left(𝐭\right)}={\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}p_n\left(𝐭\right)\psi _{in},e^{H\left(𝐭\right)}\psi _i^{}e^{H\left(𝐭\right)}={\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}p_n\left(𝐭\right)\psi _{i+n}^{},`$ (2.2.7)
where $`p_n`$ is the elementary Schur polynomial defined by the Taylor’s expansion:
$$e^{\xi (𝐭,z)}=\mathrm{exp}\left(\underset{k=1}{\overset{+\mathrm{}}{}}t_kz^k\right)=\underset{n=0}{\overset{+\mathrm{}}{}}z^np_n\left(𝐭\right).$$
(2.2.8)
The action on fermionic fields is especially simple:
$$e^{H\left(𝐭\right)}\psi \left(z\right)e^{H\left(𝐭\right)}=\psi \left(z\right)e^{\xi (𝐭,z)},e^{H\left(𝐭\right)}\psi ^{}\left(z\right)e^{H\left(𝐭\right)}=\psi ^{}\left(z\right)e^{\xi (𝐭,z)},$$
(2.2.9)
$$e^{H^{}\left(𝐭^{}\right)}\psi \left(z\right)e^{H^{}\left(𝐭^{}\right)}=\psi \left(z\right)e^{\xi (𝐭^{},z^1)},e^{H^{}\left(𝐭^{}\right)}\psi ^{}\left(z\right)e^{H^{}\left(𝐭^{}\right)}=\psi ^{}\left(z\right)e^{\xi (𝐭^{},z^1)}.$$
(2.2.10)
In the KP theory it is suitable to use another definition of the Schur function corresponding to the partition $`𝐧=(n_1,\mathrm{},n_r)`$:
$$s_𝐧\left(𝐭\right)=det\left(p_{n_ii+j}\left(𝐭\right)\right)_{1i,jr},$$
(2.2.11)
where $`p_m\left(𝐭\right)`$ is the elementary Schur polynomial defined by the Taylor’s expansion:
$$e^{\xi (𝐭,z)}=\mathrm{exp}\left(\underset{k=1}{\overset{+\mathrm{}}{}}t_kz^k\right)=\underset{n=0}{\overset{+\mathrm{}}{}}z^np_n\left(𝐭\right).$$
(2.2.12)
It is related to $`s_𝐧\left(𝐱_{\left(N\right)}\right)`$ and $`s_𝐧^{}\left(𝐱_{\left(N\right)}\right)`$, where a partition $`𝐧^{}`$ is conjugated to $`𝐧`$, as follows :
$$s_𝐧\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)=s_𝐧\left(𝐱_{\left(N\right)}\right),s_𝐧\left(𝐭^{}\left(𝐱_{\left(N\right)}\right)\right)=s_𝐧^{}\left(𝐱_{\left(N\right)}\right)$$
(2.2.13)
via the changes of variables (which is known as Miwa change of variables in the literature on the integrable systems):
$$t_m^+\left(𝐱_{\left(N\right)}\right)=\underset{i=1}{\overset{N}{}}\frac{x_i^m}{m},$$
(2.2.14)
$$t_m^{}\left(𝐱_{\left(N\right)}\right)=\underset{i=1}{\overset{N}{}}\frac{x_i^m}{m}.$$
(2.2.15)
Let us notice that $`s_𝐧\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)=0`$ for $`l\left(𝐧\right)>N`$, and $`s_𝐧\left(𝐭^{}\left(𝐱_{\left(N\right)}\right)\right)=0`$ for $`l\left(𝐧^{}\right)>N`$.
Lemma 1
For $`j_1<\mathrm{}<j_k<0i_s<\mathrm{}<i_1`$, $`sk0`$ the next formula is valid:
$$sk\left|e^{H\left(𝐭\right)}\psi _{j_1}^{}\mathrm{}\psi _{j_k}^{}\psi _{i_s}\mathrm{}\psi _{i_1}\right|0=\left(1\right)^{j_1+\mathrm{}+j_k+\left(ks\right)\left(ks+1\right)/2}s_𝐧\left(𝐭\right),$$
(2.2.16)
where the partition $`𝐧=(n_1,\mathrm{},n_{sk},n_{sk+1},\mathrm{},n_{sk+j_1})`$ is defined by the pair of partitions:
$`(n_1,\mathrm{},n_{sk})=(i_1\left(sk\right)+1,i_2\left(sk\right)+2,\mathrm{},i_{sk}),`$ (2.2.17)
$`(n_{sk+1},\mathrm{},n_{sk+j_1})=(i_{sk+1},\mathrm{},i_s|j_11,\mathrm{},j_k1).`$ (2.2.18)
The proof is achieved by direct calculation. Here $`\left(\mathrm{}|\mathrm{}\right)`$ is another notation for a partition due to Frobenius (see ).
#### 2.3 Baker-Akhiezer functions and bilinear identities
Vertex operators $`V_{\mathrm{}}\left(z\right)`$, $`V_{\mathrm{}}^{}\left(z\right)`$ and $`V_0\left(z\right)`$, $`V_0^{}\left(z\right)`$ act on the space $`C[t_1,t_2,\mathrm{}]`$ of polynomials in infinitely many variables, and are defined by the formulae:
$$V_{\mathrm{}}\left(z\right)=z^Me^{\xi (𝐭,z)}e^{\xi (\stackrel{~}{},z^1)},V_{\mathrm{}}^{}\left(z\right)=z^Me^{\xi (𝐭,z)}e^{\xi (\stackrel{~}{},z^1)},$$
(2.3.1)
$$V_0\left(z\right)=z^Me^{\xi (𝐭^{},z^1)}e^{\xi (\stackrel{~}{}^{},z)},V_0^{}\left(z\right)=z^Me^{\xi (𝐭^{},z^1)}e^{\xi (\stackrel{~}{}^{},z)},$$
(2.3.2)
where $`\stackrel{~}{}=(\frac{}{t_1},\frac{1}{2}\frac{}{t_2},\frac{1}{3}\frac{}{t_3},\mathrm{})`$, $`\stackrel{~}{}^{}=(\frac{}{t_1^{}},\frac{1}{2}\frac{}{t_2^{}},\frac{1}{3}\frac{}{t_3^{}},\mathrm{})`$.
We have the rules of the bosonization:
$`M+1|e^{H\left(𝐭\right)}\psi \left(z\right)=V_{\mathrm{}}\left(z\right)M|e^{H\left(𝐭\right)},M1|e^{H\left(𝐭\right)}\psi ^{}\left(z\right)=V_{\mathrm{}}^{}\left(z\right)M|e^{H\left(𝐭\right)},`$ (2.3.3)
$`\psi ^{}\left(z\right)e^{H^{}\left(𝐭^{}\right)}|M=V_0^{}\left(z\right)e^{H^{}\left(𝐭^{}\right)}|M+1,\psi \left(z\right)e^{H^{}\left(𝐭^{}\right)}|M=V_0\left(z\right)e^{H^{}\left(𝐭^{}\right)}|M1.`$ (2.3.4)
The Baker-Akhiezer functions and conjugated Baker-Akhiezer functions are:
$`w_{\mathrm{}}(M,𝐭,𝐭^{},z)={\displaystyle \frac{V_{\mathrm{}}\left(z\right)\tau }{\tau }},w_{\mathrm{}}^{}(M,𝐭,𝐭^{},z)={\displaystyle \frac{V_{\mathrm{}}^{}\left(z\right)\tau }{\tau }},`$ (2.3.5)
$`w_0(M,𝐭,𝐭^{},z)={\displaystyle \frac{V_0\left(z\right)\tau \left(M+1\right)}{\tau \left(M\right)}},w_0^{}(M,𝐭,𝐭^{},z)={\displaystyle \frac{V_0^{}\left(z\right)\tau \left(M1\right)}{\tau \left(M\right)}},`$ (2.3.6)
where
$$\tau (M,𝐭,𝐭^{})=M\left|e^{H\left(𝐭\right)}ge^{H^{}\left(𝐭^{}\right)}\right|M.$$
(2.3.7)
Both KP and TL hierarchies are described by the bilinear identity:
$$w_{\mathrm{}}(M,𝐭,𝐭^{},z)w_{\mathrm{}}^{}(M^{},𝐭^{},𝐭_{}^{}{}_{}{}^{},z)𝑑z=w_0(M,𝐭,𝐭^{},z^1)w_0^{}(M^{},𝐭^{},𝐭_{}^{}{}_{}{}^{},z^1)z^2𝑑z,$$
(2.3.8)
which holds for any $`𝐭,𝐭^{},𝐭^{},𝐭_{}^{}{}_{}{}^{}`$ for any integers $`M,M^{}`$.
The Schur functions $`s_𝐧\left(𝐭\right)`$ are well-known examples of tau-functions which correspond to rational solutions of the KP hierarchy. It is known that not any linear combination of Schur functions turns to be a KP tau-function, in order to find these combinations one should solve bilinear difference equation, see , which is actually a version of discrete Hirota equation. Below we shall present KP tau-functions which are infinite series of Schur polynomials, and which turn to be known hypergeometric functions (1.2.1),(1.3.2). We shall use the fermionic representation of tau-function .
### 3 Hypergeometric functions related to Schur functions
#### 3.1 KP tau-function $`\tau _r(M,𝐭,\beta )`$
Let $`r`$ be a function of one variable. Let $`D=z\frac{d}{dz}`$ acts on the basis $`\left\{z^n;nZ\right\}`$ of functions holomorphic in the punctured disk $`0<\left|z\right|<1`$ . Then we put $`r\left(D\right)z^n=r\left(n\right)z^n`$. All functions of operator $`D`$ which we consider below are given via their eigenvalues on this basis.
Let us consider an abelian subalgebra in $`\widehat{gl}\left(\mathrm{}\right)`$ formed by the set of fermionic operators
$$A_k=\frac{1}{2\pi i}\psi ^{}\left(z\right)\left(\frac{1}{z}r\left(D\right)\right)^k\psi \left(z\right),k=1,2,\mathrm{},$$
(3.1.1)
where the operator $`r\left(D\right)`$ acts on all functions of $`z`$ from the right hand side. In other terms
$$A_k=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\psi _{nk}^{}\psi _nr\left(n\right)r\left(n1\right)\mathrm{}r\left(nk+1\right),k=1,2,\mathrm{}.$$
(3.1.2)
We have $`[A_m,A_k]=0`$ for each $`m,k`$. Fermionic operators (3.1.2) resemble Toda lattice Hamiltonians $`H_k^{}`$ (2.2.5), and coincide with them if $`r\left(n\right)=1,nZ`$.
For the collection of independent variables $`\beta =(\beta _1,\beta _2,\mathrm{})`$ we denote
$$A\left(\beta \right)=\underset{n=1}{\overset{\mathrm{}}{}}\beta _nA_n.$$
(3.1.3)
For the partition $`𝐧=(n_1,\mathrm{},n_k)`$ and a function of one variable $`r`$, let us introduce the notation
$$r_𝐧\left(M\right)=\underset{i=1}{\overset{k}{}}r\left(1i+M\right)r\left(2i+M\right)\mathrm{}r\left(n_ii+M\right).$$
(3.1.4)
We set $`r_\mathrm{𝟎}\left(M\right)=1`$. Using the notation from (2.2.16) we have
Lemma 2 The following formula holds
$$0\left|\psi _{i_1}^{}\mathrm{}\psi _{i_s}^{}\psi _{j_s}\mathrm{}\psi _{j_1}e^{A\left(\beta \right)}\right|0=\left(1\right)^{j_1+\mathrm{}+j_s}r_𝐧\left(0\right)s_𝐧\left(\beta \right).$$
(3.1.5)
The proof is achieved by a direct calculation using $`e^A=1+A+\frac{1}{2}A^2+\mathrm{}`$, (3.1.2), the hook decomposition of $`𝐧`$ and (2.2.16).
Let us consider the tau-function (2.2.3) of the KP hierarchy
$$\tau _r(M,𝐭,\beta ):=M\left|e^{H\left(𝐭\right)}e^{A\left(\beta \right)}\right|M.$$
(3.1.6)
Using Taylor expanding $`e^H=1+H+\mathrm{}`$ and Lemma 1, Lemma 2 we easily get
###### Proposition 1
We have the expansion:
$$\tau _r(M,𝐭,\beta )=\underset{𝐧}{}r_𝐧\left(M\right)s_𝐧\left(𝐭\right)s_𝐧\left(\beta \right).$$
(3.1.7)
We shall not consider the problem of convergence of this series. The variables $`M,𝐭`$ play the role of KP higher times, $`\beta `$ is a collection of group times for a commuting subalgebra of additional symmetries of KP (see and Remark 7 in ). From different point of view (3.1.7) is a tau-function of two-dimensional Toda lattice with two sets of continuous variables $`𝐭`$, $`\beta `$ and one discrete variable $`M`$. Formula (3.1.7) is symmetric with respect to $`𝐭\beta `$. This ’duality’ supplies us with the string equations which characterize a tau-function of hypergeometric type (see below). In the similar expansions to (3.1.7) were considered, without specifying the coefficients and in a different context.
For given $`r`$ we define the function $`r^{}`$:
$$r^{}\left(n\right):=r\left(n\right).$$
(3.1.8)
###### Proposition 2
We have the involution:
$$\tau _r^{}(M,𝐭,\beta )=\tau _r(M,𝐭,\beta ).$$
(3.1.9)
The proof follows from the relations
$$r_𝐧^{}\left(M\right)=r_𝐧^{}\left(M\right),s_𝐧\left(𝐭\right)=s_𝐧^{}\left(𝐭\right).$$
(3.1.10)
Now let us introduce
$$\stackrel{~}{A}_k=\frac{1}{2\pi i}\psi ^{}\left(z\right)\left(\stackrel{~}{r}\left(D\right)z\right)^k\psi \left(z\right),\left(k=1,2,\mathrm{}\right),\stackrel{~}{A}\left(\beta \right)=\underset{n=1}{\overset{\mathrm{}}{}}\stackrel{~}{\beta }_n\stackrel{~}{A}_n.$$
(3.1.11)
Then we have the following generalization of Proposition 1:
###### Proposition 3
$$M\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A\left(\beta \right)}\right|M=\underset{𝐧}{}\left(\stackrel{~}{r}r\right)_𝐧\left(M\right)s_𝐧\left(\stackrel{~}{\beta }\right)s_𝐧\left(\beta \right).$$
(3.1.12)
###### Remark 1
This expansion has the following interpretation. If one uses vacuum vectors $`M||M`$ for normal ordering $`:A:=AM|A|M`$ in the formula for $`\widehat{gl}(\mathrm{})`$ commutation relation (2.1.8), he gets different $`\widehat{gl}(\mathrm{})`$ 2-cocycles $`c_M`$ which are cohomological to $`c_0`$ (2.1.10). The value of $`c_M`$ on the elements $`\stackrel{~}{A}_1,A_1`$ is $`\stackrel{~}{r}(M)r(M)`$:
$$c_M(\stackrel{~}{A}_1,A_1)=\left(\stackrel{~}{r}r\right)\left(M\right).$$
(3.1.13)
Formula (3.1.12) is an expansion of $`\widehat{GL}(\mathrm{})`$ group 2-cocycle, evaluated on the elements $`e^{\stackrel{~}{A}(\stackrel{~}{\beta })}`$,$`e^{A(\beta )}`$, in terms of $`r(M)`$, see also Appendix “Gauss factorization problem etc.”
In what follows we put $`\stackrel{~}{r}=1`$, since (3.1.12) depends only on $`\stackrel{~}{r}r`$.
#### 3.2 $`H_0(𝐓)`$, twisted fermions $`\psi (𝐓,z),\psi ^{}(𝐓,z)`$ and bosonization rules
Let $`r0`$, and put $`r\left(n\right)=e^{T_{n1}T_n}`$, where the variables $`T_n`$ are defined up to a constant independent of $`n`$. We define a Hamiltonian $`H_0\left(𝐓\right)\widehat{gl}\left(\mathrm{}\right)`$ (all $`T_nC`$ are finite):
$$H_0\left(𝐓\right):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}T_n:\psi _n^{}\psi _n:,$$
(3.2.1)
which produces the transformation (2.1.7):
$$e^{H_0\left(𝐓\right)}\psi _ne^{\pm H_0\left(𝐓\right)}=e^{\pm T_n}\psi _n,e^{H_0\left(𝐓\right)}\psi _n^{}e^{\pm H_0\left(𝐓\right)}=e^{T_n}\psi _n^{},$$
(3.2.2)
$$e^{H_0\left(𝐓\right)}\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)e^{H_0\left(𝐓\right)}=H\left(\stackrel{~}{\beta }\right),e^{H_0\left(𝐓\right)}A\left(\beta \right)e^{H_0\left(𝐓\right)}=H^{}\left(\beta \right).$$
(3.2.3)
Let $`r0`$. It is convenient to consider the fermionic operators:
$`\psi (𝐓,z)=e^{H_0\left(𝐓\right)}\psi \left(z\right)e^{H_0\left(𝐓\right)}={\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}e^{T_n}z^n\psi _n,`$ (3.2.4)
$`\psi ^{}(𝐓,z)=e^{H_0\left(𝐓\right)}\psi ^{}\left(z\right)e^{H_0\left(𝐓\right)}={\displaystyle \underset{n=\mathrm{}}{\overset{n=+\mathrm{}}{}}}e^{T_n}z^n\psi _n^{}{\displaystyle \frac{dz}{z}}.`$ (3.2.5)
For the variables $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ and $`𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)`$, and for the “Hamiltonians” $`A`$ and $`\stackrel{~}{A}`$ defined by (3.1.1), (3.1.11), one can derive the bosonization rules:
$`e^{A\left(𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)\right)}|M={\displaystyle \frac{\psi (𝐓,y_1)\mathrm{}\psi (𝐓,y_N)|MN}{\mathrm{\Delta }^+(M,N,𝐓,𝐲_{\left(N\right)})}},`$ (3.2.6)
$`e^{A\left(𝐭_{}^{}{}_{}{}^{}\left(𝐲_{\left(N\right)}\right)\right)}|M={\displaystyle \frac{\psi ^{}(𝐓,y_1)\mathrm{}\psi ^{}(𝐓,y_N)|M+N}{\mathrm{\Delta }^{}(M,N,𝐓,𝐲_{\left(N\right)})}},`$ (3.2.7)
$`M|e^{\stackrel{~}{A}\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)}={\displaystyle \frac{MN|\psi ^{}(\stackrel{~}{𝐓},\frac{1}{x_N})\mathrm{}\psi ^{}(\stackrel{~}{𝐓},\frac{1}{x_1})}{\stackrel{~}{\mathrm{\Delta }}^+(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})}},`$ (3.2.8)
$`M|e^{\stackrel{~}{A}\left(𝐭^{}\left(𝐱_{\left(N\right)}\right)\right)}={\displaystyle \frac{M+N|\psi (\stackrel{~}{𝐓},\frac{1}{x_N})\mathrm{}\psi (\stackrel{~}{𝐓},\frac{1}{x_1})}{\stackrel{~}{\mathrm{\Delta }}^{}(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})}}.`$ (3.2.9)
Here $`\stackrel{~}{T}_n`$ are related to $`\stackrel{~}{A}`$ via (3.1.11) and $`\stackrel{~}{r}\left(n\right)=e^{\stackrel{~}{T}_{n1}\stackrel{~}{T}_n}`$. Vandermond coefficients are
$`\mathrm{\Delta }^+(M,N,𝐓,𝐲_{\left(N\right)})={\displaystyle \frac{_{i<j}\left(y_iy_j\right)}{\left(y_1\mathrm{}y_N\right)^{NM}}}{\displaystyle \frac{\tau (M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}{\tau (MN,\mathrm{𝟎},𝐓,\mathrm{𝟎})}},`$ (3.2.10)
$`\mathrm{\Delta }^{}(M,N,𝐓,𝐲_{\left(N\right)})={\displaystyle \frac{_{i<j}\left(y_iy_j\right)}{\left(y_1\mathrm{}y_N\right)^{M+N}}}{\displaystyle \frac{\tau (M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}{\tau (M+N,\mathrm{𝟎},𝐓,\mathrm{𝟎})}},`$ (3.2.11)
$`\stackrel{~}{\mathrm{\Delta }}^+(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})={\displaystyle \frac{_{i<j}\left(x_ix_j\right)}{\left(x_1\mathrm{}x_N\right)^{NM1}}}{\displaystyle \frac{\tau (M,\mathrm{𝟎},\stackrel{~}{𝐓},\mathrm{𝟎})}{\tau (MN,\mathrm{𝟎},\stackrel{~}{𝐓},\mathrm{𝟎})}},`$ (3.2.12)
$`\stackrel{~}{\mathrm{\Delta }}^{}(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})={\displaystyle \frac{_{i<j}\left(x_ix_j\right)}{\left(x_1\mathrm{}x_N\right)^{N+M1}}}{\displaystyle \frac{\tau (M,\mathrm{𝟎},\stackrel{~}{𝐓},\mathrm{𝟎})}{\tau (M+N,\mathrm{𝟎},\stackrel{~}{𝐓},\mathrm{𝟎})}}.`$ (3.2.13)
The notation $`\tau (M,\mathrm{𝟎},𝐓,\mathrm{𝟎})`$ is explained in the next Subsection, see (3.3.5),(3.3.6).
Therefore in Miwa variables one can rewrite correlators (3.1.12):
$`M\left|e^{\stackrel{~}{A}\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)}e^{A\left(𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)\right)}\right|M=`$
$`{\displaystyle \frac{MN\left|\psi ^{}(\stackrel{~}{𝐓},\frac{1}{x_N})\mathrm{}\psi ^{}(\stackrel{~}{𝐓},\frac{1}{x_1})\psi (𝐓,y_1)\mathrm{}\psi (𝐓,y_N)\right|MN}{\stackrel{~}{\mathrm{\Delta }}^+(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})\mathrm{\Delta }^+(M,N,𝐓,𝐲_{\left(N\right)})}},`$ (3.2.14)
$`M\left|e^{\stackrel{~}{A}\left(𝐭^{}\left(𝐱_{\left(N\right)}\right)\right)}e^{A\left(𝐭_{}^{}{}_{}{}^{}\left(𝐲_{\left(N\right)}\right)\right)}\right|M=`$
$`{\displaystyle \frac{M+N\left|\psi (\stackrel{~}{𝐓},\frac{1}{x_N})\mathrm{}\psi (\stackrel{~}{𝐓},\frac{1}{x_1})\psi ^{}(𝐓,y_1)\mathrm{}\psi ^{}(𝐓,y_N)\right|M+N}{\stackrel{~}{\mathrm{\Delta }}^{}(M,N,\stackrel{~}{𝐓},𝐱_{\left(N\right)})\mathrm{\Delta }^{}(M,N,𝐓,𝐲_{\left(N\right)})}}.`$ (3.2.15)
#### 3.3 Toda lattice tau-function $`\tau (M,𝐭,𝐓,𝐭^{})`$
Now let us consider the Toda lattice tau-function (2.2.4), which depends on the three sets of variables $`𝐭,𝐓,𝐭^{}`$ and on $`MZ`$:
$$\tau (M,𝐭,𝐓,𝐭^{})=M\left|e^{H\left(𝐭\right)}\mathrm{exp}(\underset{\mathrm{}}{\overset{\mathrm{}}{}}T_n:\psi _n^{}\psi _n:)e^{H^{}\left(𝐭^{}\right)}\right|M,$$
(3.3.1)
where $`:\psi _n^{}\psi _n:=\psi _n^{}\psi _n0\left|\psi _n^{}\psi _n\right|0`$. Since the operator $`_{\mathrm{}}^{\mathrm{}}:\psi _n^{}\psi _n:`$ commutes with all elements of the $`\widehat{gl}\left(\mathrm{}\right)`$ algebra, one can put $`T_1=0`$ in (3.3.1). With respect to the KP and the TL dynamics the times $`T_n`$ have a meaning of integrals of motion. With respect to each pair of times $`(t_m,T_n)`$ one can consider the Liouville equation related to (3.3.1), see Appendix “Equations with respect to $`𝐓`$ variables” (the variables $`𝐭^{}`$ plays the role of integrals of motion for these Liouville equations).
As we shall see the hypergeometric functions (1.1.1),(1.1.3),(1.2.1),(1.2.6) listed in the Introduction are ratios of tau-functions (3.3.1) evaluated at special values of times $`M,𝐭,𝐓,𝐭^{}`$. It is true only in the case when all parameters $`a_k`$ of the hypergeometric functions are nonintegers. For the case when at least one of the indices $`a_k`$ is an integer, we will need a tau-function of an open Toda chain which will be considered in the next Sections.
Tau-function (3.3.1) is linear in each $`e^{T_n}`$. It is described by the Proposition
###### Proposition 4
$$\frac{\tau (M,𝐭,𝐓,𝐭^{})}{\tau (M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}=1+\underset{𝐧\mathrm{𝟎}}{}e^{\left(T_{M1}T_{n_1+M1}\right)+\left(T_{M2}T_{n_2+M2}\right)+\mathrm{}+\left(T_{Ml}T_{n_l+Ml}\right)}s_𝐧\left(𝐭\right)s_𝐧\left(𝐭^{}\right).$$
(3.3.2)
The sum is going over all different partitions
$$𝐧=(n_1,n_2,\mathrm{},n_l),l=1,2,3,\mathrm{},$$
(3.3.3)
excluding the partition $`\mathrm{𝟎}`$.
Let $`r0`$,$`\stackrel{~}{r}0`$. Then we put
$$r\left(n\right)=e^{T_{n1}T_n},\stackrel{~}{r}\left(n\right)=e^{\stackrel{~}{T}_{n1}\stackrel{~}{T}_n}.$$
(3.3.4)
Let us show the equivalence of (3.1.12) and (3.3.2) in this case. We have $`\tau (0,\mathrm{𝟎},𝐓,\mathrm{𝟎})=1`$ and
$$\tau (n,\mathrm{𝟎},𝐓,\mathrm{𝟎})=e^{T_{n1}\mathrm{}T_1T_0},n>0,$$
(3.3.5)
$$\tau (n,\mathrm{𝟎},𝐓,\mathrm{𝟎})=e^{T_n+\mathrm{}+T_2+T_1},n<0.$$
(3.3.6)
###### Proposition 5
Let $`\tau (n,\stackrel{~}{\beta },𝐓,\beta )`$ is Toda lattice tau-function (3.3.1) and $`\tau _{\stackrel{~}{r}r}(n,\stackrel{~}{\beta },\beta )`$ is defined by (3.1.12), where $`\stackrel{~}{r},r,𝐓,\stackrel{~}{𝐓}`$ are related by (3.3.4), the functions $`r,\stackrel{~}{r}`$ have no zeroes at integer values of argument then
$$\frac{\tau (n,\stackrel{~}{\beta },\stackrel{~}{𝐓}+𝐓,\beta )}{\tau (n,\mathrm{𝟎},\stackrel{~}{𝐓}+𝐓,\mathrm{𝟎})}=n\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A\left(\beta \right)}\right|n=\tau _{\stackrel{~}{r}r}(n,\stackrel{~}{\beta },\beta ).$$
(3.3.7)
This proposition follows from formulas (3.2.2)-(3.2.3).
For $`\stackrel{~}{r}=1`$ we can put $`\stackrel{~}{\beta }=𝐭`$. Then the next equations hold
$`_{t_1}_{\beta _1}\varphi _n=r\left(n\right)e^{\varphi _{n1}\varphi _n}r\left(n+1\right)e^{\varphi _n\varphi _{n+1}},e^{\varphi _n}={\displaystyle \frac{\tau _r(n+1,𝐭,\beta )}{\tau _r(n,𝐭,\beta )}},`$ (3.3.8)
$`\left(\tau \left(n\right):=\tau _r(n,𝐭,\beta )\right)\tau \left(n\right)_{\beta _1}_{t_1}\tau \left(n\right)_{t_1}\tau \left(n\right)_{\beta _1}\tau \left(n\right)=r\left(n\right)\tau \left(n1\right)\tau \left(n+1\right).`$ (3.3.9)
As we shall see eqs. (3.3.8) and (3.3.9) are still true in case $`r\left(n\right)`$ has zeroes.
If the function $`r`$ has no integer zeroes, using the change of variables
$$\phi _n=\varphi _nT_n,$$
(3.3.10)
we obtain Toda lattice equation in standard form :
$$_{t_1}_{t_1^{}}\phi _n=e^{\phi _{n+1}\phi _n}e^{\phi _n\phi _{n1}}.$$
(3.3.11)
As we see the variables $`T_n`$ might have the meaning of asymptotic values of the fields $`\varphi _n`$ for the class of tau-functions (3.3.1) which is characterized by the property $`\phi _n0`$ as $`t_10`$.
#### 3.4 Toda lattice consisted of open parts
Now we consider tau-function (3.3.1) with the modification of the definition of flows.
Definition. Let us introduce the function $`\delta `$ which is equal to zero when $`r`$ is equal to zero and is equal to unity otherwise:
$$\delta \left(n\right)=0ifr\left(n\right)=0,\delta \left(n\right)=1ifr\left(n\right)0.$$
(3.4.1)
Given collection of zeroes $`𝐦`$ of $`r`$:
$$𝐦=\left\{M_iZ\right\},M_{i+1}>M_i,r\left(M_i\right)=0.$$
(3.4.2)
we construct Hamiltonians labeled by $`𝐦`$:
$`H_k\left(𝐦\right)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\delta \left(n\right)\delta \left(n1\right)\mathrm{}\delta \left(nk+1\right)\psi _n\psi _{nk}^{},H^{}(𝐦;\beta )={\displaystyle H_k\left(𝐦\right)\beta _k}.`$ (3.4.3)
$`H_k\left(𝐦\right)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\delta \left(n+1\right)\delta \left(n+2\right)\mathrm{}\delta \left(n+k\right)\psi _n\psi _{n+k}^{},H(𝐦;𝐭)={\displaystyle H_k\left(𝐦\right)t_k}.`$ (3.4.4)
The tau-function of the open TL we are interested in, see (3.4.9) below, can be written in the three equivalent forms:
$`\tau _{op}(𝐦;M,𝐭,\beta )=M\left|e^{H\left(𝐭\right)}\mathrm{exp}({\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}T_n:\psi _n^{}\psi _n:)e^{H^{}(𝐦;\beta )}\right|M=`$ (3.4.5)
$`M\left|e^{H(𝐦;𝐭)}\mathrm{exp}({\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}T_n:\psi _n^{}\psi _n:)e^{H\left(\beta \right)}\right|M=M\left|e^{H(𝐦;𝐭)}\mathrm{exp}({\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}T_n:\psi _n^{}\psi _n:)e^{H(𝐦;\beta )}\right|M,`$ (3.4.6)
where $`T_nC`$ are some constants (times). This tau-function has the property:
$$M_i𝐦\tau _{op}(𝐦;M_i,𝐭,\beta )=1$$
(3.4.7)
for all values of times $`𝐭,\beta `$. If one consider
$$\phi _n=\mathrm{log}\frac{\tau _{op}(𝐦;n+1,𝐭,𝐓,\beta )}{\tau _{op}(𝐦;n,𝐭,𝐓,\beta )}$$
(3.4.8)
he comes to the equation of an open TL:
$$_{t_1}_{\beta _1}\phi _n=\delta \left(n\right)e^{\phi _{n1}\phi _n}\delta \left(n+1\right)e^{\phi _n\phi _{n+1}}.$$
(3.4.9)
The set of fields $`\phi _n`$ solves a number of open lattice problems in the set of intervals:
$$\left\{\phi _n,n<M_1\right\},$$
(3.4.10)
$$\left\{\phi _n,M_in<M_{i+1},M_{i+1}M_i>1\right\},$$
(3.4.11)
$$\left\{\phi _n,n>M_s\right\}.$$
(3.4.12)
The tau-function describes a set of open Toda lattices between each pair of neighbor zeroes (between neighbor zeroes $`M_{i+1}`$,$`M_i`$ there is an open chain with $`M_{i+1}M_i`$ number of sites), and two semiinfinite Toda lattices, one of them ends on the smallest zero and the other on the largest zero.
#### 3.5 Properties of the tau function $`\tau _r`$ when function r(n) has zeroes
Now let us introduce a set of $`T_n`$ variables with the help of relations
$$r\left(n\right)=e^{T_{n1}T_n}$$
(3.5.1)
for all $`n`$ where $`r\left(n\right)0`$. Equation (3.5.1) define variables $`T_n`$ uniquely up to an integration constant in each of the intervals (the number of the constants is equal to the number of intervals)
$$\left\{T_n,n<M_1\right\},$$
(3.5.2)
$$\left\{T_n,M_in<M_{i+1},M_{i+1}M_i>1\right\},$$
(3.5.3)
$$\left\{T_n,M_sn\right\}$$
(3.5.4)
separately. In case of there are zeroes such that $`M_{i+1}=M_i+1`$ one can define variables $`T_{M_i}`$, however we will not need them.
We introduce the Hamiltonian
$$H_0\left(𝐓\right)=\underset{n}{}T_n:\psi _n^{}\psi _n:,$$
(3.5.5)
where sum is over all $`n`$ satisfying one of the equations (3.5.2), (3.5.3) or (3.5.4).
We have
$$A_k=e^{H_0\left(𝐓\right)}H_k\left(𝐦\right)e^{H_0\left(𝐓\right)},A\left(\beta \right)=e^{H_0\left(𝐓\right)}H^{}(𝐦,\beta )e^{H_0\left(𝐓\right)}.$$
(3.5.6)
###### Proposition 6
Let $`\tau _{op}(𝐦,n,𝐭,𝐓,\beta )`$ is Toda lattice tau-function (3.4.5), and $`\tau _r(n,𝐭,\beta )`$ is defined by (3.1.6), where $`r`$ and $`𝐓`$ are related by (3.5.1), the functions $`r`$ has zeroes at integer values of argument described by (3.4.2) and for $`r0`$ the set of variables $`𝐓`$ is related to $`r`$ by (3.5.1)
$$\frac{\tau _{op}(𝐦,n,𝐭,𝐓,\beta )}{\tau _{op}(𝐦,n,\mathrm{𝟎},𝐓,\mathrm{𝟎})}=n\left|e^{H\left(𝐭\right)}e^{A\left(\beta \right)}\right|n=\tau _r(n,𝐭,\beta ).$$
(3.5.7)
Equations (3.3.8) and (3.3.9) are still true in case $`r(n)`$ has zeroes.
Hirota equation (3.3.9) can be viewed as recurrent relation which expresses tau-function with discrete Toda lattice variable $`n`$ via $`\tau _r(M_i\pm 1,𝐭,\beta ),\tau _r(M_i,𝐭,\beta )=1`$. It follows from (3.1.4),(3.1.7) that
$$r\left(M_k\right)=0\tau _r(M_k,𝐭,\beta )=1.$$
(3.5.8)
Then from (3.1.4),(3.1.7) we see the following. In the region (3.4.11) the series (3.1.7) has only a finite number of nonvanishing terms. For the region (3.4.12) the sum is only over the Young diagrams $`𝐧`$ of the length $`l\left(𝐧\right)<MM_1`$. For the region (3.4.10) only those diagrams $`𝐧`$ for which the conjugated diagrams $`𝐧^{}`$ have length $`l\left(𝐧^{}\right)M_sM`$ contribute the series (3.1.7).
In Appendix we shall write down a system of orthogonal polynomials related to $`𝐦`$.
###### Remark 2
There are two different ways to restrict the sum (3.1.7) to a sum over partitions of length $`l(𝐧)N`$ (or over $`l(𝐧^{})N`$). The second way is to use so-called Miwa’s change of variables.
#### 3.6 Notations
In order to simplify notations, we shall omit additional argument $`𝐦`$ and subindex which distinguish TL tau-function (3.3.1) and open TL tau-functions (3.4.5). Instead of $`\tau _{op}(𝐦,n,𝐭,𝐓,𝐭^{})`$ we shall write $`\tau (n,𝐭,𝐓,𝐭^{})`$. The notation $`\tau _r(M,𝐭,\beta )`$ will be used only for the KP tau-function (3.1.6). Also $`\beta =𝐭^{}`$. When TL higher times are expressed via Miwa change (2.2.14) or (2.2.15), sometimes we shall put the argument $`𝐱_{\left(N\right)}`$ at the place of the argument $`𝐭`$ and the argument $`𝐲_{\left(N\right)}`$ at the place of $`𝐭^{}`$, for instance $`\tau _r(M,𝐱_{\left(N\right)},𝐭^{})`$, $`\tau _r(M,𝐭,𝐲_{\left(N\right)})`$, $`\tau _r(M,𝐱_{\left(N\right)},𝐲_{\left(N\right)})`$.
#### 3.7 Linear equations for the tau-function $`\tau _r`$
Here we shall write down linear equations, which follow from the explicit fermionic representation of the tau-function (3.1.6) via the bosonization formulae (3.2) and .(3.2) These equations may be also viewed as the constraints which result in the string equations. For the variables $`𝐭^{}\left(𝐱_{\left(N\right)}\right)`$, using $`M|A=0`$ and making profit of the relation $`A_k=e^{H_0}H_ke^{H_0}`$ inside the fermionic correlator (3.2), we get the partial differential equations for the tau-function (3.1.7):
$$\frac{\tau _r(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭^{})}{t_k^{}}=\frac{1}{\stackrel{~}{\mathrm{\Delta }}}\left(\underset{i=1}{\overset{N}{}}\left(x_ir\left(D_{x_i}\right)\right)^k\right)\stackrel{~}{\mathrm{\Delta }}\tau _r(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭^{}),$$
(3.7.1)
where $`\stackrel{~}{\mathrm{\Delta }}=\stackrel{~}{\mathrm{\Delta }}^{}(M,N,\mathrm{𝟎},𝐱_{\left(N\right)})`$. These equations have the meaning of string constraint equations for the tau-function (3.1.7). In variables $`𝐭_{}^{}{}_{}{}^{}\left(𝐲_{\left(\mathrm{}\right)}\right)`$ we can rewrite (3.7.1):
$`\left(1\right)^k{\displaystyle \underset{i=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{e_{k1}(\frac{1}{y_1},\mathrm{},\frac{1}{y_{i1}},\frac{1}{y_{i+1}},\mathrm{})}{_{ji}\left(1\frac{y_i}{y_j}\right)}}{\displaystyle \frac{\tau _r(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭_{}^{}{}_{}{}^{}\left(𝐲_{\left(\mathrm{}\right)}\right))}{y_i}}=`$
$`{\displaystyle \frac{1}{\stackrel{~}{\mathrm{\Delta }}}}\left({\displaystyle \underset{i=1}{\overset{N}{}}}\left(x_ir\left(D_{x_i}\right)\right)^k\right)\stackrel{~}{\mathrm{\Delta }}\tau _r(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭_{}^{}{}_{}{}^{}\left(𝐲_{\left(\mathrm{}\right)}\right)),`$ (3.7.2)
where $`e_k\left(𝐲\right)`$ is a symmetric function defined through the relation $`_{i=1}^+\mathrm{}\left(1+ty_i\right)=_{k=0}^+\mathrm{}t^ke_k\left(𝐲\right)`$. Also we have
$`\left({\displaystyle \underset{k=1}{\overset{M+N1}{}}}k{\displaystyle \underset{i=1}{\overset{N}{}}}D_{x_i}\right)\stackrel{~}{\mathrm{\Delta }}^{}\tau (M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐓,𝐭^{}\left(𝐲_{\left(N^{}\right)}\right))\mathrm{\Delta }^{}=`$
$`\left({\displaystyle \underset{k=1}{\overset{M+N^{}1}{}}}k{\displaystyle \underset{i=1}{\overset{N^{}}{}}}\left({\displaystyle \frac{1}{y_i}}D_{y_i}y_i\right)\right)\stackrel{~}{\mathrm{\Delta }}^{}\tau (M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐓,𝐭^{}\left(𝐲_{\left(N^{}\right)}\right))\mathrm{\Delta }^{},`$ (3.7.3)
where $`\stackrel{~}{\mathrm{\Delta }}^{}=\stackrel{~}{\mathrm{\Delta }}^{}(M,N,\mathrm{𝟎},𝐱_{\left(N\right)})`$ and $`\mathrm{\Delta }^{}=\mathrm{\Delta }^{}(M,N,\mathrm{𝟎},𝐲_{\left(N^{}\right)})`$. This formula is obtained by the insertion of the fermionic operator $`res_z:\psi ^{}\left(z\right)z\frac{d}{dz}\psi \left(z\right):`$ inside the fermionic correlator. These formulae can be also written in terms of higher KP and TL times, with the help of vertex operator action, see the Appendix “Vertex operator action”. Then the relation (3.7.1) is the infinitesimal version of (A3.6), while the relation (3.7) is the infinitesimal version of (A3.7).
#### 3.8 Determinant formulae
With the help of Wick theorem one obtains the formulae.
###### Proposition 7
A generalization of Milne’s determinant formula
$$\tau _r(M,𝐭^+\left(𝐱_{\left(N\right)}\right),\beta )=\frac{det\left(x_i^{Nk}\tau _r(Mk+1,𝐭^+\left(x_i\right),\beta )\right)_{i,k=1}^N}{det\left(x_i^{Nk}\right)_{i,k=1}^N}.$$
(3.8.1)
Proof
$`\tau _r(M,𝐭^+\left(𝐱_{\left(N\right)}\right),\beta )=M\left|e^{H(𝐭^+\left(𝐱_{\left(N\right)}\right)}e^{A\left(\beta \right)}\right|M=`$ (3.8.2)
$`{\displaystyle \frac{x_1^{NM1}\mathrm{}x_N^{NM1}}{_{i<j}\left(x_ix_j\right)}}M\left|\psi _{M1}\mathrm{}\psi _{MN}\psi ^{}\left({\displaystyle \frac{1}{x_N}}\right)\mathrm{}\psi ^{}\left({\displaystyle \frac{1}{x_1}}\right)e^{A\left(\beta \right)}\right|M=`$ (3.8.3)
$`{\displaystyle \frac{\left(x_1\mathrm{}x_N\right)^{NM1}}{_{i<j}\left(x_ix_j\right)}}det\left(M\left|\psi _{Mk}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\right|M\right)_{i,k=1}^N=`$ (3.8.4)
$`{\displaystyle \frac{det\left(x_i^{Nk}\tau _r(Mk+1,𝐭^+\left(x_i\right),\beta )\right)_{i,k=1}^N}{det\left(x_i^{Nk}\right)_{i,k=1}^N}}`$ (3.8.5)
Last equality follows from:
$`M\left|\psi _{Mk}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\right|M=`$ (3.8.6)
$`=M\left|\psi _{M1}\mathrm{}\psi _{Mk+1}\psi _{Mk}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\psi _{Mk+1}^{}\mathrm{}\psi _{M1}^{}\right|M+`$ (3.8.7)
$`+{\displaystyle \underset{j=1}{\overset{k1}{}}}a_j^k\left(\beta \right)M\left|\psi _{M1}\mathrm{}\psi _{Mk+j}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\psi _{Mk+j+1}^{}\mathrm{}\psi _{M1}^{}\right|M=`$ (3.8.8)
$`=Mk+1\left|\psi _{Mk}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\right|Mk+1+`$ (3.8.9)
$`+{\displaystyle \underset{j=1}{\overset{k1}{}}}a_j^k\left(\beta \right)Mk+1+j\left|\psi _{Mk+j}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\right|Mk+1+j=`$ (3.8.10)
$`=x_i^{Mk+1}\tau _r(Mk+1,𝐭^+\left(x_i\right),\beta )+{\displaystyle \underset{j=1}{\overset{k1}{}}}a_j^k\left(\beta \right)x_i^{Mk+1+j}\tau _r(Mk+1+j,𝐭^+\left(x_i\right),\beta )`$ (3.8.11)
Where the functions $`a_j^k(\beta )`$ must be derived as the results of action of operator $`e^{A(\beta )}`$ on the fermions $`\psi _{M1},\mathrm{},\psi _{Mk}`$. Thus we have:
$`x_i^{NM1}M\left|\psi _{Mk}\psi ^{}\left({\displaystyle \frac{1}{x_i}}\right)e^{A\left(\beta \right)}\right|M=`$ (3.8.12)
$`=x_i^{Nk}\tau _r(Mk+1,𝐭^+\left(x_i\right),\beta )+{\displaystyle \underset{l=1}{\overset{k1}{}}}a_{kl}^k\left(\beta \right)x_i^{Nl}\tau _r(Ml+1,𝐭^+\left(x_i\right),\beta )`$ (3.8.13)
###### Proposition 8
For $`r0`$ we take a tau function $`\tau _r(M,𝐭^+(𝐱_{(N)}),𝐭_{}^{}{}_{}{}^{+}(𝐲_{(N)}))`$ and apply Wick’s theorem. We get the determinant formula:
$$\tau _r(M,𝐭^+\left(𝐱_{\left(N\right)}\right),𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right))=\frac{det\left(F\left(x_iy_j\right)\right)_{i,j=1}^N}{\stackrel{~}{\mathrm{\Delta }}^+(M,N,\mathrm{𝟎},𝐱_{\left(N\right)})\mathrm{\Delta }^+(M,N,𝐓,𝐲_{\left(N\right)})},$$
(3.8.14)
$$F\left(x_iy_j\right)=MN\left|\psi ^{}\left(\frac{1}{x_i}\right)\psi (𝐓,y_j)\right|MN.$$
(3.8.15)
#### 3.9 Integral representations
For the fermions (3.2.4) we easily get the relations:
$$\psi (𝐓,\alpha z)𝑑\mu \left(\alpha \right)=\psi (𝐓+𝐓\left(\mu \right),z),\psi ^{}(𝐓,\frac{1}{\alpha z})𝑑\stackrel{~}{\mu }\left(\alpha \right)=\psi ^{}(𝐓𝐓\left(\stackrel{~}{\mu }\right),\frac{1}{z})$$
(3.9.1)
where $`\mu ,\stackrel{~}{\mu }`$ are some integration measures, and shifts of times $`T_n`$ are defined in terms of the moments:
$$\alpha ^n𝑑\mu \left(\alpha \right)=e^{T_n\left(\mu \right)},\alpha ^n𝑑\stackrel{~}{\mu }\left(\alpha \right)=e^{T_n\left(\stackrel{~}{\mu }\right)}.$$
(3.9.2)
Therefore thanks to the bosonization formulae (3.2) we have the relations for the tau-function (below $`𝐭^{}`$ is defined via (2.2.14))
###### Proposition 9
Integral representation formula holds
$`{\displaystyle \stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{𝐓}}\left(\stackrel{~}{\alpha }𝐱_{\left(N\right)}\right)\frac{\tau (M,𝐭^+\left(\stackrel{~}{\alpha }𝐱_{\left(N\right)}\right),𝐓+\stackrel{~}{𝐓},𝐭_{}^{}{}_{}{}^{+}\left(\alpha 𝐲_{\left(N\right)}\right))}{\tau (M,\mathrm{𝟎},𝐓+\stackrel{~}{𝐓},\mathrm{𝟎})}\mathrm{\Delta }_𝐓\left(\alpha 𝐲_{\left(N\right)}\right)\underset{i=1}{\overset{N}{}}d\stackrel{~}{\mu }\left(\stackrel{~}{\alpha }_i\right)\underset{i=1}{\overset{N}{}}d\mu \left(\alpha _i\right)}`$
$`=\stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{𝐓}+\stackrel{~}{𝐓}\left(\stackrel{~}{\mu }\right)}\left(𝐱_{\left(N\right)}\right){\displaystyle \frac{\tau (M,𝐭^+\left(𝐱_{\left(N\right)}\right),𝐓+\stackrel{~}{𝐓}+\stackrel{~}{𝐓}\left(\stackrel{~}{\mu }\right)+𝐓\left(\mu \right),𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right))}{\tau (M,\mathrm{𝟎},𝐓+\stackrel{~}{𝐓}+\stackrel{~}{𝐓}\left(\stackrel{~}{\mu }\right)+𝐓\left(\mu \right),\mathrm{𝟎})}}\mathrm{\Delta }_{𝐓+𝐓\left(\mu \right)}\left(𝐲_{\left(N\right)}\right).`$ (3.9.3)
where $`\mathrm{\Delta }_𝐓(\alpha 𝐲_{(N)})=\mathrm{\Delta }^+(M,N,𝐓,\alpha 𝐲_{(N)})`$, $`\stackrel{~}{\mathrm{\Delta }}_{\stackrel{~}{𝐓}}(\stackrel{~}{\alpha }𝐱_{(N)})=\stackrel{~}{\mathrm{\Delta }}^+(M,N,\stackrel{~}{𝐓},\stackrel{~}{\alpha }𝐱_{(N)})`$,
$`\alpha 𝐲_{(N)}=(\alpha _1y_1,\alpha _2y_2,\mathrm{},\alpha _Ny_N)`$ and $`\stackrel{~}{\alpha }𝐱_{(N)}=(\stackrel{~}{\alpha }_1x_1,\stackrel{~}{\alpha }_2x_2,\mathrm{},\stackrel{~}{\alpha }_Nx_N)`$. In particular
$`{\displaystyle \frac{\tau (M,𝐭,𝐓,𝐭_{}^{}{}_{}{}^{+}\left(\alpha 𝐲_{\left(N\right)}\right))}{\tau (M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}\mathrm{\Delta }_𝐓\left(\alpha 𝐲_{\left(N\right)}\right)\underset{i=1}{\overset{N}{}}d\mu \left(\alpha _i\right)}=`$
$`{\displaystyle \frac{\tau (M,𝐭,𝐓+𝐓\left(\mu \right),𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right))}{\tau (M,\mathrm{𝟎},𝐓+𝐓\left(\mu \right),\mathrm{𝟎})}}\mathrm{\Delta }_{𝐓+𝐓\left(\mu \right)}\left(𝐲_{\left(N\right)}\right).`$ (3.9.4)
Remember that arbitrary linear combination of tau-functions is not a tau-function. Formulae (9) and also (9) give the integral representations for the tau-function (3.1.6). It may help to express a tau-function with the help of a more simple one. If we choose the integration measures:
$$\frac{i}{2\pi }_C\psi (𝐓,\alpha z)e^\alpha \left(\alpha \right)^{b1}𝑑\alpha =\psi (𝐓+𝐓^b,z),$$
(3.9.5)
$$_0^{\mathrm{}}\psi (𝐓,\alpha z)e^\alpha \alpha ^a𝑑\alpha =\psi (𝐓+𝐓^a,z),$$
(3.9.6)
$$_0^1\psi (𝐓,\alpha z)\alpha ^a\left(1\alpha \right)^{ba1}𝑑\alpha =\psi (𝐓+𝐓^c,z),$$
(3.9.7)
where $`C`$ starts at $`+\mathrm{}`$ on the real axis, circles the origin in the counterclockwise direction and returns to the starting point. Then
$$T_n^b=\mathrm{ln}\mathrm{\Gamma }\left(b+n+1\right),T_n^a=\mathrm{ln}\mathrm{\Gamma }\left(a+n+1\right),T_n^c=\mathrm{ln}\frac{\mathrm{\Gamma }\left(b+n+1\right)}{\mathrm{\Gamma }\left(a+n+1\right)\mathrm{\Gamma }\left(ba\right)}.$$
(3.9.8)
Also consider the $`q`$-integrals :
$$q^{\left(a+n\right)\left(a+n+1\right)}_0^{\mathrm{}}\psi (𝐓,\alpha \left(1q\right)z)E_q\left(\alpha \right)\alpha ^ad_q\alpha =\psi (𝐓+𝐓(a,q),z),$$
(3.9.9)
$$\frac{1}{\mathrm{\Gamma }_q\left(ba\right)}_0^1\psi (𝐓,\alpha z)\alpha ^a\frac{(\alpha q;q)_{\mathrm{}}}{(\alpha q^{ba};q)_{\mathrm{}}}d_q\alpha =\psi (𝐓+𝐓(a,b,q),z).$$
(3.9.10)
Then
$$T_n(a,q)=\mathrm{ln}\frac{1}{\left(1q\right)^n\mathrm{\Gamma }_q\left(a+n+1\right)},T_n(a,b,q)=\mathrm{ln}\frac{\mathrm{\Gamma }_q\left(b+n+1\right)}{\mathrm{\Gamma }_q\left(a+n+1\right)}.$$
(3.9.11)
In the same way one can consider Miwa change (2.2.15). In the Examples below we shall present hypergeometric functions listed in the Subsections 1.2 and 1.3 as tau-functions of the type (3.2). Then we are able to write down integration formulae, namely (9), which express $`{}_{p+1}{}^{}\mathrm{\Phi }_{s}^{}`$ and $`{}_{p+1}{}^{}\mathrm{\Phi }_{s+1}^{}`$ in terms of $`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}`$ with the help of (3.9.9), (3.9.10) and (9). In different integral representation formula was presented, which was based on the $`q`$-analog of Selberg’s integral of Askey and Kadell. By taking the limit $`q1`$ one can consider functions $`{}_{p}{}^{}F_{s}^{}`$. Using (3.9.8), one can express $`{}_{p+1}{}^{}_{s}^{}`$, $`{}_{p+1}{}^{}_{s+1}^{}`$ and $`{}_{p}{}^{}_{s+1}^{}`$ as integrals of $`{}_{p}{}^{}_{s}^{}`$ with the help of (3.9.6), (3.9.7) and (3.9.5) respectively.
#### 3.10 Examples
The main point of the paper is the observation that if $`r\left(D\right)`$ is a rational function of $`D`$ then $`\tau _r`$ is a hypergeometric series. If $`r\left(D\right)`$ is a rational function of $`q^D`$ we obtain $`q`$-deformed hypergeometric series. Now let us consider various $`r\left(D\right)`$.
Example 1 Let $`r=1`$. One can put $`𝐓=0`$. Then one gets
$$\tau _{r=1}(M,𝐭,𝐭^{})=\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}nt_nt_n^{}\right),$$
(3.10.1)
which is vacuum tau-function for the two-dimensional Toda lattice. Formula (3.10.1) is a manifestation of summation formulas for Schur functions . Let us note that this is also an example of function $`{}_{1}{}^{}_{0}^{}`$(1.2.6).
Example 2 Let $`r\left(n\right)=n`$, that is $`T_n=\mathrm{ln}\frac{1}{n!},n0`$ and $`T_n=\mathrm{ln}(1)^n(n1)!,n<0`$. Also let us put $`𝐭^{}=(t_1^{},0,0,\mathrm{})`$. For $`M=0,\pm 1`$ we get
$$\tau _r(0,𝐭,t_1^{})=1,\tau _r(1,𝐭,t_1^{})=e^{\xi (𝐭,t_1^{})},\tau _r(1,𝐭,t_1^{})=e^{\xi (𝐭,t_1^{})}.$$
(3.10.2)
Here $`t_1^{}`$ plays the role of spectral parameter for the vacuum Baker-Akhiezer function. This fact is in accordance to the meaning of $`t_1^{}`$ as a group time for the Galilean transformation .
Similar answers $`\tau =e^{\pm \xi (𝐭,z^{})}`$ one obtains if he substitutes $`t_n^{}=\pm n^1\left(z^{}\right)^n`$ to (3.10.1). Let us note that (3.10.2) are the functions $`{}_{1}{}^{}F_{0}^{}(0;t_1,t_2,\mathrm{})`$, $`{}_{1}{}^{}F_{0}^{}(1;t_1,t_2,\mathrm{})`$ and $`{}_{1}{}^{}F_{0}^{}(1;t_1,t_2,\mathrm{})`$ (1.3.2) which will be described below in the Example 3 (3.10).
Example 3 Let all parameters $`b_k`$ be nonintegers.
$${}_{p}{}^{}r_{s}^{}\left(D\right)=\frac{\left(D+a_1\right)\left(D+a_2\right)\mathrm{}\left(D+a_p\right)}{\left(D+b_1\right)\left(D+b_2\right)\mathrm{}\left(D+b_s\right)}.$$
(3.10.3)
If all $`a_k`$ are also nonintegers the relevant $`𝐓`$ is:
$${}_{}{}^{p}T_{n}^{s}=\mathrm{ln}\frac{\mathrm{\Gamma }\left(n+a_1+1\right)\mathrm{\Gamma }\left(n+a_2+1\right)\mathrm{}\mathrm{\Gamma }\left(n+a_p+1\right)}{\mathrm{\Gamma }\left(n+b_1+1\right)\mathrm{\Gamma }\left(n+b_2+1\right)\mathrm{}\mathrm{\Gamma }\left(n+b_s+1\right)}.$$
(3.10.4)
For the correlator (3.1.7) we have:
$$\frac{{}_{}{}^{p}\tau _{}^{s}(M,𝐭,𝐓,𝐭^{})}{{}_{}{}^{p}\tau _{}^{s}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}={}_{}{}^{p}\tau _{r}^{s}(M,𝐭,𝐭^{})=\underset{𝐧}{}s_𝐧\left(𝐭\right)s_𝐧\left(𝐭^{}\right)\frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}.$$
(3.10.5)
If in formula (3.10.5) we put
$$t_1^{}=1,t_i^{}=0,i>1$$
(3.10.6)
then $`s_𝐧\left(𝐭^{}\right)=H_𝐧^1`$, and we obtain the hypergeometric function related to Schur functions (see for help):
$`{\displaystyle \frac{{}_{}{}^{p}\tau _{}^{s}(M,𝐭,𝐓,𝐭^{})}{{}_{}{}^{p}\tau _{}^{s}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}}={}_{}{}^{p}\tau _{r}^{s}(M,𝐭,𝐭^{})={}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|t_1,t_2,\mathrm{}\right)=`$
$`{\displaystyle \underset{𝐧}{}}{\displaystyle \frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}}{\displaystyle \frac{s_𝐧\left(𝐭\right)}{H_𝐧}}.`$ (3.10.7)
In the last formula $`H_𝐧`$ is the following hook polynomial (compare with (1.2.5)):
$`H_𝐧={\displaystyle \underset{(i,j)𝐧}{}}h_{ij},h_{ij}=\left(n_i+n_j^{}ij+1\right).`$ (3.10.8)
We obtain ordinary hypergeometric function of one variable of type
$${}_{p1}{}^{}F_{s}^{}(a_2,\mathrm{},a_p;b_1,\mathrm{},b_s;\pm t_1t_1^{})=\tau _r(\pm 1,𝐭,𝐓,𝐭^{}),$$
(3.10.9)
if we take $`a_1=0`$, $`𝐭=(t_1,0,0,\mathrm{})`$, $`𝐭^{}=(t_1^{},0,0,\mathrm{})`$ .
For variables $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ the formula (3.10) turns out to be
$`{}_{}{}^{p}\tau _{r}^{s}(M,𝐭^+\left(𝐱_{\left(N\right)}\right),𝐭^{})={}_{p}{}^{}F_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|𝐱_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)}{H_𝐧}}.`$ (3.10.10)
We got the hypergeometric function (1.3.2) related to zonal polynomials for the symmetric space $`GL(N,C)/U\left(N\right)`$ . Here $`x_i=z_i^1,i=1,\mathrm{},N`$ are the eigenvalues of the matrix $`𝐗`$, and for zonal spherical polynomials there is the following matrix integral representation
$$Z_𝐧\left(𝐗\right)=Z_𝐧\left(𝐈_N\right)_{U(N,C)}\mathrm{\Delta }^𝐧\left(U^{}𝐗U\right)d_{}U,$$
(3.10.11)
where $`\mathrm{\Delta }^𝐧\left(𝐗\right)=\mathrm{\Delta }_1^{n_1n_2}\mathrm{\Delta }_2^{n_2n_3}\mathrm{}\mathrm{\Delta }_N^{n_N}`$ and $`\mathrm{\Delta }_1,\mathrm{}\mathrm{\Delta }_N`$ are main minors of the matrix $`𝐗`$, $`d_{}U`$ is the invariant measure on $`U(N,C)`$, see for the details.
Taking $`N=1`$ we obtain the ordinary hypergeometric function of one variable, which is $`x=x_1`$ now (compare with (3.10.9)). The ordinary hypergeometric function satisfies well-known hypergeometric equation
$$\left(_x{}_{p}{}^{}r_{s}^{}\left(D\right)\right){}_{p}{}^{}F_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;x)=0,D:=x_x.$$
(3.10.12)
This relation helps us to understand the meaning of function $`r`$.
It is known that the series (3.10) diverges if $`p>s+1`$ (untill any of $`a_i+M`$ is nonpositive integer). In case $`p=s+1`$ it converges in certain domain in the vicinity of $`𝐱_{\left(N\right)}=\mathrm{𝟎}`$. For $`p<s+1`$ the series (3.10) converges for all $`𝐱_{\left(N\right)}`$. These known facts (see ) can be also obtained with the help of the determinant representation (3.8.1) and properties of (1.1.1).
Example 4. Hypergeometric function of two sets of variables $`𝐱_{\left(N\right)},𝐲_{\left(N\right)}`$ we put
$${}_{p}{}^{}r_{s}^{}\left(D\right)=\frac{_{i=1}^p\left(a_i+D\right)}{_{i=1}^s\left(b_i+D\right)}\frac{1}{NM+D},$$
(3.10.13)
$$e^{T_n}=\frac{1}{\mathrm{\Gamma }\left(NM+n+1\right)}\frac{_{i=1}^p\mathrm{\Gamma }\left(a_i+n+1\right)}{_{i=1}^s\mathrm{\Gamma }\left(b_i+n+1\right)},$$
(3.10.14)
For the variables $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ and $`𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)`$ we obtain (see Section 3 of for help) the formula (1.3.2)
$`M\left|e^{H\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)}e^{A\left(𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)\right)}\right|M={}_{p}{}^{}_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{\left(N\right)_𝐧}}{\displaystyle \frac{\left(a_1+M\right)_𝐧\mathrm{}\left(a_p+M\right)_𝐧}{\left(b_1+M\right)_𝐧\mathrm{}\left(b_s+M\right)_𝐧}}.`$ (3.10.15)
Example 5 The $`q`$-generalization of the Example 3:
$${}_{p}{}^{}r_{s}^{\left(q\right)}\left(D\right)=\frac{_{i=1}^p\left(1q^{a_i+D}\right)}{_{i=1}^s\left(1q^{b_i+D}\right)}.$$
(3.10.16)
For the variables $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ and
$$y_k=q^{k1},k=1,2,\mathrm{},t_m^{}=\underset{k=1}{\overset{+\mathrm{}}{}}\frac{y_k^m}{m}=\frac{1}{m\left(1q^m\right)},m=1,2,\mathrm{}$$
(3.10.17)
we get Milne’s hypergeometric function (1.2.1):
$`M\left|e^{H\left(𝐭\right)}e^{A\left(𝐭^{}\right)}\right|M={}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{(q^{a_1+M};q)_𝐧\mathrm{}(q^{a_p+M};q)_𝐧}{(q^{b_1+M};q)_𝐧\mathrm{}(q^{b_s+M};q)_𝐧}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}s_𝐧\left(𝐱_{\left(N\right)}\right).`$ (3.10.18)
Example 6. To obtain Milne’s hypergeometric function of two sets of variables $`𝐱_{\left(N\right)},𝐲_{\left(N\right)}`$ we use $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ and $`𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)`$. This choice restricts the sum over partitions $`𝐧`$ with $`l𝐧N`$. We put
$${}_{p}{}^{}r_{s}^{\left(q\right)}\left(n\right)=\frac{_{i=1}^p\left(1q^{a_i+n}\right)}{_{i=1}^s\left(1q^{b_i+n}\right)}\frac{1}{1q^{NM+n}},$$
(3.10.19)
$`e^{T_n}={\displaystyle \frac{1}{\left(1q\right)^n\mathrm{\Gamma }_q\left(n+NM+1\right)}}{\displaystyle \frac{_{i=1}^p\left(1q\right)^n\mathrm{\Gamma }_q\left(a_i+n+1\right)}{_{i=1}^s\left(1q\right)^n\mathrm{\Gamma }_q\left(b_i+n+1\right)}},`$ (3.10.20)
$`\mathrm{\Gamma }_q\left(a\right)=\left(1q\right)^{1a}{\displaystyle \frac{(q;q)_{\mathrm{}}}{(q^a,q)_{\mathrm{}}}},(q^a,q)_n=\left(1q\right)^n{\displaystyle \frac{\mathrm{\Gamma }_q\left(a+n\right)}{\mathrm{\Gamma }_q\left(a\right)}}.`$ (3.10.21)
Here $`\mathrm{\Gamma }_q\left(a\right)`$ is a $`q`$-deformed Gamma-function
$$\mathrm{\Gamma }_q\left(a\right)=\left(1q\right)^{1a}\frac{(q;q)_{\mathrm{}}}{(q^a,q)_{\mathrm{}}},(q^a,q)_n=\left(1q\right)^n\frac{\mathrm{\Gamma }_q\left(a+n\right)}{\mathrm{\Gamma }_q\left(a\right)}.$$
(3.10.22)
We obtain (see Section 3 of for help) the Milne’s formula (1.3.1)
$`\tau _r(M,𝐭^+\left(𝐱_{\left(N\right)}\right),𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right))={}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{s_𝐧(1,q,\mathrm{},q^{N1})}}{\displaystyle \frac{(q^{a_1+M};q)_𝐧\mathrm{}(q^{a_p+M};q)_𝐧}{(q^{b_1+M};q)_𝐧\mathrm{}(q^{b_s+M};q)_𝐧}}.`$ (3.10.23)
This is the KP tau-function (but not the TL one because (3.10.19) depends on TL variable $`M`$).
To receive the basic hypergeometric function of one set of variables we must put indeterminates $`𝐲_{\left(N\right)}`$ in (3.10.18) as $`y_i=q^{i1},i=(1,\mathrm{},N)`$. Thus we have
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{(q^{a_1+M};q)_𝐧\mathrm{}(q^{a_p+M};q)_𝐧}{(q^{b_1+M};q)_𝐧\mathrm{}(q^{b_s+M};q)_𝐧}}{\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}}s_𝐧\left(𝐱_{\left(N\right)}\right).`$ (3.10.24)
And for $`N=1`$ we have the ordinary $`q`$-deformed hypergeometrical function:
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,x\right)={\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(q^{a_1+M};q)_n\mathrm{}(q^{a_p+M};q)_n}{(q^{b_1+M};q)_n\mathrm{}(q^{b_s+M};q)_n}}{\displaystyle \frac{x^n}{(q;q)_n}},x=x_1`$ (3.10.25)
which satisfies the $`q`$-difference equation (compare it with (3.10.12)
$$\left(\frac{1}{x}\left(1q^D\right){}_{p}{}^{}r_{s}^{\left(q\right)}\left(D\right)\right){}_{p}{}^{}\mathrm{\Phi }_{s}^{}(a_1,\mathrm{},a_p;b_1,\mathrm{},b_s;q,x)=0,D:=x_x,$$
(3.10.26)
where $`{}_{p}{}^{}r_{s}^{\left(q\right)}\left(D\right)`$ is defined by (3.10.16).
For the bosonic representation of hypergeometric function (3.10.25) see .
There are various applications for series (3.10.25), for instance see , and . Bosonic representation of (3.10.25) was found in . Let us note that operator $`q^D`$ which acts on fermions $`\psi \left(z\right)`$ was used in in different context.
Example 7
Notations and notions for this Example we borrowed from . Let $`r`$ be a rational function of Jackoby theta-functions $`\theta \left(2x\eta |\tau ^{}\right)`$, where $`\tau ^{}`$ is an elliptic modulus:
$`{}_{p}{}^{}r_{s}^{\left(\eta \right)}\left(n\right)={\displaystyle \frac{_{i=1}^p\theta \left(2\eta \left(a_i+n\right)|\tau ^{}\right)}{\theta \left(2\eta \left(NM+n\right)|\tau ^{}\right)_{i=1}^s\theta \left(2\eta \left(b_i+n\right)|\tau ^{}\right)}},e^{T_{n1}}={\displaystyle \frac{_{i=1}^p\left[a_i\right]_n}{\left[NM\right]_n_{i=1}^s\left[b_i\right]_n}}.`$ (3.10.27)
Here elliptic Pochhammer’s symbol $`\left[a\right]_n`$ is defined in terms of the elliptic number \[a\]
$$\left[a\right]=\theta \left(2a\eta |\tau ^{}\right),\left[a\right]_k=\left[a\right]\left[a+1\right]\left[a+2\right]\mathrm{}\left[a+k1\right].$$
(3.10.28)
One can associate the elliptic Pochhammer’s symbol with a given partition $`𝐧`$:
$$\left[a\right]_𝐧=\left[a\right]_{n_1}\left[a1\right]_{n_2}\mathrm{}\left[al+1\right]_{n_l}.$$
(3.10.29)
For the variables $`𝐭^+\left(𝐱_{\left(N\right)}\right)`$ and $`𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)`$ we can introduce the hypergeometric function
$`M\left|e^{H\left(𝐭^+\left(𝐱_{\left(N\right)}\right)\right)}e^{A\left(𝐭_{}^{}{}_{}{}^{+}\left(𝐲_{\left(N\right)}\right)\right)}\right|M={}_{p}{}^{}F_{s}^{\left(\eta \right)}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|\eta ,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{\left[N\right]_𝐧}}{\displaystyle \frac{\left[a_1+M\right]_𝐧\mathrm{}\left[a_p+M\right]_𝐧}{\left[b_1+M\right]_𝐧\mathrm{}\left[b_s+M\right]_𝐧}}.`$ (3.10.30)
As in the case of (1.2.1) this the KP tau-function which is not the TL tau-function because the factor $`\left[N\right]_𝐧`$ in the denominator does not depend on $`M`$. For $`N=1`$ we get elliptic hypergeometric function of one variable . For instance to obtain the elliptic very-well-poised hypergeometric function
$${}_{p+1}{}^{}W_{p}^{}(\alpha _1;\alpha _4,\alpha _5,\mathrm{},\alpha _{p+1};z|\eta ,\tau ^{})=\underset{n=0}{\overset{\mathrm{}}{}}z^n\frac{\left[\alpha _1+2n\right]\left[\alpha _1\right]_n}{\left[\alpha _1\right]\left[n\right]!}\underset{m=1}{\overset{p2}{}}\frac{\left[\alpha _{m+3}\right]_n}{\left[\alpha _1\alpha _{m+3}+1\right]_n},$$
(3.10.31)
we choose
$$t_n=\frac{z^n}{n},t_n^{}=\frac{1}{n},e^{T_{n1}}=\frac{\left[\alpha _1+2n\right]\left[\alpha _1\right]_n}{\left[\alpha _1\right]\left[n\right]!}\underset{m=1}{\overset{p2}{}}\frac{\left[\alpha _{m+3}\right]_n}{\left[\alpha _1\alpha _{m+3}+1\right]_n}.$$
(3.10.32)
Example 8 The hypergeometric function (1.3.1),(3.10) $`{}_{1}{}^{}\mathrm{\Phi }_{1}^{}\left(\genfrac{}{}{0pt}{}{a}{b}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)`$ can be degenerated to $`{}_{1}{}^{}\mathrm{\Phi }_{0}^{}\left(a|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)`$ by taking $`b+\mathrm{}`$ (remember that $`\left|q\right|<1`$). The limit $`b\mathrm{}`$ (with the rescaling of times $`x_i,y_iq^{\frac{b}{2}}x_i,q^{\frac{b}{2}}x_i`$) is also of interest. Consider this limit and put $`a=NM`$.
Now we get an example of KP tau function (3.1.7) which is not a hypergeometric function. Take $`T_n=\frac{\gamma }{2}\left(n+\frac{1}{2}\right)^2`$, or the same
$$r\left(D\right)=q^D,q=e^\gamma ,$$
(3.10.33)
and rescale the times once more: $`t_k=\alpha ^kp_k,t_k^{}=\alpha ^kp_k^{}`$. We get the series
$`M\left|e^{H\left(𝐭\right)}e^{A\left(𝐭^{}\right)}\right|M={\displaystyle \underset{𝐧}{}}\alpha ^{\left|𝐧\right|}e^{\gamma f_2\left(𝐧\right)}s_𝐧\left(𝐩\right)s_𝐧\left(𝐩^{}\right),`$ (3.10.34)
$`f_2\left(𝐧\right)={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}\left[\left(n_ii+M{\displaystyle \frac{1}{2}}\right)^2\left(i+M{\displaystyle \frac{1}{2}}\right)^2\right],`$ (3.10.35)
which was recently considered in (our notations $`𝐧,\alpha ,\gamma `$ are related to $`\lambda ,q,\beta `$ in respectively). This series is a generating function for double Hurwitz numbers $`Hur_{d,b}(n,m)`$ introduced in as follows. $`Hur_{d,b}(n,m)`$ is a weighted number of connected degree $`d`$ covering of $`P^1`$ with monodromy around $`0,\mathrm{}P^1`$ being $`n`$ and $`m`$, respectively, and $`b`$ additional simple ramifications. The genus of each covering is $`g=\left(b+2l\left(n\right)l\left(m\right)\right)/2`$, where $`l\left(n\right)`$ is the number of parts of $`n`$. The weight of each covering is the reciprocal of the order of its automorphism group. The formula presented in in our terms reads as
$$\mathrm{log}M|e^{H\left(𝐭\right)}e^{A\left(𝐭^{}\right)}M=\underset{d,b,n,m}{}\alpha ^d\gamma ^bp_np_m^{}Hur_{d,b}(n,m)/b!,$$
(3.10.36)
for $`A`$ see (3.1.3),(3.1.1),(3.10.33). Therefore the generating function for the double Hurwitz numbers is expressed in terms of group cocycle of the $`\mathrm{\Psi }DO`$ on the circle (see Appendix “Gauss factorization problem, additional symmetries, string equations and $`\mathrm{\Psi }DO`$ on the circle), which is the correlator under the logarithm.
Example 9. Take $`r`$ be a step function: $`r\left(n\right)=0,n<k`$ and $`r\left(n\right)=1,nk`$, and let $`k<N`$. Then
$$\tau _r(𝐱_{\left(N\right)},𝐲_{\left(N\right)})=\underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)k<N}}{}s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)$$
(3.10.37)
is equal the determinant of a Toeplitz matrix - it is a subject of Gessel’s theorem. We have the determinant representation of (3.10.37) due to the formula (3.8.1).
#### 3.11 Baker-Akhiezer functions and Sato Grassmannian
Let us write down the expression for Baker-Akhiezer functions (2.3.5) in terms of Miwa variables (2.2.15):
$$w_{\mathrm{}}(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭^{},\frac{1}{z})=\frac{\tau _r(M,𝐭^{}\left(𝐱_{\left(N+1\right)}\right),𝐭^{})}{\tau _r(M,𝐭^{}\left(𝐱_{\left(N\right)}\right),𝐭^{})}\underset{i=1}{\overset{N}{}}\left(1\frac{x_i}{z}\right),𝐱_{\left(N+1\right)}=(x_1,\mathrm{},x_N,z),$$
(3.11.1)
$$w_{\mathrm{}}^{}(M,𝐭,𝐭^{},\frac{1}{z})=\frac{\tau _r(M,𝐭+\left[𝐳\right],𝐭^{})}{\tau _r(M,𝐱_{\left(N\right)},𝐭^{})}\underset{i=1}{\overset{N}{}}\frac{1}{\left(1\frac{x_i}{z}\right)}\frac{dz}{z},\left[𝐳\right]=(z,\frac{z^2}{2},\mathrm{}).$$
(3.11.2)
We see that the variables $`x_k,k=1,\mathrm{},N`$ are zeroes of $`w_{\mathrm{}}\left(z\right)`$ and poles of $`w_{\mathrm{}}^{}\left(z\right)`$.
###### Remark 3
The associated linear problems for Baker-Akhiezer functions are read as
$$\left(_{t_1}_{t_1}\varphi _n\right)w(n,𝐭,𝐭^{},z)=w(n+1,𝐭,𝐭^{},z),$$
(3.11.3)
$$_{t_1^{}}w(n,𝐭,𝐭^{},z)=r\left(n\right)e^{\varphi _{n1}\varphi _n}w(n1,𝐭,𝐭^{},z).$$
(3.11.4)
where $`w`$ is either $`w_{\mathrm{}}`$ or $`w_0`$. The compatibility of these equations gives rise to the equation (3.3.8). Taking into account the second eq.(3.3.8), equations (3.11.3), (3.11.4) may be also viewed as the recurrent equations for the tau-functions which depend on different number of variables $`𝐱_{(N)}`$.
Let us write down a plane of Baker-Akhiezer functions (2.3.5), which characterizes Sato Grassmannian related to the tau-function (3.3.1) $`\tau _r(M,𝐭,𝐭^{})`$. We take $`x_k=0,k=1,\mathrm{},N`$ in (3.11.1) and obtain:
$$w_{\mathrm{}}(n,\mathrm{𝟎},𝐭^{},z)=z^n\left(1+\underset{m=1}{\overset{\mathrm{}}{}}r\left(n\right)r\left(n1\right)\mathrm{}r\left(nm+1\right)p_m\left(𝐭^{}\right)z^m\right),n=M,M+1,M+2,\mathrm{}.$$
(3.11.5)
The dual plane is
$$w_{\mathrm{}}^{}(n,\mathrm{𝟎},𝐭^{},z)=z^n\left(1+\underset{m=1}{\overset{\mathrm{}}{}}r\left(n\right)r\left(n+1\right)\mathrm{}r\left(n+m1\right)p_m\left(𝐭^{}\right)z^m\right)dz,n=M,M+1,M+2,\mathrm{}.$$
(3.11.6)
About these formulae see also (A4.22),(A4.23).
We see that when $`r`$ has zeroes, then in the regions (3.5.3) the Grassmannian is the finite-dimensional one. The corresponding tau-function is a particular case of the one found in , .
#### 3.12 Different representations
Let us rewrite hypergeometric series in different way representing all Pochhammer’s coefficients $`(q^a;q)_𝐧`$ and $`\left(a\right)_𝐧`$ through Schur functions. This gives us the opportunity to interchange the role of Pochhammer’s coefficients and Schur functions in (1.3.1),(3.10), and to present different fermionic representations of the hypergeometric functions. We have the relations (see ):
$$\underset{(i,j)𝐧}{}\left(1q^{a+ji}\right)=\frac{s_𝐧\left(𝐭(a,q)\right)}{s_𝐧\left(𝐭(+\mathrm{},q)\right)},\underset{(i,j)𝐧}{}\left(a+ji\right)=\frac{s_𝐧\left(𝐭\left(a\right)\right)}{s_𝐧(𝐭(+\mathrm{})},$$
(3.12.1)
where parameters $`t_m(a,q)`$ and $`t_m\left(a\right)`$ are chosen via generalized Miwa transform with multiplicity $`a`$ (remember that $`\left|q\right|<1`$)
$`t_m(a,q)={\displaystyle \frac{1\left(q^a\right)^m}{m\left(1q^m\right)}},t_m\left(a\right)={\displaystyle \frac{a}{m}},m=1,2,\mathrm{},`$ (3.12.2)
$`s_𝐧\left(𝐭(+\mathrm{},q)\right)=\underset{a+\mathrm{}}{lim}s_𝐧\left(𝐭(a,q)\right)={\displaystyle \frac{q^{n\left(𝐧\right)}}{H_𝐧\left(q\right)}},`$ (3.12.3)
$`s_𝐧\left(𝐭\left(+\mathrm{}\right)\right)=\underset{a+\mathrm{}}{lim}s_𝐧({\displaystyle \frac{t_1\left(a\right)}{a}},{\displaystyle \frac{t_2\left(a\right)}{a^2}},\mathrm{})=\underset{a+\mathrm{}}{lim}{\displaystyle \frac{1}{a^{\left|𝐧\right|}}}s_𝐧\left(𝐭\left(a\right)\right)={\displaystyle \frac{1}{H_𝐧}}.`$ (3.12.4)
Now we rewrite the series (3.10.18) and (3.10) only in terms of Schur functions:
$`{}_{p}{}^{}\mathrm{\Phi }_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|q,𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=\tau _r(M,𝐭(+\mathrm{},q),𝐭^{})`$
$`={\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{_{k=1}^ps_𝐧\left(𝐭(a_k+M,q)\right)}{_{k=1}^ss_𝐧\left(𝐭(b_k+M,q)\right)}}\left(s_𝐧\left(𝐭(+\mathrm{},q)\right)\right)^{sp+1}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{s_𝐧\left(𝐭(N,q)\right)}},`$ (3.12.5)
$`{}_{p}{}^{}_{s}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+M,\mathrm{},a_p+M}{b_1+M,\mathrm{},b_s+M}}|𝐱_{\left(N\right)},𝐲_{\left(N\right)}\right)=\tau _r(M,𝐭\left(+\mathrm{}\right),𝐭)=`$
$`{\displaystyle \underset{\genfrac{}{}{0pt}{}{𝐧}{l\left(𝐧\right)N}}{}}{\displaystyle \frac{_{k=1}^ps_𝐧\left(𝐭\left(a_k+M\right)\right)}{_{k=1}^ss_𝐧\left(𝐭\left(b_k+M\right)\right)}}\left(s_𝐧\left(𝐭\left(+\mathrm{}\right)\right)\right)^{sp+1}{\displaystyle \frac{s_𝐧\left(𝐱_{\left(N\right)}\right)s_𝐧\left(𝐲_{\left(N\right)}\right)}{s_𝐧\left(𝐭\left(N\right)\right)}}.`$ (3.12.6)
A nice feature of this formulae is that they do not contain number coefficients at all, it is a sum of ratios of Schur functions only.
We obtain different fermionic representations of hypergeometric functions (3.12), (3.12), and they are parametrized by a complex noninteger number $`b`$:
###### Proposition 10
For $`bC`$ and for $`r={}_{p}{}^{}r_{s}^{}`$ (see (3.10.3)) we have
$$\tau _r(M,𝐭\left(+\mathrm{}\right),𝐭^{})=\tau _{r_b}(M,𝐭\left(b+M\right),𝐭^{}),r_b=\frac{r}{b+D}.$$
(3.12.7)
For $`r={}_{p}{}^{}r_{s}^{(q)}`$ (see (3.10.16)) we have
$$\tau _r(M,𝐭(+\mathrm{},q),𝐭^{})=\tau _{r_b}(M,𝐭(b+M,q),𝐭^{}),r_b=\frac{r}{1q^{b+D}}.$$
(3.12.8)
###### Remark 4
There are two ways to restrict the sum (3.1.7) to the sum over partitions of length $`l(𝐧)N`$. First, if we use Miwa’s change (2.2.14), then $`s_𝐧(𝐱_{(N)})=0`$, for $`𝐧`$ with length $`l(𝐧)>N`$. The second way is to restrict the Pochhammer’s coefficients: if we put $`a_i=N`$ for one $`i`$ from (3.10.16) equal to $`N`$, then the coefficient $`(q^{a_i},q)_𝐧`$ vanishes for $`l(𝐧)>N`$. Since we expressed Pochhammer’s coefficients in terms of Schur functions in (3.12.1) both ways have the same explanation. Indeed
$$t_m(N,q)=\frac{1}{m}\frac{1\left(q^N\right)^m}{1q^m}=\frac{1}{m}\left(1+\left(q\right)^m+\left(q^2\right)^m+\mathrm{}+\left(q^{N1}\right)^m\right).$$
(3.12.9)
Therefore we obtain for Miwa’s change: $`x_1=1,x_2=q,\mathrm{},x_N=q^{N1}`$ and
$$s_𝐧\left(𝐭(N,q)\right)=s_𝐧(1,q,\mathrm{},q^{N1})=0,l\left(𝐧\right)>N.$$
(3.12.10)
The same we have for the sum over partitions $`𝐧`$ such that $`l(𝐧^{})<K`$. Again the first way has to be realized through the following Miwa’s change of variables:
$$t_m=\underset{i=1}{\overset{K}{}}\frac{x_i^m}{m},s_𝐧\left(𝐭\right)=s_𝐧^{}\left(𝐱_{\left(K\right)}\right).$$
(3.12.11)
The second way is to make one of the parameters, for example $`a_j`$ from (3.10.16) equal to $`(K)`$. In this case
$$s_𝐧\left(𝐭(K,q)\right)=s_𝐧^{}(\frac{1}{q},\frac{1}{q^2},\mathrm{},\frac{1}{q^K})=0,l\left(𝐧^{}\right)>K.$$
(3.12.12)
### 4 Further generalization. Examples of Gelfand-Graev hypergeometric functions
#### 4.1 Generalization
Formula (3.3.7) is related to ’Gauss decomposition’ of operators inside vacuums $`0\left|\mathrm{}\right|0`$ into diagonal operator $`e^{H_0\left(𝐓\right)}`$ and upper triangular operator $`e^{H\left(𝐭\right)}`$ and lower triangular operator $`e^{H^{}\left(𝐭^{}\right)}`$ the last two have the Toeplitz form. Now let us consider more general two-dimensional Toda chain tau-function
$$\tau =M\left|e^{H\left(𝐭\right)}ge^{A\left(𝐭^{}\right)}\right|M,$$
(4.1.1)
where we decompose $`g`$ in the following way:
$$g(\stackrel{~}{\gamma },\gamma )=e^{\stackrel{~}{A}_1\left(\stackrel{~}{\gamma }_1\right)}\mathrm{}e^{\stackrel{~}{A}_k\left(\stackrel{~}{\gamma }_k\right)}e^{A_l\left(\gamma _l\right)}\mathrm{}e^{A_1\left(\gamma _1\right)},$$
(4.1.2)
where each of $`\stackrel{~}{\gamma _i},\gamma _i,\stackrel{~}{A}_i,A_i`$ has an additional index: $`\stackrel{~}{\gamma }_{in},\gamma _{in},\stackrel{~}{A}_{in},A_{in}`$, ($`n=1,2,\mathrm{}`$). Here each of $`A_i\left(\gamma _i\right)`$ has a form as in (3.1.1),(3.1.3) and corresponds to operator $`r^i\left(D\right)`$, while each of $`\stackrel{~}{A}_i\left(\stackrel{~}{\gamma }_i\right)`$ has a form of (3.1.11) and corresponds to operator $`\stackrel{~}{r}^i\left(D\right)`$. Collections of variables $`\stackrel{~}{\gamma }=\left\{\stackrel{~}{\gamma }_{in}\right\},\gamma =\left\{\gamma _{in}\right\}`$ play the role of coordinates for some wide enough class of Clifford group elements $`g`$. This tau-function is related to rather involved generalization of the hypergeometric functions we considered above. Tau-function (4.1.1),(4.1.2) may be considered as the result of applying of the additional symmetries to the vacuum tau function, which is 1, see Appendix “The vertex operator action”.
Let us calculate this tau-function. First of all we introduce a set consisting of $`m+1`$ partitions:
$$\left(𝐧_\mathrm{𝟏},\mathrm{},𝐧_𝐦,𝐧_{𝐦+\mathrm{𝟏}}=𝐧\right),0𝐧_\mathrm{𝟏}𝐧_\mathrm{𝟐}\mathrm{}𝐧_𝐦𝐧_{𝐦+\mathrm{𝟏}}=𝐧,$$
(4.1.3)
see for the notation $``$ for the partitions. The corresponding set
$$\mathrm{\Theta }_𝐧^m=(𝐧_\mathrm{𝟏},\theta _1,\mathrm{},\theta _m),\theta _i=𝐧_{𝐢+\mathrm{𝟏}}𝐧_𝐢,i=1,\mathrm{},m$$
(4.1.4)
depends on the partition $`𝐧`$ and the number $`m+1`$ of the partitions. We take as $`s_\mathrm{\Theta }(𝐭^{},\gamma )`$ the product which is relevant to the set $`\mathrm{\Theta }_𝐧^m`$ and depending on the set of variables $`\mu _i=\left\{\mu _{ij}\right\}`$ ($`i=(1,\mathrm{},m+1)`$, $`j=(1,2\mathrm{})`$)
$`s_{\mathrm{\Theta }_𝐧^m}\left(\mu \right)=s_{𝐧_\mathrm{𝟏}}\left(\mu _1\right)s_{\theta _1}\left(\mu _2\right)\mathrm{}s_{\theta _m}\left(\mu _{m+1}\right).`$ (4.1.5)
Here $`s_{\theta _i}`$ is a skew Schur function (see ). Further we define function $`r_{\mathrm{\Theta }_𝐧^m}\left(M\right)`$:
$$r_{\mathrm{\Theta }_𝐧^m}\left(M\right)=r_{𝐧_\mathrm{𝟏}}\left(M\right)r_{\theta _1}^1\left(M\right)\mathrm{}r_{\theta _m}^m\left(M\right),$$
(4.1.6)
where the function $`r_{\theta _i}^i\left(M\right)`$ , a skew analogy of $`r_𝐧\left(M\right)`$ from (3.1.4), is
$$r_{\theta _i}\left(M\right)=\underset{j=1}{\overset{s}{}}r\left(n_j^{\left(i\right)}j+1+M\right)\mathrm{}r\left(n_j^{\left(i+1\right)}j+M\right),$$
(4.1.7)
where $`𝐧_{𝐢+\mathrm{𝟏}}=(n_1^{\left(i+1\right)},\mathrm{},n_s^{\left(i+1\right)})`$. If the function $`r^i\left(m\right)`$ has no poles and zeroes at integer points then the relation
$$r_{\theta _i}^i\left(M\right)=\frac{r_{𝐧_{𝐢+\mathrm{𝟏}}}^i\left(M\right)}{r_{𝐧_𝐢}^i\left(M\right)},i=1,\mathrm{},m$$
(4.1.8)
is correct. To calculate the tau function we need the Lemma
Lemma 3 Let partitions $`𝐧=(i_1,\mathrm{},i_s|j_11,\mathrm{},j_s1)`$ and $`\stackrel{~}{𝐧}=(\stackrel{~}{i}_1,\mathrm{},\stackrel{~}{i}_r|\stackrel{~}{j}_11,\mathrm{},\stackrel{~}{j}_r1)`$ satisfy the relation $`𝐧\stackrel{~}{𝐧}`$. The following is valid:
$`0\left|\psi _{\stackrel{~}{i}_1}^{}\mathrm{}\psi _{\stackrel{~}{i}_r}^{}\psi _{\stackrel{~}{j}_r}\mathrm{}\psi _{\stackrel{~}{j}_1}e^{A^i\left(\gamma _i\right)}\psi _{j_1}^{}\mathrm{}\psi _{j_s}^{}\psi _{i_s}\mathrm{}\psi _{i_1}\right|0=`$
$`=\left(1\right)^{\stackrel{~}{j}_1+\mathrm{}+\stackrel{~}{j}_r+j_1+\mathrm{}+j_s}s_\theta \left(\gamma _i\right)r_\theta \left(0\right),\theta =𝐧\stackrel{~}{𝐧}.`$ (4.1.9)
Proof: the proof is achieved by direct calculation (see Example 22 in Sec 5 of for help).
Then we obtain the generalization of Proposition 1:
###### Proposition 11
$$\tau _M(𝐭,𝐭^{};\gamma ,\stackrel{~}{\gamma })=\underset{𝐧}{}\underset{\mathrm{\Theta }_𝐧^k}{}\underset{\mathrm{\Theta }_𝐧^l}{}\stackrel{~}{r}_{\mathrm{\Theta }_𝐧^k}\left(M\right)r_{\mathrm{\Theta }_𝐧^l}\left(M\right)s_{\mathrm{\Theta }_𝐧^k}(𝐭,\stackrel{~}{\gamma })s_{\mathrm{\Theta }_𝐧^l}(𝐭^{},\gamma ),$$
(4.1.10)
where $`\stackrel{~}{r}_{\mathrm{\Theta }_𝐧^k}(M)`$ and $`r_{\mathrm{\Theta }_𝐧^l}(M)`$ are given by (4.1.7).
With the help of this series one can obtain different hypergeometric functions.
#### 4.2 The example of Gelfand,Graev and Retakh hypergeometric series
Let us consider the tau function:
$$\tau (M,\stackrel{~}{\beta },\beta ;\gamma )=M\left|e^{\stackrel{~}{A}\left(\stackrel{~}{\beta }\right)}e^{A_l\left(\gamma _l\right)}\mathrm{}e^{A_1\left(\gamma _1\right)}e^{A\left(\beta \right)}\right|M.$$
(4.2.1)
We put
$$\stackrel{~}{\beta }=(x,\frac{x^2}{2},\frac{x^3}{3},\mathrm{}),\beta =(y_1,0,0,\mathrm{}),\gamma _i=(y_{i+1},0,0,\mathrm{})i=(1,\mathrm{},l).$$
(4.2.2)
We obtain the series
$`\tau (M,x,y_1,\mathrm{},y_{l+1})={\displaystyle \underset{n_1,\mathrm{},n_{l+1}=0}{\overset{+\mathrm{}}{}}}\stackrel{~}{r}_{\left(n_1+\mathrm{}+n_{l+1}\right)}\left(M\right)r_{\mathrm{\Theta }_𝐧^l}\left(M\right){\displaystyle \frac{\left(xy_1\right)^{n_1}\mathrm{}\left(xy_{l+1}\right)^{n_{l+1}}}{n_1!\mathrm{}n_{l+1}!}}=`$ (4.2.3)
$`{\displaystyle \underset{n_1,\mathrm{},n_{l+1}Z}{}}c(n_1,\mathrm{},n_{l+1})\left(xy_1\right)^{n_1}\mathrm{}\left(xy_{l+1}\right)^{n_{l+1}},c(n_1,\mathrm{},n_{l+1})={\displaystyle \frac{\stackrel{~}{r}_{\left(n_1+\mathrm{}+n_{l+1}\right)}\left(M\right)r_{\mathrm{\Theta }_𝐧^l}\left(M\right)}{\mathrm{\Gamma }\left(n_1+1\right)\mathrm{}\mathrm{\Gamma }\left(n_{l+1}+1\right)}},`$ (4.2.4)
where $`\mathrm{\Theta }_𝐧^l`$ corresponds to the set of simple partitions-rows
$$𝐧_\mathrm{𝟏}=\left(n_1\right),𝐧_\mathrm{𝟐}=\left(n_1+n_2\right),\mathrm{},𝐧_{l+1}=\left(n_1+\mathrm{}+n_{l+1}\right)$$
(4.2.5)
When functions $`b_i(n_1,\mathrm{},n_{l+1})`$ defined as
$$b_i(n_1,\mathrm{},n_{l+1})=\frac{c(n_1,\mathrm{},n_i+1,\mathrm{},n_{l+1})}{c(n_1,\mathrm{},n_{l+1})},i=1,\mathrm{},l+1$$
(4.2.6)
are rational functions of $`(n_1,\mathrm{},n_{l+1})`$, then tau function (4.2.3) is a Horn hypergeometric series .
Above series for the special choice of functions $`r^i\left(D\right)`$ can be deduced from the Gelfand, Graev and Retakh series defined on the special lattice and corresponding to the special set of parameters. Let us take the rational functions $`r^i\left(D\right)`$:
$`r^i\left(D\right)={\displaystyle \frac{_{j=1}^{p^{\left(i\right)}}\left(D+a_j^{\left(i\right)}\right)}{_{m=1}^{s^{\left(i\right)}}\left(D+b_m^{\left(i\right)}\right)}},\left(i=0,\mathrm{},l\right),r^0\left(D\right)=r\left(D\right)`$ (4.2.7)
$`\stackrel{~}{r}\left(D\right)={\displaystyle \frac{_{j=1}^{p^{\left(l+1\right)}}\left(D+a_j^{\left(l+1\right)}\right)}{_{m=1}^{s^{\left(l+1\right)}}\left(D+b_m^{\left(l+1\right)}\right)}}`$ (4.2.8)
Let define $`N=p^{\left(0\right)}+s^{\left(0\right)}+2_{j=1}^l\left(p^{\left(j\right)}+2s^{\left(j\right)}\right)+p^{\left(l+1\right)}+s^{\left(l+1\right)}+l+1`$ and consider complex space $`C^N`$. In this space we consider the $`l+1`$-dimensional basis $`B`$ and the vector $`\upsilon `$ consisting of parameters.
$`p_0=s_0=0,p_i=p^{\left(i1\right)}+s^{\left(i\right)},s_i=s^{\left(i1\right)}+p^{\left(i\right)},i=(1,\mathrm{},l)`$
$`p_{l+1}=p^{\left(l\right)}+p^{\left(l+1\right)},s_{l+1}=s^{\left(l\right)}+s^{\left(l+1\right)},N={\displaystyle \underset{j=1}{\overset{l+1}{}}}\left(p_j+s_j\right)+l+1`$ (4.2.9)
$`𝐟^i=\left(𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+1}+\mathrm{}+𝐞_{p_1+s_1+\mathrm{}+p_{i1}+s_{i1}+p_i}\right)+`$
$`+\left(𝐞_{p_1+s_1+\mathrm{}+p_{i1}+s_{i1}+p_i+1}+\mathrm{}+𝐞_{p_1+s_1+\mathrm{}+p_{i1}+s_{i1}+p_i+s_i}\right),i=1,\mathrm{},l+1`$ (4.2.10)
where $`𝐞_i=\underset{\text{N}}{\underset{}{(0,\mathrm{},0,\stackrel{i}{\widehat{1}},0,\mathrm{})}}`$. The lattice $`BC^N`$ is generated by the vector basis of dimension $`l+1`$:
$$𝐛^i=𝐟^i+\mathrm{}+𝐟^{l+1}+𝐞_{Nl1+i},i=1,\mathrm{},l+1$$
(4.2.11)
Vector $`\upsilon C^N`$ is defined as follows (compare with (4.2)):
$`\upsilon ^i=(a_1^{\left(i1\right)}𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+1}+\mathrm{}+a_{p^{\left(i1\right)}}^{\left(i1\right)}𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p^{\left(i1\right)}}+`$
$`+b_1^{\left(i\right)}𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p^{\left(i1\right)}+1}+\mathrm{}+b_{s^{\left(i\right)}}^{\left(i\right)}𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i})+`$
$`+((b_1^{\left(i1\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+1}+\mathrm{}+(b_{s^{\left(i1\right)}}^{\left(i1\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s^{\left(i1\right)}}+`$
$`+(a_1^{\left(i\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s^{\left(i1\right)}+1}+\mathrm{}+(a_{s^{\left(i\right)}}^{\left(i\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s_i})`$ (4.2.12)
for $`i=(1,\mathrm{}l)`$, and
$`\upsilon ^{l+1}=(a_1^{\left(l\right)}𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+1}+\mathrm{}+a_{p^{\left(l\right)}}^{\left(l\right)}𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p^{\left(l\right)}}+`$
$`+a_1^{\left(l+1\right)}𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p^{\left(l\right)}+1}+\mathrm{}+a_{s^{\left(l+1\right)}}^{\left(l+1\right)}𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p_{l+1}})+`$
$`+((b_1^{\left(l\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p_{l+1}+1}+\mathrm{}+(b_{s^{\left(l\right)}}^{\left(l\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p_{l+1}+s^{\left(l\right)}}+`$
$`+(b_1^{\left(l+1\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p_{l+1}+s^{\left(l\right)}+1}+\mathrm{}+(b_{s^{\left(l+1\right)}}^{\left(l+1\right)}1)𝐞_{p_0+s_0+\mathrm{}+p_l+s_l+p_{l+1}+s_{l+1}})`$ (4.2.13)
Vector $`\upsilon `$ is:
$$\upsilon =\upsilon ^1+\mathrm{}+\upsilon ^{l+1}$$
(4.2.14)
Now we can write down Gelfand, Graev and Retakh hypergeometric series corresponding to the lattice $`B`$ and vector $`\upsilon `$:
$`F_B(\upsilon ;z)={\displaystyle \underset{𝐛B}{}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{z_j^{\upsilon _j+b_j}}{\mathrm{\Gamma }\left(\upsilon _j+b_j+1\right)}}+`$
$`{\displaystyle \underset{n_1,\mathrm{},n_{l+1}Z}{}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{z_j^{\upsilon _j+n_1b_j^1+\mathrm{}+n_{l+1}b_j^{l+1}}}{\mathrm{\Gamma }\left(\upsilon _j+n_1b_j^1+\mathrm{}+n_{l+1}b_j^{l+1}+1\right)}}`$ (4.2.15)
Let us compare this series with tau function (4.2.3):
$$F_B(\upsilon ;𝐳)=c_1(a,b)g_1\left(𝐳\right)\mathrm{}c_{l+1}(a,b)g_{l+1}\left(𝐳\right)\tau (M,x,y_1,\mathrm{},y_{l+1})$$
(4.2.16)
where
$`c_i^1(a,b)=\mathrm{\Gamma }(1a_1^{\left(i1\right)})\mathrm{}\mathrm{\Gamma }(1a_{p^{\left(i1\right)}}^{\left(i1\right)})\mathrm{\Gamma }(1b_1^{\left(i\right)})\mathrm{}\mathrm{\Gamma }(1b_{s^{\left(i\right)}}^{\left(i\right)})\times `$
$`\times \mathrm{\Gamma }\left(b_1^{\left(i1\right)}\right)\mathrm{}\mathrm{\Gamma }\left(b_{s^{\left(i1\right)}}^{\left(i1\right)}\right)\mathrm{\Gamma }\left(a_1^{\left(i\right)}\right)\mathrm{}\mathrm{\Gamma }\left(a_{s^{\left(i\right)}}^{\left(i\right)}\right),i=1,\mathrm{},l`$ (4.2.17)
$`c_{l+1}^1(a,b)=\mathrm{\Gamma }(1a_1^{\left(l\right)})\mathrm{}\mathrm{\Gamma }(1a_{p^{\left(l\right)}}^{\left(l\right)})\mathrm{\Gamma }(1a_1^{\left(l+1\right)})\mathrm{}\mathrm{\Gamma }(1a_{s^{\left(l+1\right)}}^{\left(l+1\right)})\times `$
$`\times \mathrm{\Gamma }\left(b_1^{\left(l\right)}\right)\mathrm{}\mathrm{\Gamma }\left(b_{s^{\left(l\right)}}^{\left(l\right)}\right)\mathrm{\Gamma }\left(b_1^{\left(l+1\right)}\right)\mathrm{}\mathrm{\Gamma }\left(b_{s^{\left(l+1\right)}}^{\left(l+1\right)}\right)`$ (4.2.18)
$`{\displaystyle \frac{z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+1}\mathrm{}z_{p_1+s_1+\mathrm{}+p_{i1}+s_{i1}+p_i+s_i}}{\left(z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+1}\right)\mathrm{}\left(z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i}\right)}}=1,i=2,\mathrm{},l`$ (4.2.19)
$`{\displaystyle \frac{z_{p_1+1}\mathrm{}z_{p_1+s_1}z_{Nl}}{\left(z_1\right)\mathrm{}\left(z_{p_1}\right)}}=y_1`$ (4.2.20)
$$y_i=z_{Nl1+i},i=2,\mathrm{},l+1$$
(4.2.21)
$`{\displaystyle \frac{z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+1}\mathrm{}z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+s_{l+1}}}{\left(z_{p_1+s_1+\mathrm{}+p_l+s_l+1}\right)\mathrm{}\left(z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}}\right)}}=x`$ (4.2.22)
$`g_i\left(𝐳\right)=z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+1}^{\left(a_1^{\left(i1\right)}\right)}\mathrm{}z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p^{\left(i1\right)}}^{\left(a_{p^{\left(i1\right)}}^{\left(i1\right)}\right)}\times `$
$`\times z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p^{\left(i1\right)}+1}^{\left(b_1^{\left(i\right)}\right)}\mathrm{}z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i}^{\left(b_{s^{\left(i\right)}}^{\left(i\right)}\right)}\times `$
$`\times z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+1}^{\left(b_1^{\left(i1\right)}1\right)}\mathrm{}z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s^{\left(i1\right)}}^{\left(b_{s^{\left(i1\right)}}^{\left(i1\right)}1\right)}\times `$
$`\times z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s^{\left(i1\right)}+1}^{\left(a_1^{\left(i\right)}1\right)}\mathrm{}z_{p_0+s_0+\mathrm{}+p_{i1}+s_{i1}+p_i+s_i1}^{\left(a_{s^{\left(i\right)}}^{\left(i\right)}1\right)}`$ (4.2.23)
for $`i=(1,\mathrm{}l)`$, and
$`g_{l+1}\left(𝐳\right)=z_{p_1+s_1+\mathrm{}+p_l+s_l+1}^{\left(a_1^{\left(l\right)}\right)}\mathrm{}z_{p_1+s_1+\mathrm{}+p_l+s_l+p^{\left(l\right)}}^{\left(a_{p^{\left(l\right)}}^{\left(l\right)}\right)}\times `$
$`\times z_{p_1+s_1+\mathrm{}+p_l+s_l+p^{\left(l\right)}+1}^{\left(a_1^{\left(l+1\right)}\right)}\mathrm{}z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}}^{\left(a_{s^{\left(l+1\right)}}^{\left(l+1\right)}\right)}\times `$
$`\times z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+1}^{\left(b_1^{\left(l\right)}1\right)}\mathrm{}z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+s^{\left(l\right)}}^{\left(b_{s^{\left(l\right)}}^{\left(l\right)}1\right)}\times `$
$`\times z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+s^{\left(l\right)}+1}^{\left(b_1^{\left(l+1\right)}1\right)}\mathrm{}z_{p_1+s_1+\mathrm{}+p_l+s_l+p_{l+1}+s_{l+1}1}^{\left(b_{s^{\left(l+1\right)}}^{\left(l+1\right)}1\right)}`$ (4.2.24)
### Conclusion
We get multivariable hypergeometric functions as certain tau-functions of the KP hierarchy. It means that we have a set of new relations on the multivariable hypergeometric functions. For instance all hypergeometric functions of the form (3.1.6) or of the form (4.1.10) satisfy bilinear Hirota equations of the KP hierarchy. Hypergeometric functions may be also considered as ratios of the tau-functions of the two-dimensional Toda lattice evaluated at special values of Toda lattice times. To get hypergeometric functions of a single set of arguments (1.2.6),(1.2.1) one should take $`𝐭^{}`$ as in (3.10.6), (3.10.17). To get hypergeometric functions of a double set of arguments (1.3.1), (1.3.2) one should keep the descrete time $`M`$ constant. One can get the fermionic representations for different special functions and polynomials related to these hypergeometric functions. Using integral representation one can express hypergeometric functions as the integral of rather simple hypergeometric function. We also get determinant representation of (3.1.6), which may allow to analize analytical properties of multivariable tau-functions in terms of functions of only one variable. We wrote down the system of linear equations on tau-function (3.1.6), which may allow to find applications to quantum mechanical problems. It is quite unexpected that we obtain $`q`$-deformed version of these hypergeometric functions as tau-functions not of a $`q`$-deformed KP hierarchy ,,,, but of the usual KP hierarchy. It is now an interesting problem to establish links between these results and group-theoretic approach to the $`q`$-special functions and matrix integrals. Let us note a certain similarity of some of our formulas and formulas from . We expect to work out connections with matrix models of Kontsevich type and two-matrix models related to 2D Toda lattice . We consider (in ) the multicomponent KP as a basis for obtaining new examples of hypergeometric functions.
Let us note the paper which turns to be different example of the tau function (3.3.1). This tau-function is not of hypergeometric type, but is closely related to it. This tau-function has an interesting meaning in the algebraic geometry. In the interpretation of the Schur functions as a measure on partitions is explained, thus we have an additional interpretation of hypergeometric functions as generating functions for certain probabilities described in .
In the present version of the paper we add references to the , where different examples of special functions of one and of two variables were considered as solutions of Hirota equations with variable coefficients. It will be interesting to get the fermionic representations for these examples.
## Appendix Appendix 1 Formulae involving $`r(D)`$
Let $`r\left(n\right)=e^{T_{n1}T_n}`$, $`h\left(n\right)=e^{T_n}`$. Then the operators $`r\left(D\right),h\left(D\right)`$, where $`D=z\frac{d}{dz}`$, have the simple properties
$$F_r(M,zz^{}):=\left(1+\underset{n=1}{\overset{\mathrm{}}{}}e^{T_{Mn}T_M}\frac{\left(zz^{}\right)^n}{n!}\right)\left(zz^{}\right)^M=$$
(A1.1)
$$=\left(1\frac{1}{zz^{}}\frac{r\left(D\right)}{D+a}\right)^a\left(zz^{}\right)^M=\left(1\frac{1}{zz^{}}\frac{r\left(D\right)}{D+b}\right)^b\left(zz^{}\right)^M,a,bM,$$
(A1.2)
$$\xi _r^{\left(\mathrm{}\right)}(\beta ,z,D):=\underset{n=1}{\overset{\mathrm{}}{}}\beta _n\left(\frac{1}{z}r\left(D\right)\right)^n,\xi _r^{\left(0\right)}(𝐭,z,D):=\underset{n=1}{\overset{\mathrm{}}{}}t_n\left(r\left(D\right)z\right)^n.$$
(A1.3)
Let us consider
$`f_r(M,\beta ,z):=\left(zz^{}\right)^M\left(1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}e^{T_{Mn}T_M}z^np_n\left(\beta \right)\right)=e^{\xi _r^{\left(\mathrm{}\right)}(\beta ,z,D)}\left(zz^{}\right)^M.`$ (A1.4)
Let $`h\left(n\right)=e^{T_n}`$. Then
$`{\displaystyle \frac{h\left(D\right)}{h\left(0\right)}}f_r(M,\beta ,z)=\left(zz^{}\right)^Me^{\xi (\beta ,z^1)},\xi (\beta ,z^1):={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\beta _mz^m.`$ (A1.5)
Remark. Let us put $`m\beta _m=a\left(z^{}\right)^m,m=1,2,\mathrm{}`$ with some $`a0`$, then $`f_r(M,\beta ,z)=F_r^{}(M,zz^{})`$, where $`r^{}\left(D\right)=\left(D+a\right)r\left(D\right)`$, since the elementary Schur polynomials (2.2.12) $`p_n\left(\beta \right)=\frac{\left(a\right)_n}{n!}\left(z^{}\right)^n`$. Eq.(3.12.7) generalizes this property for all partitions.
## Appendix Appendix 2 Orthogonal polynomials and matrix integrals
It is known that the hypergeometric functions (1.3.2) appear in the group representation theory and are connected with the so-called matrix integrals . On the other hand the set of examples reveals a connection between the matrix integrals and the soliton theory. To establish this connection in our case, it is useful to consider the related systems of the orthogonal polynomials. Let us briefly describe how to write down these polynomials.
Let $`M_+`$ be the largest integer zero of $`r`$. Then the function
$$f_r^+\left(zz^{}\right)=\underset{n=0}{\overset{+\mathrm{}}{}}\left(zz^{}\right)^{n+M_+}e^{T_{n+M_+}T_{M_+}}$$
(A2.1)
is the eigenfunction of the operator $`\frac{1}{z}r\left(D\right)`$ with the eigenvalue $`z^{}`$. Since operator $`r\left(D\right)`$ is invertible on the functions $`\left\{z^M,M>M_+\right\}`$ we write
$$f_r^+\left(zz^{}\right)=\left(1\frac{1}{r\left(D\right)}zz^{}\right)^1\left(zz^{}\right)^{M_+}.$$
(A2.2)
For example if we take $`r\left(n\right)=n`$ we obtain
$$f_r^+\left(zz^{}\right)=e^{zz^{}}.$$
(A2.3)
We use this function as weight function for a system of orthogonal polynomials $`\left\{\pi _n^\pm ,n=0,1,2,\mathrm{}\right\}`$, related to the hypergeometric solution of KP:
$$_\gamma \pi _n^{}(𝐭,\beta ,z)e^{\xi (𝐭,z)}f_r^+\left(zz^{}\right)e^{\xi (\beta ,z^{})}\pi _m^+(𝐭,\beta ,z^{})𝑑z𝑑z^{}=e^{\varphi _{M_++n}(𝐭,\beta )}\delta _{n,m}.$$
(A2.4)
The corresponding two matrix integral is the following one
$$\tau (M,𝐭,\beta )=e^{Tr\xi (𝐭,Z)}f_r^+\left(Tr\left(ZZ^{}\right)\right)e^{Tr\xi (\beta ,Z^{})}𝑑Z𝑑Z^{}.$$
(A2.5)
Here $`Z,Z^{}`$ are Hermitian $`M\times M`$ matrices.
## Appendix Appendix 3 The vertex operator action
Now we present relations between hypergeometric functions which follow from the soliton theory, for instance see ,. Let us introduce the operators which act on functions of $`𝐭`$ variables:
$$\mathrm{\Omega }_r^{\left(\mathrm{}\right)}\left(𝐭^{}\right):=\frac{1}{2\pi i}\underset{ϵ0}{lim}V_{\mathrm{}}^{}\left(z+ϵ\right)\xi _r^{\left(\mathrm{}\right)}(𝐭^{},z,D)V_{\mathrm{}}\left(z\right)𝑑z,$$
(A3.1)
$$\mathrm{\Omega }_r^{\left(0\right)}\left(𝐭\right):=\frac{1}{2\pi i}\underset{ϵ0}{lim}V_0\left(z+ϵ\right)\xi _r^{\left(0\right)}(𝐭,z,D)V_0^{}\left(z\right)z^2𝑑z,$$
(A3.2)
where $`V_{\mathrm{}}\left(z\right),V_{\mathrm{}}^{}\left(z\right),V_0\left(z\right),V_0^{}\left(z\right)`$ are defined by (2.3.1). For instance
$$\mathrm{\Omega }_r^{\left(\mathrm{}\right)}\left(𝐭^{}\right)=\mathrm{\Omega }_r^{\left(0\right)}\left(𝐭\right)=\underset{n>0}{}nt_nt_n^{},r=1.$$
(A3.3)
The bosonization formulae , give the relations which connect hypergeometric functions (3.1.7) and hypergeometric functions (4.1.10):
$$e^{\mathrm{\Omega }_{\stackrel{~}{r}^1}^{\left(\mathrm{}\right)}\left(\stackrel{~}{\gamma }_1\right)}\mathrm{}e^{\mathrm{\Omega }_{\stackrel{~}{r}^k}^{\left(\mathrm{}\right)}\left(\stackrel{~}{\gamma }_k\right)}\mathrm{}e^{\mathrm{\Omega }_{r^l}^{\left(\mathrm{}\right)}\left(\gamma _l\right)}\mathrm{}e^{\mathrm{\Omega }_{r^l}^{\left(\mathrm{}\right)}\left(\gamma _l\right)}\tau _r(M,𝐭,𝐭^{})=\tau (M,𝐭,𝐭^{};\stackrel{~}{\gamma },\gamma ),$$
(A3.4)
$$e^{\mathrm{\Omega }_{r^l}^{\left(0\right)}\left(\gamma _l\right)}\mathrm{}e^{\mathrm{\Omega }_{r^l}^{\left(0\right)}\left(\gamma _l\right)}\mathrm{}e^{\mathrm{\Omega }_{\stackrel{~}{r}^k}^{\left(0\right)}\left(\stackrel{~}{\gamma }_k\right)}\mathrm{}e^{\mathrm{\Omega }_{\stackrel{~}{r}^1}^{\left(0\right)}\left(\stackrel{~}{\gamma }_1\right)}\tau _r(M,𝐭,𝐭^{})=\tau (M,𝐭,𝐭^{};\stackrel{~}{\gamma },\gamma ).$$
(A3.5)
In particular we have shift argument formulae for the tau function (3.1.7):
$$e^{\mathrm{\Omega }_r^{\left(\mathrm{}\right)}\left(\gamma \right)}\tau _r(M,𝐭,𝐭^{})=\tau _r(M,𝐭,𝐭^{}+\gamma ),e^{\mathrm{\Omega }_r^{\left(0\right)}\left(\gamma \right)}\tau _r(M,𝐭,𝐭^{})=\tau _r(M,𝐭+\gamma ,𝐭^{}).$$
(A3.6)
Also we have
$$e^{_{\mathrm{}}^{\mathrm{}}\gamma _nZ_{nn}}\tau (M,𝐭,𝐓,𝐭^{})=e^{_{\mathrm{}}^{\mathrm{}}\gamma _nZ_{nn}^{}}\tau (M,𝐭,𝐓,𝐭^{})=\tau (M,𝐭,𝐓+\gamma ,𝐭^{}),$$
(A3.7)
where
$$Z_{nn}=\frac{1}{4\pi ^2}\frac{z^n}{z_{}^{}{}_{}{}^{n}}V_{\mathrm{}}^{}\left(z^{}\right)V_{\mathrm{}}\left(z\right)𝑑z𝑑z^{},Z_{nn}^{}=\frac{1}{4\pi ^2}\frac{z^n}{z_{}^{}{}_{}{}^{n}}V_0^{}\left(z^{}\right)V_0\left(z\right)𝑑z𝑑z^{}.$$
(A3.8)
For instance
$$e^{_{\mathrm{}}^{\mathrm{}}T_nZ_{nn}}\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}nt_nt_n^{}\right)=e^{_{\mathrm{}}^{\mathrm{}}T_nZ_{nn}^{}}\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}nt_nt_n^{}\right)=\tau (M,𝐭,𝐓,𝐭^{}).$$
(A3.9)
## Appendix Appendix 4 Gauss factorization problem, additional symmetries, string equations and $`\mathrm{\Psi }DO`$ on the circle
Let us describe relevant string equations following Takasaki and Takebe ,. We shall also consider this topic in a more detailed paper.
Let us introduce infinite matrices to describe KP and TL flows and symmetries, see . Zakharov-Shabat dressing matrices are $`K`$ and $`\overline{K}`$. $`K`$ is a lower triangular matrix with unit main diagonal: $`\left(K\right)_{ii}=1`$. $`\overline{K}`$ is an upper triangular matrix. The matrices $`K,\overline{K}`$ depend on parameters $`M,𝐭,𝐓,𝐭^{}`$. The matrices $`\left(\mathrm{\Lambda }\right)_{ik}=\delta _{i,k1}`$, $`\left(\overline{\mathrm{\Lambda }}\right)_{ik}=\delta _{i,k+1}`$. For each value of $`𝐭,𝐓,𝐭^{}`$ and $`MZ`$ they solve Gauss (Riemann-Hilbert) factorization problem for infinite matrices:
$`\overline{K}=KG(M,𝐭,𝐓,𝐭^{}),G(M,𝐭,𝐓,𝐭^{})=\mathrm{exp}\left(\xi (𝐭,\mathrm{\Lambda })\right)\mathrm{\Lambda }^MG(\mathrm{𝟎},𝐓,\mathrm{𝟎})\overline{\mathrm{\Lambda }}^M\mathrm{exp}\left(\xi (𝐭^{},\overline{\mathrm{\Lambda }})\right).`$ (A4.1)
We put $`\mathrm{log}\left(\overline{K}_{ii}\right)=\varphi _{i+M}`$, and a set of fields $`\varphi _i(𝐭,𝐭^{}),\left(\mathrm{}<i<+\mathrm{}\right)`$ solves the hierarchy of higher two-dimensional TL equations.
Take $`L=K\mathrm{\Lambda }K^1`$, $`\overline{L}=\overline{K}\overline{\mathrm{\Lambda }}\overline{K}^1`$, and $`\left(\mathrm{\Delta }\right)_{ik}=i\delta _{i,k}`$, $`\widehat{M}=K\mathrm{\Delta }K^1+M+nt_nL^n`$, $`\widehat{\overline{M}}=\overline{K}\mathrm{\Delta }\overline{K}^1+M+nt_n^{}\overline{L}^n`$. Then the KP additional symmetries ,,,, and higher TL flows are written as
$$_{\beta _n}K=\left(\left(r\left(\widehat{M}\right)L^1\right)^n\right)_{}K,_{\beta _n}\overline{K}=\left(\left(r\left(\widehat{M}\right)L^1\right)^n\right)_+\overline{K},$$
(A4.2)
$$_{t_n^{}}K=\left(\overline{L}^n\right)_{}K,_{t_n^{}}\overline{K}=\left(\overline{L}^n\right)_+\overline{K}.$$
(A4.3)
Then the string equations are
$$\overline{L}L=r\left(\widehat{M}\right),$$
(A4.4)
$$\widehat{\overline{M}}=\widehat{M}.$$
(A4.5)
The first equation is a manifestation of the fact that the group time $`\beta _1`$ of the additional symmetry of KP can be identified with the Toda lattice time $`t_1^{}`$. In terms of tau-function we have the equation (A3.6) in terms of vertex operator action ,, or the equations (3.7.1) in case the tau-function is written in Miwa variables.
The second string equation (A4.5) is related to the symmetry of our tau-functions with respect to $`𝐭\beta `$.
When
$$r\left(M\right)=M+a,$$
(A4.6)
the equations (A4.4),(A4.5) describe $`c=1`$ string , see ,. In this case we easily get the relation
$$[\overline{L},L]=1.$$
(A4.7)
The string equations in the form of Takasaki allows us to notice the similarity to the different problem. The dispersionless limit of (A4.7) (and also of (A4.5), (A4.4), where $`r\left(M\right)=M^n`$, and of (A4.6)) will be written as
$$\overline{\lambda }\lambda =\mu ^n,nZ,$$
(A4.8)
$$\overline{\mu }=\mu .$$
(A4.9)
The case when $`\overline{\lambda }`$ and $`\overline{\mu }`$ are complex conjugate of $`\lambda `$ and of $`\mu `$ respectively, is of interest. These string equations (mainly the case $`n=1`$) were recently investigated to solve the so-called Laplacian growth problem, see . We are grateful to A.Zabrodin for the discussion on this problem. For the dispersionless limit of the KP and TL hierarchies see .
In case the function $`r\left(n\right)`$ has zeroes (described by divisor $`𝐦=(M_1,\mathrm{})`$: $`r\left(M_k\right)=0`$), one needs to produce the replacement:
$$\mathrm{\Lambda }\mathrm{\Lambda }\left(𝐦\right),\overline{\mathrm{\Lambda }}\overline{\mathrm{\Lambda }}\left(𝐦\right),$$
(A4.10)
where new matrices $`\mathrm{\Lambda }\left(𝐦\right)`$,$`\overline{\mathrm{\Lambda }}\left(𝐦\right)`$ are defined as
$$\left(\mathrm{\Lambda }\left(𝐦\right)\right)_{i,j}=\delta _{i,j1},jM_k,\left(\mathrm{\Lambda }\left(𝐦\right)\right)_{i,j}=0,j=M_k,$$
(A4.11)
$$\left(\overline{\mathrm{\Lambda }}\left(𝐦\right)\right)_{i,j}=\delta _{i,j+1},iM_k,\left(\overline{\mathrm{\Lambda }}\left(𝐦\right)\right)_{i,j}=0,i=M_k.$$
(A4.12)
This modification describes the open TL equation (3.4.9):
$$_{t_1}_{\beta _1}\phi _n=\delta \left(n\right)e^{\phi _{n1}\phi _n}\delta \left(n+1\right)e^{\phi _n\phi _{n+1}}.$$
(A4.13)
The set of fields $`\varphi _{\mathrm{}},\mathrm{},\varphi _{M_1},\varphi _{M_1+1}\mathrm{}`$ consists of the following
parts due to the conditions
$$\underset{n=\mathrm{}}{\overset{M_s}{}}\varphi _n=0,\underset{n=M_k+1}{\overset{M_{k+1}}{}}\varphi _n=0,\underset{n=\mathrm{}}{\overset{M_1+1}{}}\varphi _n=0,$$
(A4.14)
which result from $`\tau (M_k,𝐭,𝐭^{})=1`$.
Each internal part (3.5.3) of this Toda chain is a $`M_{k+1}M_k`$ sites chain. There are two semiinfinite parts which correspond to (3.5.2) and (3.5.4).
###### Remark 5
The matrix $`r(\widehat{M})`$ contains $`(\widehat{M}b_i)`$ in the denominator. The matrix $`(\widehat{M}b_i)^1`$ is $`K(\mathrm{\Delta }b_i)^1K^1(1+O(𝐭))`$ (compare the consideration of the inverse operators with ).
The KP tau-function (3.1.7) can be obtained as follows.
$$G(M,𝐭,𝐓,𝐭^{})=G(\mathrm{𝟎},𝐓,\mathrm{𝟎})U(M,𝐭,\beta ),U(M,𝐭,\beta )=U^+\left(𝐭\right)U^{}(M,\beta ).$$
(A4.15)
$$U^+\left(𝐭\right)=\mathrm{exp}\left(\xi (𝐭,\mathrm{\Lambda })\right),U^{}(M,\beta )=\mathrm{exp}\left(\xi (\beta ,\mathrm{\Lambda }^1r\left(\mathrm{\Delta }+M\right))\right),$$
(A4.16)
The matrix $`G(\mathrm{𝟎},𝐓,\mathrm{𝟎})`$ is related to the transformation of the eq.(3.3.11) to the eq.(3.3.8). By taking the projection $`UU_{}`$ for nonpositive values of matrix indices we obtain a determinant representation of the tau-function (3.1.7):
$$\tau _r(M,𝐭,\beta )=\frac{detU_{}(M,𝐭,\beta )}{det\left(U_{}^+\left(𝐭\right)\right)det\left(U_{}^{}(M,\beta )\right)}=detU_{}(M,𝐭,\beta ),$$
(A4.17)
since both determinants in the denominator are equal to one. Formula (A4.17) is also a Segal-Wilson formula for $`GL\left(\mathrm{}\right)`$ 2-cocycle $`C_M(U^+\left(𝐭\right),U^{}\left(\beta \right))`$. Choosing the function $`r`$ as in Section 3.2 we obtain hypergeometric functions listed in the Introduction.
###### Remark 6
Therefore the hypergeometric functions which were considered above have the meaning of $`GL(\mathrm{})`$ two-cocycle on the two multiparametrical group elements $`U^+(𝐭)`$ and $`U^{}(M,\beta )`$. Both elements $`U^+(𝐭)`$ and $`U^{}(M,\beta )`$ can be considered as elements of group of pseudodifferential operators on the circle. The corresponding Lie algebras consist of the multiplication operators $`\{z^n;nN_0\}`$ and of the pseudodifferential operators $`\{\left(\frac{1}{z}r(z\frac{d}{dz}+M)\right)^n;nN_0\}`$. Two sets of group times $`𝐭`$ and $`\beta `$ play the role of indeterminates of the hypergeometric functions (3.10.5). Formulas (3.1.7) and (3.1.12) mean the expansion of $`GL(\mathrm{})`$ group 2-cocycle in terms of corresponding Lie algebra 2-cocycle
$$c_M(z,\frac{1}{z}r\left(D\right))=r\left(M\right),c_M(\stackrel{~}{r}\left(D\right)z,\frac{1}{z}r\left(D+M\right))=\stackrel{~}{r}\left(M\right)r\left(M\right).$$
(A4.18)
Japanese cocycle is cohomological to Khesin-Kravchenko cocycle for the $`\mathrm{\Psi }DO`$ on the circle:
$$c_Mc_0\omega _M,$$
(A4.19)
which is
$$\omega _M(A,B)=res_{}A[\mathrm{log}\left(D+M\right),B]𝑑z,A,B\mathrm{\Psi }DO.$$
(A4.20)
For the group cocycle we have
$$C_M(e^{{\scriptscriptstyle z^nt_n}},e^{{\scriptscriptstyle \left(z^1r\left(D\right)\right)^n\beta _n}})=\tau _r(M,𝐭,\beta ),$$
(A4.21)
where we imply that the order of $`\mathrm{\Psi }`$DO $`r(D)`$ is 1 or less. About properties of $`e^{{\scriptscriptstyle (z^1r(D))^n\beta _n}}`$ see (A1.3),(A1.4).
###### Remark 7
It is interesting to note that in case of hypergeometric functions $`{}_{p}{}^{}F_{s}^{}`$ (1.2.6) the order of $`r`$ is $`ps`$ (see Example 3), and the condition $`ps1`$ is the condition of the convergence of this hypergeometric series, see . Namely the radius of convergence is finite in case $`ps=1`$, it is infinite when $`ps<1`$ and it is zero for $`ps>1`$ (this is true for the case when no one of $`a_k`$ in (1.2.6) is nonnegative integer).
###### Remark 8
The set of functions $`\{w(n,z),n=M,M+1,M+2,\mathrm{}\}`$, where
$$w(n,z)=\mathrm{exp}\left(\underset{m=0}{\overset{\mathrm{}}{}}t_m^{}\left(\frac{1}{z}r\left(D\right)\right)^m\right)z^n,$$
(A4.22)
may be identified with Sato Grassmannian (3.11.5) related to the cocycle (A4.21). The dual Grassmannian (3.11.6) is the set of one forms $`\{w^{}(n,z),n=M,M+1,M+2,\mathrm{}\}`$,
$$w^{}(n,z)=\mathrm{exp}\left(\underset{m=0}{\overset{\mathrm{}}{}}t_m^{}\left(\frac{1}{z}r\left(D\right)\right)^m\right)z^ndz.$$
(A4.23)
## Appendix Appendix 5 Equations with respect to $`𝐓`$ variables
Let us show that variables $`𝐓`$ play the role of time variables. We shall write down some equation involving the differentiation with respect to $`𝐓`$. The relevant Liouville equations is constructed in terms of
$$e^{v_k}=\frac{M+1\left|e^{H\left(𝐭\right)}\psi _ke^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|MM1\left|e^{H\left(𝐭\right)}\psi _k^{}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M}{M\left|e^{H\left(𝐭\right)}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M^2}.$$
(A5.1)
The equations are
$$\frac{^2v_k}{t_1T_k}=e^{v_k},kZ.$$
(A5.2)
To get three-wave equations take
$$\beta _{nm}=\frac{M\left|e^{H\left(𝐭\right)}\psi _m\psi _n^{}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M}{M\left|e^{H\left(𝐭\right)}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M},$$
(A5.3)
$$\beta _{1^{}n}=\frac{M+1\left|e^{H\left(𝐭\right)}\psi _ne^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M}{M\left|e^{H\left(𝐭\right)}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M},\beta _{n1^{}}=\frac{M1\left|e^{H\left(𝐭\right)}\psi _n^{}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M}{M\left|e^{H\left(𝐭\right)}e^{H_0\left(𝐓\right)}e^{H^{}\left(𝐭^{}\right)}\right|M}.$$
(A5.4)
Then we obtain
$$_{T_n}\beta _{1^{}m}=\beta _{1^{}n}\beta _{nm},_{T_n}\beta _{m1^{}}=\beta _{mn}\beta _{n1^{}},_{t_1}\beta _{mn}=\beta _{m1^{}}\beta _{1^{}n},mn,n,mZ.$$
(A5.5)
One obtains these equations with the help of the Lax type representation :
$$[_{T_m}+\beta _{1^{}m}_{t_1}^1\beta _{m1^{}},_{T_n}+\beta _{1^{}n}_{t_1}^1\beta _{n1^{}}]=0.$$
(A5.6)
It is possible to write down the discrete versions of these equations, and (the discrete and the continues) equations involving only $`𝐓`$ variables.
## Appendix Appendix 6 Orthogonal $`q`$-polynomials
Now we present the fermionic representation of polynomials listed in the Introduction. These polynomials are obtained by the specification of (3.10.25).
$`q`$-Askey-Wilson polynomials are defined as
$`p_n(x;a,b,c,d|q)=a^n(ab;q)_n(ac;q)_n(ad;q)_n`$
$`\times {}_{4}{}^{}\phi _{3}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{q^n,q^{n1}abcd,ae^{i\eta },ae^{i\eta }}{ab,ac,ad}}\right|q,q),x=\mathrm{cos}\eta .`$ (A6.1)
Let operator $`{}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)`$ be
$${}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{n+D}\right)\left(1abcdq^{n1+D}\right)\left(1ae^{i\eta }q^D\right)\left(1ae^{i\eta }q^D\right)}{\left(1abq^D\right)\left(1acq^D\right)\left(1adq^D\right)}.$$
(A6.2)
For this operator we have:
$$\frac{{}_{}{}^{4}\tau _{}^{3}(M,𝐭,𝐓,𝐭^{})}{{}_{}{}^{4}\tau _{}^{3}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}=M\left|e^{H\left(𝐭\right)}e^{A\left(𝐭^{}\right)}\right|M,$$
(A6.3)
where
$$A_k=\frac{1}{2\pi i}\psi ^{}\left(z\right)\left(\frac{1}{z}{}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)\right)^k\psi \left(z\right),k=1,2,\mathrm{}.$$
(A6.4)
If we take the variables $`𝐭^{}`$ as in (3.10.17) and
$$𝐭=(q,\frac{q^2}{2},\frac{q^3}{3},\mathrm{}),$$
(A6.5)
$`{\displaystyle \frac{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(M,𝐭,𝐓,𝐭^{})}{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}}={}_{4}{}^{}\phi _{3}^{}\left({\displaystyle \genfrac{}{}{0pt}{}{q^{Mn},q^{M+n1}abcd,aq^Me^{i\eta },aq^Me^{i\eta }}{q^Mab,q^Mac,q^Mad}}|q,q\right)=`$
$`{\displaystyle \underset{m=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{(q^{Mn};q)_m(q^{M+n1}abcd;q)_m(aq^Me^{i\eta };q)_m(aqe^{i\eta };q)_m}{(abq^M;q)_m(acq^M;q)_m(adq^M;q)_m}}{\displaystyle \frac{q^m}{(q;q)_m}}.`$ (A6.6)
For $`M=0`$ we obtain $`q`$-Askey-Wilson polynomials:
$`p_n(x;a,b,c,d|q)=aq^n(ab;q)_n(ac;q)_n(ad;q)_n{\displaystyle \frac{{}_{}{}^{4}\tau _{}^{3}(0,𝐭,𝐓,𝐭^{})}{{}_{}{}^{4}\tau _{}^{3}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})}}.`$ (A6.7)
If parameters in (A6.2) are
$$a=q^{\frac{2\alpha +1}{4}},b=q^{\frac{2\nu +1}{4}},c=q^{\frac{2\alpha +3}{4}},d=q^{\frac{2\nu +3}{4}},$$
(A6.8)
we get a fermionic representation for continuous $`q`$-Jacobi polynomials
$$P_n^{(\alpha ,\nu )}\left(x|q\right)=\frac{(q^{\alpha +1};q)_n}{(q;q)_n)}\frac{{}_{}{}^{4}\tau _{}^{3}(0,𝐭,𝐓,𝐭^{})}{{}_{}{}^{4}\tau _{}^{3}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})}.$$
(A6.9)
For $`c=a,b=d=q^{\frac{1}{2}}a`$ we have $`q`$-Gegenbauer polynomials
$$C_n(\mathrm{cos}\eta ;\mu |q)=\frac{(\mu ^2;q)_n}{\mu ^{\frac{n}{2}}(q;q)_n}\frac{{}_{}{}^{4}\tau _{}^{3}(0,𝐭,𝐓,𝐭^{})}{{}_{}{}^{4}\tau _{}^{3}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})},\mu =a^2.$$
(A6.10)
Clebsch-Gordan coefficients $`C_q(𝐥,𝐣)`$ see .
Let
$${}_{3}{}^{}r_{2}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{jl_1+D}\right)\left(1q^{l_1+j+1+D}\right)\left(1q^{l+m+D}\right)}{\left(1q^{l_2l+j+1+D}\right)\left(1q^{ll_2+j+D}\right)}.$$
(A6.11)
For variables from (3.10.17) and (A6.5) we have
$$\frac{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(M,𝐭,𝐓,𝐭^{})}{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}={}_{3}{}^{}\mathrm{\Phi }_{2}^{}\left(\genfrac{}{}{0pt}{}{jl_1+M,l_1+j+1+M,l+m+M}{l_2l+1+M,ll_2+j+M}|q,q\right).$$
(A6.12)
Thus we have the fermionic representation
$`C_q(𝐥,𝐣)={\displaystyle \frac{\left(1\right)^{l_1j}q^B\mathrm{\Delta }\left(𝐥\right)\left[l+l_2j\right]!\left([𝐥,𝐣]\left[2l+1\right]\right)^{\frac{1}{2}}}{\left[l_1l_2+l\right]!\left[l+l_2l_1\right]!\left[l_2l+j\right]!\left[l_1j\right]!\left[l_2+k\right]!\left[lm\right]!}}{\displaystyle \frac{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,𝐭,𝐓,𝐭^{})}{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})}}.`$ (A6.13)
$`q`$-Hahn polynomials
Let us take the operator:
$${}_{3}{}^{}r_{2}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{n+D}\right)\left(1abq^{n+1+D}\right)\left(1q^{x+D}\right)}{\left(1aq^{D+1}\right)\left(1q^{DN}\right)},nN.$$
(A6.14)
The corresponding tau function whose variables are defined in (3.10.17) and (A6.5):
$$\frac{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(M,𝐭,𝐓,𝐭^{})}{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}={}_{3}{}^{}\phi _{2}^{}\left(\genfrac{}{}{0pt}{}{q^{Mn},abq^{M+n+1},q^{Mx}}{aq^{M+1},q^{MN}}|q,q\right).$$
(A6.15)
Therefore the fermionic representation of $`q`$-Hahn polynomials is
$$Q_n(q^x;a,b;N|q)=\frac{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,𝐭,𝐓,𝐭^{})}{{}_{3}{}^{}\tau _{2}^{\left(q\right)}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})}.$$
(A6.16)
$`q`$-Racah polynomials Now we take as $`r\left(D\right)`$
$${}_{4}{}^{}r_{3}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{n+D}\right)\left(1abq^{n+1+D}\right)\left(1q^{x+D}\right)\left(1cdq^{x+1+D}\right)}{\left(1aq^{D+1}\right)\left(1bdq^{D+1}\right)\left(1cq^{D+1}\right)}.$$
(A6.17)
The tau function is:
$$\frac{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(M,𝐭,𝐓,𝐭^{})}{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}={}_{4}{}^{}\phi _{3}^{}\left(\genfrac{}{}{0pt}{}{q^{Mn},abq^{M+n+1},q^{Mx},cdq^{M+x+1}}{aq^{M+1},bdq^{M+1},cq^{M+1}}|q,q\right),$$
(A6.18)
where $`𝐭^{}`$ and $`𝐭`$ as in (3.10.17) and (A6.5). Thus we have an expression
$$R_n(\mu \left(x\right);a,b,c,d|q)=\frac{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(0,𝐭,𝐓,𝐭^{})}{{}_{4}{}^{}\tau _{3}^{\left(q\right)}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})},\mu \left(x\right)=q^x+cdq^{x+1}.$$
(A6.19)
Little $`q`$-Jacobi polynomials
Setting operator $`r\left(D\right)`$:
$${}_{2}{}^{}r_{1}^{\left(q\right)}\left(D\right)=\frac{\left(1q^{n+D}\right)\left(1abq^{n+1+D}\right)}{\left(1aq^{D+1}\right)},$$
(A6.20)
we get tau function:
$$\frac{{}_{2}{}^{}\tau _{1}^{\left(q\right)}(M,𝐭,𝐓,𝐭^{})}{{}_{2}{}^{}\tau _{1}^{\left(q\right)}(M,\mathrm{𝟎},𝐓,\mathrm{𝟎})}={}_{2}{}^{}\phi _{1}^{}\left(\genfrac{}{}{0pt}{}{q^{Mn},abq^{M+n+1}}{aq^{M+1}}|q,qx\right),$$
(A6.21)
where $`t_m=\frac{\left(qx\right)^m}{m}`$ and $`𝐭^{}`$ as in (3.10.17).
$$p_n(x;a,b|q)=\frac{{}_{2}{}^{}\tau _{1}^{\left(q\right)}(0,𝐭,𝐓,𝐭^{})}{{}_{2}{}^{}\tau _{1}^{\left(q\right)}(0,\mathrm{𝟎},𝐓,\mathrm{𝟎})}.$$
(A6.22)
## Appendix Appendix 7
Since we have $`\mathrm{exp}(T_n:\psi _n^{}\psi _n:)=1+(e^{T_n}1):\psi _n^{}\psi _n:`$ the tau-function (3.3.1) is a linear function of each $`e^{T_m}1`$. Take in the series (3.3.1) a coefficient before the monomial $`_i\left(e^{T_{m_i}}1\right)`$, where $`\left\{m_i\right\}`$ form a set of indices $`X`$. Denote this coefficient as $`\rho (X,𝐭,𝐭^{})\mathrm{exp}\left(_{k=1}^{\mathrm{}}kt_kt_k^{}\right)`$. The function $`\rho (X,𝐭,𝐭^{})`$ has the following meaning in the probability theory . The Schur measure on partitions is defined via the formula:
$$\mu (𝐧,𝐭,𝐭^{})=s_𝐧\left(𝐭\right)s_𝐧\left(𝐭^{}\right)\mathrm{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}kt_kt_k^{}\right),$$
(A7.1)
where $`𝐧`$ is a partition, and $`𝐭,𝐭^{}`$ are parameters. Then the function $`\rho (X,𝐭,𝐭^{})`$ describes the probability that the set $`\left\{n_ii\right\}`$ contains a given set $`XZ`$. The Schur measure has a natural group-theoretical interpretation, see . In this context it may be of interest to consider the following measure on partitions, which generalizes (A7.1):
$$\mu _{\stackrel{~}{r}r}(𝐧,𝐭,𝐭^{})=\frac{\stackrel{~}{r}_𝐧\left(M\right)r_𝐧\left(M\right)s_𝐧\left(𝐭\right)s_𝐧\left(𝐭^{}\right)}{M\left|e^{\stackrel{~}{A}\left(𝐭\right)}e^{A\left(𝐭^{}\right)}\right|M},$$
(A7.2)
see (3.1.12). Let us note that in the case of hypergeometric tau-functions, function $`r_𝐧`$ is a rational function of the Schur polynomials, see (3.12),(3.12) below.
### Acknowledgements
One of the authors (A.O.) thanks Vl.Dragovich and T.Shiota for the helpful numerous discussions. D.S. would like to thank S.Senchenko for helpful discussions. We thank S.Milne for sending , L.A.Dickey for , J.Harnad and A.Its for the interesting remarks after the first author’s talk at the CRM workshop “Isomonodromy deformations”, May 2000. We thank also A.Zabrodin for the discussion of the paper . We thank A. Nakamura for sending . |
warning/0001/cond-mat0001247.html | ar5iv | text | # On the problem of neural network decomposition into some subnets
## I Hopfield’s model of a neural network
A neural network of size $`n`$ is a set of $`n`$ connected spin variables (spins) $`\sigma _i`$; each $`\sigma _i`$ can be either $`1`$ or $`1`$:
$$\sigma _i=\{\pm 1\},i=1,2,\mathrm{}n.$$
$`(1)`$
The interaction between spins is described by a connection matrix. Let $`J_{ii^{}}`$ be the connection strength between the spins $`\sigma _i`$ and $`\sigma _i^{}`$<sup>*</sup><sup>*</sup>*For the sake of simplicity we suppose that there is no self-interaction in the system: $`J_{ii}=0i`$., and let $`\sigma _i(t)`$ be the value of $`i`$th spin at time $`t`$, then
$$h_i(t)=\underset{i^{}=1}{\overset{n}{}}J_{ii^{}}\sigma _i^{}(t)$$
$`(2)`$
represents the local field that the spin $`\sigma _i`$ experiences at time $`t`$. Under the action of this field the new value of the spin $`\sigma _i`$ at the next moment $`t+1`$ is:
$$\sigma _i(t+1)=\{\begin{array}{cc}\hfill \sigma _i(t),& \text{ if }h_i(t)\sigma _i(t)0\hfill \\ \hfill \sigma _i(t),& \text{ if }h_i(t)\sigma _i(t)<0\hfill \end{array}$$
$`(3)`$
The vectors which coordinates are $`\{\pm 1\}`$ only is called the configuration vectors. We denote the configuration vectors by small Greek letters.
It is convenient to describe the state of the network at time $`t`$ by $`n`$-dimensional configuration vector
$$\stackrel{}{\sigma }(t)=(\sigma _1(t),\sigma _2(t),\mathrm{},\sigma _n(t)).$$
If we introduce the connection matrix $`𝐉=\left(J_{ii^{}}\right)_1^n`$ and define the quadratic form
$$E(t)=\underset{i,i^{}=1}{\overset{n}{}}J_{ii^{}}\sigma _i(t)\sigma _i^{}(t)=(𝐉\stackrel{}{\sigma }(t),\stackrel{}{\sigma }(t)),$$
$`(4)`$
then it is easy to show that for any symmetrical connection matrix $`𝐉`$ the overturn of a spin $`\sigma _i(t)`$, which value does not coincide with the sign of $`h_i(t)`$, leads to the decrease of $`E(t)`$:
$$E(t+1)=E(t)+4\sigma _i(t)h_i(t).$$
$`(5)`$
$`E(t)`$ can be interpreted as the energy of the state $`\stackrel{}{\sigma }(t)`$. As the number of network states is finite and the $`i`$th spin does not turn over if $`h_i(t)=0`$, it is obvious that the final state of the network would be a state which corresponds to a minimum (may be local) of the energy $`E(t)`$. In such a state every spin $`\sigma _i`$ will be align with its local field $`h_i`$ and there will be no further evolution of the network. These states are called the fixed points of the network. Consequently, if the configuration vector $`\stackrel{}{\sigma }^{}=(\sigma _1^{},\sigma _2^{},\mathrm{},\sigma _n^{})`$ is a fixed points, then
$$\sigma _i^{}=sgn\left(\underset{i^{}=1}{\overset{n}{}}J_{ii^{}}\sigma _i^{}^{}\right),i=1,2,\mathrm{},n.$$
$`(6)`$
In what follows the configuration vectors which are fixed points will be marked by superscripts ”\*”.
Let’s define a neural network which is called Hopfield’s network. Let $`p`$ be a number of preassigned configuration vectors $`\stackrel{}{\xi }^{(l)}`$, which are called the memorized patterns:
$$\stackrel{}{\xi }^{(l)}=(\xi _1^{(l)},\xi _2^{(l)},\mathrm{},\xi _n^{(l)}),l=1,2,\mathrm{},p.$$
$`(7)`$
(The superscripts numerate the vectors from $`\mathrm{R}^\mathrm{n}`$ and the subscripts numerate their coordinates. Usually it is assumed that $`p<n`$ or even $`p<<n`$.) J.Hopfield proposed to use the connection matrix of the form:
$$J_{ii^{}}=\{\begin{array}{cc}_{l=1}^p\xi _i^{(l)}\xi _i^{}^{(l)},\hfill & ii^{}\hfill \\ 0,\hfill & i=i^{},i,i^{}=1,2,\mathrm{},n.\hfill \end{array}$$
$`(8)`$
The matrix $`𝐉`$ (8) is a symmetric matrix with zero diagonal elements. Then, the fixed points are the minima of the energy $`E`$ given by Eq.(4). If we define $`(p\times n)`$-matrix $`𝚵`$ with $`p`$ memorized patterns (7) as the rows,
$$𝚵=\left(\begin{array}{c}\stackrel{}{\xi }^{(1)}\\ \stackrel{}{\xi }^{(2)}\\ \mathrm{}\\ \stackrel{}{\xi }^{(p)}\end{array}\right)=\left(\begin{array}{cccc}\xi _1^{(1)}& \xi _2^{(1)}& \mathrm{}& \xi _n^{(1)}\\ \xi _1^{(2)}& \xi _2^{(2)}& \mathrm{}& \xi _n^{(2)}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \xi _1^{(p)}& \xi _2^{(p)}& \mathrm{}& \xi _n^{(p)}\end{array}\right)$$
$`(9)`$
then the expression for the connection matrix takes the form
$$𝐉=𝚵^T𝚵p𝐈,$$
$`(10)`$
where $`(n\times p)`$-matrix $`𝚵^T`$ is the transpose of matrix $`𝚵`$ and $`𝐈`$ is the unit matrix in the space $`\mathrm{R}^\mathrm{n}`$. Therefore the searching of the fixed points of Hopfield’s network reduces to the maximization of the functional
$$(𝚵^T𝚵\stackrel{}{\sigma },\stackrel{}{\sigma })=𝚵\stackrel{}{\sigma }^2.$$
But this problem can be reformulated, if $`n`$ $`p`$-dimensional vectors $`\stackrel{}{\xi }_i`$, which are the columns of matrix $`𝚵`$ are introduced:
$$\stackrel{}{\xi }_i=\left(\begin{array}{c}\xi _i^{(1)}\\ \xi _i^{(2)}\\ \mathrm{}\\ \xi _i^{(p)}\end{array}\right)\mathrm{R}^p,i=1,2,\mathrm{},n.$$
$`(11)`$
In contrast to $`n`$-dimensional vectors $`\stackrel{}{\xi }^{(l)}`$ defined by Eq.(7), here the subscripts numerate the vectors $`\stackrel{}{\xi }_i`$ from $`\mathrm{R}^\mathrm{p}`$ and the superscripts numerate their coordinates.
It is easy to see, that the problem of maximization of the functional $`𝚵\stackrel{}{\sigma }^2`$ takes the form:
$$\underset{i=1}{\overset{n}{}}\sigma _i\stackrel{}{\xi }_i\text{ max},\text{ where }\sigma _i=\{\pm 1\}i.$$
$`(12)`$
In other words, we have to find out such a weighted sum of the $`p`$-dimensional vectors $`\stackrel{}{\xi }_i`$ with the weights are equal $`\{\pm 1\}`$, which length would be maximal. In what follows the expression (12) would be the start point of our consideration.
## II Factor analysis and extremal grouping of parameters
The problem (12) is a special case of the problem which is well-known for the centroid method of the factor analysis. The basic idea of the factor analysis is to replace the great number of the parameters, which describe the objects under investigation, by a considerably lesser set of specially constructed characteristics provided that such replacement would not lead to the loss of the essential information about these objects.
The formalization of this idea can be done in the following way. Let us have $`p`$ objects which are represented by the vectors $`\stackrel{}{x}^{(l)}=(x_1^{(l)},x_2^{(l)},\mathrm{},x_n^{(l)}),l=1,2,\mathrm{},p`$ in the space $`\mathrm{R}^\mathrm{n}`$. Let’s consider the $`(p\times n)`$-matrix $`𝐗`$, which rows are the object-vectors $`\stackrel{}{x}^{(l)}`$. (This matrix is an analog of the matrix $`𝚵`$ (9), but now the matrix elements can be an arbitrary real numbers, and not $`\pm 1`$ only.) On the other hand the matrix $`𝐗`$ can be described as the matrix which columns are the parameter-vectors $`\stackrel{}{x}_i`$:
$$\stackrel{}{x}_i=\left(\begin{array}{c}x_i^{(1)}\\ x_i^{(2)}\\ \mathrm{}\\ x_i^{(p)}\end{array}\right),i=1,2,..n.$$
(We recall that the vectors from the space $`\mathrm{R}^\mathrm{n}`$ are numerated by superscripts: $`l=1,\mathrm{},p`$, and the vectors from the space $`\mathrm{R}^\mathrm{p}`$ by subscripts: $`i=1,\mathrm{},n`$.)
If a relatively small number $`t`$ ($`t<<n`$) of such $`p`$-dimensional vectors $`\stackrel{}{f}_1,\stackrel{}{f}_2,\mathrm{},\stackrel{}{f}_t`$ can be found, that the papameter-vectors $`\stackrel{}{x}_i`$ can be represented in the form
$$\stackrel{}{x}_i=\underset{s=1}{\overset{t}{}}a_{is}\stackrel{}{f}_s+\stackrel{}{a}_i,i=1,2,\mathrm{},n,$$
where the remainders $`\stackrel{}{a}_i`$ are small in some sense and can be omitted, then the objects can be described by the characteristics $`\stackrel{}{f}_s`$ instead of the initially used parameters $`\stackrel{}{x}_i`$. Indeed, due to the smallness of the remainders $`\stackrel{}{a}_i`$, characteristics $`\stackrel{}{f}_s`$ adequately describe the investigated phenomenon. But it is much more convenient to work if the number of the parameters is considerable reduced. The characteristics $`\stackrel{}{f}_s`$ are called the essential factors.
The various models of the factor analysis differ in the forms in which the factors $`\stackrel{}{f}_s`$ are sought and the sense in which the smallness of $`\stackrel{}{a}_i`$ is understood. In the centroid method the first factor $`\stackrel{}{f}_1`$ is sought as a linear combination $`_{i=1}^n\sigma _i\stackrel{}{x}_i`$ of the parameters $`\stackrel{}{x}_i`$ with the weights $`\sigma _i=\{\pm 1\}`$, that have a maximal length
$$\stackrel{}{f}_1\underset{i=1}{\overset{n}{}}\sigma _i^{}\stackrel{}{x}_i,\text{ where }\underset{i=1}{\overset{n}{}}\sigma _i^{}\stackrel{}{x}_i=\underset{\sigma _i=\{\pm 1\}}{\mathrm{max}}\underset{i=1}{\overset{n}{}}\sigma _i\stackrel{}{x}_i.$$
$`(13)`$
The comparison of Eq.(12) and Eq.(13) shows that the problem of the network fixed points searching is equivalent to the construction of the first centroid factor for the set of the $`p`$-dimensional vectors $`\stackrel{}{\xi }_i`$ (11).
In the centroid method after the construction of the first factor $`\stackrel{}{f}_1`$, the vectors $`b_{i1}\stackrel{}{f}_1`$, where $`b_{i1},i=1,2,.,n`$ are some coefficients, are subtracted from each parameter-vector $`\stackrel{}{x}_i`$. In such a way we obtain a new set of vectors $`\stackrel{}{x}_i^{}=\stackrel{}{x}_ib_{i1}\stackrel{}{f}_1`$ for which their own factor is constructed by analogy. This factor would be the second factor for the initial parameters $`\stackrel{}{x}_i`$ . This process will be repeated till the vectors which are obtained after the next step would be small enough. For details see .
An important generalization of the factor analysis was the idea of the extremal grouping of the parameters suggested by E.M.Braverman in 1970. Braverman introduced a model of the factor analysis where an essentially nonuniform distribution of the vectors $`\stackrel{}{x}_i`$ in the space $`\mathrm{R}^\mathrm{p}`$ was taken into account.
Indeed, if the number $`n`$ of the parameter-vectors is very large, it is possible that they can be divided into some compact groups such that the vectors joined into one group are ”strongly correlated” with each other and are ”weakly correlated” with the parameters included into other groups. Then it is reasonable to construct the factors not for the full set of the parameter-vectors, but for every compact group separately. If these groups are compact enough, we can restrict ourselves with the first factor of each group only. To divide the parameter-vectors into these compact groups, Braverman suggested an approach connected with the maximization of a certain functional depending both on the grouping of the parameters and on the choice of the factors.
Let’s write down Braverman’s functional. Let $`p`$-dimensional vectors $`\stackrel{}{x}_1,\stackrel{}{x}_2,\mathrm{},\stackrel{}{x}_n`$ be divided into current disjoint groups $`A_1,A_2,\mathrm{},A_t`$:
$$A_1A_2\mathrm{}A_t=\{1,2,\mathrm{},n\}.$$
For every group $`A_s`$ the first centroid factor can be constructed as the solution of the problem:
$$\underset{iA_s}{}\sigma _i^{}\stackrel{}{x}_i=\underset{\sigma _i=\{\pm 1\}}{\mathrm{max}}\underset{iA_s}{}\sigma _i\stackrel{}{x}_i.$$
$`(14)`$
Then, the partition into $`t`$ the most compact groups is obtained as a result of the maximization of the functional:
$$M(A_1,A_2,\mathrm{},A_t)=\underset{iA_1}{}\sigma _i^{}\stackrel{}{x}_i+\underset{iA_2}{}\sigma _i^{}\stackrel{}{x}_i+\mathrm{}+\underset{iA_t}{}\sigma _i^{}\stackrel{}{x}_i\mathrm{max}$$
$`(15)`$
where $`\sigma _i^{}`$ are the solutions of the problem (14) for every group $`A_s,s=1,2,..,t`$. We want to notice, that, though the problem of maximization of the functional (15) is very hard, the method for the extremal grouping of parameters was successfully used for various problems in engineering, economics, sociology, psychology and other fields .
## III Neural networks decomposition into some subnets
Let in Eqs.(14),(15) the vectors $`\stackrel{}{x}_i`$ be replaced by the vectors $`\stackrel{}{\xi }_i`$ from Eq.(11), i.e. only the vectors with the coordinates $`\{\pm 1\}`$ are under consideration. Then, in the framework of the neural network paradigm, the problem (14),(15) can be interpreted as the problem of the grouping of the network neurons into some connected groups.
Indeed, natural networks have evident differential structure: different neuron groups have different functions, they respond for the regulation/analysis of different aspects of a complicate pattern which is worked over by the network. To some extent every such neuron group can be treated as an autonomous neural network of the smaller size which is dealing with some specific features of the pattern.
Let a network be consisted of some groups of neurons (subnets) $`A_1,A_2,\mathrm{},A_t`$. There is one universal mechanism for the functioning of all network neurons: a spin $`\sigma _i`$ turns over if its sign does not coincide with the sign of the field $`h_i`$ acting on this spin. However, it is reasonable to assume that the incoming excitations from the neurons belonging to the same group as the neuron $`\sigma _i`$ affect this neuron stronger then the excitations from the neurons of other groups (those, which analyze the same pattern from other points of view). This hierarchy of excitations can be modelled in different ways. As an initial model it can be assumed that:
$$h_i(t)=\underset{i^{}A_s}{\overset{n}{}}J_{ii^{}}\sigma _i^{}(t)$$
$`(16)`$
where the summation is taken over all neurons belonging to the same group $`A_s`$ as the $`i`$th neuron.
The subnet consisting of the neurons from the group $`A_s`$ is evolving to one of its fixed points. This leads us to the problem (14). And the network as a whole is acting so, that the composite functional
$$M(\{\sigma _i\}_1^n)=\underset{iA_1}{}\sigma _i\stackrel{}{\xi }_i+\underset{iA_2}{}\sigma _i\stackrel{}{\xi }_i+\mathrm{}+\underset{iA_t}{}\sigma _i\stackrel{}{\xi }_i\mathrm{max}$$
$`(17)`$
would be maximized.
We have discussed the situation when the neurons are already decomposed into groups $`A_1,A_2,\mathrm{},A_t`$. If the structure of the groups is unknown, but their number $`t`$ is fixed, it is necessary to maximize the functional (17) with respect to the structure of the groups $`A_s`$ as well as with respect to all the weights $`\sigma _i`$ inside every group. In this case Eqs.(14),(15), where the vectors $`\stackrel{}{\xi }_i`$ have to be substituted instead of the vectors $`\stackrel{}{x}_i`$, describe the optimal decomposition of the network into $`t`$ autonomous subnets.
Here some remarks must be done. Firstly, it is easy to see, that when the number of the groups $`t`$ increases, the functional $`M(A_1,..A_t)`$ (15) is nondecreasing (it follows from the triangle inequality). This functional attains it’s global maximum when the number of the groups $`t`$ is equal $`n`$. However, it is a trivial decomposition. Simple geometric arguments show that when a group of strongly correlated vectors $`\stackrel{}{x}_i`$ is divided into two subgroups, the functional (15) increases negligibly. So, the problem is not to get the global maximum of the functional (15), but to obtain such a number $`t^{}`$ of the groups beginning with which the further increase of the number of the groups would not lead to the substantial increase of this functional. About these $`t^{}`$ groups we can speak as about the proper number of the subnets which constitute the initial $`n`$-network.
Furthermore, it is reasonable to try to interpret the specific characteristics of each obtained subnet in meaning terms. In other words, we can try to understand what kind of pattern’s characteristics are analyzed by each particular subnet, i.e. we must determine what kind of neurons are joined in the group. On this step the monograph , which reflects the accumulated experience in this field, can be useful.
Secondly, the above mentioned program can be fulfilled only if we are able to solve two problems: A) to find out the compact groups of the vectors $`\stackrel{}{\xi }_i`$; B) to determine the optimal configuration $`\{\sigma _i^{}\},iA_s`$, for each group. What concerns the problem B, actually all the attempts to create an effective algorithm for the maximization of the functional $`(𝐉\stackrel{}{\sigma },\stackrel{}{\sigma })`$ are devoted to this problem.
The problem A is much less studied and seems to be more complicated. Usually, it is solved by step by step transferring of $`p`$-dimensional vectors from one group to another, and the comparison of the values of the functional (15) for the consequently obtained grouping. When $`n`$ is rather large, in such a way only the determination of the local maximum of the functional $`M`$ is guaranteed. We know not so much papers , devoted exactly to the problem of finding of the global maximum of a functional of type (15). In these papers the general case of vectors $`\stackrel{}{x}_i`$ with real coordinates is studied. As for neural networks, the vectors $`\stackrel{}{\xi }_i`$ are specific: their coordinates are $`\{\pm 1\}`$. It can be hope that the specific character of the vectors $`\stackrel{}{\xi }_i`$ would make it possible to present effective method for the searching of the compact groups.
And the last remark. Although the proposed approach was formulated for Hopfield’s model, it can be generalized for the case of an arbitrary symmetric connection matrix: it is sufficient to replace in Eqs. (14), (15) and (17) the term
$$\underset{iA_s}{}\sigma _i^{}\stackrel{}{x}_i$$
by
$$\left(\underset{i,i^{}A_s}{}J_{ii^{}}\sigma _i\sigma _i^{}\right)^{1/2},$$
and all reasoning are valid.
This project was supported partially by Russian Basic Research Foundation. |
warning/0001/gr-qc0001082.html | ar5iv | text | # Dynamical study of the empty Bianchi type I model in generalised scalar-tensor theory
## 1 Introduction.
The scalar-tensor theories of gravitation allow to the gravitational constant to vary. Such a phenomenon happens in a large number of theories which try to unify gravitation with the other interaction forces. In the vacuum case, the most general form of the action of the scalar-tensor theories is written :
$$S=\left[F(\varphi )R1/2(\phi )^2U(\phi )\right]\sqrt{g}d^4x$$
(1)
where $`\phi `$ is a scalar field, $`U(\phi )`$ a potential. We get General Relativity with $`F(\phi )=cte`$ and Brans-Dicke theory with $`U=0`$, $`F(\phi )=\phi ^2/8\omega `$ and $`\omega =cte`$. When $`F(\phi )`$ is anatically invertible this action can always be written with a Brans-Dicke scalar field. Putting $`\varphi =F(\phi )`$ and $`\omega (\phi )=F/\left[2(dF/d\phi )^2\right])`$, we get :
$$S=\left[\varphi R\frac{\omega (\varphi )}{\varphi }(\varphi )^2U(\varphi )\right]\sqrt{g}d^4x$$
(2)
We will take $`U(\varphi )=0`$ so that we can obtain a Newtonian limit for the weak fields . Techniques to find exact or asymptotic solutions to the field equations derived from action (2), with or without matter, in an anisotropic Universe, by means of a conformal transformation, have been described in . Exact solutions and asymptotic behaviours of the scale factor have been analysed for the generalised scalar-tensor theory in FLRW model with matter in . Dynamical studies have been made for Brans-Dicke theory in a FLRW model in . Here, we will work in an empty Bianchi type I Universe. We will introduce new variables, write the field equations with their first derivatives and then perform an analysis to get analytically the sign of the first and second derivatives of the metric functions, without asymptotic methods, whatever $`\omega (\varphi )`$. Hence we will get the qualitative form of these functions in the Brans-Dicke frame for any time: are they increasing or decreasing, do extrema exist and if so, how many, is there inflation, do they tend towards a power law type, etc.
In section 2, we write the field equations of the vacuum Bianchi type I model with the new variables. In section 3, we study particular values of these variables and in section 4 we describe the method which gives the sign of the first derivatives of the metric functions, depending on the form of $`\omega (\varphi )`$. In section 5, we apply our method to three different forms of the coupling function which are all such that $`\omega \mathrm{}`$ and $`\omega _\varphi \omega ^30`$ if we adjust some of their parameters. These two limits ensures that the PPN parameters converge towards values in agreement with the observational data . Thus the different theories, corresponding to different choices of the coupling $`\omega (\varphi )`$, converge towards relativistic behaviours. In section 6, we examine the three metric functions and under what conditions they are increasing or decreasing together, etc. In the section (7), we describe the method giving the sign of the second derivatives of the metric functions and examine in which conditions they can be decelerated at late time. We apply our results to the coupling functions of section 5.
## 2 The field equations
The metric is:
$$ds^2=dt^2+a^2(\omega ^1)^2+b^2(\omega ^2)^2+c^2(\omega ^3)^2$$
(3)
where the $`\omega ^i`$ are the 1-forms of the Bianchi type I model, $`t`$ the proper time and $`a(t)`$, $`b(t)`$, $`c(t)`$ the metric functions depending on $`t`$. We define the $`\tau `$ time as:
$$d\tau =abcdt$$
(4)
and then, the field equations and the Klein-Gordon equation are written:
$`{\displaystyle \frac{a^{,,}}{a}}{\displaystyle \frac{a^{,2}}{a^2}}+{\displaystyle \frac{a^,}{a}}{\displaystyle \frac{\varphi ^,}{\varphi }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\omega ^,}{3+2\omega }}{\displaystyle \frac{\varphi ^,}{\varphi }}`$ $`=`$ $`0`$
$`{\displaystyle \frac{b^{,,}}{b}}{\displaystyle \frac{b^{,2}}{b^2}}+{\displaystyle \frac{b^,}{b}}{\displaystyle \frac{\varphi ^,}{\varphi }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\omega ^,}{3+2\omega }}{\displaystyle \frac{\varphi ^,}{\varphi }}`$ $`=`$ $`0`$ (5)
$`{\displaystyle \frac{c^{,,}}{c}}{\displaystyle \frac{c^{,2}}{c^2}}+{\displaystyle \frac{c^,}{c}}{\displaystyle \frac{\varphi ^,}{\varphi }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\omega ^,}{3+2\omega }}{\displaystyle \frac{\varphi ^,}{\varphi }}`$ $`=`$ $`0`$
$`{\displaystyle \frac{a^,}{a}}{\displaystyle \frac{b^,}{b}}+{\displaystyle \frac{a^,}{a}}{\displaystyle \frac{c^,}{c}}+{\displaystyle \frac{b^,}{b}}{\displaystyle \frac{c^,}{c}}+{\displaystyle \frac{\varphi ^,}{\varphi }}({\displaystyle \frac{a^,}{a}}+{\displaystyle \frac{b^,}{b}}+{\displaystyle \frac{c^,}{c}}){\displaystyle \frac{\omega }{2}}({\displaystyle \frac{\varphi ^,}{\varphi }})^2`$ $`=`$ $`0`$ (6)
$`\varphi ^{,,}={\displaystyle \frac{\omega ^,\varphi ^,}{3+2\omega }}`$ (7)
We integrate (7) and get :
$$A\varphi ^,\sqrt{3+2\omega }=1$$
(8)
$`A`$ being an integration constant. We see in this last expression that the coupling function must be superior to -3/2 so that the square root is real. We use (7) to introduce the second derivative of the scalar field in (5) and put:
$$\alpha =\frac{a^,}{a}\varphi ,\text{ }\beta =\frac{b^,}{b}\varphi ,\text{ }\gamma =\frac{c^,}{c}\varphi ,\text{ }\varphi ^,=\mathrm{\Phi }$$
(9)
After integrating, the field equations become:
$`\alpha +{\displaystyle \frac{1}{2}}\mathrm{\Phi }=\alpha _0`$
$`\beta +{\displaystyle \frac{1}{2}}\mathrm{\Phi }=\beta _0`$
$`\gamma +{\displaystyle \frac{1}{2}}\mathrm{\Phi }=\gamma _0`$
$`\alpha \beta +\alpha \gamma +\beta \gamma +\mathrm{\Phi }(\alpha +\beta +\gamma ){\displaystyle \frac{1}{4}}(A^23\mathrm{\Phi }^2)=0`$ (10)
$`\alpha _0`$, $`\beta _0`$, $`\gamma _0`$ being integration constants. The constraint imposes the condition :
$$\alpha _0\beta _0+\alpha _0\gamma _0+\beta _0\gamma _0=(4A^2)^1$$
(11)
The physical solutions are such that the metric functions and the scalar field are positive. Hence, the sign of the variables $`\alpha `$, $`\beta `$, $`\gamma `$ will be the same as the sign of the first derivative of the metric functions. The sign of $`\mathrm{\Phi }`$ will be the same as $`\varphi ^,`$. Negative scalar fields have already been considered in but it means that, in the Einstein frame, the gravitational constant will be negative. For this reason, many authors deal with positive scalar fields. We will do the same, but the method can easily be extended to negative ones. In what it follows, we will consider only the metric function $`a`$. What we write for $`a`$ will be valid for $`b`$ and $`c`$. Let us say a few words about exact solutions of the field equations. From (2), we can easily show that:
$$a=exp(\frac{\alpha _0}{\varphi }𝑑\tau +cte)\varphi ^{1/2}$$
(12)
The scalar field can be calculated by integrating and inverting (13):
$$d\tau =A\sqrt{3+2\omega }𝑑\varphi $$
(13)
Therefore, we can obtain exact solutions of the metric functions for the simple form of the coupling function.
What is the link between the results we will obtain in the $`\tau `$ time and the behaviours of the metric functions in the $`t`$ time. Since $`a(\tau )=a(\tau (t))=a(t)`$, the amplitudes of the metric functions will be the same in both $`\tau `$ and $`t`$ times. Moreover as:
$$da/dt=da/d\tau d\tau /dt=da/d\tau (abc)^1$$
with $`abc>0`$, the sign of the first derivatives of the metric functions will not be different in $`\tau `$ or $`t`$ time. Of course the amplitudes of all the derivatives will be different. While it will always be possible to determine asymptotically the amplitudes of $`a^,`$, this will not be the same for $`\dot{a}`$.
Therefore, as we are mainly interested by the sign of $`a^,`$, $`a^{,,}`$ and $`\ddot{a}`$, this is not important. The sign of the second derivatives will be different in both times since
$$d^2a/dt^2=\ddot{a}=\left[a^{,,}a^,(a^,/a+b^,/b+c^,/c)\right](abc)^2$$
an overdot denoting differentiation with respect to $`t`$.
For these reasons, all that we will say about the sign of the first derivatives will apply to both $`t`$ and $`\tau `$ time. Hence, results of section 3, 4, 5, 6 and in particular table 1 (except the sign of the second derivative of the scalar field which will be different by $`\varphi ^{,,}`$ in the $`t`$ time) will not change in $`t`$ time since they depends on the sign of constants or first derivative of $`\omega `$ with respect to $`\varphi `$. In section 7, where we will deal with the sign of the second derivatives, we will study separately the sign of $`a^{,,}`$ and $`\ddot{a}`$.
Another difference between $`\tau `$ and $`t`$ time is that, for instance, $`t`$ can diverge for a finite value of $`\tau `$. It can, for instance, transform a Universe that exists during a finite $`\tau `$ time into a Universe which would exist in an infinite $`t`$ time. But we will not pay attention to this type of phenomenon in our study. In fact, in most cases, we will use $`\varphi `$ as a time coordinate, particularly in section 5 and 7, and so we will have no need to know the intervals of $`\tau `$ or $`t`$.
## 3 Study of the first derivative of a metric function.
We consider the first equation of (2). The solution of this equation in the $`(\alpha ,\mathrm{\Phi })`$ plane is represented by a straight line. We have two cases depending on the sign of $`\alpha _0`$, which are represented on graph 1. To describe the variations of the metric function $`a`$, we have to study the dynamic of a point $`(\alpha ,\mathrm{\Phi })`$ on this straight line so that we know the sign of $`\alpha `$ and hence, this of $`a^,`$ during the time evolution. The straight line cuts the $`\mathrm{\Phi }`$ axe at $`(\alpha ,\mathrm{\Phi })=(0,2\alpha _0)`$ and the $`\alpha `$ axe at $`(\alpha ,\mathrm{\Phi })=(\alpha _0,0)`$. In $`(0,2\alpha _0)`$, we have $`\alpha =0`$. This means that :
\- the metric function $`a`$ reaches an extrema if the motion of the point $`(\alpha ,\mathrm{\Phi })`$ on the straight line is such that the sign of $`\alpha `$ change. It is an inflexion point for the metric function, if the motion of the point $`(\alpha ,\mathrm{\Phi })`$ on the straight line changes direction when it reaches $`(0,2\alpha _0)`$.
\- If the motion of the $`(\alpha ,\mathrm{\Phi })`$ point on the straight line is such that it tend asymptotically towards $`(0,2\alpha _0)`$ then a possible explanation is that the scalar field vanishes or that $`a\tau `$.
In $`(\alpha _0,0)`$, the first derivative of the scalar field disappears. We will show below that the scalar field is a monotone function of $`\tau `$. Hence , $`\varphi ^,=0`$ can be an inflexion point for $`\varphi `$ in the $`\tau `$ time if the motion of the point $`(\alpha ,\mathrm{\Phi })`$ changes direction after reaching $`(\alpha _0,0)`$. Otherwise it means that the scalar field tends towards a constant. In this last case, we have $`\varphi \varphi ^{}=cte`$ and (8) shows that $`\omega \mathrm{}`$. If we put $`\varphi =\varphi ^{}`$ in the field equations (2), the metric functions are written $`a=e^{\alpha _0\varphi ^1(\tau \tau _0)}`$, $`b=e^{\beta _0\varphi ^1(\tau \tau _0)}`$, $`c=e^{\gamma _0\varphi ^1(\tau \tau _0)}`$, and become in the proper time $`a=a_0t^{p_1}`$ , $`b=b_0t^{p_2}`$ , $`c=c_0t^{p_3}`$ with $`p_i=1`$ and $`p_i^2=12^1A^2(\alpha _0+\beta _0+\gamma _0)^2`$. Hence, when $`(\alpha ,\mathrm{\Phi })(\alpha _0,0)`$, the metric functions tend towards a Kasnerian behaviour.
We can make the following general observations valid in the $`\tau `$ time: when $`\mathrm{\Phi }[2\alpha _0,0]`$, the more increasing (decreasing) the scalar field is, the more decreasing (increasing) the metric function will be. When $`\mathrm{\Phi }[2\alpha _0,0]`$, the scalar field and the metric function increase (decrease) if $`\alpha _0>0`$ ($`\alpha _0<0`$).
The last remark will concern the representation, in the $`(\alpha ,\mathrm{\Phi })`$ plane, of the solutions of the first equation in (7). If we take as a convention that $`\sqrt{3+2\omega }>0`$, equation (8) shows that the sign of $`\varphi ^,=\mathrm{\Phi }`$ depends on the sign of the integration constant $`A`$. Hence the solution represented in figure 1 by the straight line is physically composed of two separate solutions represented by two half-line, one corresponding to $`A>0`$ and then $`\mathrm{\Phi }>0`$ and the other to $`A<0`$ and then $`\mathrm{\Phi }<0`$. So, to the first equation of (2) correspond four types of behaviours for the metric function and the scalar field, depending on the sign of $`\alpha _0`$ and $`A`$. We will see below that each of them can be split again in two cases depending on the sign of $`\mathrm{\Phi }^,=\varphi ^{,,}`$. These four solutions are illustrated in figure 2. In this figure, {1}, {2}, {3}, {4} correspond to the four half-lines which represent the four physically different solutions of the first equation of (2). $`(\tau _1)`$ and $`(\tau _2)`$ represent the finite or infinite values of the time $`\tau `$ for which $`(\alpha ,\mathrm{\Phi })`$ is equal to $`(0,2\alpha _0)`$ and $`(\alpha _0,0)`$. In what it follows, we will consider the motion of a point $`(\alpha ,\mathrm{\Phi })`$ on each of the four half-lines. It depends on the form of the coupling function $`\omega (\varphi )`$. To determine it, we need an equation to know how and under which conditions $`\mathrm{\Phi }`$ varies.
## 4 Study of the metric functions and scalar field variations depending on the form of $`\omega (\varphi )`$
We have $`d\tau =abcdt`$ with $`abc>0`$. Hence $`\tau `$ is an increasing function of $`t`$ and the variations of the metric functions in the $`\tau `$ time will be the same in the $`t`$ time. From (7), we deduce the equation which gives the variation of $`\mathrm{\Phi }`$ depending on $`\omega (\varphi )`$ :
$$\mathrm{\Phi }^,=\frac{\omega _\varphi (\varphi ^,)^2}{3+2\omega }$$
(14)
with $`\omega _\varphi =\omega ^,/\varphi ^,=d\omega /d\varphi `$. $`3+2\omega `$ is positive since $`\omega >3/2`$. Then, the sign of $`\mathrm{\Phi }^,`$ depends on the sign of $`\omega _\varphi `$ which is independent of the time we consider, namely $`\tau `$ or $`t`$ (of course $`\mathrm{\Phi }^,=\varphi ^{,,}`$ and the sign of $`\ddot{\varphi }`$ will be different in the $`t`$ time. But this is not important here since our final aim is to determine the sign of the first derivatives of the metric functions which does not change in $`t`$ time). So the results we will find and which depend on the sign of the variations of $`\mathrm{\Phi }`$ will be valid in both $`t`$ and $`\tau `$ times. Hence, if $`\omega _\varphi `$ has a constant sign, the motion of the point $`(\alpha ,\mathrm{\Phi })`$ on each half-line will be monotone otherwise its direction will change depending on the sign of $`\omega _\varphi `$. We now study the case where $`\omega (\varphi )`$ is a monotone function and get eight different behaviours for the scalar field and the metric function corresponding to the split of each of the 4 previous cases in two cases. First, we consider that the coupling function is an increasing function of the scalar field. Then, $`\omega _\varphi >0`$ and from (14) we deduce that $`\mathrm{\Phi }^,=\varphi ^{,,}<0`$. Consequently, the motion of the point $`(\alpha ,\mathrm{\Phi })`$ on the half-lines will be such that $`\mathrm{\Phi }`$ decrease. Then, if we are on the half-line {1}, the point $`(\alpha ,\mathrm{\Phi })`$ moves from the left to the right. In the same time, $`\tau `$ increases and then we deduce that $`\tau _1<\tau _2`$. On {1} we have $`\mathrm{\Phi }=\varphi ^,>0`$ : the scalar field is an increasing function of $`\tau `$. When $`\mathrm{\Phi }+\mathrm{}`$, $`\alpha <0`$. $`\alpha `$ remains negative until $`(\alpha ,\mathrm{\Phi })=(0,2\alpha _0)`$, which means $`\tau =\tau _1`$, and when $`\mathrm{\Phi }[0,2\alpha _0]`$, $`\alpha `$ becomes positive. So, we deduce that the metric function is first decreasing until $`\tau =\tau _1`$ and then increases when $`\tau >\tau _1`$ until $`\tau =\tau _2`$, the value of $`\tau `$ for which the scalar field becomes a constant: the metric function can have a minimum (but it is not necessarily true as we will see below). The same type of reasoning can be applied when we consider the half-lines {2}, {3} and {4}.
If now we consider that the coupling function is a decreasing function of the scalar field, we have $`\omega _\varphi <0`$ and $`\mathrm{\Phi }^,=\varphi ^{,,}>0`$. The point $`(\alpha ,\mathrm{\Phi })`$ moves from the right to the left on each of the four half-lines and we have $`\tau _2<\tau _1`$. The same reasoning as in the case $`\omega _\varphi >0`$ will hold. Hence we get four more cases. Table 1 summarises these eight cases : we give the sign of the triplet $`(\omega _\varphi ,A,\alpha _0)`$, independent of the time we consider ($`t`$ or $`\tau `$), the scalar field and metric function variations, the direction of the motion of the point on each half-line and we allocate a number for each behaviour.
Another condition has to be fulfilled in the cases {1}, {1’}, {4}, {4’}, to have necessarily an extremum: we have to check if the value $`\mathrm{\Phi }=2\alpha _0`$ belongs to the interval in which $`\mathrm{\Phi }`$ varies. For this purpose, we rewrite the equation (8) :
$$A\mathrm{\Phi }\sqrt{3+2\omega }=1$$
(15)
We determine the interval in which the scalar field $`\varphi `$ varies by imposing the conditions $`\sqrt{3+2\omega }>0`$ and $`\varphi >0`$. Then from (15) we deduce the interval for $`\mathrm{\Phi }`$. The condition for an extremum to exist for the behaviours of type {1}, {1’}, {4} and {4’} will be that this last interval contains the value $`2\alpha _0`$. One can also check if the value of the scalar field corresponding to $`3+2\omega =(2\alpha _0A)^2`$ beholds to the interval in which $`\varphi `$ varies.
Now, we consider the case where the coupling function $`\omega (\varphi )`$ is not a monotone function of the scalar field. It means that the sign of $`\omega _\varphi `$ will change during the evolution of the dynamic. In the interval of time where $`\omega _\varphi `$ will be positive, we will have behaviours of type {1}, {2}, {3} or {4} and when it becomes negative the metric function and the scalar field will behave respectively as {1’}, {2’}, {3’} or {4’}. Hence, the behaviours of the metric function when the coupling function is not monotone will be a succession of behaviours of type {i}+{i’}+{i}+{i’}…, the repetitions of the scheme {i}+{i’} depending on the number of zero of $`\omega _\varphi `$.
Note that to achieve our goal, that is determine the variation (sign of the first derivative) of the metric function, we used quantities such that the second derivative of the scalar field or the amplitude of its first derivative are not invariant when we change time coordinate from $`\tau `$ to $`t`$. But these two quantities can always be written as function of $`\omega _\varphi `$ or $`\omega `$ which are independent of time coordinate. Therefore our method is in agreement with the fact that the sign of the first derivative of the metric function is the same in $`\tau `$ or $`t`$ time.
In the next section we will consider several forms of the coupling function with a decreasing scalar field, i.e. $`A<0`$.
## 5 Applications.
We are going to examine the variations of the metric functions with three different forms of the coupling function. The couplings we will consider are interesting for the following reasons. The first coupling is $`3+2\omega =\varphi _c^2\varphi ^{2m}`$. When $`m>0`$ and $`\varphi \mathrm{}`$ or $`m<0`$ and $`\varphi 0`$, $`\omega \varphi ^{2m}\mathrm{}`$. When $`m<1/4`$ and $`\varphi 0`$ or when $`m>1/4`$ and $`\varphi \mathrm{}`$, $`\omega _\varphi \omega ^30`$. Hence, asymptotically, the theory tends towards relativistic behaviours at late time $`(\varphi 0)`$ when $`m<1/4`$. When the scalar field becomes infinite, $`\omega (\varphi )`$ tends towards a power law that corresponds to a power or exponential law for $`F(\phi )`$ (see (1)). Power Law for $`\omega (\varphi )`$ have been studied in . This class of theories is also in agreement with the constraints imposed by the slow logarithmic decrease of the gravitational constant $`(dG/dt)G^1`$. The two other laws, $`2\omega +3=m\mathrm{ln}(\varphi /\varphi _0)^n`$ and $`2\omega +3=m1(\varphi /\varphi _0)^n^1`$ have been studied in in a FLRW Universe. For the first one, we recover the values of the PPN parameters in General Relativity when $`\varphi \varphi _0`$ if $`n>1/2`$, whereas for the second one there is no restriction on the value of the exponent $`n`$.
### 5.1 The theory $`3+2\omega =\varphi _c^2\varphi ^{2m}`$
We have :
$$\omega _\varphi =\varphi _c^2m\varphi ^{2m1}$$
(16)
The expression $`3+2\omega `$ is positive for all positive values of the scalar field. Hence $`\varphi `$ varies in $`[0,+\mathrm{}[`$. From (15) we deduce that $`\mathrm{\Phi }`$ varies in $`]\mathrm{},0]`$. If $`m`$ is positive, $`\omega _\varphi >0`$ and the metric function behaves as {2} and {4} whereas if $`m`$ is negative, $`\omega _\varphi <0`$, and it behaves as {2’} and {4’}. In the Cases {2} and {2’}, the metric function increases. In the case {4} and {4’}, from (15) we deduce that the metric function has an extremum when the scalar field is equal to $`(2\alpha _0A\varphi _c)^{1/m}`$. This last value is always positive and then belongs to the interval in which the scalar field varies. We conclude that for the types {4} or {4’}, the metric function will always have respectively a minimum or a maximum.
### 5.2 The theory $`2\omega +3=mln\varphi /\varphi _0^n`$.
We restrict the parameters to $`n>0`$, $`m>0`$ so that $`2\omega +3`$ is positive. We will first consider the case where $`\varphi >\varphi _0`$ . Then, we can write:
$$2\omega +3=m(ln\varphi /\varphi _0)^n$$
(17)
$`\omega _\varphi `$ is always negative and $`\mathrm{\Phi }]\mathrm{},0]`$. Hence, if $`\alpha _0>0`$, the metric function is increasing. If $`\alpha _0<0`$, the metric function will always have a maximum since $`\mathrm{\Phi }=2\alpha _0`$ belongs to the interval where $`\mathrm{\Phi }`$ varies.
If we chose for $`\varphi `$ the interval $`[0,\varphi _0]`$, the metric function has a minimum if $`\alpha _0<(2A\sqrt{m})^1`$. Otherwise, it is increasing.
### 5.3 The theory $`2\omega +3=m1(\varphi /\varphi _0)^n^1`$.
We restrict the parameters to $`n>0`$, $`m>0`$ and will take first $`\varphi >\varphi _0`$. Hence we have:
$$2\omega +3=m\left[(\varphi /\varphi _0)^n1\right]^1$$
(18)
$`\omega _\varphi `$ is always negative. If the integration constant $`\alpha _0`$ is positive, the metric function is increasing, whereas if $`\alpha _0`$ is negative, since $`\mathrm{\Phi }]\mathrm{},0]`$, the metric function will always have a maximum. If we choose $`\varphi [0,\varphi _0]`$, the metric function is still increasing when $`\alpha _0>0`$ but have a minimum if $`\alpha _0<0`$.
## 6 Behaviour of the three metric functions.
The graph 3 represents the solutions of the system equations (2) on the plane $`((\alpha ,\beta ,\gamma ),\mathrm{\Phi })`$. We choose without loss of generality $`\alpha _0<\beta _0<\gamma _0`$. We distinguish four cases :
1. If $`\mathrm{\Phi }>2\gamma _0`$, all the metric functions are decreasing.
2. If $`\mathrm{\Phi }[2\gamma _0,2\beta _0]`$, the metric function associated with the largest of the integration constants is increasing whereas the two others are still decreasing.
3. If $`\mathrm{\Phi }[2\beta _0,2\alpha _0]`$, the metric function associated with the smallest of the integration constants is the only one to be decreasing.
4. If $`\mathrm{\Phi }<2\alpha _0`$, the three metric functions are increasing.
If i constants among $`\alpha _0`$, $`\beta _0`$ and $`\gamma _0`$ are positive, we deduce from figure 3 that when $`\varphi `$ is increasing, whatever the form of $`\omega (\varphi )`$, only the i+1 first cases can exist, when $`\varphi `$ is decreasing, whatever the form of $`\omega (\varphi )`$, only the i+1 last cases can exist. Hence, in the case where $`\alpha _0`$, $`\beta _0`$, $`\gamma _0`$ are positive constant and $`A`$ is a negative one, all the metric functions will be increasing whatever the form of $`\omega (\varphi )`$. But, if $`\alpha _0`$, $`\beta _0`$, $`\gamma _0`$ are negative and $`A`$ positive, all the metric functions will be decreasing. We deduce also that to get three increasing metric functions which tend towards a power law, that is $`((\alpha ,\beta ,\gamma ),\mathrm{\Phi })((\alpha _0,\beta _0,\gamma _0),0)`$, when $`\tau `$(and thus $`t`$) increases, a necessary condition will be that $`\alpha _0`$, $`\beta _0`$, $`\gamma _0`$ be positive , $`A`$ and $`\omega _\varphi `$ have the same sign.
## 7 Study of the second-derivative of the metric function
In the FLRW models, a positive sign of the first and second derivatives of the scale factor with respect to the cosmic time is the sign of inflation: the expansion in the $`t`$ time is accelerated. Inflation in generalised scalar-tensor theory and in FLRW models has been studied in and . It seems to be noteworthy that it happens without a cosmological constant or potential. One can talk about inflation only when the second derivatives of the metric functions with respect to $`t`$ are positives. First, we are going to describe a method giving the sign of the second derivative of the metric function with respect to $`\tau `$ from the knowledge of $`\omega `$ and $`\omega _\varphi `$. Hence, we will be able to completely determine the qualitative form of the metric function in the $`\tau `$ time. Second, we apply it and finally we will study the sign of $`\ddot{a}`$ and obtain conditions to have inflation in Bianchi type I model.
### 7.1 Study of $`a^{,,}`$
The first spatial component of the field equations is written :
$$\frac{a^{,,}}{a}=\frac{a^{,2}}{a^2}\frac{a^,}{a}\frac{\varphi ^,}{\varphi }+\frac{1}{2}\frac{\omega ^,}{3+2\omega }\frac{\varphi ^,}{\varphi }$$
$$\varphi ^2\frac{a^{,,}}{a}=\alpha ^2\alpha \varphi ^,+\frac{1}{2}\frac{\omega _\varphi }{3+2\omega }\varphi ^{,2}\varphi $$
But $`\varphi ^,=1/(A\sqrt{3+2\omega })`$, so we get :
$$\varphi ^2\frac{a^{,,}}{a}=\alpha ^2\frac{\alpha }{A\sqrt{3+2\omega }}+\frac{1}{2}\frac{\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}$$
(19)
The sign of the left hand side of (19) is the same as $`a^{,,}`$. The right hand side of equation (19) is an equation of degree two in $`\alpha `$. Hence, we have to know the sign of this equation in order to obtain the sign of $`a^{,,}`$, i.e. to determine its roots. It is important to recall that $`\alpha `$ can be expressed as a function of the scalar field. We get :
$$\alpha =\alpha _0\frac{1}{2}\varphi ^,=\alpha _0\frac{1}{2}\frac{1}{A\sqrt{3+2\omega }}$$
(20)
Now we calculate the determinant of the second degree equation (19) :
$$\mathrm{\Delta }=\frac{1}{A^2(3+2\omega )}2\frac{\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}$$
(21)
If $`\mathrm{\Delta }`$ is negative, the second degree equation is positive for all value of $`\alpha `$ and $`a^{,,}`$ is positive. Then the dynamic of the metric function is accelerated (this is not inflation since the sign of $`a^{,,}`$ and $`\ddot{a}`$ are not necessarily the same). If $`\mathrm{\Delta }`$ is positive, the second degree equation has two real roots $`\alpha _1`$ and $`\alpha _2`$. From (21), we deduce that $`\mathrm{\Delta }<0`$ if :
$$\omega _\varphi >\frac{3+2\omega }{2\varphi }$$
(22)
The condition (22) will be true for the three metric functions. It does not depend on a specific parameter of one of these functions. Hence, when (22) is true, the dynamic of the three metric functions in the $`\tau `$ time is accelerated. If now we consider $`\mathrm{\Delta }>0`$, we find two roots :
$$\alpha _{1,2}=(\frac{1}{A\sqrt{3+2\omega }}\pm \sqrt{\frac{1}{A^2(3+2\omega )}\frac{2\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}})/2$$
(23)
With the form of the coupling function, one can deduce the conditions so that $`a^{,,}`$ be positive or negative. By conditions we mean the values of the scalar field and of the different parameters defining the form of the coupling function, which rule the sign of $`a^{,,}`$. To get this sign, we have to know the sign of:
$$\alpha _{1,2}(\varphi )\alpha (\varphi )=\alpha _0+(2A\sqrt{3+2\omega })^1\left[2\pm \sqrt{12\omega _\varphi \varphi (3+2\omega )^1}\right]$$
(24)
When $`\alpha _1\alpha `$ and $`\alpha _2\alpha `$ have the same sign, equation (19) is positive and thus $`a^{,,}`$ is positive; otherwise, it means that $`\alpha [\alpha _2,\alpha _1]`$ and then $`a^{,,}`$ is negative. At late time, if $`\varphi _{RG}`$ is the value of the scalar field for which $`\omega \mathrm{}`$ and $`\omega _\varphi \omega ^30`$ (which ensures the theory is compatible with the observation) we deduce from (24) that a necessary and sufficient condition for the dynamic of the metric function to be decelerated in the $`\tau `$ time, will be:
$$\underset{\varphi \varphi _{RG}}{lim}\omega _\varphi <2\alpha _0^2A^2(3+2\omega )^2\varphi ^1$$
(25)
### 7.2 Applications.
#### 7.2.1 Theory $`3+2\omega =\varphi _c^2\varphi ^{2m}`$
Remember that for this form of $`3+2\omega `$ we have $`\varphi [0,+\mathrm{}[`$. We continue to choose $`A<0`$ in order to have a decreasing scalar field. We get :
$$\alpha =\alpha _0\frac{1}{2}\frac{\varphi ^m}{A\varphi _c}$$
(26)
$$\alpha _{1,2}=\frac{\varphi ^m(1\pm \sqrt{12m})}{2A\varphi _c}$$
(27)
The condition (22) is satisfied when $`m>1/2`$ : in this case we always have $`a^{,,}`$, $`b^{,,}`$ and $`c^{,,}`$ positive. When $`m<1/2`$, we have to determine the sign of :
$$\alpha _{1,2}\alpha =\frac{\varphi ^m(2\pm \sqrt{12m})}{2A\varphi _c}\alpha _0$$
(28)
We will always have $`\alpha _1<\alpha _2`$.
\- If $`\alpha _0=0`$, we have $`\alpha >\alpha _1`$ for all values of the scalar field. If $`m<3/2`$, from equation (28) we deduce that $`\alpha _2<\alpha <\alpha _1`$ and thus $`\alpha ^{,,}<0`$. If $`m[3/2,1/2]`$, we get $`\alpha >\alpha _{1,2}`$ and then $`a^{,,}>0`$. Now we consider general case where $`\alpha _00`$.
* If $`m<0`$,
+ if $`\alpha _0>0`$, when $`\varphi \mathrm{}`$, $`\alpha >\alpha _1`$. If $`m[3/2,0]`$, $`\alpha >\alpha _{1,2}`$ and if $`m<3/2`$, $`\alpha [\alpha _1,\alpha _2]`$. Then the scalar field decreases and when $`\varphi 0`$, $`\alpha >\alpha _{1,2}`$.
+ If $`\alpha _0<0`$, when $`\varphi \mathrm{}`$, if $`m[3/2,0]`$, $`\alpha >\alpha _{1,2}`$, if $`m<3/2`$, $`\alpha [\alpha _1,\alpha _2]`$. When the scalar field decreases and $`\varphi 0`$, $`\alpha <\alpha _{1,2}`$.
Hence, we deduce that :
+ If $`\alpha _0>0`$,
- if $`m[3/2,0]`$, we have $`a^{,,}>0`$,
- if $`m<3/2`$, we have first $`a^{,,}<0`$ and then $`a^{,,}>0`$.
+ If $`\alpha _0<0`$,
- if $`m[3/2,0]`$, we have $`a^{,,}>0`$, then $`a^{,,}<0`$ and finally $`a^{,,}>0`$,
- if $`m<3/2`$, we have $`a^{,,}<0`$ and $`a^{,,}>0`$.
* If $`m[0,1/2]`$,
We will always have $`\varphi ^m(21\sqrt{12m})>0`$. When $`\varphi \mathrm{}`$, $`\alpha `$ is larger than $`\alpha _{1,2}`$ if $`\alpha _0>0`$ or smaller if $`\alpha _0<0`$. For all value of $`\alpha _0`$, when $`\varphi `$ decreases and tends towards 0, we have $`\alpha >\alpha _{1,2}`$.
Hence, we deduce that if $`\alpha _0<0`$, first we have $`a^{,,}>0`$, then $`a^{,,}<0`$ and at last $`a^{,,}>0`$. If $`\alpha _0>0`$, we always have $`a^{,,}>0`$.
From the knowledge of $`a^,`$ (see 5.2) and $`a^{,,}`$ it is now easy to know qualitatively the behaviours of the metric function $`a`$, depending on its different parameters $`\alpha _0`$ and $`m`$. We deduce from our qualitative analysis that:
* When $`m[0,1/2]`$ and $`\alpha _0>0`$, the metric function is increasing and accelerated. When $`\alpha _0<0`$, the metric function has a minimum. The branch before the minimum is accelerated whereas the branch after the minimum has an inflexion point and is accelerated in late time.
* When $`m>1/2`$, the dynamic of the metric function is always accelerated.
* When $`m<0`$ and $`\alpha _0>0`$, the metric function increases. It is accelerated if $`m[3/2,0]`$. If $`m<3/2`$, it is first decelerated and then accelerated: the metric function has an inflexion point. If $`\alpha _0<0`$, the metric function has a maximum. If $`m[3/2,0]`$, the dynamic is accelerated in both late and early times whereas if $`m<3/2`$, it is decelerated in early time and accelerated in late time.
Note that one can always obtain the value of the scalar field for which the sign of $`a^{,,}`$ changes by writing $`\alpha _{1,2}\alpha =0`$. We see that the theory $`3+2\omega =\varphi _c^2\varphi ^{2m}`$ is always accelerated in late time in accordance with the relation (25).
#### 7.2.2 The theory $`2\omega +3=m\mathrm{ln}\varphi /\varphi _0^n`$.
Here, we consider only the interval $`[\varphi _0,\mathrm{}[`$ for the scalar field, $`\omega _\varphi `$ is always negative and then $`\mathrm{\Delta }`$ is always positive. We have:
$$\alpha _{1,2}\alpha =\alpha _0+(2A\sqrt{m})^1(\mathrm{ln}\varphi /\varphi _0)^{n/2}(2\pm \sqrt{1+n\varphi _0\mathrm{ln}(\varphi /\varphi _0)^1})$$
(29)
When $`\alpha _0>0`$, in early time, $`\varphi \mathrm{}`$ and $`\alpha >\alpha _{1,2}`$. Then, at late time, when $`\varphi \varphi _0`$, if $`n>1`$, we have again $`\alpha >\alpha _{1,2}`$ and then the metric function increases and is accelerated whereas if $`n[0,1]`$, we have $`\alpha [\alpha _1,\alpha _2]`$. Then, the metric function increases but have an inflexion point. It is decelerated at late time.
When $`\alpha _0<0`$, the metric function has a maximum. If $`n>1`$, the dynamic is both accelerated in early and late time whereas if $`n[0,1]`$, it is just accelerated in early time.
#### 7.2.3 The theory $`2\omega +3=m1(\varphi /\varphi _0)^n^1`$.
Here again we consider the same interval for $`\varphi `$ and $`\mathrm{\Delta }`$ will be always positive. We have:
$$\alpha _{1,2}\alpha =\alpha _0+(2A\sqrt{m})^1\sqrt{(\varphi /\varphi _0)^n1}(2\pm \sqrt{1+n(\varphi /\varphi _0)^n/\left[(\varphi /\varphi _0)^n1\right]})$$
(30)
We get two important values for $`n`$: $`n=3`$ or $`n=4A^2\alpha _0^2m`$.
* When $`\alpha _0>0`$, the metric function is increasing and its behaviour is accelerated if $`n<(3,4A^2\alpha _0^2m)`$ or decelerated if $`n>(3,4A^2\alpha _0^2m)`$. If the value of $`n`$ is between $`n=3`$ and $`n=4A^2\alpha _0^2m`$, the metric function has an inflexion point and the dynamic will be accelerated at late time if $`3<4A^2\alpha _0^2m`$ or decelerated if $`3>4A^2\alpha _0^2m`$.
* When $`\alpha _0<0`$, the metric function has a maximum. Its behaviour is decelerated if $`n>(3,4A^2\alpha _0^2m)`$. If $`n<(3,4A^2\alpha _0^2m)`$ , the dynamic is accelerated at both late and early times. If the value of $`n`$ is between $`n=3`$ and $`n=4A^2\alpha _0^2m`$, the dynamic is decelerated at early time when $`3<4A^2\alpha _0^2m`$ and becomes accelerated whereas when $`3>4A^2\alpha _0^2m`$, it is first accelerated and then decelerated at late time.
In all the applications one can prove that the behaviours of $`a^{,,}`$ at early and late times are continuous. The sign of $`a^{,,}`$ does not change between the late and early times because $`(\alpha _{1,2}\alpha )^,`$ vanish for only one value of $`\varphi `$ in the intervals in which the parameters of the three theories and the scalar field are allowed to vary. If it was not the case, the sign of this last expression would vanish for, at least, two values of the scalar field.
In the next subsection we will talk about the second derivative of the metric function in $`t`$ time. For the sake of simplicity (the sign of the second derivative can change more than twice in $`t`$ time) we will not study the behaviour of these theories in the $`t`$ time (qualitatively, only the sign of the second derivative changes). Moreover, to do this we must carry out numerical computations as we will see, that seems diverge from our goal, i.e. make a general study of the dynamic whatever the coupling function.
### 7.3 Study of $`\ddot{a}`$.
Here, when $`\ddot{a}`$ and the first derivative are positives one can speak about inflation. We have:
$$\frac{\ddot{a}}{a}=\left[\frac{a^{,,}}{a}\frac{a^{,2}}{a^2}\frac{a^,}{a}(\frac{b^,}{b}+\frac{c^,}{c})\right](abc)^2$$
The relations (8) and (19) imply:
$$\frac{\ddot{a}}{a}(abc)^2\varphi ^2=\frac{1}{2}\frac{\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}\alpha (\beta _0+\gamma _0)$$
(31)
This is an equation of first degree for $`\alpha `$. Its solution is:
$$\alpha _3=\frac{1}{2}\frac{\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}(\beta _0+\gamma _0)^1$$
We use equation (20) to write:
$$\alpha \alpha _3=\alpha _0\frac{1}{2}\frac{1}{A\sqrt{3+2\omega }}\frac{1}{2}\frac{\omega _\varphi }{(3+2\omega )^2}\frac{\varphi }{A^2}(\beta _0+\gamma _0)^1$$
(32)
Then, one has to solve $`\alpha \alpha _3=0`$ for $`\varphi `$ so that we can determine the sign of this last expression for different intervals of the scalar field. This is not an easy task and to study the theories of the last subsection, we would need numerical computation. In a general manner, to simplify the resolution, one can notice that equation (32) is a third degree equation for $`(3+2\omega )^{1/2}`$. Then, $`\ddot{a}`$ is positive when $`\beta _0+\gamma _0>0`$ ($`<0`$) if $`\alpha \alpha _3>0`$ ($`<0`$) and negative when $`\beta _0+\gamma _0>0`$ ($`<0`$) if $`\alpha \alpha _3<0`$ (\>0). When a theory tends toward General Relativity, i.e. $`\varphi \varphi _{RG}`$, the dynamic of the metric function will be decelerated if:
$$\underset{\varphi \varphi _{RG}}{lim}\omega _\varphi <2A^2\alpha _0(\beta _0+\gamma _0)(3+2\omega )^2\varphi ^1$$
(33)
Under this condition one can not get inflation at late time. Note that (33) has the same form as (25) except the introduction of the constant $`\beta _0+\gamma _0`$. This comes from the fact that in the $`t`$ time, all the metric functions appear in each field equations. If we use the three coupling functions of subsection 7.2 with equation (32), one obtain complex expressions which need numerical investigations to find their zeros.
Since the presence of matter tends to slow down the expansion, one can hypothesize that (33) could be a sufficient (but not necessary) condition so that model with matter has a decelerated behaviour in the same circumstances, that is at late time when the theory tends towards a relativistic behaviour.
## 8 Conclusions.
From the form of the coupling function $`\omega (\varphi )`$, we can deduce the qualitative behaviour of the metric functions. It depends on the sign of $`d\varphi /d\tau `$, $`d\omega /d\varphi `$ and the integration constants $`\alpha _0`$, $`\beta _0`$, $`\gamma _0`$. We have studied two things : sign of the first and second derivatives of the metric functions.
For the first derivative, the main difficulty is to find the zeros of $`\omega _\varphi `$. When $`\omega (\varphi )`$ is a monotonous function of the scalar field, we have eight basic possible behaviours ({1}, {2}, {3}, {4}, {1’}, {2’}, {3’}, {4’}) for a metric function because $`d\varphi /d\tau `$, $`d\omega /d\varphi `$ and the corresponding integration constants can be positive or negative (2\*2\*2=8). When $`\omega (\varphi )`$ has one or several extrema, the behaviour of the metric function is a succession of behaviours of types {i} + {i’}, {i} and {i’} being the number of two of the eight basic behaviours, one with $`\omega _\varphi >0`$ and the other with $`\omega _\varphi <0`$. For the behaviours of type {1}, {1’}, {4} and {4’}, a complementary condition has to be fulfilled so that the metric function $`a`$ ($`b`$, $`c`$) has an extremum : the value $`2\alpha _0`$ ($`2\beta _0`$, $`2\gamma _0`$) has to be in the interval in which $`d\varphi /d\tau `$ varies otherwise the metric function is monotone. Or equivalently, a time independent formulation of this condition will be that the value of the scalar field corresponding to $`3+2\omega =(2\alpha _0A)^2`$ ($`(2\beta _0A)^2`$, $`(2\gamma _0A)^2`$) have to belong to the interval in which $`\varphi `$ varies.
For the second derivative of the metric functions in the $`\tau `$ time, if the condition (22) is fulfilled, the dynamic of the metric functions is always accelerated. If it is not the case, we have to examine, for the metric function $`a`$ for instance, the sign of $`\alpha _1\alpha `$ and $`\alpha _2\alpha `$. If these expressions have the same sign, the second derivative of $`a`$ is positive otherwise it is negative.
In the $`t`$ time, the dynamic is accelerated if (31) is positive and decelerated otherwise. If moreover, the first derivative is positive, we have inflation.
With this method we have been able to completely determine, whatever $`\tau `$, the qualitative form of the metric functions for three different theories. Each of them can be compatible with the value of the PPN parameters at late time if we adjust their parameters. By using the results of subsection 7.3 concerning the sign of the second derivative in the cosmic time and numerical calculations, it is also possible to obtain the qualitative form of the metric functions in the $`t`$ time.
Moreover, if with $`\omega +\mathrm{}`$ and $`\omega _\varphi \omega ^30`$, we want the three metric functions to be increasing and decelerated at late time in the cosmic time, we deduce of the study that we must have: $`(\alpha _0,\beta _0,\gamma _0)>0`$ and $`A`$ and $`\omega _\varphi `$ must have the same sign, which is positive since $`\omega +\mathrm{}`$ and $`\omega _\varphi <2A^2inf\text{[}\alpha _0(\beta _0+\gamma _0)`$, $`\beta _0(\alpha _0+\gamma _0)`$, $`\gamma _0(\alpha _0+\beta _0)\text{]}`$$`(3+2\omega )^2\varphi ^1`$ when $`\varphi `$ tends towards $`\varphi _{RG}`$, $`\varphi _{RG}`$ being the smallest value of the scalar field. In these conditions the metric functions have a power law form.
In section 6, we have determined the conditions to have 1, 2 or 3 increasing metric functions; in fact, this is a graphic translating of some information contained in the constraint equation of the field equations.
We have studied the simplest anisotropic cosmological model but we hope to extend this method to more complicated ones such as Bianchi types II and V and in more complex situations, i.e. with cosmological constant or potential. The main advantage of such study is to reveal completely the dynamic of the metric functions whatever the form of the coupling function and not only for a particular one or for asymptotic behaviour.
Figure 1 : solution of the first equation of (2) in the $`\alpha ,\mathrm{\Phi }`$ plane depending on the sign of $`\alpha _0`$.
Figure 2 : the four different physically solutions of the first equation of (2).
Figure 3 : representation of all the solutions of the equations (2) in the $`((\alpha ,\beta ,\gamma ),\mathrm{\Phi })`$ plane. |
warning/0001/hep-ph0001318.html | ar5iv | text | # The Nonperturbative Color Meissner Effect in a Two-Flavor Color Superconductor
## 1 Introduction
It has recently been demonstrated that the color symmetry of Quantum Chromodynamics (QCD) might be spontaneously broken through the formation of a diquark condensate. This phase is metastable at zero and small matter density but becomes stable and replaces the usual chiral-broken phase at some critical density .
Condensation of diquarks is analogous to Cooper pairing of electrons in BCS theory, in which the massive photon is the microscopic embodiment of the Meissner effect. In QCD the gauge group is $`SU(3)`$, but the consequences are similar: gauge bosons (gluons) become massive due to their interaction with the correlated fermions (quarks). Not every gluon species becomes massive, since the symmetry breaking pattern is $`SU(3)SU(2)`$ and the three gluons of the residual $`SU(2)`$ remain massless. Among the five massive modes, four are degenerate and the mass of the fifth is $`\sqrt{4/3}`$ times that of the others. This is comparable to the symmetry-broken phase of the electroweak theory. Thus the direct analogy is to the Higgs mechanism, with the diquark field forming a composite scalar and its nonzero VEV generating a dynamical Higgs effect.
In this paper this mechanism is analyzed microscopically. As the source of diquark condensation we will use the effective quark action obtained by averaging over instanton configurations . In this approach to the QCD vacuum, all information from the large background gauge fields (instantons) is encoded in the non-local form of the ’t Hooft interaction and the two instanton parameters, average size ($`\rho `$) and number density ($`N/V`$). This procedure replaces dynamical gluons with classical, nonperturbative field solutions, ignoring quantum gluonic fluctuations. The reintroduction of gluonic fluctuations will require a gauge-invariant perturbative modification of the ’t Hooft vertex. A direct consequence of gauge invariance is the transverse polarization operator, which we will compute explicitly. This is the main result of the paper as from it we obtain the gluon masses.
Since the dynamical Higgs effect is interesting regardless of density, we will first consider the vacuum alone (chemical potential $`\mu =0`$), although in this case the superconductor is only metastable . The gluon masses are found to be proportional to $`g\mathrm{\Delta }`$ where $`\mathrm{\Delta }`$ is the superconducting gap, the scale of which is a nonperturbative one set by the strength of the instanton background. Extending the formalism to finite chemical potential $`\mu `$ is straightforward, given that the density dependence of the propagators and effective action are known . However, in this paper we avoid exact calculations at $`\mu 0`$. Instead, for simplicity we assume that the in-medium gluon masses are determined by the behavior of $`g\mathrm{\Delta }(\mu )`$ up to and somewhat above the critical value of $`\mu 300`$ MeV. At this density, the point at which the color superconductor becomes the stable phase, we find that the gluon masses are about 120 MeV.
For gluons the color-breaking Meissner masses are a primary ramification of the superconducting state, however this is not the only effect of a quark medium. There will be additional contributions, among them a Debye screening mass of the order of $`g\mu `$, which do not break color symmetry. When the density becomes large, the Debye mass increases and instantons are screened out of the picture. Yet diquark condensation persists, now due to perturbative gluon exchange , and at asymptotically large density the Meissner masses are also proportional to $`g\mu `$ . So while one still has quark pairing, any matching between low and high density mechanisms remains unclear.
In this paper we consider the case of two massless flavors, which is expected to be relevant at finite-density chiral restoration given a relatively large strange quark mass . Instantons are Euclidean pseudoparticles, and thus all calculations will be in Euclidean space.
In Section 2 we recall those features of the ordinary Higgs mechanism (based on elementary scalar fields) which will also be relevant for composite scalars. In Section 3 the instanton-induced action including perturbative gluons will be formulated. In Section 4 the diquark gap equation is reviewed, and in Section 5 the color current is determined. In Section 6 we sum a set of quark-quark correlation functions to recover the Nambu-Goldstone modes of the theory, a necessary exercise as these will mix with the longitudinal gluons. These ingredients are assembled into the gluon polarization operator in Section 7, which taken in the static limit corresponds to a mass. Section 8 compares the result to chiral symmetry breaking, and in Section 9 this result is discussed and conclusions are drawn.
## 2 Higgs Mechanism
The dynamical Higgs mechanism, though technically more involved, does not differ much from the ordinary one. Let us denote the quark bilinear combination
$$\varphi ^\alpha =ϵ_{\alpha \beta \gamma }ϵ_{fg}ϵ_{ij}\left(\psi _L^{\beta fi}\psi _L^{\gamma gj}+(LR)\right),$$
(1)
where Greek letters denote color, $`f,g=1,2`$ flavor, and $`i,j=1,2`$ are spinor indices of the left ($`L`$) and right ($`R`$) components of the quark field $`\psi `$. The resulting complex field $`\varphi ^\alpha `$ is a Lorentz scalar isoscalar field belonging to the fundamental representation of the $`SU(3)`$ group. To support gauge invariance it must couple to the gauge potential via the covariant derivative, $`(_\mu )_\beta ^\alpha =_\mu \delta _\beta ^\alpha igA_\mu ^a(\lambda ^a/2)_\beta ^\alpha `$, in which the $`\lambda ^a`$ denote the eight Gell-Mann matrices. The kinetic energy term for the composite diquark field $`\varphi ^\alpha `$ is the usual
$``$ $`=`$ $`Z_\varphi ^1\left|_\mu \varphi \right|^2`$ (2)
$`=`$ $`Z_\varphi ^1[_\mu \varphi _\alpha ^{}_\mu \varphi ^\alpha +i{\displaystyle \frac{g}{2}}A_\mu ^a(\varphi _\alpha ^{}(\lambda ^a)_\beta ^\alpha _\mu \varphi ^\beta _\mu \varphi _\alpha ^{}(\lambda ^a)_\beta ^\alpha \varphi ^\beta )`$
$`+{\displaystyle \frac{g^2}{4}}A_\mu ^aA_\mu ^b\varphi _\alpha ^{}(\lambda ^a\lambda ^b)_\beta ^\alpha \varphi ^\beta ].`$
The only difference with the standard case of the elementary field is that there is, in principle, a common ‘wave function renormalization’ factor, $`Z_\varphi `$. For elementary fields $`Z_\varphi =1`$ at the tree level, however, it deviates from unity even for the elementary field when one takes into account the perturbative virtual emission of particles. For the fields that are composite from the start there is no reason for $`Z_\varphi `$ to be unity. This quantity is a priori unknown and should be determined from a dynamical calculation, as will follow below. We stress, however, that the relative weights of the three terms in Eq. (2) are fixed by gauge invariance.
If the scalar field $`\varphi ^\alpha `$ develops a nonzero VEV signaling the diquark condensation,
$$\varphi ^\alpha =2\mathrm{\Delta }_0\delta ^{\alpha 3},$$
(3)
(it can be always arranged along the third color axis and made real), then gluons obtain a mass matrix
$$M_{ab}^2=2g^2Z_\varphi ^1\mathrm{\Delta }_0^2(\lambda ^a\lambda ^b)_3^3=\{\begin{array}{cc}0\hfill & a=b=1,2,3\hfill \\ 2g^2Z_\varphi ^1\mathrm{\Delta }_0^2\hfill & a=b=4,5,6,7\hfill \\ \frac{8}{3}g^2Z_\varphi ^1\mathrm{\Delta }_0^2\hfill & a=b=8.\hfill \end{array}$$
(4)
The symmetry breaking pattern is, thus, $`SU(3)SU(2)`$; three gluons corresponding to the unbroken $`SU(2)`$ subgroup remain massless, four gluons obtain masses proportional the Higgs VEV, and the fifth gluon is $`\sqrt{4/3}`$ times as heavy. This relation will be, of course, reproduced for composite Higgs fields as well, as it follows from symmetry considerations alone.
This simple elementary-Higgs model also provides an alternative way to find the gluon masses, which we will generalize to the case of the composite Higgs field. One can compute the gluon polarization operator $`\mathrm{\Pi }_{\mu \nu }^{ab}(q)`$, as seen in Fig. 1. In Fig. 1a we take the linear coupling of the gauge potential $`A_\mu ^a`$, as given by the second term in Eq. (2), and iterate it twice. The intermediate state in this diagram is the would-be Nambu-Goldstone boson, which contributes to the polarization operator
$`\mathrm{\Pi }_{\mu \nu }^{(\mathrm{NG})ab}(q)`$ $`=`$ $`2\left[Z_\varphi ^1g\varphi _\alpha ^{}(\lambda ^a)_\gamma ^\alpha q_\mu \right]\left[{\displaystyle \frac{Z_\varphi }{q^2}}\right]\left[q_\nu Z_\varphi ^1g(\lambda ^b)_\beta ^\gamma \varphi ^\beta \right]`$ (5)
$`=`$ $`2g^2Z_\varphi ^1\mathrm{\Delta }_0^2{\displaystyle \frac{q_\mu q_\nu }{q^2}}(\lambda ^a\lambda ^b)_3^3.`$
This contribution is purely longitudinal.
Fig. 1b is the contact term; it gives a Kronecker delta contribution:
$$\mathrm{\Pi }_{\mu \nu }^{(\mathrm{contact})ab}(q)=2g^2Z_\varphi ^1\mathrm{\Delta }_0^2\delta _{\mu \nu }(\lambda ^a\lambda ^b)_3^3.$$
(6)
Combining the two we get a transverse polarization operator,
$$\mathrm{\Pi }_{\mu \nu }^{(\mathrm{full})ab}(q)=2g^2Z_\varphi ^1\mathrm{\Delta }_0^2\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)(\lambda ^a\lambda ^b)_3^3.$$
(7)
The fact that it is transverse reflects the color current conservation, even in case of a broken symmetry.
In the case of a composite Higgs field the tree diagrams of Fig. 1 will be replaced by loop diagrams, formed by quarks, however, the general setting will be rather similar to this simple case: there will be a pole contribution from the would-be Goldstone bosons in the intermediate state, and there will be a ‘contact’ term. The resulting polarization operator will be transverse, provided one takes a conserved color current.
## 3 Gauge Invariant Effective Action
Following a long tradition of instanton phenomenology, we assume the following:
* Instantons are the dominant nonperturbative contribution to low-energy QCD, and
* Low and high momentum scales are safely separable for quark and gluon fields.
These assumptions are supported a posteriori by instanton-based phenomenology of the vacuum, as reviewed in Ref. . Instantons break chiral symmetry spontaneously and axial $`U(1)`$ anomalously, both at a satisfactory magnitude. Instanton-induced interactions also provide the necessary $`qq`$ attraction for a diquark condensate, nonperturbatively.
Here we use a formulation which not only relies on the separation of quark zero and free modes, but also explicitly includes both instantons and perturbative gluons. Gluons are separated into classical instantons and quantum corrections, which we write as
$$A_\mu (x)=A_\mu ^{(I)}(x)+A_\mu ^{(\overline{I})}(x)+a_\mu (x),$$
(8)
where the sums are over all instantons and anti-instantons, each given by the ’t Hooft solution in the singular gauge . Although there is a certain distribution of instanton sizes (see, for example, ) we simply use the average value, $`\rho 1/3`$ fm, in all calculations.
The main idea of this paper is to examine the effects of color superconductivity on the gluonic excitations above the instanton vacuum, which means we must retain and analyze the $`a_\mu ^a`$. Hereafter, the term ‘gluon’ will refer to the perturbative gauge fluctuation above the instanton background, $`a_\mu ^a`$.
The origin the ’t Hooft interaction is explained in the literature and will not be repeated here. To summarize, the low momentum quarks are approximated by zero mode solutions in the presence of one instanton. Averaging over the instanton ensemble generates a vertex for dynamical quarks, mediated by the zero modes, which can be treated perturbatively when the instanton liquid is reasonably dilute. The relevant small parameter is the ratio of average instanton size ($`\rho `$) to the average inter-instanton spacing ($`R`$). From phenomenological , variational , and lattice calculations one finds
$$\frac{\rho }{R}\frac{1}{3}.$$
(9)
More details follow in Section 4.
We now consider the effective action itself. In the quark sector it is a non-local one of four-fermion operators, since quarks are connected to instantons via the quark zero modes whose spatial extent is of the order of the instanton size, $`\rho `$. Specifically, one has :
$$S_{INT}=\lambda 𝑑Ud^4z\underset{f}{\overset{N_f}{}}\left[d^4x_fd^4y_f\psi _f^{}(x_f)/\mathrm{\Phi }(x_fz,U)\stackrel{~}{\mathrm{\Phi }}(y_fz,U)/\psi ^f(y_f)\right],$$
(10)
where $`U`$ is the instanton’s $`2\times N_c`$ color/spin orientation matrix and $`z`$ is its position. The $`\mathrm{\Phi }(x)`$ is the zero mode solution for fermions in the field of one instanton; in general it depends on the chemical potential. At $`\mu =0`$ its Fourier transform is the form factor,
$$f(p)=2x\left[I_1(x)K_0(x)I_0(x)K_1(x)+\frac{1}{x}I_1(x)K_1(x)\right]_{x=p\rho /2},$$
(11)
such that $`f(0)=1`$. At $`\mu 0`$ the form factor is also known explicitly .
Because of non-locality Eq. (10) is not gauge invariant. When calculating bulk vacuum properties one needn’t worry about quantum corrections as observables are seldom sensitive to them, however in this case the dependence is crucial. A limited literature does exist in which non-local interactions are modified to be gauge invariant and we follow the same procedure here. In particular, we will minimally modify the effective action (10) to suit the present needs.
We will take the non-local four-fermion interaction as the starting point. It then becomes a matter of multiplying each quark operator by a path-ordered exponential in the background of the perturbative gluon field $`a_\mu `$, replacing as:
$`\psi (x)`$ $``$ $`\psi (x)W(x,z),`$
$`W(x,y)`$ $`=`$ $`𝒫\mathrm{exp}\left(i{\displaystyle \frac{g}{2}}{\displaystyle _y^x}𝑑s_\mu a_\mu ^a\lambda ^a\right),`$ (12)
where $`a_\mu ^a`$ is the perturbative field. As has been pointed out in the cited works, the choice of path integrated over is not unique. Yet as long as these factors transform as
$$W(x,y)U(x)W(x,y)U^{}(y)$$
(13)
by virtue of path ordering, the action will be gauge invariant. This remains true when, as in this case, the interaction involves an explicit color average since overall color is conserved. We are also concerned only with the static limit, $`q^20`$, in which results are independent of any particular choice of path.
Thus we write the modified interaction as a product over flavors,
$`S_{INT}`$ $`=`$ $`\lambda {\displaystyle }dUd^4z{\displaystyle \underset{f}{\overset{N_f}{}}}[d^4x_fd^4y_f`$ (14)
$`\times \psi _f^{}(x_f)W(x_f,z)/\mathrm{\Phi }(x_fz)\mathrm{\Phi }^{}(y_fz)/W(y_f,z)\psi ^f(y_f)],`$
noting that this includes terms of all orders in $`ga_\mu `$. Calculating the gluon polarization operator will require the linear and quadratic contributions.
To this interaction term one must add the usual quark kinetic term, minimally modified to preserve gauge invariance:
$$S_{KIN}=d^4x\psi ^{}\gamma _\mu (i_\mu +\frac{g}{2}\lambda ^aa_\mu ^a)\psi .$$
(15)
Consequently, the color current obtained from the variation of the action in respect to $`a_\mu ^a`$ will have two contributions: one is the ordinary one arising from the minimal coupling (15) and the other arising from the non-local interaction term (14).
## 4 Diquark Condensation
In this section some details of the color superconductor phase are reviewed. All expressions are Euclidean and the notation generally follows that of Ref. ; the reader is referred to this reference for details specific to this particular approach. The spontaneous breaking of chiral symmetry was also a central concern in that work, whereas here we are considering restored chiral symmetry. The main result of Refs. is a competition between chiral and diquark condensates at zero temperature and nonzero quark chemical potential. In all cases, there is a phase transition from a low-density phase of spontaneously broken chiral symmetry to a high-density one of color superconductivity.
While technically involved, our previous results arise from an effective four-quark interaction which allows for pairing of quarks as in the BCS theory. The instanton model retains more of QCD’s features than ad hoc models, notably the anomalous breaking of axial $`U(1)`$ and transmutation of dimensions, and is constructed in a more systematic way from that underlying theory.
A prominent feature of the instanton approach used here is the determination of the effective four-fermion coupling constant $`\lambda `$. This constant is determined by a saddle-point evaluation and proves to be nonlinearly dependent on the background instanton density, $`N/V`$. A gap equation obtained from a set of Schwinger-Dyson-Gorkov equations are solved to first order in $`\lambda `$, done self-consistently to determine the quark pairing gap, $`\mathrm{\Delta }_0`$. Standard quark propagators are necessarily split by a nonzero $`\mathrm{\Delta }_0`$, since this introduces a color bias and quarks are no longer color degenerate. This detail will be crucial when quark loops are computed in Section 7.
Since this paper deals only with a phase of chiral-symmetric, diquark condensation the Schwinger-Dyson-Gorkov equations and corresponding diagrams are the simple ones of Fig. 2. We define the quark propagators as
$$\psi ^{f\alpha i}(p)\psi _{g\beta j}^{}(p)=\{\begin{array}{cc}\delta _g^f\delta _\beta ^\alpha S_1(p)_j^i\hfill & \alpha ,\beta =1,2\hfill \\ \delta _g^f\delta _\beta ^\alpha S_2(p)_j^i\hfill & \alpha ,\beta =3\hfill \end{array},$$
(16)
and the anomalous Gorkov propagator as
$$\psi _L^{f\alpha i}(p)\psi _L^{g\beta j}(p)=\psi _R^{f\alpha i}(p)\psi _R^{g\beta j}(p)=ϵ^{fg}ϵ^{\alpha \beta 3}ϵ^{ij}F(p).$$
(17)
In these expressions indices $`f`$ and $`g`$ refer to flavor, $`i`$ and $`j`$ to spin, the Greek letters to color, and $`\psi _{L,R}`$ are chiral spinors. Written in the chiral $`L,R`$ basis, the $`4\times 4`$ propagator $`S_1(p)`$ is of the form:
$$S_1(p)=\left[\begin{array}{cc}0& Z(p)𝐒_0(p)^+\\ Z(p)𝐒_0(p)^{}& 0\end{array}\right],$$
(18)
while $`S_2(p)`$ is the free propagator for color 3 and hence identical to the above with the function $`Z(p)`$ absent. The notation is $`x^\pm =x_\mu \sigma _\mu ^\pm `$, where the $`2\times 2`$ matrices $`\sigma _\mu ^\pm =(\pm i\stackrel{}{\sigma },1)`$ decompose the Dirac matrices into chiral components. The off-diagonal, bare propagator is therefore written $`𝐒_0(p)^\pm =\left[p^\pm \right]^1`$.
With these definitions, we have the scalar Schwinger-Dyson-Gorkov equations
$`Z(p)`$ $`=`$ $`1\mathrm{\Delta }(p)F(p)`$
$`F(p)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }(p)Z(p)}{p^2}}.`$ (19)
The momentum dependence of the gap,
$$\mathrm{\Delta }(p)=\mathrm{\Delta }_0f(p)^2,$$
(20)
is given by the instanton-induced form factor, Eq. (11), which suppresses the interaction beyond $`p>1/\rho 600`$ MeV. This pair of equations leads to a gap equation for $`\mathrm{\Delta }_0`$,
$`\mathrm{\Delta }_0^2`$ $`=`$ $`{\displaystyle \frac{2\lambda }{N_c(N_c1)}}{\displaystyle \frac{d^4p}{(2\pi )^4}\mathrm{\Delta }(p)F(p)}`$ (21)
$`=`$ $`{\displaystyle \frac{2\lambda }{N_c(N_c1)}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)^2}{p^2+\mathrm{\Delta }(p)^2}},`$
which must be self-consistently solved with $`\lambda `$. This coupling constant is in turn determined by a saddle-point integration (exactly in the thermodynamic limit $`N,V\mathrm{}`$),
$$\lambda =\frac{4N_c(N_c1)}{N/V}\mathrm{\Delta }_0^2.$$
(22)
Combining these two equations $`\lambda `$ may be eliminated and we obtain
$$1=8\left(\frac{N}{V}\right)^1\frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)^2}{p^2+\mathrm{\Delta }(p)^2}.$$
(23)
The scale of the gap, therefore, is set by the instanton density $`N/V`$. To this order $`N/V`$ remains at its vacuum value, and although it will be affected by the finite density of quarks this is an $`𝒪(\lambda )`$ correction to the instanton weight and not considered here.
Solved numerically, the gap is $`\mathrm{\Delta }_0400`$ MeV in vacuum and drops to $`\mathrm{\Delta }_0200`$ MeV at the phase transition from chiral broken to color superconducting matter, as detailed in Ref. <sup>3</sup><sup>3</sup>3 In Ref. $`\mathrm{\Delta }_0`$ was defined as half this paper’s (and the more standard) definition..
## 5 Color Current
A color-conserving Noether current naturally follows from the modified interaction, including contributions from both $`S_{KIN}`$ (15) and $`S_{INT}`$ (14). Along with the standard quark-gluon coupling piece,
$$j_\mu ^a(q)=\frac{d^4p}{(2\pi )^4}\left[\psi _R^{}(p)\frac{\lambda ^a}{2}\sigma _\mu ^{}\psi _L(p+q)+\psi _L^{}(p)\frac{\lambda ^a}{2}\sigma _\mu ^+\psi _R(p+q)\right],$$
(24)
the current now includes a four-quark coupling to the gluon as shown in Fig. 3 and written in terms of four-momentum:
$`\stackrel{~}{ȷ}_\mu ^a(q)`$ $`=`$ $`{\displaystyle \frac{\delta S_{INT}}{\delta a_\mu ^a(q)}}|_{a_\mu ^a=0}`$ (25)
$`=`$ $`\lambda {\displaystyle \frac{ϵ^{f_1f_2}ϵ_{g_1g_2}}{4}}{\displaystyle 𝒟ie^{ix_f(p_fp_f^{})+iy_f(k_f+k_f^{})iz(p_f^{}+k_f^{})}}`$
$`\times [{\displaystyle _{x_1}^z}ds_\mu e^{iqs}\psi _{Lf_1}^{}(p_1)\lambda ^a(p_1^{},k_1^{})\psi _L^{g_1}(k_1)\psi _{Lf_2}^{}(p_2)(p_2^{},k_2^{})\psi _L^{g_2}(k_2)`$
$`+{\displaystyle _z^{y_1}}ds_\mu e^{iqs}\psi _{Lf_1}^{}(p_1)(p_1^{},k_1^{})\lambda ^a\psi _L^{g_1}(k_1)\psi _{Lf_2}^{}(p_2)(p_2^{},k_2^{})\psi _L^{g_2}(k_2)]`$
$`+(LR).`$
Indices on quarks denote chirality and flavor. The measure is
$$𝒟=dUd^4z\underset{f}{\overset{N_f=2}{}}d^4x_fd^4y_f\frac{d^4p_fd^4p_f^{}d^4k_fd^4k_f^{}}{(2\pi )^{16}}$$
(26)
and the form factors lie in the color/spin matrices
$$(p,k)_{\beta j}^{\alpha i}=U_k^\alpha ϵ^{ki}ϵ_{jl}U_\beta ^lf(p)f(k).$$
(27)
The $`U`$ are $`2\times N_c`$ color orientation matrices, averaged in each vertex. Through use of the Dirac equation for $`\psi `$ , one can explicitly verify that $`q_\mu J_\mu (q)=0`$ for $`J_\mu =j_\mu +\stackrel{~}{ȷ}_\mu `$. While this condition would remain satisfied with a transverse addition to the color current, no such addition is motivated here.
In practice we are interested in pairing off two of the four quark legs of the second vertex. Since we are considering a phase where chiral symmetry is unbroken, a chirality-violating $`\psi ^{}\psi `$ loop cannot contribute. Thus we pair either $`\psi \psi `$ or $`\psi ^{}\psi ^{}`$ and obtain the effective current,
$`\stackrel{~}{ȷ}_\mu ^a(q)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_0}{2}}{\displaystyle d^4zd^4x\frac{d^4p_fd^4p_f^{}d^4p^{}}{(2\pi )^{12}}}`$ (28)
$`\times [e^{ix(p_1p^{})iz(p^{}+p_2)}{\displaystyle _x^z}ds_\mu e^{iqs}\psi _L^{}(p_1)\lambda ^a𝒩(p_2,p^{})\psi _L^{}(p_2)`$
$`+e^{ix(p_1+p^{})iz(p^{}p_2)}{\displaystyle _z^x}ds_\mu e^{iqs}\psi _L(p_1)𝒩(p_2,p^{})^{}\lambda ^a\psi _L(p_2)]`$
$`+(LR).`$
where here the flavor/color/spin structure is
$$𝒩(p,p^{})^{ff^{},\alpha \alpha ^{},ii^{}}=iϵ^{ff^{}}ϵ^{\alpha \alpha ^{}3}ϵ^{ii^{}}f(p)f(p^{}).$$
(29)
The two terms in Eq. (28) correspond to pairing either two incoming or two outgoing quark legs of Fig. 3.
Gluon mass terms are found in the $`q^20`$ limit. When the gluon couples to nonsingular composite quark modes one can set $`q^2`$ strictly to zero and the path integrals simplify to
$`e^{ip_i^{}(zx_i)}{\displaystyle _{x_i}^z}𝑑s_\mu e^{iqs}|_{q=0}`$ $`=`$ $`e^{ip_i^{}(zx_i)}(zx_i)_\mu `$ (30)
$`=`$ $`ie^{ip_i^{}(zx_i)}{\displaystyle \frac{}{p_{i\mu }^{}}}.`$
This substitution, the result of which clearly depends only on the endpoints of the path $`s_\mu `$ rather than any particular choice of path, leads to differentiation of the form factor. On the other hand, the vertex can also couple gluons to Nambu-Goldstone modes which are singular as $`q^20`$. Their $`1/q^2`$ behavior must be countered by expanding $`\stackrel{~}{ȷ}_\mu ^a(q)`$ to order $`q_\mu `$, but before this complication arises the Nambu-Goldstone modes must be specified.
## 6 Nambu-Goldstone Modes
A symmetry has been spontaneously broken and Nambu-Goldstone modes are certain to follow. With $`SU(3)`$ being broken to $`SU(2)`$ there are five massless modes and, directly analogous to the Higgs mechanism, they do not become additional degrees of freedom. Instead they mix with and are incorporated into the five massive gluons, thereby relevant in the color Meissner effect.
Since the associated condensate is $`qq`$, these massless modes will be quark bilinears. To determine their quantum numbers one need only perform a gauge rotation on the diquark condensate and catalog the five orthogonal correlators which appear. Recalling the condensate direction as chosen in Eq. (17),
$$iϵ_{ij}ϵ_{fg}\psi _L^{fiT}\lambda ^2\psi _L^{gj}\mathrm{\Delta }_0,$$
(31)
we can list the five diquark operators which couple to the Nambu-Goldstone excitations:
$`ϵ_{fg}ϵ_{kl}\psi _L^{fkT}\lambda ^7\psi _L^{gl},iϵ_{fg}ϵ_{kl}\psi _L^{fkT}\lambda ^7\psi _L^{gl},ϵ_{fg}ϵ_{kl}\psi _L^{fkT}\lambda ^5\psi _L^{gl},`$
$`iϵ_{fg}ϵ_{kl}\psi _L^{fkT}\lambda ^5\psi _L^{gl},ϵ_{fg}ϵ_{kl}\psi _L^{fkT}\lambda ^2\psi _L^{gl}.`$ (32)
Propagators for these quark-quark modes will exhibit a simple pole at $`q^2=0`$, behavior recovered by computing the corresponding correlation functions. In order to obtain this pole we must sum $`s`$ and $`u`$ channel contributions to all orders in the four-quark coupling $`\lambda `$ <sup>4</sup><sup>4</sup>4A related summation could be done for $`t`$ channel diquark exchange as well, however these do not contribute the dominant $`1/q^2`$ behavior in the limit of zero transfer momentum.. Since connecting these modes to the gluon propagator will be itself of order $`g^2`$, we do not include any gluonic corrections to the instanton vertex. With both standard and anomalous quark propagators at our disposal we obtain the set of coupled Bethe-Salpeter-Gorkov equations diagramed in Fig. 4. The diagram shows not only standard two-body propagators, denoted $`\mathrm{\Gamma }`$, but also its anomalous analog, $`\mathrm{\Omega }`$. It is easy to see that $`\mathrm{\Omega }`$ will vanish when any of the external lines are quarks of color 3, since the vertices conserve color and the internal (Gorkov) propagators involve only colors 1 and 2.
Knowing the quantum numbers of the Nambu-Goldstone modes, we can immediately write down ansätze for the four-point functions which will be required when we later compute the gluon polarization operator:
$`\psi _{\chi f\alpha i}^{}(p)\psi _{\chi f^{}\alpha ^{}i^{}}^{}(p^{})\psi _\chi ^{}^{g\beta j}(k)\psi _\chi ^{}^{g^{}\beta ^{}j^{}}(k^{})={\displaystyle \frac{ϵ_{ff^{}}ϵ^{gg^{}}ϵ_{ii^{}}ϵ^{jj^{}}}{N_f^2(N_c1)}}[ϵ_3\lambda ^a]_{\alpha \alpha ^{}}[ϵ^3\lambda ^b]^{\beta \beta ^{}}`$
$`\times (2\pi )^4\delta ^4\left(p+p^{}kk^{}\right)f(p)f(p^{})f(k)f(k^{})\mathrm{\Gamma }_{\chi \chi ^{}}^{ab}(p+p^{})`$ (33)
$`\psi _\chi ^{f\alpha i}(p)\psi _\chi ^{f^{}\alpha ^{}i^{}}(p^{})\psi _\chi ^{}^{g\beta j}(k)\psi _\chi ^{}^{g^{}\beta ^{}j^{}}(k^{})={\displaystyle \frac{ϵ^{ff^{}}ϵ^{gg^{}}ϵ^{ii^{}}ϵ^{jj^{}}}{N_f^2(N_c1)}}[ϵ_3\lambda ^a]^{\alpha \alpha ^{}}[ϵ_3\lambda ^b]^{\beta \beta ^{}}`$
$`\times (2\pi )^4\delta ^4\left(p+p^{}+k+k^{}\right)f(p)f(p^{})f(k)f(k^{})\mathrm{\Omega }_{\chi \chi ^{}}^{ab}(p+p^{}).`$ (34)
Here, $`\left[ϵ_3\right]_{\alpha \beta }=ϵ_{3\alpha \beta }`$, and the $`\chi `$ refer to $`L`$ or $`R`$ chirality with the chiral substructures defined as
$$\mathrm{\Gamma }=\left[\begin{array}{cc}\mathrm{\Gamma }_{LL}& \mathrm{\Gamma }_{LR}\\ \mathrm{\Gamma }_{RL}& \mathrm{\Gamma }_{RR}\end{array}\right]\mathrm{\Omega }=\left[\begin{array}{cc}\mathrm{\Omega }_{LL}& \mathrm{\Omega }_{LR}\\ \mathrm{\Omega }_{RL}& \mathrm{\Omega }_{RR}\end{array}\right],$$
(35)
in which color indices have been suppressed. It is clear from their definitions that $`\mathrm{\Gamma }_{LR}=\mathrm{\Gamma }_{RL}`$, $`\mathrm{\Gamma }_{LL}=\mathrm{\Gamma }_{RR}`$, $`\mathrm{\Omega }_{LR}=\mathrm{\Omega }_{RL}`$, and $`\mathrm{\Omega }_{LL}=\mathrm{\Omega }_{RR}`$.
Written in terms of these chiral elements, the two diagrams of Fig. 4 correspond to the following set of equations:
$`\mathrm{\Gamma }_{LL}^{ab}(q)`$ $`=`$ $`\overline{\lambda }\left[1+\widehat{c}_1^a\mathrm{\Gamma }_{LR}^{ab}(q)_1(q)+\widehat{c}_2^a\mathrm{\Gamma }_{LR}^{ab}(q)_2(q)+\widehat{c}_3^a\mathrm{\Omega }_{LL}^{ab}(q)_3(q)\right]`$
$`\mathrm{\Gamma }_{LR}^{ab}(q)`$ $`=`$ $`\overline{\lambda }\left[\widehat{c}_1^a\mathrm{\Gamma }_{LL}^{ab}(q)_1(q)+\widehat{c}_2^a\mathrm{\Gamma }_{LL}^{ab}(q)_2(q)+\widehat{c}_3^a\mathrm{\Omega }_{LR}^{ab}(q)_3(q)\right]`$
$`\mathrm{\Omega }_{LL}^{ab}(q)`$ $`=`$ $`\overline{\lambda }\left[\widehat{c}_3^a\mathrm{\Gamma }_{LL}^{ab}(q)_3(q)+\widehat{c}_2^a\mathrm{\Omega }_{LR}^{ab}(q)_2(q)\right]`$
$`\mathrm{\Omega }_{LR}^{ab}(q)`$ $`=`$ $`\overline{\lambda }\left[\widehat{c}_3^a\mathrm{\Gamma }_{LR}^{ab}(q)_3(q)+\widehat{c}_2^a\mathrm{\Omega }_{LL}^{ab}(q)_2(q)\right],`$ (36)
where no sums are implied on color indices $`a,b`$. The adjoint color vectors $`\{\widehat{c}_i\}`$ determine the internal propagators and are found to be
$`\widehat{c}_1`$ $`=`$ $`(0,0,0,1,1,1,1,0)`$
$`\widehat{c}_2`$ $`=`$ $`(1,1,1,0,0,0,0,1)`$
$`\widehat{c}_3`$ $`=`$ $`(1,1,1,0,0,0,0,1)`$ (37)
The integrals which result from these loops are:
$`_1(q^2)`$ $``$ $`{\displaystyle \frac{1}{\mathrm{\Delta }_0^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{p_+p_{}}{p_+^2p_{}^2}[\mathrm{\Delta }(p_+)\mathrm{\Delta }(p_{})]\left[Z(p_+)+Z(p_{})\right]}`$
$`_2(q^2)`$ $``$ $`{\displaystyle \frac{2}{\mathrm{\Delta }_0^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}[\mathrm{\Delta }(p_+)\mathrm{\Delta }(p_{})]F(p_+)F(p_{})}`$
$`_3(q^2)`$ $``$ $`{\displaystyle \frac{2}{\mathrm{\Delta }_0^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{p_+p_{}}{p_+^2p_{}^2}[\mathrm{\Delta }(p_+)\mathrm{\Delta }(p_{})]Z(p_+)Z(p_{})},`$ (38)
where we have defined $`(p_\pm )_\mu =p_\mu \pm \frac{1}{2}q_\mu `$ and
$$\overline{\lambda }\frac{\lambda }{N_c(N_c1)}.$$
(39)
The set of Bethe-Salpeter-Gorkov equations (36) can easily be solved for correlations functions of all adjoint colors. For colors $`a=b=4\mathrm{}8`$, in direct analogy to similar calculations of pions in the instanton vacuum , the $`q^2=0`$ pole arises due to a cancellation in the denominator which is a direct consequence of the gap equation. The calculation follows identically for $`a=4,5,6,`$ and 7, where we find
$$\mathrm{\Gamma }_{LR}^{44}(q^2)=\mathrm{\Gamma }_{LL}^{44}(q^2)=\frac{\overline{\lambda }}{1\overline{\lambda }^2_1(q^2)^2}.$$
(40)
After writing the gap equation (23) in these terms,
$$1=\overline{\lambda }_1(0),$$
(41)
we have
$$\mathrm{\Gamma }_{LR}^{44}(q^2)=\overline{\lambda }\left(1[\overline{\lambda }_1(0)]^2q^2\overline{\lambda }^2\frac{}{q^2}_1(q^2)^2\right)^1\frac{Z_\varphi }{q^2},$$
(42)
where
$`Z_\varphi ^1`$ $`=`$ $`2{\displaystyle \frac{_1(q^2)}{q^2}}|_{q^2=0}`$ (43)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }_0^2}}{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{{\displaystyle \frac{\mathrm{\Delta }(p)^2\frac{1}{2}\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)+\frac{1}{4}p^2\mathrm{\Delta }^{}(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}}`$
$`{\displaystyle \frac{\frac{1}{2}\mathrm{\Delta }(p)^2\left[\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)\right]^2}{p^2[p^2+\mathrm{\Delta }(p)^2]^2}}\}.`$
The primes denote differentiation with respect to $`p`$. The value of $`Z_\varphi `$ will change with density, since not only does it depend on $`\mathrm{\Delta }_0`$ but the integrand will exhibit functional dependence on $`\mu `$. In vacuum numerical evaluation gives $`Z_\varphi ^1=8.07\times 10^3`$.
For the $`\mathrm{\Gamma }^{88}`$ and $`\mathrm{\Omega }^{88}`$ solving the coupled equations becomes more involved, since all diagrams in Fig. 4 are present. One finds
$`\mathrm{\Gamma }_{LR}^{88}(q^2)=\mathrm{\Omega }_{LR}^{88}(q^2)`$ $`=`$ $`{\displaystyle \frac{\overline{\lambda }}{1\overline{\lambda }^2\left[_2(q^2)+_3(q^2)\right]^2}}`$ (44)
$`=`$ $`{\displaystyle \frac{Z_\varphi }{q^2}},`$
where the second line is obtain by making use of the Schwinger-Dyson-Gorkov equations, Eqs. (19), in the limit of small $`q^2`$. This completes the set of Nambu-Goldstone modes. One can naturally compute the correlation function for the additional diquark correlators with $`a=b=1,2,3`$, and find a crucial sign difference in the combination of integrals:
$`\mathrm{\Gamma }_{LR}^{11}(q^2)=\mathrm{\Omega }_{LR}^{11}(q^2)`$ $`=`$ $`{\displaystyle \frac{\overline{\lambda }}{1\overline{\lambda }^2\left[_2(q^2)_3(q^2)\right]^2}}`$ (45)
$`=`$ $`{\displaystyle \frac{1}{q^2+m(q)^2}}.`$
The mass can be determined from an $`𝒪(q^2)`$ expansion of the denominator, however this will not be done here as for our purposes it is enough to verify the $`a=`$ 1, 2, and 3 modes are indeed massive.
## 7 Gluon Polarization Operator
With the gapped quark propagators, conserved current interactions, and composite modes defined in the previous sections, we now compute the leading modification to the gluon polarization operator. All diagrams prove to be color diagonal, and so we write
$$\mathrm{\Pi }_{\mu \nu }^{ab}(q^2)=\delta ^{ab}\mathrm{\Pi }_{\mu \nu }(q^2).$$
(46)
Our interest is the static limit, $`q^20`$, which may be considered an effective mass.
In the presence of a color-3, scalar diquark the gluons are divided into three classes. Gluons of adjoint colors 1, 2, and 3 belong to the residual $`SU(2)`$ gauge group and as such remain massless. Gluons 4, 5, 6, and 7 couple one gapped quark (of fundamental color 1 or 2) with the ungapped species and share a degenerate mass. Gluon 8, diagonal in fundamental color, obtains a mass proportional to the previous four. One polarization operator from each class will be explained here to avoid unnecessary repetition. All possible contributions to order $`g^2`$ are diagramed in Figs. 5 and 6 and, depending on the gluonic species, some of these diagrams vanish and others combine to cancel in the static limit.
We begin with the case of gluons 4–7, considering corrections to $`a_\mu ^4a_\nu ^4`$. Diagrams (5b), (5f), and those involving (6c) require pairs of Gorkov propagators and thus vanish since the quarks of color 3 (to which these gluons couple) cannot propagate anomalously. There are four graphs which constitute the ‘contact’ term proportional to $`\delta _{\mu \nu }`$.
Diagram (5a) is a standard loop, where one of the quarks is gapped (color 1 in the case of gluon 4) and the other not (color 3). After subtracting off the vacuum part of this diagram, which remains a concern of gluon renormalization and is not relevant to the Meissner mass, we find
$$\mathrm{\Pi }_{\mu \nu }^{(\text{5}a)}(q^20)=g^2N_f\delta _{\mu \nu }\frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}.$$
(47)
The remaining integral is finite due to an power-law cut-off in the the function $`\mathrm{\Delta }(p)`$ arising from the finite size of instantons.
Additional diagrams arise from the modified $`S_{INT}`$. Not only does this generate additional current interactions (29), but the interaction (14) itself contains a contact contribution to the gluon two-point function. This, diagram (5c), is the second variation of the action (14) with respect to the fourth gluon field,
$$\mathrm{\Pi }_{\mu \nu }^{(\text{5}c)}(q^20)=\frac{\delta ^2S_{INT}}{\delta a_\mu ^4(q)\delta a_\nu ^4(q)}|_{q^2=0}.$$
(48)
Evaluated to order $`g^2`$ it is
$$\mathrm{\Pi }_{\mu \nu }^{(\text{5}c)}(q^20)=\frac{1}{4}g^2N_f\delta _{\mu \nu }\frac{d^4p}{(2\pi )^4}\frac{\frac{1}{2}p^2\mathrm{\Delta }^{}(p)^2+\mathrm{\Delta }(p)p^2\mathrm{\Delta }^{\prime \prime }(p)}{p^2[p^2+\mathrm{\Delta }(p)^2]}.$$
(49)
Diagrams (5d) and (5e), constructed with the additional current piece, are
$`\mathrm{\Pi }_{\mu \nu }^{(\text{5}d)}(q^20)`$ $`=`$ $`{\displaystyle \frac{1}{4}}g^2N_f\delta _{\mu \nu }{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)}{p^2[p^2+\mathrm{\Delta }(p)^2]}}`$
$`\mathrm{\Pi }_{\mu \nu }^{(\text{5}e)}(q^20)`$ $`=`$ $`{\displaystyle \frac{1}{8}}g^2N_f\delta _{\mu \nu }{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{p^2\mathrm{\Delta }^{}(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}}.`$ (50)
Eqs. (47), (49), and (50) comprise the microscopic equivalent of the Higgs contact term (6), and their sum is
$`\mathrm{\Pi }_{\mu \nu }^{(\text{5})}(q^20)=`$
$`g^2N_f\delta _{\mu \nu }{\displaystyle \frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)^2\frac{1}{4}\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)\frac{1}{4}\mathrm{\Delta }(p)p^2\mathrm{\Delta }^{\prime \prime }(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}}.`$ (51)
The integral can be trivially rewritten and then manipulated to yield
$`{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{{\displaystyle \frac{\mathrm{\Delta }(p)^2\frac{1}{2}\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)}{p^2[p^2+\mathrm{\Delta }(p)^2]}}`$
$`+{\displaystyle \frac{\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)}{4p^2[p^2+\mathrm{\Delta }(p)^2]}}{\displaystyle \frac{\mathrm{\Delta }(p)p^2\mathrm{\Delta }^{\prime \prime }(p)}{4p^2[p^2+\mathrm{\Delta }(p)^2]}}\}=Z_\varphi ^1\mathrm{\Delta }_0^2.`$ (52)
The final equality is achieved by integrating the second and third terms by parts.
The Nambu-Goldstone modes couple to the gluons as in Fig. 6. The construction \[(6a)+(6b)\]$`\mathrm{\Gamma }^{44}(q^2)`$ \[(6a)+(6b)\] supplies the $`q_\mu q_\nu /q^2`$ piece to ensure transversality:
$`\mathrm{\Pi }_{\mu \nu }^{(\text{6})}(q^20)`$ $`=`$ $`g^2N_fq_\mu q_\nu {\displaystyle \frac{Z_\varphi }{q^2}}[{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{{\displaystyle \frac{\frac{1}{2}\mathrm{\Delta }(p)^2\left[\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)\right]^2}{p^2\left[p^2+\mathrm{\Delta }(p)^2\right]^2}}`$ (53)
$`{\displaystyle \frac{\mathrm{\Delta }(p)^2\frac{1}{2}\mathrm{\Delta }(p)p\mathrm{\Delta }^{}(p)+\frac{1}{4}p^2\mathrm{\Delta }^{}(p)^2}{p^2\left[p^2+\mathrm{\Delta }(p)^2\right]}}\}]^2`$
$`=`$ $`g^2N_f\mathrm{\Delta }_0^2Z_\varphi ^1{\displaystyle \frac{q_\mu q_\nu }{q^2}}.`$
We have now accounted for all contributions to gluons 4, 5, 6, and 7.
Analysis of the eighth gluon follows in a similar fashion, although with terms from every diagram. A superficial difference lies in the factors arising from the elements $`\lambda ^8`$ which lead to a polarization $`4/3`$ times the previous result. More subtle is the combination of diagrams (5a) and (5b). For $`a_\mu ^8`$ this sums to $`4/3`$ times Eqs. (47), whereas for $`a_\mu ^{1,2,3}`$ they cancel one another. This cancellation is only manifest to $`𝒪(\mathrm{\Delta }_0^4)`$ here, due to the fact that the diagrams are constructed with a pair of one-loop, resummed quark propagators as determined in Section 3. This generates a $`\mathrm{\Delta }(p)^2`$ term in each integrand denominator, in essence including higher-order terms in the perturbative expansion (in $`\lambda `$) which violate gauge invariance. After the vacuum pieces are subtracted, these diagrams contribute the following to $`a_\mu ^8`$:
$`\mathrm{\Pi }_{\mu \nu }^{(\text{5}a+\text{5}b)}(q^20)`$ $`=`$ $`{\displaystyle \frac{1}{3}}g^2N_f\mathrm{\Delta }_0^2\delta _{\mu \nu }{\displaystyle }{\displaystyle \frac{d^4p}{(2\pi )^4}}\{{\displaystyle \frac{p^2}{[p^2+\mathrm{\Delta }(p)^2]^2}}`$
$`{\displaystyle \frac{2\mathrm{\Delta }(p)^2}{[p^2+\mathrm{\Delta }(p)^2]^2}}{\displaystyle \frac{1}{p^2}}\}`$
$`=`$ $`{\displaystyle \frac{4}{3}}g^2N_f\mathrm{\Delta }_0^2\delta _{\mu \nu }{\displaystyle \frac{d^4p}{(2\pi )^4}\left\{\frac{\mathrm{\Delta }(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}+𝒪(\mathrm{\Delta }_0^4)\right\}}.`$
For gluons 1, 2, and 3, the sign of the second term is changed and the corresponding integral is of order $`\mathrm{\Delta }_0^4`$, and thus this quantity vanishes to the order to which we know quark propagators, $`\mathrm{\Delta }_0^2`$ . The same follows for pairs of diagrams with similar construction, such as (5e) and (5f).
By combining Eqs. (46), (51), (52), and (53) we arrive at the satisfyingly compact expressions
$$\mathrm{\Pi }_{\mu \nu }^{ab}(q^20)=\{\begin{array}{cc}0\hfill & a,b=1,2,3\hfill \\ g^2N_f\mathrm{\Delta }_0^2Z_\varphi ^1\delta ^{ab}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)\hfill & a,b=4,5,6,7\hfill \\ \frac{4}{3}g^2N_f\mathrm{\Delta }_0^2Z_\varphi ^1\delta ^{ab}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right)\hfill & a,b=8.\hfill \end{array}$$
(55)
Transversality requires that the contact term be proportional to the wave function renormalization of the Nambu-Goldstone modes and this result, a Ward Identity for color superconductivity, is recovered here.
Finally, to determine a numerical value for the masses we must fix the coupling constant $`g`$. Evaluating in the instanton vacuum, one finds the large finite action
$$S_0=\frac{8\pi ^2}{g^2}12,$$
(56)
or $`g2.6`$ and the perturbative expansion parameter $`\alpha _s=g^2/4\pi =0.54`$.
The gluon masses squared are thus
$$M_a^2=\{\begin{array}{cc}0\hfill & a=1,2,3\hfill \\ 2g^2\mathrm{\Delta }_0^2Z_\varphi ^1(150\mathrm{MeV})^2\hfill & a=4,5,6,7\hfill \\ \frac{8}{3}g^2\mathrm{\Delta }_0^2Z_\varphi ^1(175\mathrm{MeV})^2\hfill & a=8.\hfill \end{array}$$
(57)
These masses apply to the vacuum, $`\mu =0`$.
In order to estimate the finite-density behavior of the Meissner mass, we can simply take the values of $`\mathrm{\Delta }_0`$ for a given $`\mu `$ from the results of Ref. . As detailed in that paper, the instanton form factor (11) becomes density dependent and thus $`\mathrm{\Delta }(p)`$ should be replaced by a complicated $`\mathrm{\Delta }(p_4,|\stackrel{}{p}|,\mu )`$. However, the changes in $`Z_\varphi `$ arising from the finite-$`\mu `$ modifications of the form factors are minor compared compared to the changes in the gap magnitude, $`\mathrm{\Delta }_0`$; we numerically estimate this correction to be about 3%. For simplicity we therefore considered only the scaling from $`\mathrm{\Delta }_0(\mu )`$ (taken from previous work ). Each non-zero gluon mass is proportional to $`\mathrm{\Delta }_0Z_\varphi ^{1/2}`$ (see Eq. (57)) and therefore all will scale identically with density. The resulting gluon mass $`M(\mu )`$, as well as that of the renormalization constant $`Z_\varphi (\mu )^{1/2}`$, is shown in units of its vacuum value in Fig. 7. Note that although $`Z_\varphi `$ rises with increasing quark density, this effect is not sufficient to overcome the falling gap and the masses continuously decrease. The first point of physical relevance would occur around $`\mu 300`$ MeV, the common prediction for chiral restoration to a color superconductor. Results for matter at lower densities correspond to an unstable solution and are only of academic interest .
## 8 Similarity with Chiral Symmetry Breaking
Apart from the gauge coupling $`g`$ the gluon masses (57) are determined by the combination $`F_{qq}^22\mathrm{\Delta }_0^2Z_\varphi ^1`$. This quantity is the analog of the $`F_\pi ^2`$ constant in the chiral-broken phase, both in its physical meaning and algebraically.
In the chiral-broken phase the Nambu-Goldstone bosons are pions. The correlation function of the axial current, $`j_{\mu 5}^A=\psi ^{}\gamma _\mu \gamma _5\tau ^A\psi `$, in the massless quark limit has the transverse form
$$j_{\mu 5}^A(q)j_{\nu 5}^B(q)=\mathrm{\Pi }_{\mu \nu }^{AB}(q)=F_\pi ^2\delta ^{AB}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right).$$
(58)
The transversality of $`\mathrm{\Pi }_{\mu \nu }`$ is the consequence of the conservation of the axial current; the $`1/q^2`$ pole is due to the pion in the intermediate state. Were there gauge bosons coupled to the quark axial current their mass would be equal to $`F_\pi `$ multiplied by the corresponding gauge coupling. In case of diquark (vs. quark-antiquark) condensation the relevant currents are color ones, and we obtain a similar form for the correlation function of two color currents, Eq. (55). One needs only to multiply $`F_{qq}`$ by the gauge coupling to deduce the gluon mass.
If instantons are dilute, the leading contribution to the Kronecker part of the polarization operator arises from the $`\mathrm{\Pi }_{\mu \nu }^{(4a)}`$ piece, Eq. (47). It is the only contribution which diverges logarithmically if one neglects the momentum dependence of the gap $`\mathrm{\Delta }(p)`$. For the same reason the leading contribution to the $`Z_\varphi ^1`$ and $`F_{qq}^2`$ in the dilute limit comes from the pieces not containing the derivatives $`\mathrm{\Delta }^{}(p)`$. One has therefore in the dilute limit:
$$F_{qq}^22\frac{d^4p}{(2\pi )^4}\frac{\mathrm{\Delta }(p)^2}{p^2[p^2+\mathrm{\Delta }(p)^2]}\frac{2\mathrm{\Delta }_0^2}{8\pi ^2}\mathrm{log}\frac{R^2}{\rho ^2}$$
(59)
where $`\mathrm{\Delta }_0`$ is the superconducting gap at zero momentum. It follows from the gap equation (23) that $`\mathrm{\Delta }_0\pi \rho \sqrt{N/V}=\pi \rho /R^2`$. Similarly, the axial correlation function (58) computed in Ref. gives (in the same approximation)
$$F_\pi ^24N_c\frac{d^4p}{(2\pi )^4}\frac{M(p)^2}{[p^2+M(p)^2]^2}\frac{4N_cM_0^2}{8\pi ^2}\mathrm{log}\frac{R^2}{\rho ^2}$$
(60)
where $`M(p)`$ is the dynamical quark mass (the chiral gap) whose value at zero momentum is determined from a corresponding gap equation to be $`M_0\pi \rho \sqrt{N/VN_c}`$.
We see, thus, that not only have the Meissner mass and the $`F_\pi `$ constant analogous meaning, but their algebraic structure is quite similar. One expects, therefore, that the numerical value for the Meissner mass is of the order of $`F_\pi `$, and this expectation is confirmed by an exact numerical calculation of the previous section.
## 9 Discussion and Conclusions
We have analyzed the problem of spontaneous gauge symmetry breaking brought about by a diquark condensate. Since the broken symmetry is continuous and gauged the resulting Nambu-Goldstone modes do not remain in the spectra, instead mixing with the longitudinal components of the gauge fields to produce massive gauge bosons. This, a dynamical Higgs mechanism, can be called the color Meissner effect in the context of color superconductivity.
To reveal the gauge boson masses mathematically one has to compute the polarization operator, which ought to be transverse,
$$\mathrm{\Pi }_{\mu \nu }^{ab}(q)=M^2\delta ^{ab}\left(\delta _{\mu \nu }\frac{q_\mu q_\nu }{q^2}\right),$$
(61)
where the massless pole $`1/q^2`$ is the manifestation of the Nambu–Goldstone intermediate state. The coefficient, $`M^2`$ gives the mass of the gauge boson. In this paper, we explicitly solved this problem for the case of diquark condensation as induced by the instanton background, the effective action from which we found necessary to modify in order to maintain a conserved color current.
Through computing the gluon polarization operators to order $`g^2`$ we find the effective gluon masses to be on the order of the diquark gap. The three gluons comprising the residual $`SU(2)`$ group remain massless and hence a quark-gluon medium would become color-biased in such a phase. The analysis here was done in the limit of zero temperature and, initially, vanishing chemical potential $`\mu `$, though strictly speaking at zero $`\mu `$ the color superconductor is only metastable with the ground state being the usual chiral-broken phase. We then estimated finite-density dependence of the calculated quantities. At the critical density, where one expects the phase transition to the color superconducting phase, we deduce that the Meissner masses of the $`4^{th}`$, $`5^{th}`$, $`6^{th}`$, and $`7^{th}`$ gluon are about 120 MeV and the $`8^{th}`$ gluon has a mass about 140 MeV. These quantities would be of physical relevance should a low temperature, high density region become experimentally accessible. At a chemical potential low enough to leave the instanton background approximately unchanged ($`\mu <`$ 0.6 GeV), the instanton effects analyzed here would still be present and likely dominant.
In computing the coupling of the Nambu-Goldstone modes to gluons we have established that their mixing is described by an effective theory in which the composite diquark is replaced by a complex scalar in the fundamental representation of the gauge group, $`\varphi ^a`$. The effective Lagrangian coincides with that of elementary scalar field covariantly coupled to gluons, aside from an overall factor $`Z_\varphi `$ which is the ‘wave function renormalization’ of the composite scalar field. There are no a priori reasons for this factor to be close to unity (as it is in the case of a weakly-coupled elementary Higgs field), and indeed we find a substantial deviation from unity. Meanwhile, it is a crucial factor for the estimate of the Meissner mass.
If electromagnetic interactions were taken into account, the gluon $`a_\mu ^8`$ will invariably mix with the (massless) photon. This mixing, estimated from general arguments to be rather small , reorganizes the fields into a massive ‘new gluon’ and massless ‘new photon’. Given that the gluonic sector alone can be recast as a Yang-Mills-Higgs theory, this additional result would complete the analogy to the electroweak sector of the Standard Model.
Finally, we would like to comment on the phase transition between ordinary chiral-broken phase and the color superconductivity. The point is that all estimates existing in the literature (see Ref. and references therein) indicate that this phase transition happens alarmingly ‘early’: taken literary, the claim is that the interior region of a heavy nucleus is actually in a ‘boiling’ state. However, those estimates generally neglect the influence of the dense medium on the gluonic background fluctuations which induce the color superconductivity itself. The Meissner mass of about 150 MeV found here is a large quantity and, together with the Debye mass and other effects, will suppress instantons. Therefore, by taking into account back-influence effects one has a chance to ‘save’ ordinary nuclear matter from a premature phase transition by moving that transition to higher densities.
Although the instanton density is not expected to change significantly for any chemical potential below the inverse instanton size (600 MeV), at and beyond this point perturbative Meissner and screening masses will become increasingly important. At asymptotically large densities perturbative gluon exchange is the source of diquark condensation , whereas in this work we have been concerned with the low density, nonperturbative regime. Determining the behavior at a moderate density scale – which would be an interpolation between the two – is necessary if one wishes to confidently consider signals of color superconductivity in any potentially realizable situation, be it in some future heavy-ion collider or a neutron star.
## Acknowledgments
G.W.C. thanks D. Rischke for discussions and both the Leon Rosenfeld Fund and USDOE grant DE-FG02-88ER40388 for support. |
warning/0001/hep-ph0001309.html | ar5iv | text | # Further investigation of a relic neutralino as a possible origin of an annual–modulation effect in WIMP direct search
## I Introduction
The effect of annual modulation measured by the DAMA Collaboration in its WIMP direct search experiment with a NaI(Tl) detector and reported in Ref. was analyzed in terms of relic neutralinos in Refs.. In these papers, we proved that this interpretation is compatible with the DAMA data, and entails a relic neutralino which might have the role of a major component of dark matter in the Universe, especially when the uncertainties affecting the evaluation of the neutralino–nucleon cross section are taken into account ). We have also presented in detail other physical properties of such a neutralino, both in a Minimal Supersymmetric extension of the Standard Model and in supergravity schemes , and we have outlined how indirect measurements of WIMPS (low–energy antiprotons in cosmic rays and up-going muon fluxes from the center of the Earth and from the Sun) may bring further information , by way of constraints on the supersymmetric configurations derived from the DAMA annual–modulation results .
New data, collected by the DAMA Collaboration in a further two–year running of the NaI(Tl) experiment for an exposure of 38 475 kg $``$ day, and now presented in Ref. , confirm their previous finding of an annual–modulation effect, which does not appear to be related to any possible source of random systematics. Taking together all (old and new) samples of data for a total exposure of 57 986 kg $``$ day, the effect turns out to be at a 4$`\sigma `$ C.L. Performing a maximum likelihood analysis in terms of $`m_\chi `$ and $`\xi \sigma _{\mathrm{scalar}}^{(\mathrm{nucleon})}`$, where $`m_\chi `$ is the WIMP mass, $`\sigma _{\mathrm{scalar}}^{(\mathrm{nucleon})}`$ is the WIMP–nucleon scalar elastic cross–section, and $`\xi =\rho _\chi /\rho _l`$ is the WIMP fractional amount of local non-baryonic dark–matter density $`\rho _l`$, the DAMA Collaboration presents a 3$`\sigma `$ C.L. annual–modulation region, in the plane $`m_\chi `$$`\xi \sigma _{\mathrm{scalar}}^{(\mathrm{nucleon})}`$, whose actual size depends on whether or not the upper–bound constraints previously obtained by the same Collaboration are included, and on the values assigned to the galactic astrophysical velocities. For the purpose of the analysis carried out in the present paper, among the regions presented in Ref. we select the one, which is obtained from the annual–modulation data, by including the upper-bound constraints of Ref. , by setting $`\rho _l`$ at the standard reference value: $`\rho _l=0.3`$ GeV cm<sup>-3</sup>, and by taking into account uncertainties in the astrophysical velocities of the usual galactic Maxwellian distribution (170 km s$`{}_{}{}^{1}v_0`$ 270 km s<sup>-1</sup>; $`v_{esc}`$ = 450–650 km s<sup>-1</sup>; where $`v_0`$ is the rotational velocity of the local system at the position of the solar system and $`v_{esc}`$ is the galactic escape velocity). This region is the one shown in Fig. 1 (should one include also a bulk rotation of the dark matter halo , this region would elongate along the horizontal axis up to $`m_\chi `$ 230 GeV ). In this figure we also show the contour lines for the three values $`v_0=170,220,270`$ km s<sup>-1</sup>, separately . In the comparison of the experimental data with the theoretical evaluations one has to further consider the uncertainty in $`\rho _l`$: 0.1 GeV cm$`{}_{}{}^{3}\rho _l0.7`$ GeV cm<sup>-3</sup> . Fig. 2 displays how the DAMA annual–modulation region shifts along the vertical axis, as the value of $`\rho _l`$ is varied within its uncertainty range. The four panels correspond to the representative values: $`\rho _l`$ = 0.1, 0.3, 0.5, 0.7 GeV cm<sup>-3</sup>. In Fig. 2, as well as in all subsequent figures, where experimental results of direct and indirect WIMP measurements are compared with theoretical evaluations, separate panels are used for the four representative values of $`\rho _l`$.
In the present paper we investigate the implications of the DAMA data with the total exposure of 57 986 kg $``$ day in terms of relic neutralinos, along the lines previously developed in Refs. . We single out the set of the supersymmetric configurations compatible with the DAMA data, and then apply to this set, denoted as set $`S`$, the constraints derived from experimental indirect searches for WIMPs (up-going muons at neutrino telescopes and antiprotons in cosmic rays). In this analysis we incorporate recent, and quite significant, theoretical and experimental developments.
The supersymmetric theoretical framework adopted here is the Minimal Supersymmetric extension of the Standard Model (MSSM) , which conveniently describes the supersymmetric phenomenology at the electroweak scale, without too strong theoretical assumptions. This model has been extensively used by a number of authors for evaluations of the neutralino relic abundance and detection rates (a list of references may be found, for instance, in ).
The neutralino is defined as the lowest–mass linear superposition of photino ($`\stackrel{~}{\gamma }`$), zino ($`\stackrel{~}{Z}`$) and the two higgsino states ($`\stackrel{~}{H}_1^{}`$, $`\stackrel{~}{H}_2^{}`$): $`\chi a_1\stackrel{~}{\gamma }+a_2\stackrel{~}{Z}+a_3\stackrel{~}{H}_1^{}+a_4\stackrel{~}{H}_2^{}`$.
The MSSM contains three neutral Higgs fields: two of them ($`h`$, $`H`$) are scalar and one ($`A)`$ is pseudoscalar. At the tree level the Higgs sector is specified by two independent parameters: the mass of one of the physical Higgs fields, which we choose to be the mass $`m_A`$ of the neutral pseudoscalar boson, and the ratio of the two vacuum expectation values, defined as $`\mathrm{tan}\beta H_2/H_1`$. Once radiative corrections are introduced, the Higgs sector depends also on the squark masses through loop diagrams. The radiative corrections to the neutral and charged Higgs bosons, employed in the present paper, are taken from Refs. .
The other parameters of the model are defined in the superpotential, which contains all the Yukawa interactions and the Higgs–mixing term $`\mu H_1H_2`$, and in the soft–breaking Lagrangian, which contains the trilinear and bilinear breaking parameters and the soft gaugino and scalar mass terms.
To cast the MSSM, which originally contains a large number of parameters, into a form adequate for phenomenology, we follow the common procedure of introducing a set of restrictive assumptions at the electroweak scale: a) all trilinear parameters are set to zero except those of the third family, which are unified to a common value $`A`$; b) all squarks and sleptons soft–mass parameters are taken as degenerate: $`m_{\stackrel{~}{l}_i}=m_{\stackrel{~}{q}_i}m_0`$, c) the gaugino masses are assumed to unify at $`M_{GUT}`$, and this implies that the $`U(1)`$ and $`SU(2)`$ gaugino masses are related at the electroweak scale by $`M_1=(5/3)\mathrm{tan}^2\theta _WM_2`$.
Once these conditions are implemented in the model, the supersymmetric parameter space consists of six independent parameters. We choose them to be: $`M_2,\mu ,\mathrm{tan}\beta ,m_A,m_0,A`$ and vary these parameters in the following ranges: $`10\text{GeV}M_21\text{TeV},\mathrm{\hspace{0.33em}10}\text{GeV}|\mu |1\text{TeV},\mathrm{\hspace{0.33em}80}\text{GeV}m_A1\text{TeV},\mathrm{\hspace{0.33em}100}\text{GeV}m_01\text{TeV},3A+3,\mathrm{\hspace{0.33em}1}\mathrm{tan}\beta 50`$. We remark that the values taken here as upper limits of the ranges for the dimensional parameters, $`M_2,\mu ,m_0,m_A`$, are inspired by the upper bounds which may be derived for these quantities in SUGRA theories, when one requires that the electroweak symmetry breaking, radiatively induced by the soft supersymmetry breaking, does not occur with excessive fine tuning (see Ref. and references quoted therein).
We have further constrained our parameter space, by taking into account all the new experimental limits obtained from accelerators on supersymmetric and Higgs searches (LEP2 , CDF ). Notice that the new bounds from LEP2 and CDF constrain now rather severely the susy space, especially in the region of interest for direct detection (small $`m_h`$ and, partially, large $`\mathrm{tan}\beta `$ ).
Moreover, the constraints due to the $`bs+\gamma `$ process have been taken into account. In our analysis, the inclusive decay rate BR($`BX_s\gamma `$) is calculated with corrections up to the leading order. Next–to–leading order corrections are included only when they can be applied in a consistent way, i.e. both to standard–model and to susy diagrams. We require that our theoretical evaluation for BR($`BX_s\gamma `$) is within the range: 1.96 $`\times 10^4`$ BR($`BX_s\gamma `$) $``$ 4.32 $`\times 10^4`$. This range is obtained by combining the experimental data of Refs. at 95% C.L. and by adding a theoretical uncertainty of 25%, whenever the still incomplete next–to–leading order susy corrections cannot be applied.
Our parameter space has been further constrained by the request that the Lightest Supersymmetric Particle (LSP) is the neutralino, rather than the gluino or squarks or sleptons. The current upper bound for cold dark matter may be establish as $`\mathrm{\Omega }_{CDM}h^2<0.3`$ ($`h`$ is the usual Hubble parameter, defined in terms of the present–day value $`H_0`$ of the Hubble constant as $`hH_0/(100`$ km s$`^1`$Mpc$`{}_{}{}^{1})`$), on the basis of the most recent cosmological data . However, for sake of presentation of the results of the present analysis, which, anyway, never entail values of $`\mathrm{\Omega }_\chi h^2`$ in excess of 0.6 (see last section), we do not impose the bound $`\mathrm{\Omega }_{CDM}h^20.3`$ in our selection of susy configurations. The neutralino relic abundance is calculated here as illustrated in Ref.. We have checked that susy configurations which could potentially lead to coannihilation effects are marginal in our selected supersymmetric parameter space.
A few comments are in order here. The restrictive assumptions a) – c) adopted above in the framework of the MSSM are instrumental in reducing the otherwise large number of independent parameters to a handful set of them (six in our scheme), and in making the calculations of a number of crucial observables (such as relic abundances and event rates) manageable. The few independent parameters of this simplified MSSM have the role of relevant scales for some fundamental quantities, such as scalar masses and gaugino masses, which in turn determine the size of the numerical outputs. This version of MSSM is obviously the simplest scheme for a susy model, and the most natural one to start with. However, one has to be aware of the fact that new experimental data could eventually force one to adopt more involved versions of supersymmetric models, for instance by relaxing some GUT-inspired relation (such as $`M_10.5M_2`$) Ref. , or by including CP–violating phases .
As regards the distribution of relic neutralinos in our Galaxy, to start with we have assumed a standard halo population with a Maxwellian velocity distribution, whose dispersion speed is centered around 270 km s<sup>-1</sup> (i.e., $`v_0`$ = 220 km s<sup>-1</sup>). However, in the implementation of constraints from up-going muons at neutrino telescopes, we have also considered recent theoretical developments which may have quite contrasting effects on the expected signals . These different instances are examined in Sect. III.
Data on antiprotons in space, combined with recent evaluations of the secondary antiproton component in cosmic rays due to spallation processes, are employed in Sect. IV to put further constraints on the original set $`S`$ of susy configurations, singled out by the DAMA data.
We give the results of our combination of the annual–modulation data with indirect measurement constraints in Sect. V, where we also discuss the cosmological properties for our set of relic neutralinos and present our conclusions.
## II Set of supersymmetric configurations singled out by the annual–modulation data
In our papers of Ref. we proved that the DAMA annual–modulation region of Ref. is widely compatible with an interpretation in terms of relic neutralinos, by showing that a sizeable portion of that region is covered by supersymmetric configurations, satisfying all accelerator bounds. Now we show in Fig. 1 that the new, more constrained annual–modulation region of Ref. is still largely compatible with the relic neutralino interpretation, though the supersymmetric space is now more severely constrained by the current limits from accelerators .
In deriving the scatter plot shown in Fig. 1 we have used the scan of the susy parameter space defined in the previous section. The neutralino–nucleon cross section has been calculated with the formulae reported in Ref. . As discussed in Ref. , this cross section suffers from significant uncertainties in the size of Higgs–quark–quark and squark–quark–neutralino couplings. In fact, these couplings depend on quark masses $`m_q`$ and quark scalar densities in the nucleon $`\overline{q}q`$, which are still rather poorly determined. To be specific, we refer to the following quantities: the fractional strange–quark content of the nucleon $`y=2<\overline{s}s>/(<\overline{u}u+\overline{d}d>)`$, the quark mass ratio $`r=2m_s/(m_u+m_d)`$, and the products $`m_q<\overline{q}q>`$’s. In our analysis we have taken into account the uncertainties in these quantities. Thus, our scatter plots comprise representative points which have been derived by using both of the two following sets of values, cumulatively:
Set 1: $`y=0.33,r=29,`$ (2)
$`m_l<\overline{l}l>=\mathrm{\hspace{0.33em}23}\mathrm{MeV},m_s<\overline{s}s>=\mathrm{\hspace{0.33em}215}\mathrm{MeV},m_h<\overline{h}h>=\mathrm{\hspace{0.33em}50}\mathrm{MeV}.`$
Set 2: $`y=0.50,r=29,`$ (4)
$`m_l<\overline{l}l>=\mathrm{\hspace{0.33em}30}\mathrm{MeV},m_s<\overline{s}s>=\mathrm{\hspace{0.33em}435}\mathrm{MeV},m_h<\overline{h}h>=\mathrm{\hspace{0.33em}33}\mathrm{MeV}.`$
In Eqs. (24) $`l`$ stands for light quarks, $`s`$ is the strange quark and $`h=c,b,t`$ denotes heavy quarks. For the light quarks, we have defined $`m_l<\overline{l}l>`$ $``$ $`\frac{1}{2}[m_u<\overline{u}u>+m_d<\overline{d}d>]`$. Set 1 and set 2 bracket, at least partially, the present uncertainties. In Sect. V.B, in connection with neutralino cosmological properties we will also mention the consequences of using a more extreme set of values (set 3 of Ref. ). For the derivation of the values of the various sets see Ref. . It is worth noticing that a new derivation of the pion–nucleon sigma term, $`\sigma _{\pi N}`$, points to rather high values: $`\sigma _{\pi N}`$ = 73.5$`\pm `$9 MeV . By itself, this new result would increase the value of the quantity $`m_s<\overline{s}s>`$ given in Eq.(4) by $``$ 30%. We recall that the quantity $`m_s<\overline{s}s>`$ is crucial in establishing the size of $`\sigma _{scalar}^{(nucleon)}`$ .
As for the values to be assigned to the quantity $`\xi =\rho _\chi /\rho _l`$ we have adopted a standard rescaling recipe. For each point of the parameter space, we take into account the relevant value of the cosmological neutralino relic density. When $`\mathrm{\Omega }_\chi h^2`$ is larger than a minimal value $`(\mathrm{\Omega }h^2)_{\mathrm{min}}`$, compatible with observational data and with large–scale structure calculations, we simply put $`\xi =1`$. When $`\mathrm{\Omega }_\chi h^2`$ turns out to be less than $`(\mathrm{\Omega }h^2)_{\mathrm{min}}`$, and then the neutralino may only provide a fractional contribution to dark matter, we take $`\xi =\mathrm{\Omega }_\chi h^2/(\mathrm{\Omega }h^2)_{\mathrm{min}}`$. The value to be assigned to $`(\mathrm{\Omega }h^2)_{\mathrm{min}}`$ is somewhat arbitrary, in the range $`0.01<(\mathrm{\Omega }h^2)_{\mathrm{min}}<0.3`$. We use here the value $`(\mathrm{\Omega }h^2)_{\mathrm{min}}=0.01`$, which is conservatively derived from the estimate $`\mathrm{\Omega }_{\mathrm{galactic}}0.03`$.
As we mentioned above, Fig. 1 shows that the annual–modulation region (here depicted for $`\rho _l`$ = 0.3 GeV cm<sup>-3</sup>) is largely covered by the scatter plot. This turns out to be the case also for the other representative values of $`\rho _l`$, as is shown in Fig. 2. In each panel of this figure we only display the portion of the susy scatter plot which is contained in each of the relevant experimental region. In going from the generic scanning used for Fig. 1 to the one employed for Fig. 2, although keeping the overall range of variation of the susy parameter space, we have optimized the numerical scanning in order to have a number of configurations, large enough for our subsequent analyses. The covering by the scatter plots of the annual–modulation regions pertaining to different values of $`\rho _l`$ is more extended for large values of $`\rho _l`$ than for the small ones, as expected from the features of the generic plot of Fig. 1.
We define as set $`S`$ of susy configurations the set comprised of the configurations whose representative points in the plane $`m_\chi `$$`\sigma _{scalar}^{(nucleon)}`$ lie inside the annual–modulation regions displayed in Fig. 2. Only configurations of set $`S`$ are retained in the analyses presented hereafter. We remark that set $`S`$ is the union of all the subsets of susy configurations which refer to each of the following representative values for $`\rho _l`$ and $`v_0`$: $`\rho _l=0.1,0.3,0.5,0.7`$ GeV cm<sup>-3</sup>, $`v_0=170,220,270`$ km s<sup>-1</sup>, separately. At any stage, our results will be analysed and presented in our figures in terms of the chosen representative values of $`\rho _l`$ and $`v_0`$, separately.
Another experiment of WIMP direct detection, the CDMS experiment , is now entering the DAMA sensitivity region. The current CDMS upper bounds (either with or without subtractions) concern the upper left corner of the annual–modulation regions, which is anyway poorly populated by susy configurations; thus, they are currently marginal in constraining the set $`S`$. The upper left corner of the annual–modulation regions is also partially disallowed by independent constraints due to indirect measurements (see Sect. V).
Now we turn to the constraints which may be applied to the set $`S`$, using data from WIMP indirect search experiments. We set the limits for exclusion at the same C.L. to which the DAMA region is currently set, i.e. 99.7% C.L.
## III Constraints from neutralino–neutralino annihilation inside Earth and Sun
Indirect evidence for WIMPs in our halo may be obtained at neutrino telescopes by measurements of the up–going muons, which would be generated by neutrinos produced by pair annihilation of neutralinos captured and accumulated inside the Earth and the Sun . The size of the expected muon fluxes strongly depends on how these relic particles are distributed in the phase space and on the intrinsic efficiency of the celestial body in capturing the surrounding WIMPs.
In the case of the Sun the capture rate is essentially determined by its strong gravitational field and by the size of the cross section of neutralino scattering off single protons. Instead, in the case of the Earth the capture process may quite significantly be enhanced by coherent neutralino–nucleus cross sections, whose size depends on mass–matching condition between $`m_\chi `$ and the nuclear mass of the dominant chemical constituents of the Earth (O, Si, Mg, Fe) .
As for the phase-space neutralino distribution in our neighbourhood, together with the usual one based on the standard Maxwellian velocity distribution, whose dispersion speed is centered around 270 km s<sup>-1</sup>, we also consider two intriguing and conflicting models which have been recently discussed in the literature.
Damour and Krauss have proposed the existence of a solar–bound population, with velocities restricted to rather low values, $`v<50`$ km s<sup>-1</sup> (for other papers on hypothetical solar-bound WIMP populations, see Refs. ). The Damour–Krauss solar-bound population would have been produced by WIMPs which scattered off the Sun surface and were set (by perturbations from other planets) into orbits which cross the Earth orbit, but not the Sun. The ensuing velocities would be distributed in the range 25 km s$`{}_{}{}^{1}<v<`$ 50 km s<sup>-1</sup>. This population, although totally irrelevant for the direct measurements by the DAMA NaI–detector, whose electron–equivalent threshold energy is 2 keV, has been shown to be potentially important in making the capture of relic neutralinos by the Earth particularly efficient, with a consequent enhancement of the expected output of up–going muons from the Earth, as compared to the standard one . For simple kinematical reasons, the lower–speed cut off implies that this enhancement is limited to WIMPs of masses lower than $``$ 150 GeV.
On the other side, Gould and Alam , using arguments based on calculations of asteroids trajectories , have pointed out that solar-bound WIMPs could evolve in a way quite different from the one derived in Ref. , with an ensuing suppression of the up-going muon flux usually expected from the center of the Earth for a standard halo population. This suppression would be significant for WIMP masses above $``$ 65 GeV.
In the present paper we take into considerations all of these possible instances. First we consider the standard situation of a Maxwellian velocity distribution over the whole speed range, then we proceed to a critical examination of the other two cases, in which the low-speed interval is either overpopulated or de-populated , as compared to the standard one.
The neutrino spectrum and the ensuing up–going muon flux $`\mathrm{\Phi }_\mu `$ are calculated as explained in Refs. . Their normalization is set by the annihilation rate $`\mathrm{\Gamma }_A`$ of the neutralinos inside the celestial body (Earth or Sun), and $`\mathrm{\Gamma }_A`$ depends, in turn, on the capture rate $`C`$ of the relic neutralinos by the celestial body through the formula $`\mathrm{\Gamma }_A=\frac{C}{2}\mathrm{tanh}^2\left(\frac{t}{\tau _A}\right)`$ , where $`t`$ is the age of the macroscopic body ($`t4.5\mathrm{Gyr}`$ for Sun, Earth) and $`\tau _A=(CC_A)^{1/2}`$, $`C_A`$ is the annihilation rate proportional to the neutralino–neutralino annihilation cross–section and $`C`$ denotes the capture rate. In a given macroscopic body the equilibrium between capture and annihilation (i.e. $`\mathrm{\Gamma }_AC/2`$ ) is established only when $`t>\tau _A`$. Whereas, in the case of the Sun, the capture–annihilation equilibrium is usually reached, due to the much more efficient capture rate due to the stronger gravitational field, for the Earth, the equilibrium condition is not easily realized.
For the case of the standard halo population with a Maxwellian velocity distribution, $`C`$ and $`\mathrm{\Gamma }_A`$ are calculated as in Refs. , and the ensuing muon flux is denoted by $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{s}td}`$. For the Damour–Krauss population the quantities $`C`$ and $`\mathrm{\Gamma }_A`$ are evaluated according to the formulae of Ref. (the relevant muon flux is denoted by $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{D}K}`$). For the model conjectured by Gould and Alam , we have applied to the standard capture rate a suppression factor, which we have re–calculated ab initio in the scheme denoted as ultra–conservative in Ref. , to cover the whole range of masses involved in the present paper. For many susy configurations the suppression factor in the ensuing up–going muon fluxes from the center of the Earth is stronger than the reduction factor in the capture rate alone, due to the relation between $`\mathrm{\Gamma }_A`$ and $`C`$, previously mentioned. For these configurations a reduction in the capture rate induces in the muon flux an extra suppression due to a critical increase in the time required for reaching equilibrium. The muon flux calculated in the Gould–Alam model is denoted here as $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{G}A}`$.
All our neutrino fluxes include neutrino oscillations and use the procedure outlined in Ref. . Here we assume $`\nu _\mu \nu _\tau `$ oscillations, with values for the oscillation parameters which are taken from the best fit performed in Ref. over the whole set of experimental data on atmospheric neutrinos: $`\mathrm{\Delta }m^2=310^3`$ eV<sup>2</sup>, $`\mathrm{sin}\theta =1`$.
Some of our results are presented in Figs. 3–6, where we report various muon fluxes (or ratios of them) versus $`m_\chi `$, for the four representative values of $`\rho _l`$. The solid lines, depicted in Fig. 3 and Fig. 6, denote the 99.7% C.L. upper bounds, $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{lim}`$, derived from the data of the MACRO experiment from the center of the Earth and from the Sun, respectively (for similar limits from the Baksan experiment see Ref. ).
The scatter plots of Fig. 3 display some expected characteristic features, such as the peak at $`m_\chi `$ 50–60 GeV, due to the mass–matching between $`m_\chi `$ and $`m_{\mathrm{F}e}`$. We notice that a number of configurations induce up–going muon fluxes in excess of the experimental bounds. Figs. 4–5 show what would be the enhancement or the reduction effect in $`\mathrm{\Phi }_\mu ^{\mathrm{E}arth}`$ in the case of the Damour–Krauss population or in the Gould–Alam conjecture, respectively. The size of these effects agree with the evaluations in Refs. . For the Damour–Krauss population, the enhancement effect for some susy configurations appears larger here than in Ref. ; this is due to configurations (not considered in ) where rescaling in $`\rho _\chi `$ is effective. In Fig. 6 we display the scatter plots for the up–going muon flux from the Sun, expected for the standard halo population. The current experimental bound sets quite marginal constraints.
In Sect. V we use the results of this section to constrain the susy configurations of set $`S`$. The question, as of which model for the low–speed WIMP population among the two extremes of Refs. is applicable, is still open. Thus, we implement here the experimental bounds on the standard flux of up–going muons; namely, we exclude those configurations, whose $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{s}td}`$ is in excess of the 99.7% C.L. upper bound derived from the MACRO data.
## IV Constraints from cosmic–ray antiprotons
The possibility that annihilation of relic particles in the galactic halo might distort the spectrum of cosmic–ray antiprotons at low–kinetic energies ($`T_{\overline{p}}<`$1 GeV) has been considered by many authors . Indeed, in this energy range, the production of secondary antiprotons by interactions of primary cosmic–ray protons with the interstellar hydrogen has a kinematical drop off , which primary $`\overline{p}`$’s, created by relic neutralinos of appropriate mass and composition, might fill in, at least partially. The effectiveness of this argument to disentangle ordinary spallation contribution from a possible exotic component due to relic particles depends dramatically on how accurately the secondary spectrum is calculated .
This point was addressed in Ref. . In that paper we improved the evaluation of the energy losses undergone by secondary antiprotons during their diffusion inside the Galaxy, we noticed that the as–yet most recent experimental data (BESS95 ) were fitted reasonably well by the secondary spectrum alone, and we examined critically how much room was still available, in the low–energy spectrum, for a contribution from an exotic component. Now, new experimental data (BESS97 ) and improved evaluations of the secondary spectrum further constrain the room left for primary sources. These instances, instrumental in making the separation between primary and secondary antiprotons more difficult, nevertheless confer to the cosmic–ray antiproton measurements a potentially more important role in establishing stringent constraints for relic neutralinos of relatively low mass in our halo, once some of the sizeable, still persisting, uncertainties are reduced.
In the present work we have evaluated the primary antiproton flux, expected from neutralino annihilation, as in Ref. , restricting the supersymmetric configurations to those of set $`S`$. We refer to for all the details concerning the evaluation of the production of these primary antiprotons as well as for the properties related to their propagation in the halo and in the heliosphere. Here we only recall the features of the neutralino mass distribution function adopted in as well as here. This mass distribution function is taken spheroidal and parameterized as a function $`\rho _\chi (r,z)`$ of the radial distance $`r`$ from the galactic center in the galactic plane and of the vertical distance $`z`$ from the galactic plane in the form
$$\rho _\chi (r,z)=\rho _\chi \frac{a^2+r_{}^2}{a^2+r^2+z^2/f^2},$$
(5)
where $`a`$ is the core radius of the halo, $`r_{}`$ is the distance of the Sun from the galactic center and $`f`$ is a parameter which describes the flattening of the halo. Here we take the values: $`a=3.5`$ kpc, $`r_{}=8`$ kpc. In the case of a spherical halo ($`f=1`$), we use the value $`\rho _l=0.3`$ GeV cm<sup>-3</sup>. When $`f<1`$ (oblate spheroidal distribution), $`\rho _l`$ is taken as
$$\rho _l(f)=\rho _l(f=1)\frac{\sqrt{1f^2}}{f\mathrm{Arcsin}\sqrt{1f^2}}.$$
(6)
For each value of $`\rho _l`$ and of the relevant value of $`f`$: $`\rho _l/(\mathrm{GeVcm}^3)`$ = 0.1 ($`f`$ = 1), 0.3 ($`f`$ = 1), 0.5 ($`f`$ = 0.50), 0.7 ($`f`$ = 0.33), we have evaluated the top–of–atmosphere (TOA) antiproton fluxes, as the sum of the secondary flux and of the primary flux due to neutralino annihilation for the various supersymmetric configurations of set $`S`$, pertaining to that specific value of $`\rho _l`$. The secondary flux has been taken from Ref. . Re-acceleration effects in the cosmic rays propagation, which might also be relevant for the features of the secondary antiproton spectrum at low energies , are not included here. Solar modulation has been evaluated according to the procedure discussed in Ref. . We have compared our theoretical results with the combined experimental data of BESS95 and BESS97 , over the whole experimental energy–range (0.18 GeV $`T_{\overline{p}}`$ 3.56 GeV), by a $`\chi ^2`$ calculation.
The results are reported in Fig. 7. In the evaluation of the $`\chi ^2`$, in addition to the experimental errors, we have also taken into account the theoretical uncertainties, estimated according to the results in Refs. , with their appropriate energy dependence. Orientatively, they are in the following ranges: $`\pm `$ (45–55)% for the primary fluxes, $`\pm `$(60–75)% for the secondaries, depending on the energy bin.
In the following, we adopt the selection criterion of excluding from set $`S`$ the configuration whose reduced $`\chi ^2`$ is above the value $`\chi _r^2=2.44`$, which corresponds to a 99.7% C.L. for the 13 d.o.f. of the BESS 95+97 data. From Fig. 7 we notice that, especially at large values of $`\rho _l`$, this constraint disallows a number of susy configurations. The reason why the cosmic–ray antiprotons constraint is not more effective in constraining set $`S`$ is to be attributed mainly to the current large uncertainties affecting the evaluation of antiproton propagation in the galactic halo and in the heliosphere.
## V Results and conclusions
Now we apply the experimental bounds from indirect searches discussed in Sects. III–IV to constrain the supersymmetric configurations of set $`S`$.
### A Combining direct and indirect measurements
Fig. 8 displays the extent of the covering of the annual–modulation regions (one for each value of $`\rho _l`$) by the susy configurations, when the MACRO upper bounds are applied to $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{s}td}`$. A comparison of this figure with Fig. 2 shows that the implementation of these limits somewhat de-populate the covering regions, with a marked effect for the value of the neutralino mass which matches the mass of Iron, as expected. Apart from this, the extent of the regions covered by the scatter plots does not significantly change.
Fig. 9 depicts what would be the effect for a solar–bound WIMP population à la Damour–Krauss. Especially at low values of $`\rho _l`$ there would be some shrinking of the original regions of the scatter plots in their upper parts, but still the annual–modulation regions would be widely covered by physical susy configurations. At variance with this case, the Gould–Alam conjecture would relax the consequences of the constraints applied in obtaining the plots of Fig. 8.
Now, we return to the case where the experimental bounds $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{lim}`$ are applied on $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{s}td}`$. When, on top of these constraints, we also implement the constraints due to cosmic–ray antiprotons, we obtain that the scatter plots of Fig. 8 become somewhat de-populated, but without any appreciable modification in the contours of the covering regions, except for a quite marginal downward shift in their upper–left parts. Therefore Fig. 8 may be considered as the final situation of our analysis, once also the implementation of the antiprotons constraints has been applied. We denote as set $`T`$ the subset of $`S`$ which comprises the susy configurations not disallowed by bounds on the standard up-going muon fluxes and on cosmic–ray antiprotons.
We have analyzed the main properties of the configurations of set $`T`$; some of them are displayed in Figs. 10–11. We recall that the scatter plots of these figures are derived, as all previous ones, by using for the hadronic quantities, discussed in Sect. II, set 1 and set 2, cumulatively. In Fig. 10 we note that the configurations of set $`T`$ cover only a specific region of the susy parameter space not yet disallowed by accelerator constraints. The shape of the distribution of the representative points of $`T`$ in the plot of Fig. 10 is simply explained by the fact that the values of the scalar neutralino–nucleon cross section at the level of the DAMA data require either a large $`\mathrm{tan}\beta `$ or a small $`m_h`$ (or both of these two conditions). This constraint is stronger when the values of the hadronic quantities are restricted to set 1, alone. Fig. 11 displays a correlation among $`m_A`$ and $`m_0`$ which is mainly due to the interplay of these two quantities in generating a light $`m_h`$. Again, restricting the scatter plot to points belonging to set 1, this correlation becomes more pronounced. We recall that, at variance with constrained sugra–supersymmetric models, in the MSSM we are using here, $`m_A`$ and $`m_0`$ are treated as independent parameters.
### B Cosmological properties
We turn now to an analysis of the cosmological properties of relic neutralinos of the susy configurations of set $`T`$. The relevant plots $`\mathrm{\Omega }_\chi h^2`$ vs $`m_\chi `$ are displayed in Fig. 12. It is remarkable that the region of main cosmological interest: $`\mathrm{\Omega }_\chi h^2>0.03`$ turns out to be widely populated, with values of $`\mathrm{\Omega }_\chi h^2`$ which approach, and even exceed, what may be considered as the current upper bound for cold dark matter: $`\mathrm{\Omega }_{CDM}h^2<`$ 0.3 . This means that the DAMA annual–modulation data are compatible with a neutralino as a major component of dark matter. We stress that the scatter plot would even shift upward, should we use for the hadronic quantities discussed in Sect. II the following set: $`y=0.50`$, $`r=36`$, $`m_l<\overline{l}l>`$ = 33 MeV, $`m_s<\overline{s}s>`$ = 585 MeV, $`m_h<\overline{h}h>`$ = 21 MeV. This set of values, denoted as set 3 in Ref. , is more extreme as compared to set 1 and set 2, but still compatible with the current uncertainties.
Finally, we notice that a rather strong de-population in the plots of Fig. 12 is present around $`\mathrm{\Omega }_\chi h^20.01`$ and for large values of $`\rho _l`$. This effect is induced by the cosmic–ray antiproton constraint, since the calculated $`\overline{p}`$ fluxes have their maximal values for $`\mathrm{\Omega }_\chi h^2`$ close to the value below which we apply the rescaling of the local density, i.e. $`(\mathrm{\Omega }h^2)_{\mathrm{m}in}=0.01`$. This property is quite general in this class of calculations, and it was already commented upon, for instance, in Ref..
### C Conclusions
In the present paper we have examined the possibility that the annual–modulation effect, measured by the DAMA Collaboration at a 4$`\sigma `$ confidence level , may be interpreted in terms of relic neutralinos. We have examined this problem, by employing the Minimal Supersymmetric extension of the Standard Model, as a model which does not impose too strong theoretical prejudices on the phenomenological analysis. We have taken into account all experimental constraints, from accelerators and from WIMP indirect experiments.
Let us now summarize our main conclusions:
* The annual–modulation effect mentioned above turns out to be compatible with an interpretation in terms of relic neutralinos.
* The set of supersymmetric configurations selected by the annual–modulation data is only modestly reduced by current experimental data from WIMP indirect searches (up–going muons from the Earth and the Sun, and cosmic–ray antiprotons).
* The set of supersymmetric configurations, selected by the annual–modulation data and not disallowed by the indirect measurements, comprise configurations of relevant cosmological interest, with relic neutralinos playing the role of a major dark matter constituent.
The phenomenological analysis presented in this paper goes beyond the discussion of the experimental data specifically discussed here. We have tried to pin down the most relevant theoretical points, which are still at the origin of large uncertainties, and then require additional investigation. These are: i) size of the Higgs–quark–quark and the squark–quark–neutralino couplings, ii) properties of the WIMP distribution at low velocities (with the possible existence of a solar–bound WIMP population), iii) accurate determination of the propagation in the galactic halo and in the heliosphere for cosmic–ray antiprotons.
###### Acknowledgements.
This work was partially supported by the Research Grants of the Italian Ministero dell’Università e della Ricerca Scientifica e Tecnologica (MURST) within the Astroparticle Physics Project, by the Spanish DGICYT under grant number PB98–0693, and by the TMR network grant ERBFMRXCT960090 of the European Union.
FIGURE CAPTIONS
FIG. 1. Plot of $`\xi \sigma _{\mathrm{scalar}}^{(\mathrm{nucleon})}`$ versus $`m_\chi `$. The solid line delimits the 3$`\sigma `$ C.L. annual–modulation region, obtained by the DAMA NaI(Tl) experiment with a total exposure of 57 986 kg $``$ day . This region was obtained by including the upper-bound constraints of Ref. , by setting $`\rho _l`$ at the standard reference value: $`\rho _l=0.3`$ GeV cm<sup>-3</sup>, and by taking into account uncertainties in the astrophysical velocities of the usual galactic Maxwellian distribution. Also shown in the present figure are the contour lines for the three values $`v_0`$ = 170 km s<sup>-1</sup> (short-dashed (red) line), $`v_0`$ = 220 km s<sup>-1</sup> (long-dash–short-dashed (blue) line), $`v_0`$ = 270 km s<sup>-1</sup> (long-dashed (green) line), separately. The scatter plot is calculated in the MSSM with the scan described in Sect.I; the points of the scatter plot are coded according to the value of the relic abundance, $`\mathrm{\Omega }_\chi h^2`$, of the relevant susy configuration: dots denote $`\mathrm{\Omega }_\chi h^2<0.01`$, crosses denote $`0.01<\mathrm{\Omega }_\chi h^2<0.1`$ and empty circles denote $`\mathrm{\Omega }_\chi h^2>0.1`$.
FIG. 2. Location of the DAMA annual–modulation region for four representative values of $`\rho _l`$: $`\rho _l`$ = 0.1, 0.3, 0.5, 0.7 GeV cm<sup>-3</sup>. The scatter plots show only the configurations which lie inside the relevant annual–modulation region. A grey–level (color) code is used depending on the value of $`v_0`$ employed in the extraction of the annual–modulation region: medium grey (red) denotes points which lie in the annual–modulation region extracted by setting $`v_0`$ = 170 Km s<sup>-1</sup>, dark grey (blue) denotes points which lie in the annual–modulation region extracted by setting $`v_0`$ = 220 Km s<sup>-1</sup>, light grey (green) denotes points which lie in the annual–modulation region extracted by setting $`v_0`$ = 270 Km s<sup>-1</sup>. The three sets are superimposed in that sequential order.
FIG. 3. Scatter plot for the up–going muon flux from the center of the Earth for a standard Maxwellian distribution, $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{s}td}`$, versus $`m_\chi `$. The grey–level (color) code is the same as in Fig. 2. The solid line denotes the 99.7% C.L. upper bounds, $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{lim}`$, derived from the data of the MACRO experiment .
FIG. 4. Enhancement effect in the up–going muon flux from the center of the Earth in case of a solar–bound population à la Damour–Krauss . The grey–level (color) code is the same as in Fig. 2.
FIG. 5. Suppression effect in the up–going muon flux from the center of the Earth in case of the Gould–Alam conjecture . The grey–level (color) code is the same as in Fig. 2.
FIG. 6. Scatter plot for the up–going muon flux from the Sun for a standard Maxwellian distribution, $`\mathrm{\Phi }_\mu ^{Sun}`$, versus $`m_\chi `$. The grey–level (color) code is the same as in Fig. 2. The solid line denotes the 99.7% C.L. upper bounds, derived from the data of the MACRO experiment .
FIG. 7. Scatter plot for the reduced $`\chi _r^2`$’s in a comparison of the calculated cosmic–ray antiprotons fluxes with the combined experimental data of BESS95 and BESS97 . The horizontal line denotes the value $`\chi _r^2`$ = 2.44, which for 13 d.o.f. corresponds to a 99.7% C.L., above which we disallow susy configurations. The grey–level (color) code is the same as in Fig. 2.
FIG. 8. As in Fig. 2, once the constraints from the up–going muon fluxes from the center of the Earth are applied, assuming a Maxwellian halo distribution for relic neutralinos. The grey–level (color) code is the same as in Fig. 2.
FIG. 9. Covering of the annual–modulation regions, if the constraint $`(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{D}K}(\mathrm{\Phi }_\mu ^{\mathrm{E}arth})^{\mathrm{l}im}`$ were applied. The grey–level (color) code is the same as in Fig. 2.
FIG. 10. Scatter plot for set $`T`$ in the plane $`m_h`$$`\mathrm{tan}\beta `$. The grey–level (color) code is the same as in Fig. 2. For each panel, the lower dashed line denotes the frontier of the complete scatter plot; the upper dashed line denotes the frontier, when only set 1 for the hadronic quantities of Sect. II is employed. The hatched region on the right is excluded by theory. The hatched region on the left is excluded by present data from LEP and CDF . The solid line represents the 95% C.L. bound reachable at LEP2, in case of non discovery of a neutral Higgs boson.
FIG. 11. Scatter plot for set $`T`$ in the plane $`m_0`$$`m_A`$. The grey–level (color) code is the same as in Fig. 2. For each panel, the upper dashed line denotes the frontier of the complete scatter plot; the lower dashed line denotes the frontier, when only set 1 for the hadronic quantities of Sect. II is employed.
FIG. 12. Neutralino relic abundance $`\mathrm{\Omega }_\chi h^2`$ versus $`m_\chi `$, once the constraints from up–going muon fluxes and cosmic–ray antiprotons are applied. The hatched region is disallowed by the upper limit on cold dark matter $`\mathrm{\Omega }_{CDM}h^2<`$ 0.3 . |
warning/0001/gr-qc0001048.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The discovery of black hole thermodynamics represents a milestone in the search of a consistent unification of the principles of quantum mechanics and general relativiy ,. It combines quantum mechanics, general relativity and thermodynamics in a unique and fascinating picture (for a recent review and discussion, see e.g. ,). Yet, the unification between the ideas of general relativity and quantum mechanics is not done within a consistent framework, but recent developments in string–theory and its relation to M–theory may lead to a self–consistent picture. It is interesting to note that some of those developments are sparked also by investigations of black holes in the context of these models.
While the origins of the laws of black hole thermodynamics are unknown, they seem to enforce an upper bound on the number of states within a volume ,. This so–called holographic principle states that the maximal entropy within a volume $`V`$ in space is bounded by its surface area $`A`$, according to the Bekenstein–Hawking formula
$$S_{BH}=\frac{A}{4l_p^2}$$
(1)
where $`l_p`$ is the Planck length. It implies that all the physical degrees of freedom are somehow encoded on the surface $`A`$. A somewhat stronger statement of this principle would be that a theory which includes gravity, describing phenomena within a volume, can be reformulated as a theory which describes the evolution of the degrees of freedom on the boundary without gravity. We refer to this as the weak holographic principle. This has been supported at least in some special cases within the framework of the Anti–de Sitter(AdS)/CFT–correspondence and the M(atrix) model . However, the physical origin of the holographic principle remains mysterious.
In a recent interesting paper, ’t Hooft gave a dramatic interpretation of the holographic principle (using arguments from black hole physics): the fundamental degrees of freedom of nature are not quantum mechanical but rather deterministic, in a certain sense classical degrees of freedom . Due to information loss these primary degrees of freedom will evolve into a set of equivalence classes which (in his definition) evolve unitarily and span a Hilbert space. The maximal number of equivalence classes is given by the Bekenstein–Hawking formula (1). The information loss mechanism was left open in his paper but in the case of classical general relativity the information loss could be provided by black holes, as described by ’t Hooft.
If we combine his idea together with the holographic principle in a more general framework, we would conclude that
A classical theory with gravity within a volume $`V`$ can be formulated as a quantum theory with the degrees of freedom living on the boundary $`V`$. The number of quantum states on the boundary is given by (1).
We will refer to this principle as the strong holographic principle. According to ’t Hooft the quantum degrees of freedom are not fundamental. While “projecting” the classical theory onto the boundary (screen), information about the classical states would be lost. This, however, has to be formulated in a more quantitative way: how does the mapping from the bulk onto the boundary theory work? Is this mapping unique or is it different for different situations? In any case, spacetime has to know how to dissipate information.
We should mention that the principle, as formulated above, is supposed to be valid in a holographic spacetime, i.e. where the dual theory on the screen exist. We will come back to other spacetimes, such as those found in cosmology, in a later section. There we will give also the argument why the quantum degrees of freedom are not fundamenal.
Our every-day world is certainly described by quantum mechanics. According to the strong holographic principle we could somehow describe quantum degrees of freedom in our 3+1$`D`$ world with a classical theory (including information loss) in a higher–dimensional spacetime, whose “boundary is our world”. For example, the Einstein–Podolsky–Rosen correlation between states, confirmed by experiments in laboratories in our $`3+1D`$ world, should be classically explainable in a higher–dimensional space–time. Certainly gravity may play a fundamental role here as well as the concept of time.
The question we address in this paper is: Can we find a unique way to map the states of the classical theory onto the quantum states in the lower–dimensional boundary (spacetime)? It would certainly be more satisfying if there would be a unique relationship between the holographic (bulk) theory and the dual (screen) theory, rather than that this relation has to be found seperately for different spacetimes. The aim of this paper is to make the first step to find this correspondence between the holographic and the dual theory. We assume that this relation should be unique (whenever the dual theory exists) and give arguments that this is indeed the case.
Our starting point is the observation that there are similar correspondences between classical theories and quantum field theories (QFT) even without the inclusion of gravity. These are (we set $`c=G=1`$):
* The well–known correspondence between the partition function with periodic/
antiperiodic boundary conditions for a euclidean quantum field theory in $`D`$ dimensions and the partition function for a classical thermal field theory with temperature $`kT=\mathrm{}`$, where $`k`$ is the Boltzmann constant.
* Related to this is the stochastic quantization method. It relates a $`D+1`$ dimensional classical theory (which includes a stochastic noise) to a quantum theory in $`D`$ dimensions. In the higher–dimensional spacetime a stochastic noise plays a fundamental role.
The origin of these relationships is not known, but there are strikingly similar to the strong holographic principle. An immediately question we ask is therefore, if the strong holographic principle is related to these well known correspondences. As we will argue in a later section, the correspondence between QFT and classical theories mentioned before is a result of the holographic principle in the limit of vanishing gravity.
The paper is organized as follows: in Section 2 we review shortly the concept of stochastic quantization. In Section 3 we discuss the strong holographic principle in the limit of a flat spacetime. The important case for the Anti–de Sitter (AdS) spacetime is discussed in Section 4. In Section 5 we comment on a scalar particle in the AdS spacetime and its stochastic quantization. In Section 6 we extend our ideas to a spacetime in which the area of the screen is not constant. Our conclusions can be found in Section 7, as well as further questions which arise in the context of the ideas presented in the paper.
In our discussions we are mainly guided by black hole physics as well as expanding spacetimes, such as those which can be found in cosmological theories. However, our results are supported by a covariant formulation of the holographic principle ,. In this paper we take a rather heuristic view and formulate the ideas not in a mathematical language. We will use the terms “boundary” and “screen” interchangeably.
We mention that other groups are also asking for the mechanism of holography, in particular see ,.
## 2 Review of stochastic quantization
As mentioned in the introduction, the relation between the holographic theory and the dual theory should be unique. It is useful for the discussions in the later sections to review very briefly the well known relationship between classical and quantum field theories mentioned in the introduction. In what follows, we mention only the necessary points and refer to the exellent reviews ,.
The starting point is the fact that the Euclidean Green function for a scalar field $`\varphi `$ can be interpreted as a correlation function of a statistical system in equilibrium of temperature $`T=\mathrm{}/k`$. The Euclidean Green function is<sup>1</sup><sup>1</sup>1We consider here the simple example for a scalar field. Extensions to more complicated theories such as gauge field theories exists ,.
$$\varphi (x_1)\varphi (x_2)\mathrm{}\varphi (x_n)=\frac{𝒟\varphi \mathrm{exp}\left((1/\mathrm{})S_\mathrm{E}\right)\varphi (x_1)\varphi (x_2)\mathrm{}\varphi (x_n)}{𝒟\varphi \mathrm{exp}\left((1/\mathrm{})S_\mathrm{E}\right)}.$$
(2)
Here $`S_\mathrm{E}`$ is the Euclidean action. It was the highly ingenious idea by Parisi and Wu to interpret the Euclidean path integral measure as the stationary distribution of a stochastic process . This is the basic idea of the procedure known as stochastic quantization. In hindsight it is a rather natural interpretation of Feynman path integrals.
The field $`\varphi (x)`$ in Euclidean space with coordinates $`x`$ is now generalized and will be considered as a function of the Euclidean coordinates $`x`$ and a new fictitious time–coordinate $`t`$: $`\varphi (x)\varphi (x,t)`$. This field couples to a thermal bath. Let $`\eta `$ be a Markov stochastic variable, representing the coupling of the system to this thermal bath, with temperature $`T`$
$`\eta (x,t)`$ $`=`$ $`0;`$ (3)
$`\eta (x_1,t_1)\eta (x_2,t_2)`$ $`=`$ $`2\alpha \delta (x_1x_2)\delta (t_1t_2),`$
where $`\alpha `$ is the diffusion constant, connected with the temperature $`T`$ (and in general with a friction constant $`f`$) via
$`\alpha ={\displaystyle \frac{kT}{f}}`$
The reason why the new coordinate $`t`$ is called time is that one imagines that the system in $`D+1`$ dimensions evolves in this time. In order that we obtain the usual quantum mechanical expressions we have to set $`\alpha =\mathrm{}`$.
The basic equation of the stochastic quantization method is the Langevin equation
$$\frac{\varphi }{t}=\frac{S_E}{\varphi }+\eta (x,t),$$
(4)
where $`S_E`$ is the action of the field $`\varphi (x,t)`$:
$$S_E=𝑑x(\varphi ,_x\varphi ).$$
(5)
Here, $``$ is the Langrangian density. It has the form of the original Langrangian, but now one has to replace the field accordingly. There is no derivative with respect to the new time–coordinate $`t`$. Correlations are now defined as a average over the noise $`\eta `$. Then, in this framework, quantum correlation functions in Euclidean space are obtained in the limit $`t\mathrm{}`$:
$$\mathrm{}=\underset{t\mathrm{}}{lim}\mathrm{}_\eta :=\underset{t\mathrm{}}{lim}\frac{𝒟\eta \mathrm{exp}\left(\frac{1}{4}𝑑x𝑑t\eta ^2(x,t)\right)\mathrm{}}{𝒟\eta \mathrm{exp}\left(\frac{1}{4}𝑑x𝑑t\eta ^2(x,t)\right)},$$
(6)
where the dots represent solutions of the Langevin equation (4). One can show that as a result of the evolution of the system within the thermal bath the resulting equilibrium distribution is
$$𝒫(\varphi )\mathrm{exp}\left(\frac{S_E(\varphi )}{\mathrm{}}\right).$$
(7)
In this equation, the action $`S_E`$ is evaluated for those $`\varphi `$ which satisfy the Langevin equation.
Equivalently one can find a differential equation, the Fokker–Planck equation, for the probability distribution $`𝒫(\varphi )`$ of the stochastic process at the time $`t`$. Then quantum Green functions are obtained as
$$F(\varphi )=\underset{t\mathrm{}}{lim}𝒟\varphi F(\varphi )𝒫(\varphi ,t).$$
(8)
The time evolution for the probability distribution is described by the Fokker–Planck equation of the form
$$\frac{}{t}𝒫(\varphi ,t)=\frac{d}{d\varphi }\left[\frac{d}{d\varphi }+\frac{\delta S_E}{\delta \varphi }\right]𝒫(\varphi ,t)$$
(9)
In fact, the stochastic quantization method using the Langevin equation is equivalent to the approach starting from the Fokker–Planck equation. We refer to the existing literature.
It should be mentioned that the stochastic quantization method can not only be formulated in Eudlidean space but also in Minkowski spacetime. Here, the Langevin equation becomes
$$\frac{\varphi }{t}=i\frac{S}{\varphi }+\eta (x,t),$$
(10)
where $`S`$ is the action in Minkowski spacetime. We will, however, mainly work in the Euclidean formalism.
It is clear that the equilibrium limit (6) or (8) must exist in order to make sense for this quantization procedure. The reason why the stochastic quantization method works, is yet unknown. It is usually taken as a formal manipulation for field theories and it was used intensively for computer calculations. In what follows we will argue that there is a deeper reason why this procedure works.
## 3 The strong holographic principle in the limit of zero gravity
### 3.1 General considerations
The holographic principle, as originally formulated by ’t Hooft, is valid for any size and kind of black hole. In what follows we will discuss the case for a Schwarzschild black hole. It is well known that the surface gravity $`\kappa `$ for such a black hole is inversely proportional to its mass. Therefore, the larger the black hole is, the smaller its surface gravity. To be precise,
$$\kappa =\frac{1}{4M},$$
(11)
where $`M`$ is the black hole mass, see e.g. . For a huge black hole the surface gravity is very small, and in the case for $`M\mathrm{}`$, $`\kappa `$ approaches the value zero. But for a observer from the outside, all degrees of freedom are still located at the horizon of the black hole.
Inside a massive black hole, gravity is less important than it is the case for smaller black holes. Take for example the volume $`𝒱`$ within which the relation $`R_{ijkl}<b`$ holds, where $`R_{ijkl}M/r^3`$ is the curvature tensor and $`b`$ is some positive constant. The ratio of the volume $`𝒱r^3`$ and the black hole volume ($`r_s^3`$) will be smaller the larger the black hole is:
$$\frac{𝒱}{r_s^3}\frac{1}{M^2}.$$
(12)
On the other hand, let a observer sit at a position $`r=vr_s`$, where $`v`$ is some constant. Then (see e.g. ,)
$$R=\mathrm{curvature}\text{ }\mathrm{scalar}\frac{M^2}{r^6}\frac{M^2}{vM^6}\frac{1}{M^4}$$
(13)
for consant $`v`$. Therefore, for an observer, sitting at a constant ratio $`r/r_s`$, gravity will become weaker if the mass of the black hole grows. Of course, in the vicinity of the singularity (if there is any) gravity is important. However, we neglect for a moment this part of the black hole, because the singularity is not a problem for what follows<sup>2</sup><sup>2</sup>2The asymptotic flatness (see eq. (11)) is important.. (Although we will come back to the AdS spacetime in the next section, we should mention this example, too. The relation between a classical supergravity in the AdS spacetime and a quantum supersymmetric Yang–Mills is valid, even if there is a black hole embedded in the AdS–spacetime. What is important is the asymptotic form of the spacetime, which should be AdS. And also, if the radius of the spacetime goes to infinity, the space becomes (at least globally) flat, i.e. gravity/curvature goes to zero (for a discussion of this limit, see ).)
If the strong holographic principle mentioned in the indroduction makes sense, it should be independent of the size of the region, as much as the weak holographic principle should be independent of the size of the black hole and independent of the size of the AdS spacetime. We will argue therefore, that the strong holographic principle is valid even in the limit when gravity goes to zero<sup>3</sup><sup>3</sup>3It will become clear from our discussions that this should be the only way to define the Minskowski–space., i.e.
The strong holographic principle is valid also in asymptotically flat spacetimes.
If we make now the reasonable assumption that the relationship between classical theory and the quantum theory is valid independent of the spacetime and that the dual theory exists, we postulate that the relationship in general has to be given by the stochastic quantization procedure. We mention here that the conclusion above is not trivial. The strong holographic principle is a relation between a classical theory in a volume and a quantum theory on a boundary/screen and is different from the original formulation, which states that the number of quantum degrees of freedom is bounded by the area of the boundary.
Thus, the thermodynamics of spacetime maybe the origin of the well known correspondence mentioned in Section 2. That the lower dimensional theory is a quantum theory now follows from the Bekenstein–Hawking formula (1) with one degree of freedom per Planck unit. Again, our conclusion is valid only in the case of a spacetime, where the limit of the stochastic quantization procedure exist. We will return to the implications of other spacetimes below. We stress again that the asymptotic behaviour of the spacetime was important in the discussion above.
### 3.2 Stochastic quantization in Euclidean/Minkowski spacetime and the strong holographic principle
Let us shortly discuss the stochastic quantization procedure in flat Euclidean spacetime in the light of the strong holographic principle.
As emphasized by Bousso, the screen of Minkowski–space is either future or past null infinity when gravity is negligible anywhere. A spacelike projection is then allowed as well. The process of stochastic quantization procedure is in effect a projection according to the covariant holographic principle. The boundary theory should have a “time–coordinate”, so our intuition suggests that the projection should be spacelike. This is indeed possible, because the coordinate $`t`$ in Section 2 has no meaning. Nowhere was it stated that the higher–dimensional spacetime should have signature $`(D,1)`$ or $`(D,2)`$. The projection itself is irreversible when averaged along $`t`$.
However, if we take the Euclidean formulation of stochastic quantization, then the situation is as follows. Consider the $`D+1`$ dimensional Minkowski–space. The projection of the theory via stochastic quantization along the $`t`$–coordinate in this space is a dimensional reduction to a $`D`$ dimensional euclidean field theory, i.e. it gives a euclidean quantum field theory at the screen, which is in this case $`I^+`$.
## 4 The ADS/CFT correspondence and the strong holographic principle
We have argued that, whenever the dual theory exists, it is related to the holographic theory via stochastic quantization. The most impressive and explicit example, where a dual theory exists, is the relation between classical supergravity in the AdS spacetime and a supersymmetric Yang–Mills theory on the boundary . More specifically, the mathematical formulation of this correspondence is the equivalence between the partition functions of both theories ,:
$$Z_{\mathrm{AdS}}(\varphi _{\mathrm{bulk}})=Z_{\mathrm{CFT}}(\varphi _{\mathrm{bound}}).$$
(14)
Here, $`\varphi _{\mathrm{bulk}}`$ are the fields in the bulk–theory (a supergravity theory) taken at the boundary and $`\varphi _{\mathrm{bound}}`$ are the fields in the boundary–theory (a supersymmetric Yang-Mills theory). The equation above can be written as<sup>4</sup><sup>4</sup>4We note here that the original motivation was to relate superstring–theory on a AdS$`{}_{5}{}^{}\times S^5`$ to a supersymmetric Yang–Mills theory on the four–dimensional boundary of the AdS–spacetime. In the equations here the partition function of the string theory is approximated as the supergravity action. In fact, this is yet the only approximation where a mathematical formulation of the correspondence exist.:
$$\mathrm{exp}\left(_{\mathrm{AdS}}_{\mathrm{supergravity}}(\varphi _i(\varphi _i^{\mathrm{bound}}))\right)=\mathrm{exp}_{\mathrm{AdS}}𝒪^i\varphi _i^{\mathrm{bound}}_{\mathrm{CFT}}$$
(15)
From the point of view of the holographic (bulk) theory, the $`\varphi _i^{\mathrm{bound}}`$ represent the boundary values of the fields $`\varphi _i`$. The integral on the left–hand side represent here the classical action for the supergravity theory on the AdS spacetime with $`d+1`$ dimensions evaluated at the boundary. On the right–hand side we have the quantum expectation value of the primary fields $`𝒪_i`$ of a conformal theory on the boundary, where the boundary values $`\varphi _i^{\mathrm{boundary}}`$ act as an external source. We mention that the radius $`R`$ of the AdS–spacetime is related to the number of colors $`N`$ and the coupling strength $`g_{\mathrm{YM}}`$ in the Yang–Mills theory via $`R/l_s=(Ng_{\mathrm{YM}}^2)^{1/4}`$.
If our arguments are correct, the supergravity theory in the AdS spacetime and the supersymmetric Yang–Mills theory on the boundary should be related by stochastic quantization. In the presence of a gravitational field this procedure will be modified, as we will discuss in the next section. In the spirit of stochastic quantization and with the wisdom of hindsight one dimension of the AdS spacetime can be identified with the fictitious time and it is only natural to identify this coordinate with the radial coordinate $`r`$ of the AdS–spacetime<sup>5</sup><sup>5</sup>5This is agreement with the space–like holographic projection described in the work by Bousso . A similar remark has been made in the work by Lifschytz and Periwal . Their work was connected with the duality between string theories and gauge theories and the approach there was different from ours and based on the equivalence of the Fokker–Planck Hamiltonian for Yang–Mills theories and the loop operator (see also the work by Jevicki and Rodrigues ). Because the Fokker–Planck equation is at the heart of stochastic quantization, Lifschytz and Periwal speculated about the importance of this procedure in the context of the AdS/CFT correspondence. We believe that a mathematical rigerous proof of our arguments involve indeed the Fokker–Planck Hamiltonian.. In fact, if we insist that during the replacement $`\varphi (x)\varphi (x,r)`$ the action should be invariant under a supersymmetry, then the spacetime $`(x,r)`$ has to be compatible with this supersymmetry, which in this case is the AdS spacetime.
In stochastic quantization in Euclidean space with coordinates $`x`$ and fictitious time $`t`$ one calculates Green functions as an equilibrium limit of the fictitious time, that is (see eq.(6))
$$<P(\varphi (x))>=\underset{t\mathrm{}}{lim}<P(\varphi (t,x))>_\eta ,$$
(16)
where $`<\mathrm{}>_\eta `$ is the stochastic average and $`P`$ is some polynomial of the fields $`\varphi `$. Now, in the AdS–case the limit $`rr_{\mathrm{boundary}}`$, correspond to the boundary, where the dual theory lives. Here we see a geometrical picture for the stochastic quantization method emerging which is not obvious in the case for Euclidean or Minkowski–spacetime (in some coordinate–systems $`r_{\mathrm{boundary}}`$ is infinite). In the presence of gravity the fictitious coordinate can be a usual spatial coordinate. One may worry that in this case the fictitious “time–coordinate” is now physically important because fields propagate through it and gravity curves it. But because gravity (by holography) localizes the quantum degrees of freedom (on the black hole horizon for example) this is what one should expect.
What remains to be shown is that $`r`$ can indeed play the role as the fictitious coordinate, that the equilibrium limit exists and that the corresponding theory at $`rR`$ is a quantum supersymmetric Yang–Mills theory. What equation (6) (with $`t=r`$) then tells us is actually just the statement that correlation functions on the boundary are given by the thermal average in the higher–dimensional spacetime with coordinates $`r,x`$ and carrying this to the boundary to the AdS–spacetime. This, by the AdS/CFT–correspondence, has to be the correlation function of the fields on the boundary<sup>6</sup><sup>6</sup>6Here, we can couple the fields onto the boundary. We just have to add coupling terms in the Langragian..
Although we have argued that at the heart of the AdS/CFT–correspondence is the procedure of stochastic quantization, it is clear that every theory in a AdS spacetime should be related to a quantum theory on the boundary. The AdS/CFT correspondence itself is only a special case. We see no reason, why the equilibrium limit for a generalized Langevin equation would not exist and that the whole AdS/CFT–correspondence cannot be formulated as a problem of stochastic quantization. While our discussion was heuristic, we believe that a mathematical proof exists. This, however, is beyond the scope of this paper.
## 5 A scalar field in AdS
In this section we give a formal argument that in the case of a scalar field our ideas are justified and that the coordinate $`r`$ can indeed play the role of the fictitious coordinate. We show that there is a simple generalization of the procedure described in section 2. We are looking now for a stochastic process in the higher dimensional spacetime, described by a (generalized) Langevin equation. Here we emphazise the physics only, the details of the calculations can be found in the appendix.
As discussed in the appendix, we lift the theory into a higher–dimensional spacetime and assume the existence of a stochastic process there. The arguments there are general and we get a generalized Langevin–equation of the form
$$d\varphi =f(r)\frac{S_E}{\varphi }dr+dW,$$
(17)
with
$`<dW>`$ $`=`$ $`0`$ (18)
$`<dW(x,r)dW(x^{},r^{})>`$ $`=`$ $`2f(r)\delta (xx^{})\delta (rr^{})dr.`$ (19)
In the case of the AdS, the coordinate will now be interpreted as the radial coordinate $`r`$. The metric of this spacetime can be written as
$$ds^2=R^2\left[\frac{4\eta _{\mu \nu }dx^\mu dx^\nu }{(1r^2)^2}+dt^2\frac{1+r^2}{1r^2}\right],$$
(20)
with $`r=x_\mu x^\mu `$ and $`R`$ is the curvature radius of the AdS spacetime. As the field is projected along $`r`$ it is subject to a thermal bath, described by the noise–term in the Langevin equation (17). It should be noted that for this “fictitious” process $`r`$ is a time–coordinate.
$`f(r)`$ is a smooth function of the radial coordiate $`r`$ only. Here we find a new ingredient in the theory to be discussed: Whereas in the “usual” stochastical quantization procedure the diffusion constant is a real constant, gravity will cause this parameter to be different from point to point. A picture might be intuitive (see figure): The thermal bath representing the noise $`\eta `$ has a constant temperature along the fictitious coordinate $`t`$. In the case of the AdS spacetime, and in curved spacetimes in general, this temperature will be a function of the fictitious coordinate, here the radial coordinate $`r`$. As discussed in the appendix, the fluctuation–dissipation theorem should hold locally.
The probability distribution can be found to be:
$$𝒫(\varphi )=A\mathrm{exp}(S_E(\varphi (x,r))).$$
(21)
Here, $`\varphi `$ is a stochastical field for which $`\varphi (x,r_{\mathrm{boundary}})=\varphi (x)`$, where $`\varphi (x)`$ is the boundary fields. The constant $`A`$ fixes the normalization and therefore we find:
$$𝒫(\varphi )=\frac{e^{S_E(\varphi )}}{𝒟\varphi e^{S_E(\varphi )}},$$
(22)
where $`\varphi `$ are solutions of the Langevin–equation (17). In conclusion, have found that the field on the boundary of the AdS spacetime has the usual quantum mechanical expectation values calculated with the Feynman measure for the path integrals.
The example in the Appendix suggests that $`f(r)\sqrt{g(r)}`$, where $`g`$ is the determinant of the metric of the AdS spacetime, which is a function of $`r`$ only. More importantly, the example suggest further, that the metric (inside the boundary) itself is not important. The result (45) can be obtained by stochastic quantization in Minkowski space as well as in the AdS case. As pointed out by ’t Hooft, the holographic principle implies that the geometry inside a volume indeed is unimportant . In our approach this is connected to the topological origin of stochastic quantization (see the discussion in e.g. , , and ). In fact, it was shown that the stochastically quantized theory is equivalent to a topological field theory. Furthermore, we recall here the well known fact, that, at least in some cases, a classical stochastic process can be reformulated as supersymmetric quantum mechanical problem, see e.g. and and references therein.
To complete our discussion, we have to consider the entropy bound (1) in the framework of our theory. Unfortunately we were not able to find a way how our ideas lead to a derivation of the entropy bound. What we can only say is that the theory on the boundary has to be a quantum theory. Because of the covariant entropy conjecture formulated by Bousso and the modified version by Flanagan, Marolf and Wald the bound should be satisfied. All what has to be assumed is that the conditions for the holographic projection has to be fulfilled. The projection via stochastic quantization in the case of the AdS spacetime disussed here is a spacelike projection in the sense of Bousso. It is interesting to speculate that the saturation of the entropy (see eq. (1)) at the boundary is connected to the fact that the stochastic process reaches an equilibrium state where the entropy simply does not grow.
A crucial point in the discussion so far is that the “temperature” (or better: diffusion coefficient) along the boundary should be independent of the euclidean time, otherwise a sensible limit would not exist. This is important for what follows.
## 6 Holography in cosmological spacetimes: implications
Our discussion so far was based on the case of a spacetime where the holographic theory and the dual theory exist. Certainly our universe is dynamical and expanding.
That the weak holographic principle has to be modified in more general spacetimes, such as those in cosmology for example, was discussed by Fischler and Susskind , Bak and Rey , Veneziano , Easther and Lowe , and Kaloper and Linde , see also the discussions in ,,,,. Based on these earlier ideas, Bousso formulated a general covariant holographic principle . In what follows we will not be able to give a general theory. Rather we will discuss the implication for the suggestion by ’t Hooft, that the fundamental degrees of freedom are not quantum mechanical.
The duality between quantum and classical theories, as stated by the strong holographic principle, answers the question which theory is “more fundamental”, because at this point both theories are not on the same footing. We consider several gedankenexperiments in this section which makes this point clear.
Suppose that our universe was in its earlier epoch in a AdS–spacetime state and (approximately) static. All degrees of freedom live on the boundary of the AdS–spacetime and all processes in the bulk can be described by the boundary theory. Let there now be a process which turns the state of the universe into another one, say a matter dominated phase. This can happen for example if the cosmological constant becomes positive through some dynamical processes and decays into particles. What happens to the boundary in this process? An observer in the bulk will be able to see only a part of the de Sitter–space. The horizon is $`H^1`$, where $`H`$ is the expansion rate. This horizon gets dynamical when the universe becomes matter dominated. In such cases it was proposed that the holographic principle should be replaced by the (generalized) second law . This would imply that it is not longer useful to talk about a quantum boundary theory, because degrees of freedom maybe created or destroyed. It is difficult to see if such a theory is compatible with the second law as well as with unitarity in general; in short: there will be no quantum mechanical degrees of freedom.
What happens to the bulk theory? It would hardly makes sense if the theory, which is classical in the beginning becomes now “more and more” quantum. It is more likely that the theory remains classical. Of course, our discussion implies that the theory on the boundary can only exist if the holographic theory exist, because the spacetime in this theory was the starting point. This implies, however, that the quantum degrees of freedom are not the fundamental ones but the classical degrees of freedom in the bulk. This is in agreement with ’t Hooft’s proposal (and the philosophy of stochastic quantization itself).
Another example would be a black hole. We call all degrees of freedom on the boundary quantum mechanical when the black hole surface is static and non–growing, but we just have to throw matter in the black hole to destroy the quantum character of the degrees of freedom. During this process, the total entropy will grow.
Given the ideas in this paper, the fundamental difference between a spacetime with constant screen area and a general expanding spacetime seems to be that the former allows for a limit of the stochastic process, i.e. the stochastic process described by the variable $`\eta `$ is in local thermodynamical equilibrium and the procedure of stochastic quantization makes sense, i.e. the limit $`t\mathrm{}`$ exists (see eq. (6)). As explained in the last section, the temperature for $`T(rr_{\mathrm{boundary}(1)})`$ and the temperature $`T(rr_{\mathrm{boundary}(2)})`$, where the fictitious coordinate $`r`$ runs from $`r_{\mathrm{boundary}(1)}`$ to $`r_{\mathrm{boundary}(2)}`$, must be constant. However, there is no reason to believe that this is always the case. Non–equilibrium processes in the holographic theory are related to the fact, that, to use the words of ’t Hooft, the equivalence classes, which may form, do not evolve not unitarily in this case and the corresponding quantum theory does not exist. In this context it is interessting to note that Marchesini has shown that the loop equations of non–abelian gauge theories are equivalent to the equilibrium condition within the context of stochastic quantization . In fact, stochastic quantization only makes sense (and is defined in that way) as an equilibrium limit (for $`t\mathrm{}`$). In conclusion: a dual theory always exists when a sensible equilibrium limit for $`t\mathrm{}`$ exists, because then stochastic quantization is applicable, and the condition for that is connected to the spacetime structure, described by the holographic theory. The screen, on which the degrees of freedom are projected, has to allow for a sensible limit<sup>7</sup><sup>7</sup>7The screen in AdS for example allows for such a projection.. We see here a connection of our argumentation to the work by Easther and Lowe, who in particular argued that in the case of general spacetimes the holographic principle has to be replaced by the second law: if the entropy of a system grows, it is not in thermodynamical equilibrium and hence a stochastic quantization procedure makes no sense. The black hole example mentioned above is a good example for what is going on here.
It is interesting to speculate on the role of supersymmetry. Supersymmetry is connected to stochastic processes in a subtle way , because of the hidden supersymmetry in the Langevin and Fokker–Planck equation. It is well known that supersymmetry is possible only in certain spacetimes, such as AdS, but not, for example, in an expanding universe. But here, stochastic quantization breaks down, too.
We can draw an important conclusion from the discussion above: In some very strong time–varying gravitational fields, where the holographic principle makes no sense and has to be replaced with the generalized second law, we expect significant deviations from the quantum mechanical predictions. This is because the relationship between quantum–mechanical and classical theories is lost in those strong time–variating gravitational fields, similar to the process of dropping matter into a black hole. There the degrees of freedom might be destroyed or generated, which differs from ordinary quantum mechanical behaviour. One possible observational consequence would be the violation of unitarity in our $`3+1D`$ world in such strong time–varying gravitational fields.
## 7 Conclusions and outlook
The strong holographic principle mentioned in the introduction is stronlgy connected to the well known relationships between classical field theories, statistical mechanics and quantum theory in Minkowski/Euclidean space. In this paper we gave some arguments why this is the case. We have argued that whenever the dual boundary theory exists, the relationship between the holographic (bulk) theory and the dual (screen) theory should be unique, i.e. independent of the spacetime, so that the strong holographic principle is valid even in the limit of vanishing gravity (curvature). Because the Minkowski spacetime can (and should) be viewed as a limit process of vanishing gravity, we argued that one can here also find a relationship between classical and quantum mechanics. The natural candidate we propose is the stochastic quantization procedure, which relates a $`D+1`$–dimensional stochastic classical theory to a $`D`$–dimensional quantum theory. This relationship should hold whenever a dual theory exist, especially in the case of the AdS/CFT correspondence. We have argued (but not shown), that the AdS/CFT–correspondence can indeed be seen as a process of stochastic quantization, where this process has a geometrical meaning in the AdS–spacetime. We have argued, that the radial coordinate of the AdS spacetime can act as the stochastic time. Furthermore, we argued that stochastic quantization has the property that the (smooth) interior metric has no effect on the physics on the boundary. This property of holography was first discussed by ’t Hooft .
The idea by ’t Hooft of the emergence of quantum degrees of freedom was based on information loss at the classical level. While it has to be shown how exactly quantum states emerge from classical states, the information loss was here provided by the noise. What is not known at this point is the origin and the physical meaning of that noise, which has to be included for the stochastic quantization procedure<sup>8</sup><sup>8</sup>8In it was argued that chaos may play a significant role. The authors mention the important example of classical Yang–Mills theories, which are chaotic dynamical systems . In their work the Langevin–equation plays an essential role, which provides a effective description on large scales. For another approach on holography and chaos see .. We believe that the noise is not an artifical quantity in the sense that it has no physical meaning. Rather, it seems to be connected with the coupling of matter and spacetime itself. We note here that we don’t believe that the Langevin–equation is a fundamental description, but an effective one.
In a sense, if the ideas presented here have something to do with reality, the existence of quantum degrees of freedom which we observe in the laboratory are a result of the existence of extra dimensions. Our arguments suggest strongly that quantum mechanical expectation values should be seen as an average over a stochastic process and are therefore not fundamental quantities itself. This stochastic process seems to be connected to spacetime in a subtle way. Obviously, there is a connection with the interpretation of quantum mechanics by Nelson , and that spacetime itself is the source for the stochastic noise needed in this work.
A lot of work remains to be done. Most importantly, we left open if the AdS/CFT–correspondence can really be understood as a process of stochastic quantization. It is very important to investigate this case, first because it can confirm (or not confirm) our ideas. Secondly, if it turns out that the stochastic quantization is at the heart of holography (in every spacetime), the AdS spacetime is a very good example where one can learn more about the thermodynamics/statistics of this spacetime. Furthermore, one may hope to find hints about the nature of the stochastic noise. The way pioniered by Periwal and Lifshytz should tell us more. Our approach should also be discussed within the framework of M(atrix) theory: whereas we worked in the spacetime picture, the ideas presented here should have a more fundamental interpretation.
Apart from the case of the AdS spacetime, one should consider the case of a black hole in the light of the ideas presented here. Of course, a covariant formulation of the ideas presented here would be desireable.
Another important problem to be solved is to find a mathematical expression for the condition that the dual quantum theory exists. We believe that this question is deeply related to the thermodynamics of spacetime and therefore to the generalized second law.
Acknowledgements:
We thank Stephon Alexander, Robert Brandenberger, Miquel Dorca, Damien Easson, Antal Jevicki, Jerome Martin, Matthias Soika and Shan-Wen Tsai for useful discussions and critism at several stages of this project and for making the paper readable. We are grateful to Helmuth Hüffel for pointing out some useful references and comments. This work was supported by DAAD/NATO.
## Appendix A Treatment of the Langevin equation
In this appendix we justify our steps in Section 5. The field $`\varphi (x)`$ in Euclidean space (with coordinates $`x`$) is now “lifted” onto a higher–dimensional manifold with an additional coordinate $`r`$ and metric $`g_{\mu \nu }`$, which will be here the AdS–spacetime. On this manifold we postulate a stochastic process $`\eta `$ and imagine, that the generalized scalar field $`\varphi (x,r)`$ is “propagating” along the coordinate $`r`$. We set $`\mathrm{}=c=G=k=1`$.
Usually the AdS is considered as a submanifold in a higher–dimensional covering space with symmetry group $`SO(2,D)`$. However, we make use of the (euclidean) form (see e.g. )
$$ds^2=R^2\left[\frac{4}{(1r^2)^2}\left(dr^2+r^2d\mathrm{\Omega }^2\right)+d\tau ^2\frac{1+r^2}{1r^2}\right]$$
(23)
where $`r<1`$ is the AdS spacetime and $`r=1`$ is the boundary.
Before we discuss the specific example of the AdS spacetime, let us make some general considerations. Consider a point $`(x,r)`$ on the higher–dimensional spacetime. The boundary of this spacetime is our original euclidean space. We postulate a stochastical differential equation which may viewed as a “generalized” Langevin–equation<sup>9</sup><sup>9</sup>9This is actually not the Langevin–equation as in the usual stochastic quantization procedure because the cofficients depend on the coordinate $`r`$. The equation has the form of what is called It$`\widehat{o}`$’s Langevin equation. In what follows, the eq. (24) only has to allow for the correct limit, i.e. we want to recover the euclidean measure in the path integral as $`r`$ approaches the boundary.. Because we are in a curved space we have to treat the problem locally, that is, we write generally
$$d\varphi =\mathrm{\Gamma }(x,r)\frac{S_E}{\varphi (x,r)}dr+dW(x,r),$$
(24)
where we write formally “$`dW(x,r)=\eta (x,r)dr`$”. The meaning of $`\mathrm{\Gamma }`$ will become clear if one notices that this quantity “absorbes” the change of a coordinate transformation $`r\stackrel{~}{r}`$. It describes also the strength of dissipation in the Langevin equation. Because $`\varphi (x,r)`$ is a scalar function, as well as the Euclidean action, $`dW(x,r)`$ has to be a scalar function, too. For the transformation for $`\mathrm{\Gamma }`$ we find therefore
$$\mathrm{\Gamma }^{}=\mathrm{\Gamma }\frac{dr}{d\stackrel{~}{r}}.$$
(25)
In what follows, we will consider $`\mathrm{\Gamma }`$ and $`dW`$ as a function of $`r`$ only. This should be the case in the AdS, for example, reflecting the symmetries of this space<sup>10</sup><sup>10</sup>10In fact, the symmetries of the AdS spacetime motivated us to choose the form (24) for the Langevin–equation and the form of the correlations below.. The correlation for the process $`dW`$ is assumed to be
$$<dW>=0\text{ }\text{ }\mathrm{and}\text{ }<dW(r)dW(r^{})>=2\alpha (r)\delta (rr^{})dr,$$
$`\alpha (r)`$ describes the strength of the fluctuations and is therefore related to the local temperature of the bath. The transformation of $`\alpha (r)`$ is:
$$\alpha ^{}=\alpha \frac{dr}{d\stackrel{~}{r}}.$$
(26)
$`\alpha `$, the diffusion coefficient, is therefore in general not constant along the coordinate $`r`$.<sup>11</sup><sup>11</sup>11We may introduce another stochastic variable which has a constant temperature along the radial coordinate. This variable transforms then as the variable $`\mathrm{\Gamma }`$. Finally we have to specify our boundary conditions for the Langevin–equation, which is
$$\underset{rr_{\mathrm{boundary}}}{lim}\varphi (x,r)=\varphi (x),$$
(27)
i.e. the field is the original field when we approach the boundary (which could also be the event horizon in the case of a black hole). Of course, it is clear that the generalized temperature should be a smooth function as we approach the horizon/boundary. Furthermore, we assumed that the notion of temperature locally makes sense, as usual in non–equilibrium thermodynamics. However, what is the Langevin equation for this problem? Are there conditions for $`\alpha (r)`$ and $`\mathrm{\Gamma }(r)`$? And can we find the solution we want?
In order to answer these questions, we will shall now derive the Fokker–Planck equation for this problem, using Ito’s stochastic calculus. Consider a functional $`F(\varphi (x,r))`$ of $`\varphi (x,r)`$. The Taylor series is, using the Langevin equation
$`dF(\varphi )`$ $`=`$ $`{\displaystyle \frac{\delta F}{\delta \varphi }}d\varphi +{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^2F}{\delta \varphi ^2}}d\varphi ^2+\mathrm{}`$ (28)
$`=`$ $`\mathrm{\Gamma }{\displaystyle \frac{\delta F}{\delta \varphi }}{\displaystyle \frac{\delta S}{\delta \varphi }}dr+{\displaystyle \frac{\delta F}{\delta \varphi }}dW+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^2F}{\delta \varphi ^2}}\left[\mathrm{\Gamma }^2\left({\displaystyle \frac{\delta S}{\delta \varphi }}\right)^2dr^2+\mathrm{}+dW^2\right]+\mathrm{}`$
If we neglect the higher order terms and take the average of this equation we find
$$\frac{d<F(\varphi )>}{dr}=\mathrm{\Gamma }\frac{\delta F}{\delta \varphi }\frac{\delta S}{\delta \varphi }+\alpha \frac{\delta ^2F}{\delta \varphi ^2},$$
(29)
where we used the conditions for the stochastic noise. Now we introduce a probability distribution $`𝒫(\varphi ,r)`$, defined by
$$<\mathrm{}>=𝒟\varphi 𝒫(\varphi ,r)\mathrm{},$$
(30)
where the dots represent a polynom in the field $`\varphi `$. It is
$$\frac{d<F(\varphi )>}{dr}=𝒟\varphi \frac{𝒫}{r}F(\varphi ).$$
(31)
Integrating (29) by parts we then find the Fokker–Planck equation
$$\frac{𝒫}{r}=\mathrm{\Gamma }(r)\frac{\delta }{\delta \varphi }\left(\frac{\delta S}{\delta \varphi }𝒫\right)+\alpha (r)\frac{\delta ^2}{\delta \varphi ^2}𝒫.$$
(32)
We will now investigate if we can find a solution of the form
$$𝒫(\varphi ,r)=A(r)e^{S_E(\varphi (x,r))}.$$
(33)
Inserting this into the Fokker–Planck equation gives
$$\frac{A^{}(r)}{A(r)}=\left(\mathrm{\Gamma }(r)\alpha (r)\right)\left(\frac{\delta ^2S}{\delta \varphi ^2}\left(\frac{\delta S}{\delta \varphi }\right)^2\right).$$
(34)
From this equation we find with find, using seperation of variables, for $`A(r)`$:
$$A(r)=𝒞\mathrm{exp}\left(_0^r𝑑\stackrel{~}{r}\left(\mathrm{\Gamma }(\stackrel{~}{r})\alpha (\stackrel{~}{r})\right)\right),$$
(35)
and for $`\varphi `$:
$$\frac{\delta ^2S}{\delta \varphi ^2}\left(\frac{\delta S}{\delta \varphi }\right)^2=.$$
(36)
The last equation is not consistent to solve because $`S_E`$ is a free function of $`\varphi `$. If we use $`\mathrm{\Gamma }(r)=\alpha (r)`$ in the Langevin–equation (24) from the very beginning of the calculation, we would find from equation (34) that $`A(r)=𝒞=`$ constant and no restriction to $`S_E`$ would apply. The condition $`\mathrm{\Gamma }(r)=\alpha (r)`$ is a result of the well known fluctuation–dissipation theorem which should hold at every point. We fix $`𝒞`$ by the requirement
$$𝒟\varphi 𝒫(\varphi )=1.$$
(37)
In conclusion, we can formulate the problem as a stochastic process with the Langevin equation is given by eq. (24) if we assume the fluctuation–dissipation theorem (here in its local form). Furthermore, the temperature along the fictitious coordinate is not constant. The probability distribution is given by
$$𝒫(r,\varphi )=\frac{e^{S_E(\varphi (x,r))}}{𝒟\varphi e^{S_E(x,r)}}.$$
(38)
Because $`\varphi (x,r)\varphi (x)`$ as we reach the boundary of the spacetime, the probability distribution approaches the euclidean path integral measure for the field living there<sup>12</sup><sup>12</sup>12We remember again that $`S_E`$ is the euclidean action for the boundary field $`\varphi (x)`$..
As a last step we have to find an expression for the function $`\mathrm{\Gamma }`$. It is instructive to discuss a simple example in zero dimensions. Although this example is simple, it will illuminate two important points: first we will find an expression for $`\mathrm{\Gamma }`$ and second we will argue that the metric of the higher dimensional space is unimportant.
Consider the field $`\varphi (r)`$ with the action
$$S_E=\frac{1}{2}m^2\varphi ^2.$$
(39)
We want to calculate the correlation function $`<\varphi ^2>`$. With eq. (29) it follows that
$$d<\varphi ^2>=2\mathrm{\Gamma }(r)m^2<\varphi ^2>dr+2\alpha (r)dr.$$
(40)
For the solution we make the ansatz ($`C`$ is constant):
$$<\varphi ^2>=Ce^{\beta (r)}+A(r).$$
(41)
Then we find the following conditions
$`2\mathrm{\Gamma }(r)m^2`$ $`=`$ $`\beta ^{}(r),`$ (42)
$`2\alpha (r)`$ $`=`$ $`A^{}(r)+\beta ^{}(r)A(r).`$ (43)
According to the fluctuation–dissipation theorem $`\mathrm{\Gamma }(r)=\alpha (r)`$. Then a solution is
$$<\varphi ^2>=C\mathrm{exp}\left(2m^2_0^r\mathrm{\Gamma }(\stackrel{~}{r})𝑑\stackrel{~}{r}\right)+\frac{1}{m^2}.$$
(44)
Because $`\mathrm{\Gamma }(r)`$ is a positive function we find for $`r\mathrm{}`$ (which may correspond to the boundary in a certain coordinate system)
$$<\varphi ^2>\frac{1}{m^2},$$
(45)
which can also be obtained using the path integral with the measure $`\mathrm{exp}(S_E)`$.
One might wonder that in the calculation above the explicit form of $`g_{\mu \nu }`$ was not used. In fact, the only ingredient was that the diffusion parameter of the fictitious bath depends only on the fictitious coordinate $`r`$, i.e. the radial coordinate of the AdS. $`\mathrm{\Gamma }(r)`$ must have a singular behaviour there in order that eq. (45) holds. Now we observe that the transformation of $`\mathrm{\Gamma }(r)`$ is the same as a usual tensor density, such as the determinant of the metric tensor. This suggests, that we could make the ansatz ($`p`$ is a positive constant)
$$\mathrm{\Gamma }(r)=p\sqrt{g},$$
(46)
because in the case of AdS the determinant of the metric tensor is a function of $`r`$ only: $`g=g(r)`$. It is singular at the boundary. Indeed, with this ansatz we find in the coordinates (23) (reducing to one dimension) the behaviour (45) for the correlation function. Furthermore, the result (45) can be obtained if one considers a Minkowski–space instead of an AdS space (see e.g. ) . This point is disussed further in Section 5. |
warning/0001/cond-mat0001275.html | ar5iv | text | # Density-Matrix Renormalization-Group Analysis of Quantum Critical Points: I. Quantum Spin Chains
## I Introduction
Quantum critical points are characterized by fluctuations over all length and time scales and by the appearance of power law scaling. In this paper we present a simple but powerful numerical method to access quantum critical points in one-dimensional systems. The method combines the density-matrix renormalization-group (DMRG) algorithm and finite-size scaling ideas. We illustrate the method by applying it to several well-understood quantum spin chains. In a second paper to follow we apply the method to new classes of supersymmetric spin chains which describe various disordered electron systems.
The development of the density-matrix renormalization-group (DMRG) algorithm by White represented an important improvement over previous numerical methods for the study of low dimensional lattice models. It has been applied to a wide variety of systems. The DMRG approach was first used to study the ground state properties and low-energy excitations of one-dimensional chains. It has been extensively applied to the study of various spin chains. Low-lying excited states of the spin-$`1`$ and spin-$`1/2`$ Heisenberg antiferromagnets have been calculated. Likewise, spin-$`1`$ chains with quadratic and biquadratic interactions, a spin-$`2`$ antiferromagnetic chain, spin-$`1/2`$ and spin-$`1`$ chains with dimerization and/or frustration (next-nearest-neighbor coupling), and frustrated spin-$`3/2`$ and spin-$`2`$ chains have all been studied. Edge excitations at the ends of finite spin chains and the effects of perturbations such as a weak magnetic field coupled to a few sites have been considered. Randomness in the form of random transverse magnetic field in a spin-$`1/2`$ XY model, random exchange couplings, and random modulation patterns of the exchange, has been examined. Finally, alternating spin magnitudes, the presence of a constant or a staggered magnetic field in a spin-$`1`$ chain, bond doping, the effects of a local impurity, and interactions with quantum phonons have also been considered.
Most of the above work involves systems in which the first excited state is separated from the ground state by a non-zero energy gap as the DMRG works best for gapped systems. First attempts to extract critical behavior of gapless systems used the DMRG to generate renormalization transformations of the coupling constants in the Hamiltonian. Hallberg et al. studied the critical behavior of $`S=1/2`$ and $`S=3/2`$ quantum spin chains with periodic boundary conditions through extensive calculations of ground state correlation functions at different separations and different chain sizes $`L`$. Spin correlation functions in an open chain have also been calculated and compared with results calculated from low-energy field theory, showing that estimates of the amplitudes can also be obtained. The approach described in this paper was applied to the spin-$`1/2`$ Heisenberg chain and a non-Hermitian supersymmetric (SUSY) spin chain. More recently, critical behavior of classical one-dimensional reaction-diffusion models and the two-dimensional Potts model has been studied using the finite-size DMRG algorithm. Bulk and surface exponents of the Potts and Ising model have been obtained by using the DMRG to calculate correlation functions at different separations and collapsing curves obtained at different system sizes. The SUSY chain describing the spin quantum Hall effect (SQHE) plateau transition was also examined in some detail. Critical exponents were extracted and compared to exact predictions. Thermodynamic properties of other two-dimensional classical critical systems have also been studied by the DMRG method. Finally, Andersson et al. investigated the convergence of the DMRG in the thermodynamic limit for a gapless system of non-interacting fermions.
The method described in this paper combines the DMRG algorithm with finite-size scaling analysis, and yields accurate critical exponents. The main advantage of the method is its simplicity. Only the calculation of ground state correlations near the middle of chains with open boundary conditions are required. The relatively simple “infinite-size” DMRG algorithm is particularly accurate for this job. In Sec. II we describe the method. The tight-binding model can be solved exactly and in Sec. III we use it to illustrate our scaling analysis. DMRG results are presented in Sec. IV for the anisotropic $`S=1/2`$ Heisenberg antiferromagnet and several critical exponents are obtained. An analytical calculation shows that multiplicative logarithmic corrections – which complicate the extraction of accurate critical exponents – may be avoided in some instances. In Sec. V, the $`S=1`$ antiferromagnetic spin chain is studied, focusing on the critical point that separates the Haldane and the dimerized phases. We conclude with a summary in Sec. VI.
## II The DMRG / Finite-Size Scaling Approach
We first describe how critical exponents may be obtained from a finite-size scaling analysis of chains with open or fixed boundary conditions. These boundary conditions are the simplest to implement in DMRG calculations. In the next subsection the DMRG algorithm itself is briefly described.
### A Finite-Size Scaling
To illustrate the sorts of power-law scaling we wish to examine, first consider the case of a spin chain with periodic boundary conditions that is at its critical point. The system can be moved away from criticality by turning on a uniform magnetic field, say in the x-direction, at each site:
$`H_B=h{\displaystyle \underset{j=1}{\overset{L}{}}}S_j^x.`$ (1)
This perturbation makes the correlation length finite:
$`\xi _B|h|^{\nu _B}.`$ (2)
Explicit dimerization, breaking the symmetry of translation by one site, also moves the system away from criticality. For a Heisenberg antiferromagnet, this can be realized by the addition of a staggering term $`R`$ to the Hamiltonian:
$`H={\displaystyle \underset{j=1}{\overset{L1}{}}}[1+(1)^jR]\stackrel{}{S}_j\stackrel{}{S}_{j+1}.`$ (3)
The correlation length $`\xi `$ in this case scales as
$`\xi |R|^\nu .`$ (4)
Thus there are two independent exponents which correspond to these two perturbations of critical spin chains. Two-parameter scaling functions can be written for various observables and, for a finite system, these involve two dimensionless variables: the ratios $`L/\xi `$ and $`L/\xi _B`$. The induced dimerization, defined for now as the modulation of the $`xx`$ and $`yy`$ spin-spin correlations on even versus odd links,
$$\mathrm{\Delta }=(1)^j\left[S_j^xS_{j+1}^x+S_j^yS_{j+1}^yS_{j1}^xS_j^x+S_{j1}^yS_j^y\right],$$
(5)
is of course independent of the site index for periodic chains, and scales as a function of the chain length $`L`$, the field $`h`$, and the dimerization parameter $`R`$ as:
$`\mathrm{\Delta }(L,R,h)=\mathrm{sgn}(R)|R|^{\alpha _\mathrm{\Delta }}f_\mathrm{\Delta }(L|R|^\nu ,L|h|^{\nu _B}).`$ (6)
When the applied magnetic field is removed, $`h=0`$, and this expression simplifies to:
$`\mathrm{\Delta }(L,R)`$ $`=`$ $`\mathrm{sgn}(R)|R|^{\alpha _\mathrm{\Delta }}g_\mathrm{\Delta }(L|R|^\nu )`$ (7)
$``$ $`L^{x_\mathrm{\Delta }}R\mathrm{as}R0,`$ (8)
where the second line follows from the fact that when the perturbation $`R`$ is very small, or equivalently when the correlation length is larger than the system size, the net induced dimerization must be an analytic, linear, function of $`R`$. Therefore, for $`|x|1`$, the scaling function $`g_\mathrm{\Delta }(x)`$ is given by:
$`g_\mathrm{\Delta }(x)=|x|^{\alpha _\mathrm{\Delta }/\nu }(a_1|x|^{1/\nu }+a_2|x|^{2/\nu }+\mathrm{});`$ (9)
the first term yields linear dependence of $`\mathrm{\Delta }`$ in $`R`$ in the $`R0`$ limit, in agreement with Eq. 8, and the subsequent terms are higher order corrections. To recover the correct $`L`$-dependence, we must set
$$x_\mathrm{\Delta }=\frac{1\alpha _\mathrm{\Delta }}{\nu }.$$
(10)
The exponent $`x_\mathrm{\Delta }`$ and the correlation length exponent $`\nu `$ satisfy the usual relation
$$\nu =\frac{1}{2x_\mathrm{\Delta }}.$$
(11)
The applied magnetic field also polarizes the spins along the chain. The scaling form for the spin moment at each site is given by:
$`S^x=\mathrm{sgn}(h)|h|^{\alpha _B}f_B(L|R|^\nu ,L|h|^{\nu _B}).`$ (12)
With no applied dimerization, $`R=0`$, and we expect the simple power-law:
$`S^xL^{x_B}h\mathrm{as}h0.`$ (13)
Therefore, $`x_B=(1\alpha _B)/\nu _B`$.
Alternatively, dimerization can be induced by open boundary conditions, and we take advantage of this fact to extract critical exponents. As depicted in Fig. 1, open boundary conditions favor enhanced nearest-neighbor spin-spin correlations on the two outermost links. Chains of increasing length $`L=4,6,8,\mathrm{}`$ exhibit alternating patterns of dimerization on the interior bonds. Likewise, spin moments may be induced in the interior of the chain by applying a magnetic field to the ends of the chain. Strong applied edge magnetic fields completely polarize the end spins and induce non-zero and alternating spin moments along the chain. Alternatively, spin moments can be induced as before by a staggered magnetic field applied along the entire chain. Here however we consider only edge magnetic fields.
We monitor the induced dimerization and spin moments at the center of the chain as the chain length $`L`$ is enlarged via the DMRG algorithm. This scaling analysis is convenient because the relatively simple infinite-size DMRG algorithm applies to open chains and is most accurate at the center region of the chain where we focus our attention. The induced dimerization and spin moments in the interior of the chain show power-law scaling at the critical point. Igloi and Rieger demonstrated power-law scaling for a variety of open boundary conditions (free, fixed and mixed). At the critical point $`R=0`$ and $`h=0`$ the induced dimerization scales as a power-law with possible multiplicative logarithmic corrections:
$$\mathrm{\Delta }(L/2)=L^{x_\mathrm{\Delta }}(\mathrm{ln}L)^{y_\mathrm{\Delta }}\left(a+\frac{b}{L}+\mathrm{}\right).$$
(14)
A similar expression holds for the induced spin moment at the center of the chain, $`S^x(L/2)`$, with the replacement of the exponents $`x_\mathrm{\Delta }x_B`$ and $`y_\mathrm{\Delta }y_B`$.
### B Infinite-Size DMRG Algorithm
The name “density-matrix renormalization-group” is something of a misnomer as the method is most accurate away from critical points, when there is an energy gap for excitations. It is helpful to think of the DMRG algorithm as a systematic variational approximation for the calculation of the ground state and/or low-lying excitations, principally in one dimension. The Hilbert space of a quantum chain generally grows exponentially with the chain length, and eventually must exceed available computer memory. The DMRG algorithm is an efficient way to truncate the Hilbert space; as the size of the space retained can be varied (up to machine limits) it is possible to ascertain the size of errors introduced by the truncation.
For simplicity, we use the so-called “infinite-size” DMRG algorithm. As the algorithm has been described in some detail by White, we just sketch the essentials of the method. It begins with the (numerically exact) diagonalization of an open chain consisting of just four sites, each site having on-site Hilbert space of dimension $`D`$. For quantum spin chains $`D=2S+1`$, thus $`D=2`$ for the spin-1/2 Heisenberg antiferromagnet. The chain is then cut through the middle into two pieces, one half of which is interpreted as the “system” and the other half as the “environment,” the two parts combined being thought of as the entire “universe” of the problem, see Fig. 2. At this point the reduced density matrix for the system, of size $`DM\times DM`$ is constructed by performing a partial trace over the environment half of the chain. It is defined by:
$$\rho _{ij}=\underset{i^{}=1}{\overset{DM}{}}\mathrm{\Psi }_{ii^{}}\mathrm{\Psi }_{ji^{}},$$
(15)
where $`\mathrm{\Psi }_{ii^{}}=ii^{}|\mathrm{\Psi }`$ are the real-valued matrix elements of the eigenstate of interest (the “target” which is often the ground state) projected onto a basis of states labeled by unprimed Roman index $`i`$ which covers the system half of the chain and primed index $`i^{}`$ which covers the environment half of the chain. The eigenvalues of the reduced density matrix are real, positive, and sum up to one; these are interpreted as probabilities. We keep only the $`M`$ most probable eigenstates corresponding to the largest eigenvalues, and discard the remaining $`M(D1)`$ eigenstates. The retained states form a new basis for the problem. Next, two new sites are added to the middle of the chain and the pieces are connected, yielding a chain of size $`L=6`$. The process is then repeated by finding the targeted state of this chain, constructing the new reduced density matrix and again projecting onto the $`M`$ most probable states. As the chain length grows in steps of two, the total Hilbert space dimension grows by a multiplicative factor of $`D^2`$. None of the Hilbert space is thrown away until the chain grows large enough that its Hilbert space exceeds the space that is held in reserve, in other words until $`D^L>D^2M^2`$. The truncation process damages the outer regions of the chain the most, and the central region is treated most accurately.
One advantage of the method presented in this paper is that critical exponents are extracted from ground-state correlations only. Excited states are not needed for these exponents and there is no need to calculate the excitation gap. Furthermore, the finite size analysis described in the previous subsection takes advantage of the fact that the DMRG algorithm works best with open chains and treats the central region of the chain most accurately. The use of the more complicated finite-size algorithm might yield even more accurate results. However, we show below that we can calculate critical exponents to an accuracy of a few percent or better with the infinite-size algorithm.
## III Tight Binding Model at Half-Filling
As a simple first illustration of our finite-size scaling method we study the ordinary tight binding model of spinless fermions hopping from site-to-site along a chain at half-filling. Obviously, the DMRG algorithm is not needed in this case as we can solve the quadratic problem exactly via a Fourier transform. Due to particle-hole symmetry, at half-filling the chemical potential is zero. The correlation length exponent for this system is $`\nu =1`$. A direct way to see this is by introducing the staggering parameter $`R`$ to modulate the amplitude of the hopping matrix elements on even versus odd links:
$`H`$ $`=`$ $`t{\displaystyle \underset{j=0}{\overset{L1}{}}}[1+(1)^jR](c_j^{}c_{j+1}+h.c.).`$ (16)
To diagonalize the Hamiltonian, in the case of periodic boundary conditions $`c_0=c_L`$, we introduce separate fermion operators for even and odd sites as follows:
$`c_{2j}=d_{2j}`$ (17)
$`c_{2j1}=e_{2j}`$ (18)
After the Fourier transformation to momentum-space, the Hamiltonian can be written as:
$`H=t{\displaystyle \underset{k}{}}\left\{[(1R)+e^{2ik}(1+R)]d_k^{}e_k+[(1R)+e^{2ik}(1+R)]e_k^{}d_k\right\},`$ (19)
where the lattice spacing $`a=1`$. For each $`k`$, diagonalization of the $`2\times 2`$ matrix yields the dispersion relation:
$`ϵ_k=\pm 2t\sqrt{1(1R^2)\mathrm{sin}^2(k)}.`$ (20)
At half-filling the ground state has all states with $`ϵ_k<0`$ occupied. The left and right Fermi points are, respectively, $`k_F=\pm \pi /2`$. Hence the gap $`m=2t|R|`$. As the correlation length $`\xi m^1|R|^1`$ we obtain $`\nu =1`$. Since $`\nu ^1=2x_\mathrm{\Delta }=1`$, the dimerization exponent $`x_\mathrm{\Delta }=1`$.
We now reproduce this result using the finite-size scaling method applied to open chains. We consider a finite chain of length $`L`$ with open boundary conditions and calculate the induced dimerization $`\mathrm{\Delta }(j)=(1)^jc_j^{}c_{j+1}c_{j+1}^{}c_{j+2}`$ around the chain center $`j=L/2`$, and extract its leading dependence on $`L`$. Open boundary conditions are imposed by using the Fourier transform
$`c_j`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2(L+1)}}}{\displaystyle \underset{m=1}{\overset{L}{}}}c_{k_m}(e^{ik_mj}e^{ik_mj}),`$ (21)
$`k_m`$ $`=`$ $`{\displaystyle \frac{\pi }{L+1}}m,m=1,2,\mathrm{},L`$ (22)
as this enforces $`c_0=c_{L+1}=0`$. Filling all of the negative energy states at half-filling, the expectation value of the dimerization at $`L/2`$ can be found by straightforward calculation:
$$\mathrm{\Delta }(L/2)\frac{1}{L+1}\underset{m=\frac{L}{2}+1}{\overset{L}{}}\left[\mathrm{cos}[k_m(L+3)]\mathrm{cos}[k_m(L+1)]\right].$$
(23)
This sum can be evaluated numerically with the result that $`x_\mathrm{\Delta }1`$ as $`L\mathrm{}`$ as shown in Fig. 3, in agreement with the explicit calculation for the periodic chain. It is also easy to show that open chains with an odd number of sites have vanishing induced dimerization at the center of the chain, as expected by the symmetry of reflection about the central site.
The induced density moment can likewise be obtained either directly by studying the effects of a staggered chemical potential $`\mu _{stag}`$ (which doubles the size of the unit cell from one to two sites and thus generates a gap $`m=2|\mu _{stag}|`$) or by the inclusion a local chemical potential $`\mu `$ at the two ends of the chain:
$`HH\mu (c_0^{}c_0+c_{L1}^{}c_{L1}).`$ (24)
Again the system consists of $`L`$ sites, the site index running from $`0`$ to $`L1`$, and there are open boundary condition at $`j=0`$ and $`j=L1`$. For large $`\mu 0`$, the boundary condition is equivalent to enforcing unit occupancy at the chain ends, $`n_0=n_{L1}=1`$. This boundary condition is satisfied by the Fourier transform
$`c_j`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2(L1)}}}{\displaystyle \underset{m=0}{\overset{L1}{}}}c_{k_m}(e^{ik_mj}+e^{ik_mj})`$ (25)
with
$`k_m`$ $`=`$ $`{\displaystyle \frac{\pi }{L1}}m,m=0,1,\mathrm{},L1.`$ (26)
Again it is a simple exercise to calculate the occupancies. At the chain ends we obtain: $`c_0^{}c_0=c_{L1}^{}c_{L1}=1`$ in agreement with the boundary condition. At the center of the chain the occupancy can be evaluated analytically,
$`c^{}(L/2)c(L/2)={\displaystyle \frac{1}{L1}}{\displaystyle \underset{m=\frac{L}{2}}{\overset{L1}{}}}\left[1+\mathrm{cos}(k_mL)\right].`$ (27)
It scales as $`c^{}(L/2)c(L/2)1/2L^1`$. Hence $`\nu _B=x_B=1`$ in agreement with the direct calculation of these exponents.
## IV Spin-$`1/2`$ Antiferromagnet
We next turn to the study of a richer system: spin-$`1/2`$ antiferromagnetic chains. We begin with the XY model, which can be solved exactly by a Jordan-Wigner mapping to the tight binding model. We then study the anisotropic XXZ model. The isotropic Heisenberg model is treated separately as there are complicating multiplicative logarithmic corrections to scaling at the isotropic point.
### A XY model
The Hamiltonian for the spin-$`1/2`$ XY model,
$$H=J\underset{j=0}{\overset{L2}{}}\left[S_j^xS_{j+1}^x+S_j^yS_{j+1}^y\right],$$
(28)
can be written in terms of spinless fermion creation and annihilation operators $`c_j^{}`$ and $`c_j`$ via the Jordan-Wigner transformation. An up spin in the z-direction at site $`i`$ then corresponds to having the site occupied by a fermion, while spin down corresponds to an empty site. The Hamiltonian of Eq. 28 is mapped to a nearest-neighbor tight binding Hamiltonian with $`t=J`$. Based on our analysis in the previous section we can conclude that $`\nu =1`$ for the XY model.
Fig. 4 presents our DMRG results for the induced dimerization and induced spin moments, in the x- and in the z-directions, at the center of the chain as a function of the chain length, $`L`$.
The exponents are obtained from the slopes of the curves shown in Fig. 4. The induced dimerization exponent for $`\mathrm{\Delta }(L/2)`$ is close to $`1`$ ($`x_\mathrm{\Delta }=0.99\pm 0.01`$) as expected from the relation $`\nu =1/(2x_\mathrm{\Delta })`$. The slope of the log-log plot of the induced spin moment in the z-direction is also close to $`1`$ ($`x_B=1.01\pm 0.02`$). This result is also expected since it is equivalent to the exponent for the induced density moment in the tight binding model as discussed in the previous section. In the case of the induced spin moment in the x-direction, the exponent is $`0.248\pm 0.003`$. This value compares well with the exact number of $`1/4`$ as derived in the next section.
### B XXZ model
Next consider the nearest-neighbor, spin-1/2 XXZ Heisenberg antiferromagnet:
$$H=J\underset{j=0}{\overset{L2}{}}\left[S_j^xS_{j+1}^x+S_j^yS_{j+1}^y+\gamma S_j^zS_{j+1}^z\right].$$
(29)
Anisotropy in the coupling between the z-components of the spins may be varied by changing $`\gamma `$. Performing the Jordan-Wigner transformation, the XY terms again yield the tight binding Hamiltonian. Low-energy excitations therefore occur near the two Fermi points at $`k=\pm \pi /2a`$. We may treat the non-Gaussian $`\gamma `$ term as a perturbation and focus on excitations around these Fermi points by defining left and right moving low-energy quasiparticles. Taking the continuum limit and keeping only the low-energy modes, the tight binding term is then effectively described by the massless fermions. The $`S_j^zS_{j+1}^z`$ term is quartic in the fermion operators. Integrating out the high-energy modes, it will renormalize the fermion velocity and also contain interaction terms.
We then implement Abelian bosonization, with UV cutoff $`\alpha `$. The effective Hamiltonian is a sine-Gordon model (a derivation can be found in Ref. 49):
$$H=H_0\frac{y_\varphi }{2\pi \alpha ^2}𝑑x\mathrm{cos}[\sqrt{8\pi }\varphi (x)]$$
(30)
where
$$H_0=u𝑑x\left[K\mathrm{\Pi }^2+\frac{(_x\varphi )^2}{K}\right].$$
(31)
Here $`u=2Ja=2a`$ is the bare Fermi velocity and the constant $`K1+y_0/2`$ depends on the anisotropy $`\gamma `$. The XY limit correspond to $`y_\varphi =0`$.
The long distance behavior of the staggered part of $`S^z`$ and $`S^{}`$ are given in terms of the boson fields as:
$`S^z(x)`$ $``$ $`(1)^{x/\alpha }\mathrm{cos}[\varphi (x)/R]`$ (32)
$`S^{}(x)`$ $``$ $`(1)^{x/\alpha }e^{i2\pi R\stackrel{~}{\varphi }(x)}`$ (33)
where the radius $`R`$ is given by
$`R=\sqrt{{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{\mathrm{cos}^1\gamma }{2\pi ^2}}}.`$ (34)
First consider the anisotropic case $`\gamma 1`$. The isotropic case has logarithmic corrections to scaling that are dealt with in the next section. For $`\gamma >1`$ the interaction term is relevant and the system is gapped, and in the Ising universality class. Indeed, in the limit $`\gamma \mathrm{}`$ it is the Ising model. For $`\gamma <1`$ the interaction term is irrelevant, the system is gapless and $`\mathrm{\Delta }(L/2)`$ and $`S^x(L/2)`$ should exhibit power law decay, with no log corrections as there are no marginal operators. The log-log plots of Fig. 5 (a) show the induced spin moment in the x-direction at the chain center for different values of the anisotropy $`\gamma `$. The edge magnetic field in the x-direction is fixed, $`h=1.0`$. As expected, for $`\gamma >1`$ there is exponential decay and in the cases $`\gamma <1`$ the exponents $`x_B(\gamma )`$ are found by fitting the curves in Fig. 5 (a) to the form of Eq. 14. The exponents $`y_B`$ are set equal to zero, the higher order corrections are included and give very small deviations from a simple linear fit. In Fig. 5 (b) the exponents $`x_B(\gamma )`$ are compared to the exact value $`x_B(\gamma )=\pi R^2(\gamma )`$ obtained by Affleck. Agreement is found at the percent level. Affleck derived the exponent as follows. The edge magnetic field in the x-direction applied at $`j=0`$ corresponds to a term
$$H_B=hS^x(0)=\mathrm{constant}\times h\mathrm{cos}[\sqrt{2\pi }\stackrel{~}{\varphi }(0)]$$
(35)
in the Hamiltonian. For sufficiently large $`h`$ the energy is minimized by setting
$$\stackrel{~}{\varphi }(0)=0\varphi _R(0)=\varphi _L(0).$$
(36)
Regarding $`\varphi _R`$ as an analytic continuation of $`\varphi _L`$, we may identify
$$\varphi _R(x)=\varphi _L(x).$$
(37)
Using this boundary condition, the induced spin moment is given by
$`S^x(j)(1)^{j/\alpha }e^{i2\pi R\varphi _L(j)}e^{i2\pi R\varphi _L(j)}{\displaystyle \frac{(1)^{j/\alpha }}{(2j)^{\pi R^2(\gamma )}}}.`$ (38)
For the XY model ($`\gamma =0`$), the induced spin moment in the x-direction therefore decays with exponent $`\pi R^2(0)=1/4`$.
A log-log plot of the induced dimerization at the center of the chain for various values of the anisotropy $`\gamma `$ is shown in Fig. 6 (a). The free boundary condition at the chain ends corresponds to setting:
$`\stackrel{}{S}_0=\stackrel{}{S}_{L+1}=0.`$ (39)
This condition translates to $`\varphi _R(x)=\varphi _L(x)+\pi R`$ in terms of the boson fields which yields:
$`\mathrm{\Delta }(j)(1)^{j/\alpha }\mathrm{cos}[\varphi (j)/R]{\displaystyle \frac{(1)^{j/\alpha }}{(2j)^{1/4\pi R^2(\gamma )}}}.`$ (40)
In Fig. 6 (b) the exponents obtained from the slopes of the curves in Fig. 6 (a) are plotted against the exact values $`x_\mathrm{\Delta }=1/4\pi R^2(\gamma )`$. Again agreement is found at the percent level.
Another quantity of interest is the sum, instead of the difference, of the spin-spin correlation function on adjacent bonds near the center of the chain:
$`ϵ(L/2){\displaystyle \frac{1}{2}}\left(\stackrel{}{S}_{L/2}\stackrel{}{S}_{L/2+1}+\stackrel{}{S}_{L/21}\stackrel{}{S}_{L/2}\right),`$ (41)
which at $`\gamma =1`$ equals the energy density per bond and therefore does not vanish in the thermodynamic limit. Fig. 7 is a plot of $`ϵ(L/2)`$ as a function of the system size $`L`$ at the isotropic point $`\gamma =1`$.
As expected, this quantity approaches a constant value $`ϵ(\mathrm{})`$ in the thermodynamic limit. After subtracting the extrapolated value at $`L\mathrm{}`$, $`ϵ(L/2)`$ too exhibits power law decay of the form of Eq. 14. The constant $`ϵ(\mathrm{})`$ can be found by an iteration process. Starting with an initial value for $`ϵ(\mathrm{})`$ obtained from a rough extrapolation of the curve in Fig. 7, we fit the subtracted value $`ϵ(L/2)ϵ(\mathrm{})`$ to a power-law form. The extrapolated value $`ϵ(\mathrm{})`$ is then adjusted slightly until an optimal fit to a pure power law is attained. The extrapolated value found this way is $`ϵ(\mathrm{})=0.443148`$ and Fig. 8 shows the power law behavior of the subtracted quantity.
We obtain an exponent of $`2.1\pm 0.1`$ in the scaling of $`ϵ(L/2)ϵ(\mathrm{})`$. This is as expected from the linear dispersion relation of Heisenberg antiferromagnets: in a Lorentz-invariant theory the energy density operator has dimension 2.
The DMRG result for the energy per bond is extremely accurate and can be compared with the exact value obtained from the Bethe ansatz solution of $`ϵ=1/4\mathrm{ln}2=0.44314718`$. It is crucial to note that the open boundary conditions induce staggering in the strength of the bonds along the chain. To eliminate this effect, the energy per bond must be calculated as the average of the bond energy from two consecutive bonds at the center of the chain. Suggestions that infinite-size DMRG results for the center region of the chain are not very accurate appear to have failed to take this effect into account. We have also checked our results at different anisotropies. For the XY case ($`\gamma =0.0`$), we obtain $`ϵ(\mathrm{})=0.318310`$ extrapolating from chains up to $`L=200`$ and $`M=128`$ and the exact result is $`1/\pi =0.3183099`$.
### C Logarithmic Corrections to Scaling
In the isotropic XXX limit, the interaction $`\mathrm{cos}[\sqrt{8\pi }\varphi (x)]`$ in the low-energy effective Hamiltonian Eq. 30 becomes marginal and can generate multiplicative logarithmic corrections to scaling. In this section we calculate its effect on the scaling of the induced spin moment $`S^x(L/2)`$ when an edge magnetic field $`H_B`$ in the x-direction is applied. Cancellations occur and in this case there are no multiplicative $`\mathrm{ln}(L)`$ corrections. As a practical matter, the cancellation of the logarithmic corrections means that numerical calculations of the exponent $`x_B`$ are particularly precise. We note that finite-size scaling of the spin-spin correlation function has been previously calculated for a spin-$`1/2`$ chain with periodic boundary conditions.
The coupling constants in the sine-Gordon Hamiltonian (Eq. 30) renormalize under a change of the ultraviolet cutoff $`\alpha \alpha e^l`$ according to the renormalization group equations:
$`{\displaystyle \frac{dy_0}{dl}}`$ $`=`$ $`y_\varphi ^2(l),`$ (42)
$`{\displaystyle \frac{dy_\varphi }{dl}}`$ $`=`$ $`y_\varphi (l)y_0(l).`$ (43)
As noted in the previous section, a large edge magnetic field applied at $`x=0`$ in the x-direction enforces the boundary condition $`\varphi _R(x)=\varphi _L(x)`$ (Eq. 37). Thus
$$S^x(x)(1)^{x/a}\mathrm{cos}[\sqrt{2\pi }\stackrel{~}{\varphi }(x)]e^{i\sqrt{2\pi }\varphi _L(x)}e^{i\sqrt{2\pi }\varphi _L(x)}.$$
(44)
For the free theory, which corresponds to the XY model $`y_\varphi =0`$, the induced spin moment is simply
$$S^x(x)_0\mathrm{exp}\left[KU_L(2x)\right]$$
(45)
where
$$U_L(x)=\frac{1}{2}\mathrm{ln}(\frac{\alpha +ix}{\alpha }).$$
(46)
But in the general XXZ case we ascertain the effect of the marginal operator by following a procedure similar to one developed by Giamarchi and Schulz who calculated correlation functions for finite periodic chains. We first define the function:
$`F(x)e^{KU_L(2x)}S^x(x).`$ (47)
At the XY point $`y_\varphi =0`$ clearly $`F(x)=1`$. For small $`x`$, an expansion of $`F`$ in powers of $`y_\varphi `$ converges, and for sufficiently small coupling $`y_\varphi `$, $`F(x)1`$. Upon rescaling, the function $`F(x)`$ also depends on the new length scale and on the rescaled coupling constants $`y_0(l)`$ and $`y_\varphi (l)`$. By an argument similar to the one employed by Kosterlitz, the effect of rescaling $`\alpha e^l\alpha `$ is:
$`F(x,\alpha e^l,y(l))=I(dl,y(l))F(x,\alpha e^{l+dl},y(l+dl)),`$ (48)
where $`y(l)`$ denotes all the couplings as function of the scaling parameter $`l`$. The rescaled short distance cutoff is then $`\alpha (l)=e^l\alpha `$, where $`\alpha `$ is the initial cutoff. Rescaling can be repeated until $`\alpha (l)x`$, at which point we have:
$`F(x,x,y(\mathrm{ln}(x/\alpha )))=O(1).`$ (49)
The contributions to the function $`F`$ from repeated rescalings, until $`\alpha (l)`$ reaches $`x`$, can be written explicitly as:
$`F(x,\alpha ,y(\alpha ))={\displaystyle \underset{l=0}{\overset{l=\mathrm{ln}(x/\alpha )}{}}}I(dl,y(l))=\mathrm{exp}\left\{{\displaystyle _0^{\mathrm{ln}(x/\alpha }}\mathrm{ln}\left[I(dl,y(l))\right]𝑑l\right\}.`$ (50)
We proceed to calculate the function $`I`$. First we expand $`S^x(x)`$ in powers of $`y_\varphi `$, writing it in terms of averages with respect to the free Hamiltonian,
$`S^x(x)`$ $``$ $`e^{KU_L(2x)}+{\displaystyle \frac{y_\varphi }{2\pi \alpha ^2}}{\displaystyle d^2x^{}S^x(x)\mathrm{cos}[\sqrt{8\pi }\varphi (x^{})]_0}+`$ (52)
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{y_\varphi }{2\pi \alpha ^2}}\right)^2{\displaystyle d^2x_1d^2x_2S^x(x)\mathrm{cos}[\sqrt{8\pi }\varphi (x_1)]\mathrm{cos}[\sqrt{8\pi }\varphi (x_2)]_0}+\mathrm{}`$
The averages are given by
$`S^x(x)\mathrm{cos}[\sqrt{8\pi }\varphi (x^{})]_0=0`$ (53)
and
$`S^x(x)\mathrm{cos}[\sqrt{8\pi }\varphi (x_1)]`$ $`\mathrm{cos}[\sqrt{8\pi }\varphi (x_2)]_0`$ (55)
$`{\displaystyle \frac{1}{2}}\mathrm{exp}\left[KU_L(2x)+4KU_L(2x_1)+4KU_L(2x_2)4KU(x_1+x_2)4KU(x_1x_2)\right]`$
The $`O(y_\varphi ^2)`$ term can be simplified by assuming that the main contribution comes from configurations where $`x_1`$ and $`x_2`$ are very close to each other. Introducing new integration variables
$`r`$ $``$ $`x_1x_2`$ (56)
$`R`$ $``$ $`{\displaystyle \frac{x_1+x_2}{2}}`$ (57)
and expanding $`U_L`$ in powers of $`r`$, which is assumed to be small,
$`U_L(2x_1)=U_L(2R+r)=U_L(2R)+r_RU_L(2R)+\mathrm{},`$ (58)
$`U_L(2x_2)=U_L(2Rr)=U_L(2R)r_RU_L(2R)+\mathrm{},`$ (59)
we obtain the average
$`S^x(x)\mathrm{cos}[\sqrt{8\pi }\varphi (x_1)]\mathrm{cos}[\sqrt{8\pi }\varphi (x_2)]_0{\displaystyle \frac{1}{2}}\mathrm{exp}\left[KU_L(2x)4KU(r)\right].`$ (60)
The dependence on $`R`$ cancels out. The expansion Eq. 52 becomes
$`S^x(x)`$ $``$ $`e^{KU_L(2x)}\left[1+{\displaystyle \frac{y_\varphi ^2\mathrm{\Omega }}{4\alpha ^2}}{\displaystyle _\alpha }𝑑re^{4KU(r)}\right],`$ (61)
where $`\mathrm{\Omega }𝑑R`$ is a measure of the linear size of the system. Next consider the effect of rescaling $`\alpha ^{}=\alpha e^{dl}`$, where $`dl`$ is infinitesimal. Using
$`{\displaystyle _\alpha ^{\mathrm{}}}𝑑x={\displaystyle _\alpha ^\alpha ^{}}𝑑x+{\displaystyle _\alpha ^{}^{\mathrm{}}}𝑑x,`$ (62)
we obtain:
$$S^x(x)e^{KU_L(2x)}\left[1+\frac{y_\varphi ^2}{4\alpha ^2}dl+\frac{y_\varphi ^2}{4\alpha ^2}_\alpha ^{}𝑑re^{4KU(r)}\right]$$
(63)
Matching this result with Eq. 48, we find:
$$I(dl,y_0(l),y_\varphi (l))\mathrm{exp}\left[\frac{y_\varphi ^2(l)}{4\alpha ^2}dl\right],$$
(64)
hence from Eq. 50 and Eq. 47, we have
$`S^x(x)\mathrm{exp}\left\{KU_L(2x)+{\displaystyle _0^{\mathrm{ln}(x/\alpha )}}{\displaystyle \frac{y_\varphi ^2(l)}{4\alpha ^2}}𝑑l\right\}.`$ (65)
Using the RG equations (Eq. 43), the solution at large $`l`$ is $`y_\varphi (l)1/l`$ and
$`S^x(x)\left({\displaystyle \frac{x}{\alpha }}\right)^{1/2}\mathrm{exp}\left\{{\displaystyle _0^{\mathrm{ln}(x/\alpha )}}𝑑l\left[O({\displaystyle \frac{1}{l^2}})\right]\right\}.`$ (66)
There are no multiplicative $`\mathrm{ln}(x)`$ corrections, as these would require terms of order $`O(1/l)`$ in the integrand inside the exponential in Eq. 66. In our calculation, $`O(1/l)`$ terms do not appear, only $`O(1/l^2)`$ and higher-order terms. As a check, we can repeat the same procedure for $`S^z(x)`$, with the edge field now oriented in the z-direction. Of course this should give the same result since the system is isotropic, but as the Jordan-Wigner transformation picks the z-direction as the spin quantization axis, the equivalence is not obvious, and the check is non-trivial. Again, explicit calculation shows that $`O(1/l)`$ terms do not arise. This result is in reasonable agreement with our numerical results. Fitting the DMRG data (see Fig. 9) to the form Eq. 14, we obtain $`x_B=0.485\pm 0.01`$ with a small non-zero value for the log exponent $`y_B=0.06\pm 0.01`$. By contrast, in the case of the induced dimerization we obtain $`x_\mathrm{\Delta }=0.57\pm 0.01`$ and $`y_\mathrm{\Delta }=0.10\pm 0.05`$. The error was estimated from deviations obtained by fitting the $`M=128`$ data over different ranges of L ($`4L600`$) and by comparison with $`M=64`$ data ($`4L300`$). Results are systematically improved by increasing the value of $`M`$.
Finite-size scaling behavior for the XXX model with open boundary conditions and periodic boundary conditions were obtained from DMRG calculations of ground state energies and correlation functions $`S^z(x)S^z(x+r)`$ for different system sizes and separations $`r`$. In our approach, critical exponents are extracted from expectation values at the center of the chain only. The chain size is increased via the infinite-size DMRG method. It is also advantageous to extract power law exponents when there are no logarithmic corrections.
### D Conformal Anomaly
Finally, we may calculate the value of the conformal anomaly, $`c`$. We note that the central charge of the RSOS model and of the spin-$`3/2`$ Heisenberg chain, which is in the same universality class as the spin-$`1/2`$ chain, have previously been obtained using the DMRG. The conformal anomaly can be extracted by finding the coefficient of the $`1/L`$ finite-size correction to the free energy, equivalent at zero temperature to the ground state energy. We fit the ground state energy $`E_0(L)`$ to the following form:
$$E_0(L)=AL+B+C/L+\mathrm{}$$
(67)
The extensive contribution, proportional to $`A`$, and the constant term $`B`$ are non-universal. At the isotropic point $`\gamma =1`$ our results for the case of blocksize $`M=128`$ and for chain lengths in the range $`30L100`$ yield $`C0.323`$. To relate this coefficient to the conformal anomaly we must normalize it by dividing by the speed of low-lying excitations, $`v`$. The speed can be obtained by extrapolation to the thermodynamic limit of the gap to the lowest-lying excitation multiplied by the chain length:
$$v=\underset{L\mathrm{}}{lim}\frac{\mathrm{Gap}(L)\times L}{\pi }.$$
(68)
We find $`v=2.44`$. Now for open boundary conditions,
$$c=\frac{24C}{\pi v}1.01.$$
(69)
This compares well with the value of $`c=1`$ appropriate for a single boson or the pair of left and right moving fermions.
## V Spin-$`1`$ chain
As a final example we apply the DMRG / finite-size scaling method to the isotropic spin-$`1`$ antiferromagnetic chain. This problem is more challenging numerically as the on-site Hilbert space now has dimension $`D=3`$ instead of $`D=2`$. The most general nearest-neighbor Hamiltonian for the spin-$`1`$ chain includes the possibility of a biquadratic spin-spin interaction term:
$$H=\underset{j=0}{\overset{L2}{}}\left[\mathrm{cos}\theta \stackrel{}{S}_j\stackrel{}{S}_{j+1}+\mathrm{sin}\theta (\stackrel{}{S}_j\stackrel{}{S}_{j+1})^2\right]$$
(70)
The phase diagram can be represented on a circle parameterized by $`\theta `$ as depicted in Fig. 10.
Generically there is a gap to excitations in the antiferromagnetic region of the phase diagram, in accord with the Haldane conjecture. The point $`\theta =0`$ corresponds to the usual pure bilinear Heisenberg antiferromagnet. At the point $`\mathrm{tan}\theta =1/3`$ the Hamiltonian can be written as a sum of positive-definite projection operators, and the exact ground state is the AKLT valence bond solid (VBS). Negative $`\mathrm{sin}\theta `$ favors dimerization, as the energy is minimized by concentrating singlet correlations on isolated dimers. The dimerized phase also is gapped: a dimer must be broken to generate a spin excitation. The point that separates the dimerized and Haldane phases lies at $`\theta =\pi /4`$ and can be solved exactly by the Bethe ansatz. The chain is quantum critical at this integrable point. The ground state is non-degenerate here as well as in the dimerized and Haldane phases.
DMRG calculations clearly delineate the two massive phases and the critical point separating them, even for relatively small block Hilbert sizes $`M`$. In Figs. 11 and 12, the blocksize $`M=81`$ for the massive phases. Thus the results are numerically exact up only to chain lengths $`L=10`$. For chain lengths $`L>10`$ the Hilbert space is truncated via the DMRG algorithm. To increase accuracy, results at the critical point were obtained with a larger Hilbert size for the blocks, $`M=256`$.
The induced spin moment at the center of the chain decays exponentially in both the Haldane and the dimerized phases, as expected. The induced dimerization at the chain center also decays exponentially in the Haldane phase, but approaches a non-zero constant in the dimerized phase as it must. Power law decay in both observables occurs at the critical point. Fitting the $`M=256`$ data shown in Fig. 12 at the critical point $`\theta =\pi /4`$ we obtain dimerization exponents $`x_\mathrm{\Delta }=0.37\pm 0.01`$ and $`y_\mathrm{\Delta }=0.3\pm 0.05`$, reflecting the apparent presence of a marginal interaction and consequent multiplicative logarithmic corrections to scaling. Likewise, for a field of $`h=1.0`$ applied to the chain ends, the exponents for the spin operator are $`x_B=0.34\pm 0.01`$ and $`y_B=0.23\pm 0.05`$. The values of the exponents compare to the exact values $`x_\mathrm{\Delta }=3/80.375`$ and $`x_B=3/8`$. To the best of our knowledge there are no analytic results at the integrable point $`\theta =\pi /4`$ (which corresponds to a $`k=2`$ SU(2) WZW model) on the size of the logarithmic corrections $`y_\mathrm{\Delta }`$ and $`y_B`$, at least for open boundary conditions.
Finally, we may repeat the analysis of the conformal anomaly described above in subsection IV D for the case of the spin-1 chain at its critical point. For $`M=256`$ and fitting over chain lengths $`10L26`$ we find that the speed of excitations is $`v=3.69`$, $`C=0.508`$, and hence $`c=1.05`$. This value is close to its exact value of 1, demonstrating that the conformal anomaly can be reliably extracted even from relatively short chains.
## VI Conclusion
We have presented a simple method for studying critical behavior of quantum spin chains. Accurate critical exponents can be extracted. For small on-site Hilbert space sizes ($`D=2`$ for the spin-$`1/2`$ chain and $`D=3`$ for spin-$`1`$ chains) the method does not require supercomputers. Results can be systematically improved by increasing the size of $`M`$, the dimension of the Hilbert space retained in the blocks, up to limits set by machine memory and speed. The DMRG method works best for massive, non-critical, systems, but it is also quite accurate even at critical points. Critical exponents can be calculated at percent level accuracy. We showed that the leading multiplicative logarithmic correction to the scaling of the induced spin moment cancels out in the case of the isotropic spin-$`1/2`$ Heisenberg antiferromagnet. Thus accurate exponents can sometimes be found numerically despite the presence of marginal interactions.
Use of the “finite-size” DMRG algorithm might improve the method, but good results were obtained with the relatively simpler “infinite-size” DMRG algorithm. The reason for this is that the finite-size scaling method employed here focuses on the scaling of observables near the center of the chain only, where the “infinite-size” algorithm is particularly accurate. The method can be used to study new systems. For example, several non-interacting but disordered electron systems, like the integer and spin quantum Hall transitions, can be described by supersymmetric Hamiltonians. In a paper which follows, we employ the combined DMRG/finite-size method in combination with analytic calculations to understand the behavior of these supersymmetric spin chains.
Acknowledgments We thank I. Affleck, M. P. A. Fisher, V. Gurarie, J. Kondev, M. Kosterlitz, A. Ludwig and T. Senthil for useful discussions. This work was supported in part by the NSF under Grants Nos. DMR-9357613, DMR-9712391. Computations were carried out in double-precision C++ on Cray PVP machines at the Theoretical Physics Computing Facility at Brown University. |
warning/0001/hep-ph0001149.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Few years ago , combining several of the currently used philosophies, a high quality description of existing high energy elastic $`pp`$ and $`\overline{p}p`$ scattering data was obtained. The main lessons of this study performed at the Born level were:
(1) an Odderon contribution is absolutely necessary to reproduce quantitatively well the data; while its presence is not explicitly needed at $`t=0`$, its inclusion is necessary to have a good fit of the other $`|t|`$ data, specially in the dip region and in the high-$`|t|`$ domain
(2) hints are found that secondary structures (diffraction-like) develop in angular distributions with increasing energies at intermediate $`|t|`$ values in both $`pp`$ and $`p\overline{p}`$ angular distributions.
In particular, it was suggested that such structure effects should be well visible at LHC while only extremely precise data could perhaps show the effect at RHIC energies.
However, answers to some important points are still incomplete. In particular, what is a good model for the Pomeron ? What is the behavior of the scattering amplitude at high (”asymptotic”) energy ? Are large-$`|t|`$ data dominated by the Odderon ? Better, does a special criterium exist proving the Odderon presence ? Can one settle the question about the sign of $`\alpha _O(0)1\delta _O`$ concerning the intercept<sup>5</sup><sup>5</sup>5Originally it was claimed that $`\delta _O>0`$. More recently, counterarguments have been given to suggest that $`\delta _O`$ should be negative. This possibility had been anticipated in on purely phenomenological grounds and, subsequently, we have found that such a requirement is, in general, consequence of unitarity . However latest QCD calculation gives $`\delta _O=0`$. of the Odderon ? Are secondary structures always predicted at large-$`|t|`$ when $`s`$ increases, i.e. do they arise ”naturally” and are they model dependent ? What is the rôle of eikonalization ? Does an amplitude that fits well the data exhibit automatically a zero in the real part of the even component of the amplitude as required by general theorems ?
Only partial answers presently exist (see and references therein).
The Pomeron remains a most mysterious entity in spite of its resurgence from Diffractive Deep Inelastic Scattering (DDIS) data<sup>6</sup><sup>6</sup>6For an update on the subject, see e.g. .. Many models, however, exist and we are going to probe a few. Even at low-$`|t|`$ (first diffraction cone), it has been shown that existing data do not allow to select among Pomeron models. The present data are, very likely, not yet asymptotic; this (see and references therein) makes it very difficult with the existing data to establish a definite asymptotic behavior for the amplitude.
The rôle of eikonalization has not been fully clarified in spite of having been investigated by many authors but many results have been obtained recently .
The Odderon is instrumental in reproducing the large-$`|t|`$ data. While $`t=0`$ data are presumably dominated by the Pomeron, which in this region hides the Odderon, very precise data could be useful to shed light on its existence .
Predictions of secondary structures have appeared many times in the past . The large spectrum of predictions in the position of these secondary dips shows that things are actually more complicated than anticipated long ago . It is not enough that a given scheme inherently generates oscillations (like the Bessel function of an impact parameter representation); interference effects are very important in determining their position. The model dependence of these predictions, however, is not so important; it is the prediction itself of the existence of secondary structures which matters.
In this paper, four of the above points are taken into special considerations. The first is the investigation of the rôle and properties of the different varieties of eikonalization procedures one can devise. The second concerns the appearance of secondary dips and structures. These two points are strictly interconnected, being the second, to some extent, the physical counterpart of the other. The third is devoted to the behavior of the real part of the even amplitude close to zero. The fourth concerns the rôle of the Odderon in the construction of the amplitude and in the reproduction of the data.
The eikonalization procedure and its consequences is one of the principal subjects we discuss in this paper. We briefly revise (in Sec. 3) the Ordinary Eikonalization (OE) and, after (re)discovering its limits, we proceed to discuss a one-parameter generalization called Quasi Eikonalization (QE) and to propose a three-parameter extension which we term Generalized Eikonalization (GE) (see ). Although a useful tool to alleviate the violations of $`s`$-channel unitarity at some level (as emphasized in ), eikonalization does not mean unitarization.
The effects of Ordinary Eikonalization as compared to the use of the Born amplitude have been studied within a pure Pomeron model (without aiming at quantitatively reproducing the data), and also in a ”more realistic” model including Pomeron, Odderon and secondary Reggeons, fitted to the high energy data for $`pp`$ and $`\overline{p}p`$ elastic scattering . The somewhat surprising results of this ”realistic” approach were:
(i) a failure to find within the eikonalized model as high quality a fit as within its Born approximation , even when readjusting the parameters and even when confining oneself to the ISR data, limited to low-$`|t|`$ ;
(ii) a rapid numerical convergence of the rescattering series: a limited number of rescatterings (4, in addition to the Born term) is sufficient to obtain a very good approximation at present energies ;
(iii) when rescattering corrections were taken into account, a second break in the slope revealed around $`|t|`$ 4 GeV<sup>2</sup> in the angular distribution at 300-500 GeV, creating the seed of a diffraction-type pattern at higher energies ; this break becomes a shoulder and then a true dip moving down to $`|t|`$ 3 GeV<sup>2</sup> when $`\sqrt{s}`$ increases up to 14 TeV. This substructure should be seen at LHC but might even be detected at RHIC if the data are very precise.
Going one step further, the one-parameter extension (QE) and much more, the three-parameters generalization (GE) prove very useful to improve the agreement with the data and, therefore, in removing the conflict found in (i) above. In addition, it helps in understanding the appearance of secondary structures, which stirred considerable interest and which is intriguing enough that we should reconsider further both their origin and their model dependence. The variety of descriptions giving rise to these diffraction-like multiple structures may suggest them to be essentially model independent; on the other hand, this is not established in an unambiguous way and deserves further theoretical analysis<sup>7</sup><sup>7</sup>7We stress once more that several models of $`pp`$ and $`\overline{p}p`$ elastic scattering (see e.g. ) have given hints, in the past, of the possible appearance of a succession of dips or shoulders in the angular distributions, at large-$`|t|`$ values and at superhigh energies.. In the light of this, we have undertaken a most careful analysis of several models both eikonalized and in the Born approximation, trying to ascertain whether or not the predictions of secondary structures could be related to some general pattern. By-products of our investigation turn out to be the verification that the Odderon intercept $`\alpha _O(0)1\delta _O`$ obtained in the various fits is invariably non-positive and, empirically, very close to zero and that the real part of the even amplitude has the zero predicted by general theorems near $`|t|=0`$.
In Sec. 2, we report about several non-eikonalized models with some details on their specific Pomeron and Odderon components. In Sec. 3, we do the same about OE, QE and GE. The results are presented in Sec. 4, some general conclusions are given in Sec. 5.
## 2 The input Born
We focus on the (dimensionless) crossing-even and -odd amplitudes $`a_\pm (s,t)`$ of the $`pp`$ and $`\overline{p}p`$ reactions <sup>8</sup><sup>8</sup>8 Here and in the following, we denote by lower case letters the Born (or input) amplitudes and by the corresponding capital letters their eikonalized counterparts.
$$a_{pp}^{\overline{p}p}(s,t)=a_+(s,t)\pm a_{}(s,t),$$
(1)
for which we have data<sup>9</sup><sup>9</sup>9 For all versions, we fitted the adjustable parameters over a set of $`1000`$ $`pp`$ and $`\overline{p}p`$ data of both forward observables (total cross-sections $`\sigma _t`$ and $`\rho `$ratios of real to imaginary part of the amplitude) in the range $`4\sqrt{s}`$ (GeV)$`1800`$ and angular distributions ($`\frac{d\sigma }{dt}`$) in the ranges $`23\sqrt{s}`$ (GeV)$`630`$ and $`0|t|14`$ GeV<sup>2</sup>. The references to the original literature can be found in . on :
i) total cross-sections
$$\sigma _t=\frac{4\pi }{s}\mathrm{}\mathrm{m}A(s,t=0),$$
(2)
ii) differential cross-sections
$$\frac{d\sigma }{dt}=\frac{\pi }{s^2}\left|A(s,t)\right|^2,$$
(3)
iii) ratio of the real to the imaginary forward amplitudes
$$\rho =\frac{\mathrm{}\mathrm{e}A(s,t=0)}{\mathrm{}\mathrm{m}A(s,t=0)}.$$
(4)
The crossing even part in the Born amplitude is a Pomeron (to which a $`f`$Reggeon is added) while the crossing odd part is an Odderon (plus an $`\omega `$Reggeon)
$$a_+(s,t)=a_P(s,t)+a_f(s,t),a_{}(s,t)=a_O(s,t)+a_\omega (s,t).$$
(5)
For simplicity the two Reggeons have been taken in the standard form
$$a_R(s,t)=a_R\stackrel{~}{s}^{\alpha _R(t)}e^{b_Rt},\alpha _R(t)=\alpha _R(0)+\alpha _R^{}t,(R=f\mathrm{and}\omega ),$$
(6)
where $`a_f`$ ($`a_\omega `$) is real (imaginary). We begin with trajectories whose parameters are fixed as in previous works (for example ) $`\alpha _f(t)=0.69+0.84t`$, and $`\alpha _\omega (t)=0.47+0.93t`$ (with $`t`$ in GeV<sup>2</sup>), close to the values obtained in other recent fits (e.g. ). As it turns out, however, a best fit requires some variation of these parameters. Thus, at the price of economy in the parameters, we end up letting them vary.
We have investigated wide classes of choices where the input amplitude (”Born term”) for the Pomeron $`a_P(s,t)`$ and for the Odderon $`a_O(s,t)`$ is either a monopole (i.e. a simple pole in the angular momentum $`J`$plane) or a ”dipole” (i.e. a linear combination of a simple pole with a double pole).
The forms of $`a_P(s,t)`$ in the case of a monopole $`(M)`$ and of a dipole ($`D`$) are
$$a_P^{(M)}(s,t)=a_P\stackrel{~}{s}^{\alpha _P(t)}e^{b_Pt},$$
(7)
and
$$a_P^{(D)}(s,t)=a_P\stackrel{~}{s}^{\alpha _P(t)}\left[e^{b_P(\alpha _P(t)1)}(b_P+\mathrm{}n\stackrel{~}{s})+d_P\mathrm{}n\stackrel{~}{s}\right],$$
(8)
where $`a_P`$ is real.
The difference between a monopole and a dipole results in an amplitude for the second that grows with an additional power of $`\mathrm{}ns`$.
The Odderon may be constructed with the same requirements. It is, however, known that the rôle of the Odderon at $`t=0`$ is negligible but no theoretical prescription is known as how to cut it. A simple way out is to multiply the monopole or dipole form by a convenient damping factor. We choose
$$a_O(s,t)=(1\mathrm{exp}\gamma t)a_O^{(M)}(s,t),$$
(9)
or
$$a_O(s,t)=(1\mathrm{exp}\gamma t)a_O^{(D)}(s,t).$$
(10)
In (9) (or (10)), the amplitude on the r.h.s. is constructed along the same lines as in (7) (or (8)) for $`a_P^{(M,D)}(s,t)`$. $`a_P`$, however, is real while $`a_O`$ is imaginary. As usual,
$$\stackrel{~}{s}=\frac{s}{s_0}e^{i\frac{\pi }{2}},(s_0=1\mathrm{GeV}^2),$$
(11)
enforces $`su`$ crossing and $`\alpha _i(t)`$ are the trajectories taken, for simplicity, of the linear form<sup>10</sup><sup>10</sup>10Linear trajectories are an oversimplification that, strictly, violates analyticity. In addition, at large $`|t|`$ this may be dangerous in practice. We ignore this complication.
$$\alpha _i(t)=1+\delta _i+\alpha _i^{}t,(i=P,O).$$
(12)
It appears impossible to discriminate between (D) or (M), on general grounds; only the phenomenological results seem to prefer (D) over (M). For the sake of economy we confine our presentation to the dipole case, which gives somewhat better phenomenological results.
Some authors maintain that a perturbative (a large-$`|t|`$) term behaving like $`|t|^4`$ (and complying with perturbative QCD requirements according to <sup>11</sup><sup>11</sup>11 We should, however, not forget that at, even at the largest $`|t|`$ values, the ratio $`|t|/s`$ is really rather small so that we are in a domain closer to the usual Regge kinematics than to that of perturbative QCD. ) is to be added to the Odderon. When the Born amplitude is eikonalized, however, all rescattering corrections implied by eikonalization are, in principle, already taken into account. Adding another large-$`|t|`$ term at the Born level would mimic further rescattering corrections and would lead to double counting in the eikonalized models. We shall not consider this option.
We remark that most good fits require $`\delta _P>0`$ implying what is known as a supercritical Born Pomeron i.e. a Born amplitude which, taken at face value, will eventually exceed the Froissart-Martin unitarity bound even though at extremely high energies (other kinds of troubles would arise much earlier ). This special violation of unitarity is removed by all kinds of eikonalization. Nevertheless, one must verify that the unitarity constraints
$$\delta _P\delta _O,\mathrm{and}\alpha _P^{}\alpha _O^{}.$$
(13)
are satisfied (see ). The slope parameter for the Pomeron, finally, is expected to be in the vicinity of its ”world” value $`\alpha _P^{}0.25`$ GeV<sup>-2</sup> and this turns out to be, indeed, the result of the fit (see Section 4.2).
As the last comment, we recall that, in the context of the choice of the eikonalization procedure, a singular solution is, in principle, possible, whereby the Odderon dominates over the Pomeron . For the sake of completeness, we have also tried this option, unphysical as this appears but, as expected, such possibility is ruled out by the results of the fits; the fit with an Odderon dominating over the Pomeron is rather poor and unacceptable.
## 3 Eikonalization procedures
In eikonal models, the scattering amplitudes are expressed in the impact parameter (”$`b`$”) representation. First, one defines the Fourier-Bessel’s (F-B) transform of the Born amplitude
$$h_{pp}^{\overline{p}p}(s,b)=\frac{1}{2s}_0^{\mathrm{}}a_{pp}^{\overline{p}p}(s,q^2)J_0(bq)q𝑑q\mathrm{with}q=\sqrt{t}.$$
(14)
This is related to the eikonal function (”eikonal” for brevity) by
$$\chi _{pp}^{\overline{p}p}(s,b)=2h_{pp}^{\overline{p}p}(s,b).$$
(15)
The (complete) analytical forms of the Born amplitudes (both (M) and (D)) in $`b`$-space are given in Appendix A.
In all eikonalization procedures, one first derives the eikonalized amplitude $`H_{pp}^{\overline{p}p}(s,b)`$ in the $`b`$-representation; the inverse F-B transform leads then to the usual eikonalized amplitude in the $`st`$ space
$$A_{pp}^{\overline{p}p}(s,t)=2s_0^{\mathrm{}}H_{pp}^{\overline{p}p}(s,b)J_0(b\sqrt{t})b𝑑b.$$
(16)
The main technical problem of eikonalization is the derivation of $`H_{pp}^{\overline{p}p}(s,b)`$ once $`h_{pp}^{\overline{p}p}(s,b)`$ are given. In what follows we make explicit this step in, we believe, the most general form so far derived.
### 3.1 Ordinary and Quasi Eikonalization
In the ordinary eikonal (OE) formalism, $`H_{pp}^{\overline{p}p}(s,b)`$ is the sum over all rescattering diagrams in the approximation when there are only two nucleons on the mass shell in any intermediate state
$$H_{pp,QE}^{\overline{p}p}(s,b)=\frac{1}{2i}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{\left[2ih_{pp}^{\overline{p}p}(s,b)\right]^n}{n!}\right).$$
(17)
This limitation neglects the possibility to take into account multiparticle states. In the quasi eikonal (QE) procedure , the effect of these multiparticle states in the various exchange diagrams is realized introducing one additional ”weight” parameter $`\lambda `$ and the eikonalized amplitude in the $`b`$-representation (17) is replaced by
$$H_{pp,QE}^{\overline{p}p}(s,b)=\frac{1}{2i}\underset{n=1}{\overset{\mathrm{}}{}}\lambda ^{n1}\frac{\left[2ih_{pp}^{\overline{p}p}(s,b)\right]^n}{n!}.$$
(18)
The above series is meant to represent the sum of all possible multiple exchanges of Pomerons, Odderons and secondary Reggeons ($`n=1`$ corresponds to the Born approximation, $`n=2`$ to double exchanges, etc… ). Its explicit analytical form is
$$H_{pp,QE}^{\overline{p}p}(s,b)=\frac{1}{2i\lambda }\left(\mathrm{exp}\left[i\lambda \chi _{pp}^{\overline{p}p}(s,b)\right]1\right).$$
(19)
As it is obvious, the value $`\lambda =1`$ corresponds to OE, which appears, therefore, as a particular case of QE.
However, it is not clear why all intermediate states between the exchanges of two Pomerons or two Odderons (or between one Pomeron and one Odderon) could be described by just one and the same parameter $`\lambda `$ or, differently stated that all the weights for the various intermediate internal couplings (two Pomerons, two Odderons or one Pomeron and one Odderon) should be the same. It would appear more ”natural” that the various exchanges should require different weights. Differently rephrased, in the QE procedure, we do not distinguish intermediate states between $`PP,OO`$ and $`PO`$ exchanges. Giving up this assumption gives rise to a new kind of generalized eikonal (GE) procedure where all these intermediate states may have different weights.
### 3.2 Generalized Eikonalization
#### 3.2.1 with 3 $`\lambda :\lambda _\pm ,\lambda _0.`$
Consider again the separate form of the amplitude (1), and let the crossing-even and crossing-odd input in the $`b`$-representation be
$$h_\pm h_\pm (s,b)=\frac{1}{2s}_0^{\mathrm{}}𝑑qqJ_0(bq)a_\pm (s,q^2),(q^2=t).$$
(20)
Here, postponing for a moment the consideration of the most general scheme (5) when secondary Reggeons are included, we temporarily simplify the notation for the crossing even- and the crossing odd-part as if they were made by just the Pomeron and the Odderon respectively (later, we will reinstate the complete contribution)
$$a_+(s,t)=a_P(s,t),a_{}(s,t)=a_O(s,t).$$
(21)
A priori, we have three different configurations of exchanges in the intermediate states which we show diagrammatically in Fig. 1 and where the various possibilities, $`PP,OO`$ and $`PO`$ are described, phenomenologically, by three constants $`\lambda _+,\lambda _{},\lambda _0`$.
With this notation, we can deduce
$$H^{\overline{p}p}(s,b)=h_++h_{}+H[PP]+H[OO]+H[PO]+H[OP].$$
(22)
where (see for the details of the derivation)
$$\begin{array}{ccc}\hfill 2i\lambda _+H[PP]& =\underset{n=2}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{n1}{}}\underset{i=1}{\overset{m}{}}\frac{1}{(m+n)!}\left(\genfrac{}{}{0pt}{}{n1}{i}\right)\left(\genfrac{}{}{0pt}{}{m1}{i1}\right)z^ix^ny^m\hfill & \\ & +\underset{n=2}{\overset{\mathrm{}}{}}\underset{m=n}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{n1}{}}\frac{1}{(m+n)!}\left(\genfrac{}{}{0pt}{}{n1}{i}\right)\left(\genfrac{}{}{0pt}{}{m1}{i1}\right)z^ix^ny^m+\underset{n=2}{\overset{\mathrm{}}{}}\frac{1}{n!}x^n\hfill & \\ & =z\underset{n=2}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{\mathrm{}}{}}\frac{x^ny^m}{(n+m)!}(n1)_2F_1(1m,2n;2;z)\hfill & \\ & +e^xx1,\hfill & \end{array}$$
(23)
$$\begin{array}{ccc}\hfill 2i\lambda _0H[PO]& =\underset{n=1}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{n}{}}\underset{i=1}{\overset{m}{}}\frac{1}{(m+n)!}\left(\genfrac{}{}{0pt}{}{n1}{i1}\right)\left(\genfrac{}{}{0pt}{}{m1}{i1}\right)z^ix^ny^m\hfill & \\ & +\underset{n=1}{\overset{\mathrm{}}{}}\underset{m=n+1}{\overset{\mathrm{}}{}}\underset{i=1}{\overset{n}{}}\frac{1}{(m+n)!}\left(\genfrac{}{}{0pt}{}{n1}{i1}\right)\left(\genfrac{}{}{0pt}{}{m1}{i1}\right)z^ix^ny^m\hfill & \\ & =z\underset{n=1}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{\mathrm{}}{}}\frac{x^ny^m}{(n+m)!}_2F_1(1m,1n;1;z)\hfill & \end{array}$$
(24)
with
$$x=2i\lambda _+h_+,y=2i\lambda _{}h_{}z=\frac{\lambda _0^2}{\lambda _+\lambda _{}}.$$
$`H[OO]`$ is obtained from $`H[PP]`$ with the replacement $`h_+h_{}`$ and $`\lambda _+\lambda _{}`$, and $`H[OP]=H[PO]`$.
The amplitude $`H_{pp}(s,b)`$ has the same form as $`H^{\overline{p}p}(s,b)`$ with the replacement $`h_{}h_{}`$.
Unexpectedly, one can obtain a compact analytical form from (22). Omitting all details of calculations, which can be found in , the final expression for the three-parameters eikonalized amplitudes are
$$\begin{array}{ccc}\hfill H_{pp,GE}^{\overline{p}p}(s,b)& =\frac{i}{2(\lambda _0^2\lambda _+\lambda _{})}\{a+e^{i\left(\lambda _+h_+\pm \lambda _{}h_{}\right)}\hfill & \\ & \times [a\mathrm{cos}\varphi _\pm +i\frac{c_+h_+\pm c_{}h_{}}{\varphi _\pm }\mathrm{sin}\varphi _\pm ]\},\hfill & \end{array}$$
(25)
where we have introduced three constants $`a`$ and $`c_\pm `$ defined as
$$a=2\lambda _0\lambda _+\lambda _{},$$
(26)
$$c_\pm =\lambda _+\lambda _{}2\lambda _0^2\lambda _\pm ^2+2\lambda _0\lambda _\pm ,$$
(27)
in terms of the parameters of the model and the functions (of $`s`$ and $`b`$)
$$\varphi _\pm =\sqrt{(\lambda _+h_+\lambda _{}h_{})^2\pm 4\lambda _0^2h_+h_{}}.$$
(28)
Considering a general case, when there are no any special relations between $`\lambda _i`$, we have found in that the unitarity inequality
$$|H_{pp,GE}^{\overline{p}p}(s,b)|1$$
can be satisfied, in general, only if $`\delta _O0`$ <sup>12</sup><sup>12</sup>12The obvious inequalities $`|h_{}||h_+|`$ and $`|\mathrm{}\mathrm{e}h_+||\mathrm{}\mathrm{m}h_+|,\mathrm{}\mathrm{m}h_+>0`$ , which are valid at high energy, are assumed.. Two special cases, namely, $`\lambda _0^2=\lambda _{}\lambda _+`$ (see below) and $`\lambda _+=\lambda _0`$ allow $`\delta _O`$ to be positive. However in all cases unitarity requires the following restrictions
$$\delta _O\delta _P,\alpha _O^{}(0)\alpha _P^{}(0),\lambda _+1/2.$$
(29)
As anticipated above, it is easy to prove that these results, obtained in the case of 2 Reggeons (P and O), hold in the case where 4 Reggeons are grouped 2 by 2 to form a crossing even ($`P+f`$) and a crossing odd ($`O+\omega `$) contribution with the original definitions (5).
#### 3.2.2 with 2 $`\lambda :\lambda _\pm .`$
A considerable simplification is brought if the factorization $`\lambda _0=\sqrt{\lambda _+\lambda _{}}`$ is assumed (this is also treated in great details in ). In practice, the main advantage of this particular case are simplified expressions for the required amplitudes resulting in a significant gain in computer time when fitting the data. In this case, the eikonalized amplitude has the form
$$\begin{array}{ccc}\hfill H_{pp,GE}^{\overline{p}p}(s,b)& =h_+\pm h_{}+\left(\frac{h_+\sqrt{\lambda _+}\pm h_{}\sqrt{\lambda _{}}}{h_+\lambda _+\pm h_{}\lambda _{}}\right)^2\hfill & \\ & \times \left(\frac{e^{2i(h_+\lambda _+\pm h_{}\lambda _{})}1}{2i}(h_+\lambda _+\pm h_{}\lambda _{})\right).\hfill & \end{array}$$
(30)
From unitarity, either
$$\delta _O0,\lambda _+1/2\text{with}\lambda _{}\text{arbitrary}$$
(31)
or
$$\lambda _{}=\lambda _+1/2\text{with}0\delta _O\delta _P.$$
(32)
The second case (Eq.(32)) coincides with the previously considered QE method.
### 3.3 Rescattering series (in $`st`$ space).
The fact that the eikonalization procedures discussed previously lead to close analytical forms ((19) or (25)) for the amplitudes $`H(s,b)`$, allows us, in principle, to use them in the F-B transform (16) in order to derive the completely eikonalized physical amplitudes $`A(s,t)`$. The compact analytical expressions ((19) or (25)), however, require a very time-consuming numerical integration. The infinite expansions ((18) or (23),(24)), on the other hand, can be more convenient if one has a rapid convergence of the rescattering series. Fortunately, this condition is fulfilled by both the monopole and the dipole. These models are, therefore, interesting candidates to test the number and quality of exchanges necessary to give a final good accuracy in the calculation of the observables.
To be specific, we rewrite the QE amplitude isolating the Born term
$$A_{pp,QE}^{\overline{p}p}(s,t)=a_{pp}^{\overline{p}p}(s,t)+\underset{n=2}{\overset{\mathrm{}}{}}\lambda ^{n1}a_{pp;n}^{\overline{p}p}(s,t),$$
(33)
where from (14) and (16)
$$a_{pp;n}^{\overline{p}p}(s,t)=\frac{i}{n!}s_0^{\mathrm{}}\left[2ih_{pp}^{\overline{p}p}(s,b)\right]^nJ_0(b\sqrt{t})b𝑑b.$$
(34)
Each rescattering term can be calculated analytically only in some specific cases, for example again in the monopole or dipole models (see e.g. for the dipole, the monopole calculations are less involved). In practice, we find that a finite number of $`4`$ terms is sufficient to insure proper convergence of the rescattering series ($`n[2,5]`$).
In the GE case, we rewrite the amplitude as
$$A_{pp,GE}^{\overline{p}p}(s,t)=a_{pp}^{\overline{p}p}(s,t)+\underset{n_+=0}{\overset{\mathrm{}}{}}\underset{n_{}=0}{\overset{\mathrm{}}{}}a_{pp;n_+,n_{}}^{\overline{p}p}(s,t),$$
(35)
where, we have to compare (35) with (23), (24) to obtain the identification. The analytical expressions for evaluating the double series are given in Appendix B in the (most involved) case of the dipole model (Pomeron + Odderon + Reggeons).
In agreement with what we found for QE, the convergence of the rescattering series for GE is obtained by keeping only the four first terms ($`n_\pm [0,1]`$).
## 4 Results
As already mentioned, only the results for the dipole model are shown in what follows.
### 4.1 Born input amplitude
We have satisfied ourselves that the general pattern remains always the same : a wisely chosen ”Born” amplitude can reproduce the data very well but, depending on this choice, the Pomeron (and the Odderon) become supercritical and the Froissart-Martin bound is, in principle, exceeded. At the Born level, secondary structures may or may not appear; when they do, they are generally due to an additive contribution to the simple (monopole and dipole) models. For completeness, given the simplicity of the approach, we give in Table 1 the parameters of the fit. Surprisingly, the Odderon intercept equals 1, as recently claimed . The reader, however, should keep in mind that this Born approach and its parameters should not be considered as anything fundamental; they can be used as a shortcut for giving a reasonable account of the existing data but hardly to derive general properties.
| | Pomeron | Odderon |
| --- | --- | --- |
| $`\delta _i`$ | 0.071 | 0.0 |
| $`\alpha _i^{}`$ (GeV$`{}_{}{}^{2})`$ | 0.28 | 0.12 |
| $`b_i`$ | 14.56 | 28.1 |
| $`a_i`$ | -0.066 | 0.10 |
| $`d_i`$ | 0.07 | -0.06 |
| $`\gamma `$(GeV$`{}_{}{}^{2})`$ | - | 1.56 |
| | $`f`$-Reggeon | $`\omega `$-Reggeon |
| $`a_R`$ | -14.0 | 9.0 |
| $`b_R`$ (GeV$`{}_{}{}^{2})`$ | 1.64 | 0.38 |
| $`\alpha _R(0)`$ | 0.72 | 0.46 |
| $`\alpha _R^{}`$ (GeV$`{}_{}{}^{2})`$ | 0.50 | 0.50 |
Table 1. Parameters of the dipole model fitted at the Born level (dipole Pomeron $`i=P`$, dipole Odderon $`i=O`$ vanishing at $`t=0`$, secondary Reggeons $`R=f,\omega `$).
### 4.2 Eikonalized models
A general feature of all eikonalizations models is that, even when the original Born amplitude exceeds the unitarity limit (remind a good fit generally requires $`\delta _P>0`$), this violation is removed upon eikonalizing.
We remark that the OE procedure does not change the number of parameters chosen at the Born level; one parameter ($`\lambda `$) is added within the QE procedure and two ($`\lambda _\pm `$) or three ($`\lambda _\pm ,\lambda _0`$) within the GE procedure. Furthermore, we may easily reduce the GE model to the QE model by setting $`\lambda _+=\lambda _{}=\lambda _0\lambda `$ and the QE model to the OE model by setting $`\lambda =1`$.
We tested all procedures of eikonalization, either complete or partial. In the latter case, typically, one may choose not to eikonalize the Reggeons because they do not induce a unitarity violation. Whatever the procedure for eikonalizing, we find that the parameters obtained and the conclusions are qualitatively the same.
From the best fit view point some comments help the reader :
(i) the set of experimental data which are very difficult to reproduce with non vanishing eikonalized dipole Odderon are the ratios $`\rho _{pp}^{\overline{p}p}(s,t=0)`$. This justifies our choice (9-10) of a Born Odderon input vanishing at $`t=0`$;
(ii) leaving the secondary Reggeons parameters free to be adjusted improves considerably the quality of the fit to the dip in the ISR energy domain.
#### 4.2.1 Results of the OE and QE fits
Invariably (and surprisingly), ordinary eikonalization (OE) leads to a fit which is poorer than in the Born case but secondary structures emerge.
The QE version of the dipole model improved with respect to OE case is still poorer than the one obtained at the Born level but one finds a good reproduction of the data up to and including the dip for $`pp`$ and the shoulder for $`\overline{p}p`$.
In the QE version with fixed trajectories for the secondary Reggeons, we find a ”supercritical” Pomeron with $`\delta _P0.06`$ (i.e. lower than the value found in ) and a ”critical” Odderon $`\delta _O0.03`$ as expected. The slope parameter for the Pomeron $`\alpha _P^{}0.25`$ GeV<sup>-2</sup> agrees with the ”world” value, and for the Odderon we find $`\alpha _O^{}0.11`$ GeV<sup>-2</sup>. The single parameter characterizing the method of quasi eikonalization with respect to the ordinary one is found closed to its lower unitarity limit $`\lambda 0.5`$. This, in practice, tends to reduce the effect of high multiple exchanges.
Concerning the shape of the diffraction like-structures, we find significant modifications due to QE with respect to previous work in which the OE method had been used. Specifically, the dip-bump secondary structure shifts towards somewhat lower-$`|t|`$ and delays its appearance till higher energies are reached. More precisely, in the QE (i) the first dip moves down from $`|t|`$ 1.2 GeV<sup>2</sup> to 0.5 GeV<sup>2</sup> when $`\sqrt{s}`$ goes up from 60 GeV to 14 TeV; (ii) a break in the slope appears around $`|t|`$ 4.0 GeV<sup>2</sup> when the energy is around 500 GeV, becoming a shoulder and then a true dip which recedes to $`|t|`$ 1.5 GeV<sup>2</sup> when $`\sqrt{s}`$ increases to 14 TeV.
It is very instructive to compare the relative virtues of OE and QE. Generally speaking, as repeatedly stated, both eliminate conflicts with the unitarity limit and the convergence of the rescattering series is comparable (see above), but the QE method appears to cure some undesirable features of the OE, regarding the quality of the fit.
#### 4.2.2 Results of the GE fits
Following the same motivations as above for comparing the QE /OE versions, we discuss now the implications of generalizing the eikonalization with the two or three parameters $`\lambda _\pm ,\lambda _0`$ (instead of a unique parameter $`\lambda `$ for the QE case), using the same Born amplitude; again, we report only the dipole results.
The GE version with two $`\lambda `$-parameters (and with fixed Reggeon trajectories) leads to a good reproduction of the data with well structured secondary dips. The various values of the parameters are slightly different from the version with one $`\lambda `$, in particular $`\lambda _+0.5`$ and $`\lambda _{}0.44`$.
| $`\lambda _+`$ | 0.5 | |
| --- | --- | --- |
| $`\lambda _{}`$ | 0.55 | |
| $`\lambda _0`$ | 1.24 | |
| | Pomeron | Odderon |
| $`\delta _i`$ | 0.073 | -.0.005 |
| $`\alpha _i^{}`$ (GeV$`{}_{}{}^{2})`$ 0.27 | 0.054 | |
| $`b_i`$ | 9.0 | 26.6 |
| $`a_i`$ | -0.114 | -0.019 |
| $`d_i`$ | 0.165 | -0.09 |
| $`\gamma `$(GeV$`{}_{}{}^{2})`$ | $``$ | 1.37 |
| | $`f`$-Reggeon | $`\omega `$-Reggeon |
| $`a_R`$ -12.95 | 16.44 | |
| $`b_R`$ (GeV$`{}_{}{}^{2})`$ | 1.24 | 3.50 |
| $`\alpha _R(0)`$ | 0.81 | 0.47 |
| $`\alpha _R^{}`$ (GeV$`{}_{}{}^{2})`$ | 1.07 | 0.57 |
Table 2. Parameters of the dipole model fitted with the more sophisticated GE procedure (see also Table 1).
The version with three $`\lambda `$-parameters and fixed trajectories for the secondary Reggeons gives also a good reproduction of the data. The situation about the diffractive structures is partially different, now the break around $`|t|4`$ GeV<sup>2</sup> becomes a dip which moves to $`|t|3`$ GeV<sup>2</sup> at TeV energies. The various parameters are close to those of the previous case (for two $`\lambda `$-parameters); the value of $`\delta _O`$ remains negative and moves closer to zero. For the $`\lambda `$-s we find $`\lambda _+0.5,\lambda _{}0.1,\lambda _00.86`$.
The best fit is obtained if we allow some variation for the intercepts and slopes of the Reggeon trajectories. The values of free parameters are given in Table 2.
Three points are worth emphasizing:
i) the Odderon intercept $`\alpha _O(0)1\delta _O`$ consistently turns out to be negative in agreement with general arguments ; however in practice such a small value is obtained that we do not contradict the most recent QCD value
ii) the real part of the even amplitude automatically exhibits a zero at small $`|t|`$ values (typically, $`|t|.30`$ GeV<sup>2</sup> at $`\sqrt{s}=546`$ GeV) and this value recedes towards zero as $`\sqrt{s}`$ increases (typically $`|t|.27`$ GeV<sup>2</sup> at $`\sqrt{s}=1800`$ GeV) and we predict it at $`|t|.23`$ at $`\sqrt{s}=14`$ TeV. This result is in agreement with a general theorem by A. Martin . Also, a second zero appears for $`|t|1.5`$ GeV<sup>2</sup> at $`\sqrt{s}=546`$ GeV which moves to $`|t|1.25`$ GeV<sup>2</sup> at $`\sqrt{s}=1.8`$ TeV (see Table 3);
iii) the eikonalized Odderon contributes to reproduce perfectly the large-$`|t|`$ region.
| energy | 1<sup>st</sup> zero | 2<sup>d</sup> zero |
| --- | --- | --- |
| | (GeV<sup>2</sup>) | (GeV<sup>2</sup>) |
| 546 GeV | 0.30 | 1.5 |
| 1800 GeV | 0.27 | 1.25 |
| 14 TeV | 0.23 | 0.95 |
| 40 TeV | 0.17 | 0.85 |
Table 3. Positions ($`|t|`$ values) of the first two zeros of the real part of the even eikonalized amplitude.
We observe the same evolution of the structures as in the previous case.
Concerning the values of the GE parameters, collected in Table 2, we find a ”supercritical” Born Pomeron with $`\delta _P0.073`$ (i.e. greater than the QE value), while the slope parameter for the Pomeron is $`\alpha _P^{}0.27`$ GeV<sup>-2</sup> and for the Odderon $`\alpha _O^{}0.05`$ GeV<sup>-2</sup>. The tree parameters characterizing the effect of the generalized method of eikonalization are $`\lambda _+0.5`$, $`\lambda _{}=0.55`$ and $`\lambda _0=1.24`$.
Thus, the generalized eikonalization procedure gives better results than the quasi eikonalization and a fortiori than the ordinary eikonalization.
That the $`\chi ^2/d.o.f.(7.0)`$ remains pretty large is the consequence of not having made any ”wise selection” of the data. The resulting curves, however, are quite satisfactory as it is shown in Figs. 2 - 5 (the results are given for the complete set of parameters in Table 2).
The extrapolations of the total cross section and of the $`\rho `$-ratio are shown in Fig. 6. The angular distributions for the energies to be reached in the near future exhibit the secondary structure especially at LHC in Fig. 7.
Figure 2. Comparison with the data of the fit to total cross-sections for $`\overline{p}p`$ (full dots) and $`pp`$ (hollow triangles) processesfor the most sophisticated generalized eikonalization (GE) procedure.
Figure 3. Same as Fig.2 for $`\rho `$-ratios.
Figure 4. Comparison with the data of the fit to differential cross-sections for $`pp`$ process for the most sophisticated generalized eikonalization (GE) procedure. A $`10^2`$ factor between each successive curve is omitted.
Figure 5. Same as Fig.4 for $`\overline{p}p`$ process. The Tevatron data are not fitted.
Figure 6. Calculated observables within the GE dipole model, versus the energy and compared to the data (cf ) : total cross-section $`\sigma _{tot}`$ (the cosmic ray data are not fitted) and $`\rho `$-ratio.
Figure 7. Extrapolations to RHIC and LHC energies of the calculated $`pp`$ differential cross-sections.
## 5 Concluding remarks
Let us try to answer some of the questions raised in the Introduction. Of course, we do not have the final prescription for the Pomeron. Many of the forms discussed above give a good reproduction of the data; several of them (and many others in the literature) seem to work well both at the Born and at the eikonalized level (in particular, the Dipole Pomeron).
Often, the Born Pomeron is found to be supercritical ($`\delta _P>0`$) which implies an intrinsic problem with unitarity; this is removed by (all kinds of) eikonalization. Thus, the rôle of eikonalization is very important for the asymptotic behavior of all physical quantities. In all cases the eikonalization restores the correct high energy behavior of the supercritical Pomeron. While the data for total cross-sections do not contradict the $`\mathrm{}n^2s`$ behavior resulting from the eikonalization of a supercritical Born Pomeron, they are not incompatible with a $`\mathrm{}ns`$ form. The inclusion in the fit of the data at $`t0`$ is absolutely necessary to get an unambiguous conclusion on the behavior of all physical quantities.
The presence of the Odderon contribution, as repeatedly emphasized, is necessary to reproduce well the angular distributions data in the dip-region and for large-$`|t|`$ values but its contribution is required by the fit to be negligibl in the forward domain.
The problem of the Odderon intercept remains very complicated but the general agreement, in LLA, is now that the Odderon intercept is closed to 1 with $`\delta _O<0`$ or $`\delta _O=0`$ . This agrees with our findings (see also ).
A burning question concerns whether or not it is possible to get a definite prediction about the existing of secondary structures. At the Born level, the presence or absence of secondary structures rests on the specific properties of the Born amplitude (like an oscillatory component in the Pomeron amplitude). In this case, therefore, the prediction of secondary structures appear quite model dependent. The rôle of eikonalization is very important in this context. In the dipole case, structures appear in the angular distribution as soon as a double Pomeron exchange is taken into account; the trend consolidates when the number of rescattering corrections $`n`$ increases and takes a definite form when several exchanges are included. This appears to be the case in all eikonalization procedures. We conclude that secondary structures are unambiguously predicted by any eikonalization process. This reinforces previous conclusions by other authors . In fact, as emphasized by Horn and Zachariasen , oscillations in $`t`$ should be expected from the properties of Bessel functions in the F-B transforms unless some special feature of the eikonal destroys them.
Of all eikonalization procedures discussed, GE with 3 parameters leads to the best account of the data.
Finally, we emphasize that the real part of the even amplitude at high energy has a zero in the small-$`t`$ region, as anticipated by a general theorem .
In conclusion, while we believe that LHC will definitely prove (or disprove) the validity of our predictions of secondary structures and about the zero of the real part of the even amplitude, we insist on how valuable it would be to have both $`pp`$ and $`p\overline{p}`$ options available, at the same machine and at the highest energies in order to check not only our predictions but a whole host of theoretical high energy theorems.
Acknowledgements. We would like to thank G. Lamot for his help with the fortran code in particular for the hypergeometric functions. Two of us (EM and EP) would like to thank the Institut de Physique Nucléaire de Lyon for the hospitality and two of us (MG and EM) would like to thank the Theory Physics Department of the University of Torino for the hospitality. Financial support by the INFN and the MURST of Italy and from th IN2P3 of France is gratefully acknowledged.
APPENDIX A
Analytical Born amplitude in the $`b`$-space
We have now to define the analytical expressions of the Born amplitudes in b-space
$$h_{pp}^{\overline{p}p}(s,b)=h_f(s,b)+h_P(s,b)\pm \left[h_O(s,b)+h_\omega (s,b)\right]h_+\pm h_{}$$
$`(A1)`$
from which we will derive the eikonalized amplitude. With our choices of Born (s,t) amplitudes, all the analytical F-B’s transforms are readily obtained<sup>13</sup><sup>13</sup>13 Recall that the ”couplings” $`a_f,a_P`$ are real and $`a_\omega ,a_O`$ are imaginary; for the secondary Reggeons
$$h_R(s,b)=\frac{1}{2}a_R\frac{\stackrel{~}{s}^{\alpha _R(0)}}{s}\frac{\mathrm{exp}(\frac{b^2}{4B_R})}{2B_R};B_R=\alpha _R^{}\mathrm{}n\stackrel{~}{s}+b_R,R=(f,\omega ),$$
$`(A2)`$
where we have defined $`B_R`$ in terms of the slopes $`b_R`$ introduced earlier in (6)). The Pomeron part depends on our choice (7 or 8): for the monopole we would have
$$h_P^{(M)}(s,b)=\frac{1}{2}a_P\frac{\stackrel{~}{s}^{\alpha _P(0)}}{s}\frac{\mathrm{exp}(\frac{b^2}{4B_P})}{2B_P};B_P=\alpha _B^{}\mathrm{}n\stackrel{~}{s}+b_P,$$
$`(A3)`$
while for the dipole
$$h_P^{(D)}(s,b)=\frac{ia_P}{4\alpha _P^{}s_0}\left(e^{r_{1,P}\delta _P\frac{b^2}{4B_{1,P}}}+d_Pe^{r_{2,P}\delta _P\frac{b^2}{4B_{2,P}}}\right).$$
$`(A4)`$
For our Odderon monopole (9) we have
$$h_O^{(M)}(s,b)=\frac{1}{2}a_O\frac{\stackrel{~}{s}^{\alpha _O(0)}}{s}\left[\frac{\mathrm{exp}(\frac{b^2}{4B_O})}{2B_O}\frac{\mathrm{exp}(\frac{b^2}{4\stackrel{~}{B}_O})}{2\stackrel{~}{B}_O}\right],$$
$`(A5)`$
where $`B_O=\alpha _O^{}\mathrm{}n\stackrel{~}{s}+b_O`$ and $`\stackrel{~}{B}_O=\alpha _O^{}\mathrm{}n\stackrel{~}{s}+b_O+\gamma `$. Finally, for our Odderon dipole (10)
$$\begin{array}{ccc}h_O^{(D)}(s,b)\hfill & =\hfill & \frac{ia_O}{4s_0}(e^{r_{1,O}\delta _O\frac{b^2}{4D_{1,O}}}\frac{r_{1,O}}{D_{1,O}}e^{r_{1,O}\delta _O\frac{b^2}{4\stackrel{~}{D}_{1,O}}}\frac{r_{1,O}}{\stackrel{~}{D}_{1,O}}\hfill \\ & +\hfill & d_Oe^{r_{2,O}\delta _O\frac{b^2}{4D_{2,O}}}\frac{r_{2,O}}{D_{2,O}}d_Oe^{r_{2,O}\delta _O\frac{b^2}{4\stackrel{~}{D}_{2,O}}}\frac{r_{2,O}}{\stackrel{~}{D}_{2,O}}).\hfill \end{array}$$
$`(A6)`$
We have defined
$$r_{1,i}=\mathrm{}n\stackrel{~}{s}+b_i,r_{2,i}=\mathrm{}n\stackrel{~}{s},(i=P,O),$$
$`(A7)`$
and
$$B_{i,P}=\alpha _P^{}r_{i,P},D_{i,O}=\alpha _O^{}r_{i,O},\stackrel{~}{D}_{i,O}=\alpha _O^{}r_{i,O}+\gamma ,(i=1,2).$$
$`(A8)`$
APPENDIX B
GE Dipole Model and Rescattering series in $`st`$ space
As mentioned in the text, the monopole and the dipole model are useful to study various properties, such as convergence of the rescattering series expansion, together with the effect of generalizing the eikonalization since each rescattering term is tractable analytically. Here, we consider only the dipole case as an example.
The OE dipole model has been investigated in . The extension to the QE case is straightforward. We rewrite the GE amplitude as
$$A_{pp}^{\overline{p}p}(s,t)=2s_0^{\mathrm{}}H_{pp}^{\overline{p}p}(s,b)J_0(b\sqrt{t})b𝑑b,$$
$`(B1)`$
with
$$H_{pp}^{\overline{p}p}(s,b)=h_+\pm h_{}+H[PP]+H[OO]+2H[PO];$$
$`(B2)`$
the rescattering contributions $`H[PP,OO,PO]`$ are given in (23),(24). We split the Born contribution and the rescattering series of the GE dipole model (with 3 $`\lambda `$’s) which runs over the two indexes $`n_\pm `$ from 0 to infinity
$$A_{pp,GE}^{\overline{p}p}(s,t)=a_{pp}^{\overline{p}p}(s,t)+\underset{n_+=0}{\overset{\mathrm{}}{}}\underset{n_{}=0}{\overset{\mathrm{}}{}}a_{pp;n_+,n_{}}^{\overline{p}p}(s,t).$$
$`(B3)`$
Introducing the 4 partial contributions of the eikonal function $`\chi (s,b)`$
$$h_+=\frac{1}{2}\left(\chi _P(s,b)+\chi _f(s,b)\right)h_{}=\frac{1}{2}\left(\chi _O(s,b)+\chi _\omega (s,b)\right),$$
$`(B4)`$
known analytically from Appendix A and separating the three contributions, we obtain in the GE dipole case
$$a_{pp;n_+,n_{}}^{\overline{p}p}(s,t)=is\frac{(\pm i)^{n_++n_{}}(\lambda _+)^{n_+}(\lambda _{})^n_{}}{(n_++n_{}+2)!}$$
$$\times (F_{n_+,n_{}}.I+F_{n_{},n_+}.II+G_{n_+,n_{}}.III),$$
$`(B5)`$
where we have introduced the hypergeometric functions $`{}_{2}{}^{}F_{1}^{}`$ (with the real argument $`z=\frac{\lambda _0^2}{\lambda _+\lambda _{}}`$)
$$F_{n_\pm ,n_{}}=z(n_\pm +1)._2F_1(1n_{},n_\pm ;2;z).(1\delta _{n_{},0})+\delta _{n_{},0},$$
$$G_{n_+,n_{}}=_2F_1(n_{},n_+;1;z).$$
In (B5) we have also defined the inverse F-B’s transforms
$$I=\lambda _+\underset{\mathrm{}=0}{\overset{n_++2}{}}\underset{m=0}{\overset{n_{}}{}}\left(\begin{array}{c}n_++2\\ \mathrm{}\end{array}\right)\left(\begin{array}{c}n_{}\\ m\end{array}\right)\mathrm{Int}_{n_++2\mathrm{},n_{}m,\mathrm{},m}(s,t),$$
$`(B6)`$
$$II=\lambda _{}\underset{\mathrm{}=0}{\overset{n_+}{}}\underset{m=0}{\overset{n_{}+2}{}}\left(\begin{array}{c}n_+\\ \mathrm{}\end{array}\right)\left(\begin{array}{c}n_{}+2\\ m\end{array}\right)\mathrm{Int}_{n_+\mathrm{},n_{}+2m,\mathrm{},m}(s,t),$$
$`(B7)`$
$$III=\pm 2\frac{\lambda _+\lambda _{}}{\lambda _0}\underset{\mathrm{}=0}{\overset{n_++1}{}}\underset{m=0}{\overset{n_{}+1}{}}\left(\begin{array}{c}n_++1\\ \mathrm{}\end{array}\right)\left(\begin{array}{c}n_{}+1\\ m\end{array}\right)\mathrm{Int}_{n_++1\mathrm{},n_{}+1m,\mathrm{},m}(s,t).$$
$`(B8)`$
Once again, in these expressions $`+()`$ corresponds to $`\overline{p}p`$ ($`pp`$); $`\left(\begin{array}{c}n\\ k\end{array}\right)`$ is the binomial cœfficient and Int $`(s,t)`$ is the following integral over the 4 components of the eikonal function
$$\mathrm{Int}_{\lambda ,\mu ,l,m}(s,t)=_0^{\mathrm{}}\chi _P^\lambda (s,b)\chi _O^\mu (s,b)\chi _f^l(s,b)\chi _\omega ^m(s,b)J_0(b\sqrt{t})b𝑑b.$$
$`(B9)`$
An analytic expression for this integral has been written in the case when the Odderon does not contain a killing factor at $`t=0`$. It is a straightforward exercise to derive the complete analytical form from (B9). |
warning/0001/hep-th0001038.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The goal of the present paper is to investigate non-perturbative dynamics in the pure Y-M theory which is the key to the calculation of the vacuum condensates.
To understand the actual dynamics and the role of the non-perturbative effects we should have an explicit form of the non-perturbative fields. As a rule the non-perturbative effects are associated with classical fields characterized by topological charges. To see how the non- perturbative fluctuations generate physical amplitudes we shall treat the non-perturbative part of the vacuum expectation value (v.e.v.) of the energy-momentum tensor $`\theta _{\mu \mu }`$ in Y-M theory in four dimensions. In the case of theories without dimensional parameters the v.e.v. of $`\theta _{\mu \mu }`$ gives the vacuum energy density $`\epsilon _{vac}`$ and the characteristic mass scale thus defining properties of the effective theory . It has been noticed that the connection of non-perturbative effects with classical fields is not always straightforward. Indeed, there is a set of models, such as $`O\left(N\right)`$ non-linear sigma models ($`N>3`$) which do not have topological solutions, but do have non-perturbative effects which have been obtained within perturbation theory . Besides, the physical vacuum has no topological charge $`\left(Q=0\right)`$ by definition but we believe its structure has a non-perturbative nature.
There are several methods of calculation of the non-perturbative contributions into v.e.v.’s of different correlations. A solution of the problem in Y-M theory have been suggested by ’t Hooft and was a subject of intense study. However, there are several problems arising within this approach when instanton background is considered in the dilute-gas approximation . One of them is that contributions to the integrals over the instanton size come from the region of large sizes where the initial approximation is no longer valid. Therefore, the exact magnitude of the non-perturbative effects remains unknown. Recently an alternative method has been proposed in which the non-perturbative dynamics has been investigated for the example of the two dimensional $`O\left(N\right)`$ non-linear sigma model in the large $`N`$ limit . In this model any instanton effects are absent. Here we shall use the basic ideas which have been developed in these papers. For this reason we recall the main features of the approach.
It was shown that fluctuations describing the vacuum properties of a theory are subject to the requirements for potential energy to be in the minimum and for conjugated canonical momenta to be zero. Such fluctuations are not operators but a c-number function. For similar reason the temporal component of the gauge field are c-number functions.
It is evident that constant fields may obey these conditions. In such a case the quantum fluctuations around the constant background describe the perturbative properties of the theory. If non-constant fields satisfy these conditions then the quantum fluctuations around the non-constant background describe the non-perturbative properties of the theory.
The regularization procedure is essential for calculation of vacuum condensates. The regularization by separating the quantum fields in different points is used.To this end the point of the space in which the regularization procedure has to be done should be replaced by a sphere $`S^2`$ having a small radius $`r`$ which is set to zero at the very end of the calculation. The quantum field is defined on the surface of the sphere. From dimensional consideration $`\epsilon _{vac}`$ has to be proportional to $`1/r^4`$, therefore, naive $`\epsilon _{vac}`$ goes to infinity when $`r`$ tends to zero. In real fact the non-perturbative value of $`\epsilon _{vac}`$ is finite owing to quantum effects.
The non-trivial classical fields, mentioned above, are characterized by the topological charge $`Q`$. The energy funtional for static non-trivial configuration is defined as $`E=4\pi Q/g^2`$, where $`g`$ is the coupling constant. Obviously that the physical vacuum has no topological charge and may be presented as a sum of classical configurations which contribute to the topological charge with different signs. However, there is an alternative possibility. If the physical space has two boundaries, the topological charge may be equal to zero for non-trivial classical configurations. Such situation emerges due to the suggested regularization method. The point is that introduction of the sphere of small radius gives one more boundary (another boundary is at large distance) and the topological charge is given by the difference of contributions from the boundaries and is equal to zero for the infinitely small radius say, for the radius which is inverse to the mass of the ultraviolet cutoff.
In the present paper we shall adop the basic ideas which are obtained in sigma models to investigate the non-perturbative dynamics and to calculate the vacuum energy density in Y-M theory. In sect.2 the non-perturbative structure of vacuum state is discussed in general form. We describe the classical fields and discuss their properties in sect.3. In sect.4 we calculate the vacuum energy density. It is necessary to stress that the vacuum state is not obtained by solving the Schodinger equation. We belive that the constructed vacuum state is one of a coherent type.
## 2 The non-perturbative structure of the vacuum state.
In the present section we suggest the method of construting of the non-perturbative vacuum state using perturbative vacuun state. Our basic idea can be better explained in terms of the quantum mechanics.
Let us assume that $`\mathrm{\Psi }\left(x\right)`$ is the ground state of some quantum system and $`T`$ is the translation operator. We know that the average value of the quantity $`x+x_0=T^1xT`$ is determined by the following expresion
$$\overline{x}=_{\mathrm{}}^{\mathrm{}}𝑑x\mathrm{\Psi }^{}\left(x\right)\left(x+x_0\right)\mathrm{\Psi }\left(x\right)=_{\mathrm{}}^{\mathrm{}}𝑑x\left(T\mathrm{\Psi }\right)^{}x\left(T\mathrm{\Psi }\right)=x_0.$$
This example shows we do not need to know, generally speaking, the explicit ground-state $`\mathrm{\Psi }\left(xx_0\right)`$ to obtain the average value of $`\overline{x}0`$. We can get it knowing only the translation operator and the parameter of the transformation $`x_0`$. Notice that $`\mathrm{\Psi }\left(x\right)`$ and $`\mathrm{\Psi }\left(xx_0\right)`$ do not satisfy the requirement of orthogonality. If the anologous situation takes place in quantum field theory then having the non-perturbative fluctuation and some translation operator one can construct the vacuum state and to calculate vacuum condensate. However, contrary to quantum mechanics one has to know some starting ground state in field theory because of v.e.v. is obtained by another mean.It turns out that perturbative ground state may be selected as the starting state. This situation facilitates the problem.
It was shown in the framework of the sigma models in two dimensions that this idea leads to the correct result for the vacuum condensate of $`\theta _{\mu \mu }`$ . Therefore we are about to discuss how the idea is realized in quantum field theory.
Let us show how the non-perturbative state is constructed. At the beginning the non-perturbative vacuum fields fluctuations are defined and then the translation operator and the non-perturbative vacuum state is constructed.
Assume that $`L`$ is the Lagrangian of some field system and $`h`$ is the Hamiltonian density. There is a relation
$$L=\dot{\varphi }\pi _0h,$$
where $`\varphi `$ is the field and $`\pi _\mu =\delta L/\delta _\mu \varphi `$ is the canonical momentum. Here group indexes are omitted. The Hamiltonian is defined as
$$H=𝑑x^{n1}\left(\pi _\mu ^2+V\left(\varphi \right)\right).$$
Let us assume that the potential gets the minimum $`V\left(\varphi _v\right)=0`$ when $`\varphi _v^2\left(x\right)=const`$. The finite energy condition is satisfied if $`\pi _\mu =0,V\left(\varphi \right)=0`$ everywhere at large distances. The solutions satisfying this condition are a set of the vacuum fields because we believe that fields take their vacuum value at large distances in theory without sources. It is believed that vacuum is costructed in the same way at any point of the physical space. Therefore the condition defines the vacuum field in any point of the physical space except a singular point. The vacuum field can not be quantizated due to the fact that their conjugated canonical momenta are equal to zero $`\pi _0=0`$ and $`\varphi _v\left(x\right)`$ is c-number function.
If the canonical momentum is $`\pi _\mu =_\mu \varphi `$, then the solutions of the condition are trivial $`\varphi _v=const`$. In this case the theory is quantized around constant background and the quantum fluctuations describe the perturbative properties of the theory .
If the canonical momentum involves a gauge field $`A_\mu `$,i.e. $`\pi _\mu =_\mu \varphi +gA_\mu \varphi `$ ,then there are solutions of the condition $`\pi _\mu =o`$ which are topologically trivial $`A_\mu =0,\varphi _0=const`$ and topologically non-trivial $`A_\mu =1/gU_\mu U^1,\varphi _v=U\varphi _0`$ at large distances. U(x) is an arbitrary element of a gauge group and must be non-singular only on a bondary. The fields $`\varphi _v\left(x\right)`$ are c-number functions only if we quantize the theory over the non-constant background $`\pi _0=_o\varphi _v+U_\mu U^1\varphi _v=0`$. In the perturbation theory, when $`A_\mu =0`$ on the boundary and non-singular in all physical space, the canonical momentum $`\pi _\mu =_\mu \varphi 0`$ therefore the fields $`\varphi _v\left(x\right)`$ are quantum fields which describe the non-perturbative properties of the theory as will be argued in the paper. As discussed in this is the reason why some results which are obtained in the perturbative methods are related to the non-perturbative effects.
Let us suppose we quantize our theory over the constant background. If the constant is equal to zero then the vacuum of the theory may be described in a usual way as the eigenfunction of the annihilation operator $`a_0`$ with zero eigenvalue of the momentum $`a_0|0>=0`$ and v.e.v. of $`\varphi \left(x\right)`$ is $`<0\left|\varphi \left(x\right)\right|0>=0`$. If one wants to have $`<0\left|\varphi \left(x\right)\right|0>0`$ then a new vacuum state have to be constructed. The operator which carries out the transformation $`\varphi v`$ is well known to be the operator of the canonical momentum. The translation operator is
$$T=\mathrm{exp}v𝑑x_0\varphi =\mathrm{exp}v\left(a_0a_0^+\right)$$
and the new vacuum state is written as $`|v>=T|0>`$. In this case the annihilation operator $`a_0`$ has no zero eigenvalue $`a_0|v>=v|v>`$ and v.e.v. of $`\varphi \left(x\right)`$ is $`<v\left|\varphi \left(x\right)\right|v>=v`$ .
If the vacuum fluctuation is a function then the translation operator can be written as
$$T=\mathrm{exp}𝑑x^{n1}\left(\varphi _v\left(x\right)\pi _0+_0\varphi _v\left(x\right)\varphi \left(x\right)\right).$$
Here the last term in the integral generates the translation of the canonical momentum into $`_0\varphi _v\left(x\right)`$.
Considering v.e.v. of some operator in the non-perturbative vacuum state $`|\varphi _v>=T|0>`$ we obtain that all fields operators are substituted on the vacuum fluctuation and then v.e.v. have to be calculated using the same method of regularization as is given below. Notice that the vacuum state $`T|0>`$ is the total vacuum state which contains perturbative part $`|0>`$ and non-perturbative part $`\left(T1\right)|0>`$ of the vacuum state therefore there is no orthogonality between the perturbative and the total vacuum states.
Here we recall that if there is a symmetry in the theory then the symmetry generator $`Q`$ annihilates the vacuum state $`Q|0>=0`$. The condition of the symmetry breaking is the absence of the symmetry invariance of the vacuum state $`Q|0>0`$. In is necessary to stress that the symmetry generator is common one but the vacuum states are different states. If we should like to obtain the vacuum energy changes then Hamiltonian have to be precisely the same in the vacuum states.
The anologous situation takes place in this case. By choosing as the starting state the vacuum state of the perturbative theory we have to use free Hamiltonian of the theory and non-perturbative vacuum state to obtain vacuum energy changes.
However there is a problem arising within this approach in Y-M theory. The point is that free Hamiltonian of the gauge fields is not invariant under SU(2) local gauge transformation corresponding to Y-M theory. Therefore the vacuum energy which is obtained by this the method is gauge dependent.
It is known that the vacuum energy may be decomposed into non-perturbative and perturbative parts. The pertubative part is scale dependent. The non-perturbative part is scale independent. Besids the non-perturbative part of the vacuum energy is defined on the non-trivial topological fields because the physical local gauge transformations which are trivial from the point of view of topology can not break down the non-perturbative results. The gauge transformations change only the perturbative part of the vacuum energy which should be subtracted because it has no physical meaning.
Notice that the vacuum energy density is usually obtained as v.e.v. of the trace of the energy-momentum tensor which have to be defined in gauge non-invariant way according to the above arguments.
The translation operator in Y-M theory is definefed as
$$T\mathrm{exp}𝑑x^3J_i^aE_i^a$$
in the case of the static vacuum cofigurations. Using this vacuum state we can obtain v.e.v. of $`\theta _{\mu \mu }`$ by the operator method.However for pedagogical purpose the calculation of the vacuum energy density in Y-M theory is carried out using functional integral.
## 3 Classical vacuum in Yang-Mills theory
The structure of the vacuum in the pure Y-M theory was well studied. An extensive reference list can be find in . Here we are interested in static properties of the vacuum. At the begining we will remind the formulation given in paper . To set out notations we briefly recapitulate some results relevant to our work.
The Hamiltonian of the classical Y-M theory in Euclidian space is
$$H=\frac{1}{2}d^3x\left(E_k^a\right)^2+\left(B_k^a\right)^2,$$
(1)
where $`E_k^a=F_{k4}^a`$, $`B_k^a=1/2\epsilon _{klm}F_{lm}^a`$, $`F_{ij}^a=_iA_j^a_jA_i^a+g\epsilon ^{abc}A_i^aA_j^c`$; $`a,b,c`$ are the group indices; $`k,l,m=1,2,3`$. Let the gauge group be $`SU\left(2\right)`$ and $`\left(t^a\right)^{bc}=\epsilon ^{abc}`$ be generators of the group $`SU\left(2\right)`$ in the adjoint representation. In this case the fields $`A_\mu ^a`$ are real. We also imply that the gauge fields are subject to the first class constraint and there is the gauge fixing condition $`_\mu A_\mu =0`$.
The gauge fields which describe the classical properties of the vacuum must satisfy the condition $`H=0`$ in the static case. This, in its turn, yields
$$F_{4k}^a=F_{lm}^a=0,$$
(2)
i.e. the canonical momenta are equal to zero. It follows that the fields which describe the classical vacuum of the theory are pure gauge.
The space components of the potential $`A`$ are
$$A_i=\frac{1}{g}U\left(x\right)_iU^T\left(x\right),$$
(3)
where $`U\left(x\right)`$ is an arbitrary element of the gauge group which is taken independent of $`x_4`$ and $`U\left(x\right)U\left(x\right)^T=1`$. Here it is necessary to stress that the gauge field (3) may be eliminated by a gauge transformation only if $`U\left(x\right)`$ goes to zero at large distances. However we shall also deal with non-trivial topological configurations therefore the function $`U\left(x\right)`$ must be non-singular on the boundary and corresponds to topological charge. As has already been mentioned, it is possible that the total topological charge and energy functional are equal to zero even when there are non-trivial topological fields. Such possibility is discussed in the next section. In this case $`A_i`$ (3) can not be eliminated by gauge transformations .
The components of the gauge field $`A_4`$ are usually chosen equal to zero but, as shown in , $`A_40`$ defines the topological charge of the monopole and, as we shall see later, other non-perturbative properties of the theory in the static case.
The static configurations $`A_4^a`$ can be obtained from the condition (2)
$$F_{4k}^a=_kA_4^a+g\epsilon ^{abc}A_k^bA_4^c=0.$$
(4)
Multiplying this equation by $`A_4^a`$ and taking into account antisymmetricity of the tensor $`\epsilon ^{abc}`$, we get
$$_k\left(A_4^a\right)^2=0,i.e.\left(A_4^a\right)^2=const.$$
(5)
The constant in (5) is arbitrary but it is important that it may be chosen non-zero . Although the classical theory has no dimensional parameter we can introduce arbitrary quantity $`\mu `$ due to the condition (5) and define
$$A_4^a\left(x\right)=\mu \mathrm{\Phi }^a\left(x\right),$$
(6)
where $`\mathrm{\Phi }^a`$ is a dimensionless function which is subject to the condition
$$\mathrm{\Phi }^a\left(x\right)\mathrm{\Phi }^a\left(x\right)=1.$$
(7)
It follows from (7) that the fields $`\mathrm{\Phi }^a\left(x\right)`$ can be represented as
$$\mathrm{\Phi }\left(x\right)=U\left(x\right)\mathrm{\Phi }_0,$$
(8)
where $`\mathrm{\Phi }_0`$ is a constant vector in the colour space. As follows from eq.(8) the shift of $`\varphi `$ in $`x`$-space is equivalent to the rotation of $`\varphi _0`$ in the colour space.
The vector potential $`A_k^a`$ may be expressed by a combination of the fields $`\mathrm{\Phi }^a\left(x\right)`$. Multiplying eq.(4) by $`\epsilon ^{ial}\mathrm{\Phi }^l`$ we obtain
$$J_k^i=g\left(\delta ^{in}\mathrm{\Phi }^i\mathrm{\Phi }^n\right)A_k^n,$$
(9)
where $`J_k^i=\epsilon ^{ial}_k\mathrm{\Phi }^a\mathrm{\Phi }^l`$. Introducing $`P^{in}=\delta ^{in}\mathrm{\Phi }^i\mathrm{\Phi }^n`$ then $`P^{in}\mathrm{\Phi }^n=0`$ and $`P^{in}P^{nl}=P^{il}`$. Therefore multiplying (9) by $`P^{if}`$ we get
$$P^{if}\left(J_k^fgA_k^f\right)=0.$$
(10)
The solution of eq. (10) is
$$A_k^a=\frac{1}{g}\left(J_k^an_k\mathrm{\Phi }^a\right),$$
(11)
where $`n_k\left(x\right)`$ is an arbitrary function. The solution (11) also satisfies the condition $`F_{ij}=0`$ (2) if $`n_k\left(x\right)`$ is a pure gauge in the group $`U\left(1\right)`$, i.e. $`_in_j_jn_i=0`$. The quantity $`n_i`$ defines the projection of $`A_k^a`$ onto $`\mathrm{\Phi }^a\left(x\right)`$. Indeed, from (11) we get
$$n_k=gA_k^a\mathrm{\Phi }^a\left(x\right)$$
(12)
$`J_i^a`$ is a conserved current $`\left(_iJ_i^a=0\right)`$ associated with the global symmetry $`SO\left(3\right)`$ of the theory. The conservation of the current follows from the relation
$$^2\mathrm{\Phi }^a=\mathrm{\Phi }^a\left(_i\mathrm{\Phi }^a\right)^2,$$
which can be obtained from eq.(7).
The gauge fixing condition $`_\mu A_\mu =_kA_k=0`$ may be satisfied by requiring the orthogonality of $`\mathrm{\Phi }^a`$ and $`A_k^a`$ or $`n_k=0`$, i.e. $`n_k`$ is a gauge-fixing parameter. In such a case the gauge potential (11) is not fixed by the Coulomb gauge. The property of the Coulomb gauge was first discussed by V.N. Gribov .
Let us find how the fields $`\mathrm{\Phi }^a\left(x\right)`$ are expressed in terms of the current $`J_k^a`$. To do this we may use eq. (4). Substituting (6) and(11) into (4) we obtain
$$_k\mathrm{\Phi }^a\left(x\right)+\epsilon ^{abc}J_k^b\left(x\right)\mathrm{\Phi }^c\left(x\right)=0.$$
(13)
The solution of eq.(13) can be written as
$$\mathrm{\Phi }\left(x\right)=P\mathrm{exp}\left(\underset{x_0}{\overset{x}{}}𝑑z_iJ_i\right)\mathrm{\Phi }\left(x_0\right).$$
(14)
where $`J_i=J_i^at^a`$ and $`x_0`$ is an arbitrary point. As it is clear from (14) and (7) the values of the field $`\mathrm{\Phi }\left(x\right)`$ in two different fixed points $`x_1`$ and $`x_2`$ are related to each other by the global group transformation and any invariant of the global transformation is independent of choosing the point $`x_0`$. One can see from (3) and(11) that the current $`J_k`$ is definded as
$$J_k=U_kU^1,$$
(15)
when $`n_k=0`$. Therefore the integral (14) does not depend on the integration path if the function $`U\left(x\right)`$ has no singularities. However to have non-trivial topological effects we should also admit the occurrence of gauge functions which have a singularity. Then the function (14) is a pure gauge at any region which does not contain the singular points. However it is ambiguous in the whole space therefore globally the fields $`J_k`$ are not a pure gauge. We recall that the integral $`𝑑z_i\widehat{J}_i`$ taken along a closed contour encircling the singularity is not equal to zero, i.e. depends on the integration contour.
## 4 On the calculation of the vacuum energy density.
The vacuum energy density is defined as
$$\epsilon _{vac}=\frac{1}{4}<0\left|\theta _{\mu \mu }\right|0>,$$
(16)
where $`\theta _{\mu \mu }`$ is the trace of the energy-momentum tensor, the numerical coefficient is determined by the dimension of the physical space $`D=4`$.The vacuum energy density is defined as a sum of the perturbative and non-perturbative contributions. Since we are interested in the non-perturbative effects we will now assume that the perturbative contributions into the vacuum energy density are subtracted from $`\epsilon _{vac}`$ in eq. (16). The non-perturbative part of the vacuum energe density is denoted by $`\epsilon _{vac}^n`$.
It is well known that the dilatation anomaly in gauge theories is proportional to the $`\beta `$-functon and the field strength squared but the scale independent value of the dilatation anomaly is only $`g^2\left(F_{\mu \nu }^a\right)^2`$. Therefore only this part $`\epsilon _{vac}^n`$ has a physical significance. For this reason we shall calculate the vacuum condensate of $`g^2\left(F_{\mu \nu }^a\right)^2`$. Notice that this method seems to avoid the computation of the dilatation anomaly condensate if renormalization constants are known at all orders.
Let the Lagrangian be defined in terms of the renormalisation gauge fields and constants as
$$L=\frac{1}{4}\left(F_{\mu \nu }^a\right)^2\frac{1}{2\alpha }\left(_\mu A_\mu ^a\right)^2\frac{1}{4}\left(Z_31\right)\left(_\mu A_\nu ^a_\nu A_\mu ^a\right)^2.$$
(17)
The ghost are omitted because they are essential only in perturbative calculations. The gauge fields, the gauge fixing parameter and the coupling constant are given by
$$A_b=Z_3^{1/2}A,\alpha _b=Z_3\alpha ,g_b=Z_gg,$$
where $`A_b`$ ,$`g_b`$ ,$`\alpha _b`$ are bare quantities. The quantities $`Z_g`$ and $`Z_3`$ are calculated at the one-loop level. In such a case we have
$$Z_3=1+\frac{g_b^2}{16\pi ^2}\left(\frac{13}{6}N\frac{\alpha _b}{2}\right)\mathrm{ln}M^2/\mu ^2,Z_g=1+\frac{g_b^2}{16\pi ^2}\frac{13}{6}N\mathrm{ln}M^2/\mu ^2,$$
(18)
where $`\mu ^2`$ is a normalization point, $`M^2`$ is an ultraviolet cutoff.
The last term in eq.(17) is a counterterm which is not gauge invariant. As above gauge invariance of the counterterm is not required.
It is known that the conformal anomaly $`\theta _{\mu \mu }`$ can be obtained as the variation of action when the ultraviolet cutoff $`M^2`$ changes into $`\left(1+\eta \right)M^2`$ with the coupling constant kept fixed. Here $`\eta `$ is the parameter of the global scaling transformation. Then we can obtain
$$\theta _{\mu \mu }=\frac{g_b^2}{64\pi ^2}b\left(_\mu A_\nu ^a_\nu A_\mu ^a\right)^2.$$
(19)
Here $`b=13/6N\alpha _b/2`$.
In terms of the generating functional $`Z_E`$ in Euclidian space, v.e.v. of $`\theta \mu \mu `$ can be written as
$$<0\left|\theta _{\mu \mu }\right|0>=Z_E^1DA_i\theta _{\mu \mu }^Rexp^{S_E}$$
(20)
Here the quantity $`\theta _{\mu \mu }^R`$ have been obtained from eq.(19) by introducing the regularization.
Let us calculate the vacuum energy density at some point $`x_0`$. Let us assume that $`x_0`$ is enclosed by a sphere $`S^2`$ of a small radius $`r`$. To regularize the quantity $`\theta _{\mu \mu }`$ we define the quantum fields in different points on the surface of the sphere $`S^2`$. Let us put the radius to be inverse of the parameter $`\mu `$ from eq.(6).
Notice that making so we spoil the gauge symmetry. It should be restored if the operator is the gauge invariant quantity. However,since the operator $`\theta _{\mu \mu }`$ eq.(19) is not gauge invariant and from the very beginning, we may not take care of the gauge symmetry at all. Also it should be kept in mind that the quantity $`A_4^a`$ are c-number fields and therefore the quantity $`\left(_kA_4^a\right)^2`$ do not have to be regularized. As will be shown later therefore the contributions of the electrical strengh in the non-pertrubative vacuum energy are absent for the static fields.
According to V.N. Gribov , the contributions of the configurations $`J_i^a`$ can not be compensated by the ghost fields because the Faddeev-Popov determinat has zero and the standart method has to be improved. In the static case the action has a minimum on this configurations as pure gauge configurations in all physical space besides the point of the singularity. Here we keep in mind that contributions to the integral which come from the sphere rigion are ignored,but the contributions from the surface of the sphere are.
At the begining one may speculate that the singular point is $`x_0`$ where the operator $`\theta _{\mu \mu }`$ is defined. We recall that the point $`x_0`$ is surrounded by the sphere of a small radius and area bounded by the sphere surface is excluded from consideration. In such case the topological charge equals to zero. Really, the topological charge is
$$Q=_S_{\mathrm{}}𝑑\sigma _kj_k_{s_r}𝑑\sigma _kj_k=0,$$
(21)
where $`j_k=\epsilon _{kij}_iA_j^a\mathrm{\Phi }^a`$, $`S_{\mathrm{}}`$ and $`s_r`$ are the respective surfaces of the sphere at large and small distances. The result (21) is valid for any infinitely small quantity $`r`$, which has to be considered as a ”physical” zero. The result (21) is explained by the absence of the singularity in the region between two sphere surfaces where the fluctuations are defined.
At first glance it would seem that we discussed a special case because the singular point coincides with the point in which the operator $`\theta _{\mu \mu }`$ is examined. Indeed, this is not the case. As the function $`\mathrm{\Phi }^a\left(x\right)`$ are subject to constraint (7) which is independent of a position of the singular point and the quantity $`\left(_i\mathrm{\Phi }^a\right)^2`$ is a translational invariant function therefore we can give any position to the singularity for any combination being quadratic in $`\mathrm{\Phi }^a`$ or $`_i\mathrm{\Phi }^a`$. As we shall see, the v.e.v. of $`\theta _{\mu \mu }`$ may be written as one of such combinations due to $`\theta _{\mu \mu }`$ is independent of the position of the singularity.
We recall that the action has the minimum on the fields $`A_i^a=\frac{1}{g}J_i^a`$ therefore v.e.v. of $`\theta _{\mu \mu }`$ can be written as follows
$$\begin{array}{c}<0\left|\theta _{\mu \mu }\right|0>=\underset{\mathrm{\Delta }0}{lim}g_b^2b/64\pi ^2[2\mu ^2\left(_i\mathrm{\Phi }^a\right)^21/g_b^2(_iJ_j^a(x_0\mathrm{\Delta })\\ _jJ_i^a(x_0\mathrm{\Delta }))(_iJ_j^a(x_0+\mathrm{\Delta })_jJ_i^a(x_0+\mathrm{\Delta }))].\end{array}$$
(22)
Here, the first term is obtained from the term $`\left(_iA_4\right)^2=\mu ^2\left(_i\mathrm{\Phi }\right)^2`$. This term is the singular one. It can be verified that the term tends to infinity as $`1/r^2`$, when $`r`$ goes to zero. To do this one needs to use parametrization of the function $`\mathrm{\Phi }^a\left(x\right)=\left(xx_0\right)^a/r`$.The quantity $`1/r^2`$ may be considered to be of the order of the ultraviolet cuttoff squared for a small $`r^2`$. Therefore the first term gives the contribution to the perturbative part of the vacuum energy density . Since we are interested in non-perturbative part of the vacuum energy density then we shall not discuss the first term.We recall that we assumed that the perturbative conributions into v.e.v. of $`\theta _{\mu \mu }`$ are subtracted.
The non-perturbative part of the vacuum energy is associated with the second term in (22) and can be written as
$$\begin{array}{c}<0\left|\theta _{\mu \mu }\right|0>=lim_{\mathrm{\Delta }0}b/64\pi ^2\left(_iJ_j^a\left(x_0\mathrm{\Delta }\right)_jJ_i^a\left(x_0\mathrm{\Delta }\right)\right)\\ \left(_iJ_j^a\left(x_0+\mathrm{\Delta }\right)_jJ_i^a\left(x_0+\mathrm{\Delta }\right)\right).\end{array}$$
(23)
It is convinient to rewrite eq.(23) in another form, making use the eq.(2)
$$F_{ij}=_iJ_j^a_jJ_i^a+\epsilon ^{abc}J_i^bJ_j^c=0,$$
and we arrive at
$$\begin{array}{c}<0\left|\theta _{\mu \mu }\right|0>=\\ lim_{\mathrm{\Delta }0}b/64\pi ^2\{\left(_i\mathrm{\Phi }^a(x_0\mathrm{\Delta })_i\mathrm{\Phi }^a\left(x_+\mathrm{\Delta }\right)\right)^2\left(\mathrm{\Phi }^b(x_0+\mathrm{\Delta })\mathrm{\Phi }^b(x_0\mathrm{\Delta })\right)^2\\ \left(_i\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)_j\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)\right)\\ \left(_i\mathrm{\Phi }^b(x_0+\mathrm{\Delta })_j\mathrm{\Phi }^b(x_0+\mathrm{\Delta })\right)\left(\mathrm{\Phi }^c(x_0\mathrm{\Delta })\mathrm{\Phi }^c(x_0+\mathrm{\Delta })\right)^2\}.\end{array}$$
(24)
Here the terms of the type $`_i\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)`$ were omitted because they tend to zero when $`\mathrm{\Delta }`$ goes to zero. The quantity $`_i\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)_j\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)`$ is approximately equal to $`\delta _{ij}/3_k\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)_k\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)`$ and due to that we get
$$\epsilon _{vac}^n=\underset{\mathrm{\Delta }0}{lim}b/384\pi ^2\left(_i\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)_i\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)\right)^2\left(\mathrm{\Phi }^b\left(x_0\delta \right)\mathrm{\Phi }^b\left(x_0+\mathrm{\Delta }\right)\right)$$
(25)
Now it should be shown that the quantity $`\epsilon _{vac}`$ is scale independent constant. We can do it in two steps. As the fields $`\mathrm{\Phi }^a`$ have been defined on the surface of the sphere $`S^2`$ then at the first step we can place the fields $`\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)`$ and $`\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)`$ in one and the same point on the surface keeping the radius of the sphere constant. At the second step the radius tends to zero.It is convenient to introduce new variables $`\mathrm{\Delta }_i`$ and $`r^2=\mathrm{\Delta }_i^2`$ holding $`x_0`$ fixed. Making use of the parametrization of the function $`\mathrm{\Phi }^a=\gamma \mathrm{\Delta }^a/r`$ we have
$$\left(_i\mathrm{\Phi }^a\right)^2=\frac{2}{r^2}\gamma ^2.$$
(26)
Here $`\gamma `$ is the normalization constant which is given below.
To calculate the quantity $`\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)`$ one can use eq.(14) and obtain
$$\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)=\mathrm{\Phi }^b\left(x_0\right)\mathrm{\Phi }^c\left(x_0\right)\left(P\mathrm{exp}\left(\underset{x_0\mathrm{\Delta }}{\overset{x_0+\mathrm{\Delta }}{}}𝑑z_iJ_i\right)\right)_{bc}.$$
(27)
We can write eq.(27) in a more suitable form. To this end one expands of exponential function in eq.(27) in powers of small $`\mathrm{\Delta }`$ to second order and gets that the first oder term is equal to zero due to antisymmetric tensor $`\epsilon _{abc}`$ and the second order term contains $`\mathrm{\Phi }^b\mathrm{\Phi }^c`$ times $`\mathrm{\Phi }^b\mathrm{\Phi }^c`$ and thus one gets $`\left(\mathrm{\Phi }^2\right)^2`$ =1. Therefore eq.(27) can be written as follows
$$\mathrm{\Phi }^a\left(x_0\mathrm{\Delta }\right)\mathrm{\Phi }^a\left(x_0+\mathrm{\Delta }\right)=\left(P\mathrm{exp}\left(\underset{x_0\mathrm{\Delta }}{\overset{x_0+\mathrm{\Delta }}{}}𝑑z_iJ_i\right)\right)_{aa}.$$
(28)
The contour of integration should be closed around the sphere with the fixed radius. Then substituting (26) and (28) into (25) we have
$$\epsilon _{vac}^n=\underset{r0}{lim}b/96\pi ^2\left(\frac{\gamma ^2}{r^2}P\mathrm{exp}\left(𝑑z_iJ_i\right)\right)^2$$
(29)
To calculate the integral around the circle use is made of the identy $`\left(J_i^at^a\right)^{bc}=\mathrm{\Phi }^b_i\mathrm{\Phi }^c_i\mathrm{\Phi }^b\mathrm{\Phi }^c`$ which yeilds
$$𝑑z_iJ_i^{bc}=𝑑\mathrm{\Phi }^c\mathrm{\Phi }^bd\mathrm{\Phi }^b\mathrm{\Phi }^c=2S^{bc}$$
(30)
The area $`S^{bc}`$ is enclosed by the circle.The third axis may be oriented normally to the area $`S^{bc}`$ in the colour space.The field $`\mathrm{\Phi }^3`$ is fixed in such a way that fields $`\mathrm{\Phi }^1`$ and $`\mathrm{\Phi }^2`$ are normalized as $`\left(\mathrm{\Phi }^i\right)^2=\lambda ,i=1,2`$. In this case the area $`S^{12}`$ equals $`\pi \lambda `$ and the integral around the circle is $`2\pi \lambda `$. For $`\epsilon _{vac}^n`$ we have
$$\epsilon _{vac}^n=b/96\pi ^2\left(\frac{\gamma ^2}{r^2}\mathrm{exp}\left(2\pi \lambda \right)\right)^2.$$
(31)
Here $`\epsilon _{vac}`$ is obtained on the scaling mass $`\mu ^2=1/r^2`$. Now we should get the vacuum energy density corresponding to $`r^2=1/M^2`$. Since the fields were defined on the scaling mass $`\mu ^2`$ they can be rewritten in terms of the bare fields. Then we have $`J_i=Z_3^{1/2}J_i`$ and $`\mathrm{\Phi }_b=Z^{1/4}\mathrm{\Phi }_b`$. In this case the relation between $`\lambda `$ and unrenormalized constant $`\lambda _b`$ is written as $`\lambda =Z^{1/2}\lambda _b`$. Besides we have to define a relation $`<0\left|\theta _{\mu \mu }\right|0>=Z_3^1<0\left|\theta _{\mu \mu }\right|0>_b`$. Clearly the constant $`\gamma ^2`$ in eq.(26) is $`Z_3^{1/2}`$, and we get
$$\epsilon _{vac}^n=b/96\pi ^2\left(\mu ^2\mathrm{exp}\left(2\pi \lambda _bZ_3^{1/2}\left(\mu ^2\right)\right)\right).$$
(32)
The value $`\epsilon _{vac}`$ is scale independent only if
$$\mu ^2\frac{d\epsilon _{vac}^n}{d\mu ^2}=0.$$
Using (32) and (18) we obtain that if
$$\lambda _b=8\pi /\left(g_b^2b\right),$$
then $`\epsilon _{vac}^n`$ is scale independent. Now we obtain the final result
$$\epsilon _{vac}^n=b/96\pi ^2\left(M^2\mathrm{exp}\left(16\pi ^2/bg_b^2\right)\right)^2.$$
(33)
The quantity $`\epsilon _{vac}^n`$ is scale independent and rises as $`N`$ which is agremeant with the familiar result .
Notice that $`\lambda _b`$ is a large quantity ($`\lambda 1`$) due to $`g_b^2\left(M^2\right)/4\pi 1`$,therefore $`\mathrm{\Phi }_3^2=1\lambda _b`$ may be negative quantity and this, in its turn, implies that the condition $`\left(\mathrm{\Phi }^a\right)^2=1`$ is not the equation of a sphere in view of quantum effects, i.e. the monopoles are absent in the scale independent quantum theory.
## 5 Conclusions.
In the present paper we have calculated the non-perturbative part of the vacuum energy density in pure Y-M theory which is determined by the static vacuum fluctuations in one-loop level. The result was obtained by using the physical ideas which are derived in the special case of sigma models .
It was shown that the vacuum energy density in Y-M theory is the scale independent quantity, which corresponds to the topological charge equal to zero.
The method of regularization is proposed which is based on the critical assumption that at each point of the space in which a product of two operators can be replaced by a sphere $`S^2`$ having a small radius which is set to zero at the very end of the calculation and the quantum fields are separated in diffirent point on the surface of the sphere. The calculation have been carried out using functional integral but can be done by the operator method. For this aim the vacuum state has been constructed.
Acknowledements.
The author wishes to thank V.A. Novikov, A.D. Mironov, Yu.M. Makeenko for useful discussions. The paper was partially supported by the Russian Foundation of Fundamental Reseach (grant No 98-02-17316). |
warning/0001/quant-ph0001021.html | ar5iv | text | # Untitled Document
Hidden assumptions in decoherence theory
Italo Vecchi
Bahnhofstr. 33 - 8600 Duebendorf - Switzerland
email: vecchi@weirdtech.com
The present note is devoted to a critique of some aspects of current decoherence theory. Its main point is that several common claims related to decoherence theory are based on questionable hidden assumptions. We refer to and for excellent introductions to the subject of decoherence theory and for background material.We focus on Joos excellent survey of decoherence theory , whose clarity makes its relatively easy to spot the inconsistencies in the argument. According to the argument, given a system $`S`$ in a superposition of eigenstates $`|n`$ and its environment $`W`$ in a state $`\mathrm{\Phi }_o`$ the pointer states are identified as those states $`|\mathrm{\Psi }(t)`$ in $`W`$ resulting from the interaction between $`S`$ and $`W`$
$$|n|\mathrm{\Phi }_oexp(iH_{int})|n|\mathrm{\Phi }_o=:|n|\mathrm{\Phi }_n(t).$$
The states $`|\mathrm{\Phi }_n(t)`$ result from the entanglement of W with S through the interaction Hamiltonian $`H_{int}`$ and are usually referred to as the ”pointer positions”. An act of measurement on $`W`$ induces a collapse of its state vector into one of the pointer vector s, yielding information about the state of the system $`S`$. The states $`|\mathrm{\Phi }_n(t)`$ are descibed in as the states of the ”rest of the world”. The basic ambiguity underlying this description of the decoherence process may be formulated as follows. Any vector basis can be chosen as a pointer basis. The environment and any measurement device can be described using an arbitrarily chosen basis $`|\mathrm{\Psi }(t)`$. The privileged pointer basis referred to by Joos is relative to an observer, as defined by a measurement operator. The measurement device or the environment do not chose a basis. The observer does. The privileged pointer basis is determined by the set of possible outcomes of a measurement act performed by an observer. It is the intervention of the observer on the measurement apparatus in the course of the measurement process that determines the pointer basis.An example may clarify the underlying issue. We know that Planck’s radiation law in black body theory is obtained maximising entropy on discrete energy spectra. In the black body model both absorption and emission are continuous processes, but the entropy is maximised on discrete energy spectra. Entropy maximisation may be applied to other sets of observables too, but it will yield different results. If the observer is associated with continuous energy spectra then entropy maximisation yields the Jeans-Raleigh law. Other observables yield other distribution laws. This extends to decoherence, so that the result of the decoherence process is seen to depend on the observer, as defined by a set of observables or, equivalently, by a measurement operator.The role of the observer in the decoherence argument is indeed acknowledged in , as is the fact that the superpositions in the system are not destroyed but merely cease to be identifiable by local observers. However the pointer basis is implicitly treated as an intrinsic property of the interaction between the system and its environment or a measurement device. This tacit assumption is necessary for the decay of the off-diagonal interference terms of the system’s density matrix,
$$\rho _S=\underset{n,m}{}c_m^{}c_n|mn|\rho _S=\underset{n,m}{}c_m^{}c_n\mathrm{\Phi }_m|\mathrm{\Phi }_n|mn|$$
which is then interpreted as the vanishing of superpositions. The assumption however leads to inconsistencies, as shown by the following analysis. The assumption that the pointer basis is an intrinsic property of the environment would not matter if the decoherence argument was independent of the chosen pointer basis. However this is not the case. According to the argument in and , the decoherence process induces the decay of the off-diagonal elements of the systems density matrix,
$$\rho _S\underset{n}{}|c_n|^2|nn|$$
which is interpreted as the emergence of a set of stable macroscopic states. The density matrix however is defined in terms of the pointer basis. Different pointer basis lead to different density matrices for the same state vectors. It is immediate to see that the decoherence process, i.e. the decay of the off diagonal terms in the density matrix, does not commute with a change of basis. Indeed given a density matrix $`A`$ , let $`C`$ be a change of basis and , $`C^1`$ its inverse and D the operator that equates to null the off-diagonal elements. Then
$$DA(C^1DC)A$$
so that the result of the decoherence process depends on the pointer basis, which is selected by the observer and is independent of the underlying physical process. Indeed any two non-commuting operators induce pointer basis for which the above inequality holds. The states associated with a diagonal density matrix in one basis describe superpositions in the other basis. An example of different pointer basis inducing different decoherence processes is actually considered in , but the authors limits themselves to pointing out the ”right” pointer basis, without analysing its dependence on the observer. The above indicates that the result of the decoherence process depends on the observer, but it also reveals that there must be a flaw in the decoherence argument. The claim that interaction with the environment induces the diagonalisation of the system’s density matrix must be wrong, since the diagonalisation process depends on the chosen basis, which is not an intrinsic property of the environment but of the observer. Indeed if one examines the argument leading to the diagonalisation of the system’s density matrix, one discovers that it is based on the unphysical no-recoil assumption on the scattering process (), which serves the sole purpose of preventing the environment from eroding the diagonal elements of the system’s density matrix. Under the no-recoil assumption interaction with the environment action can only deplete the off-diagonal elements of the system’s density matrix. The no-recoil assumption forces the density matrix into a very singular form, where the off-diagonal terms converge rapidly to zero, while the diagonal termss remains intact . Applying the no-recoil assumption to a different basis however leads to a diagonal matrix describing a different physical state and which is not diagonal under a change of basis, as shown above. The only possibility to preserve the consistency of the decoherence argument is to acknowledge that the decoherence process induces the decay of all matrix elements, since indeed if $`DA=0`$ for any matrix $`A`$ then the equality $`DA=(C^1DC)A`$ holds for for any change of basis $`C`$. In that case the decoherence process describes the observer’s loss of information, not only on superpositions, but on the state of the system. The special status of superpositions is indeed spurious, since it depends on the measurement operator being considered, i.e. on the observer. The singling out of superpositions for special destructive treatment appears as an anthropomorphic artefact, based on unphysical assumptions. Certain results of decoherence theory’s preserve their validity in the light of the above criticism. Indeed Zurek s ”predictability sieve” () may be reformulated as the claim that if a system’s state can be tracked by an observer, it will behave as expected , so that decoherence reflects the inability of an observer subject to the second principle of thermodynamics to keep track of the system’s interaction with the environment.The pointer basis is a privileged reference system. The belief that Nature does not provide privileged reference systems, unless nudged into doing so by anthropomorphic assumptions, provided the motivation for this note.
References
A. Joos ”Decoherence Through Interaction with the Environment” in ”Decoherence and the Appearance of a Classical World in Quantun Theory” D.Giulini et al. ed,, Spinger, Berlin-Heidelberg-New York 1999. J.P. Paz , W.H. Zurek ”Quantum limits of decoherence: Environment induced superselection of energy eigenstates” Phys.Rev.Lett. 82 (1999) 5181-5185. Zurek W.H. ”Preferred Observable of Predictability, Classicality and the Environment Induced Decoherence” Progr. Theor. Phys. 89, 281-312. |
warning/0001/gr-qc0001040.html | ar5iv | text | # 1. Introduction
## 1. Introduction
The loop quantum gravity approach has yielded a number of interesting results. A mathematical arena has been defined in which the constraints of quantum gravity have been expressed as quantum operators. The complete kernel of the diffeomorphism constraints has been obtained and efforts are on to find the kernel of the Hamiltonian constraint However, contact with the classical limit (i.e. general relativity) has been elusive.
Since very little is known regarding the interpretation of the kernel of the constraint operators, the unambiguous results pertaining to the classical limit have been obtained at the kinematic level wherein the diffeomorphism and Hamiltonian constraints are ignored . By taking recourse to the arguments of Rovelli , it is, however, not inconceivable that kinematic results may be physically relevant. Moreover, in any situation with classical boundary conditions (e.g. black hole horizons, asymptotically flat spacetimes), the classical constraint vector fields leave the boundary conditions invariant. Hence, at the boundary, the smearing functions (lapse and shift) for the constraints typically vanish and kinematic results may acquire physical significance.
Even at the kinematic level, almost all work to date is restricted to an exploration of the classical limit of (functionals of) the spatial metric or densitized triad operators . <sup>1</sup><sup>1</sup>1An exception is . ‘Weave’ states have been constructed which approximate classical metrical information. It is possible that the conjugate (connection) variable fluctuates wildly in such states and, if so, these states cannot be quasi-classical.
In this work we propose a framework to analyze both the metrical as well as the connection degrees of freedom with a view towards the classical limit. The reason a new framework is required is as follows.
The connection dependent operators which have unambiguous classical counterparts are the traces of holonomies around loops. The latter are denoted by $`T_\gamma ^0(A)`$ with
$$T_\gamma ^0(A)=\frac{1}{2}TrH_\gamma (A),H_\gamma (A)=P\mathrm{exp}_\gamma A_a𝑑x^a,$$
(1)
where $`\gamma `$ is a loop embedded in the spatial manifold $`\mathrm{\Sigma }`$, $`A_a`$ is an SU(2) connection, $`H(\gamma )`$ is the holonomy and $`Tr`$ denotes the trace in the $`j=\frac{1}{2}`$ representation. It would be natural to explore the classical limit in terms of these operators. Then, quasiclassical states would be required to approximate the set of all holonomies and say, surface areas, through quantum expectation values of the corresponding operators with low fluctuations. Unfortunately, as we show in the beginning of section 2, holonomies of a classical connection along all possible loops cannot be approximated by any quantum state! More precisely, it is only on a countable set of loops (the set of all loops is, of course, uncountable) that holonomies have a chance of being approximated. However, an arbitrary choice of this countable set is in conflict with spatial covariance.
Therefore, a new framework is needed to analyze the connection degrees of freedom. In this work we propose such a framework. Our main ideas are as follows. Any quasiclassical state must approximate, in some way, the data corresponding to both the spatial metric as well as the connection. To approximate a given spatial metric we need states defined on a “large enough” graph. We require that this graph gives a latticization of the (compact) spatial manifold. Then the preferred set of loops are naturally identified as those which lie on the lattice. We make these ideas precise in section 2 in such a way that the resulting framework is spatially covariant.
Next, given the set of loops on a lattice, we would like to approximate a classical connection. The natural set of operators to consider are the holonomies along these loops. Since we are interested in approximating classical behaviour at scales much larger than the Planck length, it is enough to restrict attention to loops of size much larger than the Planck scale (the size of a loop is measured by the metric part of the classical data). Thus one natural set of connection operators for an analysis of the classical limit are the holonomies along large loops which lie on the lattice. However, the lattice structure suggests an alternative set of operators. These are the ‘magnetic flux’ operators of lattice gauge theory which measure the non-abelian magnetic flux through the plaquettes of the lattice. They are constructed in the usual way from holonomies along the plaquettes. For reasons which we spell out in section 3, we choose to base our analysis of the classical limit on these operators rather than the large loop holonomies. We devote section 3 to this change of focus from holonomy operators to flux operators.
In section 4 we work out our ideas in detail for the case of two spatial dimensions and explicitly display states which approximate aspects of both the classical spatial metric and the $`SU(2)`$ connection. We also indicate how our constructions can be extended to the case of three spatial dimensions.
Section 5 is devoted to a discussion of various issues which arise in the context of our proposal, with an emphasis on its less robust aspects. The discussion in this section indicates that some of our ideas are too simplistic whereas others posess attractive features; it thus points to ways in which the proposal may be modified. We also show, in section 5, how our proposal can be extended to the diffeomorphism invariant context of Rovelli’s work wherein the Hussain Kuchař model is coupled to a matter reference system. Section 6 contains our conclusions.
There seems to be no single viewpoint with regard to the role, within the framework of loop quantum gravity, of considerations at the purely kinematic level. Therefore it is appropriate that we spell out our viewpoint before describing our results.
The aim of loop quantum gravity is to construct a quantum theory which has general relativity as its classical limit. Since, in this approach, the Hamiltonian constraint operator is poorly understood, it would be premature to discuss the classical limit at the full dynamical level. Even at the (spatial) diffeomorphism invariant level, with the exception of the total volume operator, quantum operators corresponding to diffeomorphism invariant classical observables have not been constructed. Without these operators it is difficult to interpret the theory and discuss its classical limit. Since even the kinematic state space is very different from that of conventional flat space quantum field theory, it makes sense to understand the classical limit first at this kinematic level, where even the diffeomorphism constraints are ignored. The classical limit consists of smooth metric and connection data. The approximation of smooth metrical data by weave states is already subtle and the approximation of smooth connection data is still an open question. It is our view that an analysis of the classical limit at the kinematic level may clarify strategies for analysing the classical limit at the spatial diffeomorphism invariant level and finally, (once the Hamiltonian constraint is well understood) at the fully dynamical level.
Independent of the above ‘structural’ role of understanding the classical limit of kinematic gravity, is the question of whether results at the kinematic level have any relevance to physical predictions of full blown quantum gravity. <sup>2</sup><sup>2</sup>2Note that we are not considering those special situations mentioned earlier in this section, involving boundary conditions, where kinematic results are already ‘gauge invariant’. For example, the discreteness of the spectrum of the area operator is often cited by some workers as a physical prediction. Since the Hamiltonian constraint is not well understood, we refrain from discussing this issue in the context of full quantum gravity. Instead, we restrict our attention to the possibility of promoting kinematic results to predictions at the diffeomorphism invariant level. One way to promote kinematic results to the diffeomorphism invariant level, is to couple the gravitational variables to a matter reference system as in, for example, . This can be done only if the kinematic framework for the gravitational variables is spatially covariant. In what follows, we shall be guided by this requirement of spatial covariance.
Notation and Conventions: We assume familiarity with the loop quantum gravity approach (for example see and references therein) and use notation which is standard in the field. $`a,b..`$ are spatial indices, $`i,j..`$ are internal $`SU(2)`$ indices, $`A_a^i(x)`$ is the $`SU(2)`$ connection and $`\stackrel{~}{E}_i^a(x)`$ is the densitized triad.
$`\{A_a^i(x),\stackrel{~}{E}_i^b(y)\}=\iota G_0\delta _a^b\delta (x,y)`$ where $`\iota `$ is the (real) Immirzi parameter . We shall restrict attention to piecewise analytic loops/graphs.
$`\overline{𝒜}`$ is the completion (via a projective limit construction) of the space of smooth connections $`𝒜`$, $`\overline{𝒜/𝒢}`$ is the Gel’fand completion of the space of smooth connections modulo gauge and $`d\mu _H`$ denotes the Ashtekar-Lewandowski (or Haar) measure on $`\overline{𝒜}`$ as well as on $`\overline{𝒜/𝒢}`$.
$`\widehat{O}`$ is the operator version of the classical object $`O`$, $`\widehat{O}^{}`$ is its adjoint and $`O^{}`$ is the complex conjugate of $`O`$. We shall often denote the expectation value of $`\widehat{O}`$ in the quantum state under discussion as $`<\widehat{O}>`$. $`l_{0P}`$ is the length constructed from the dimension-full gravitational coupling $`G_0`$, $`\mathrm{}`$ and $`c`$. In 3+1 dimensions, $`l_{0P}=\sqrt{\frac{G_0\mathrm{}}{c^3}}`$. We shall use units in which $`\mathrm{}=c=1`$.
## 2. The necessity for a new framework and a sketch of our proposal
The most straightforward approach to an analysis of the classical limit of loop quantum gravity would be to construct minimum uncertainty states for the basic operators of the theory. These operators are the ‘configuration’ operators, $`\widehat{T}_\gamma ^0`$, and suitable ‘momentum’ operators. The latter may be chosen as the area operators, $`\widehat{A}_S`$ ($`A_S`$ is the area of a surface $`S`$ in $`\mathrm{\Sigma }`$).
A tentative definition of a quasi-classical state as a minimum uncertainty state for this set of operators is as follows. A kinematic quasi-classical state $`|\psi >L^2(\overline{𝒜/𝒢},d\mu _H)`$ which approximates the $`SU(2)`$ gauge equivalence class of the classical data, $`(A_{0a}^i(x),\stackrel{~}{E}_{0i}^b(x))`$, is such that, for all $`\gamma ,S`$
(i) $`|<\widehat{T}_\gamma ^0>T_\gamma ^0(A_0)|`$ and $`\mathrm{\Delta }\widehat{T}_\gamma ^0=\left(<(\widehat{T}_\gamma ^0)^2><\widehat{T}_\gamma ^0>^2\right)^{\frac{1}{2}}`$ are small.
(ii) $`|<\widehat{A}_S>A_S(E_0)|`$ and $`\mathrm{\Delta }\widehat{A}_S`$ are small compared to $`A_S(E_0)`$.
Since $`|T_\gamma ^0(A)|<1`$, we interpret ‘small’ in (i) as ‘small compared to 1’.
We now show that no quantum state exists in the kinematic Hilbert space for which (i) is true for all loops $`\gamma `$. The kinematical Hilbert space, $`L^2(\overline{𝒜/𝒢},d\mu _H)`$, is spanned by the set of cylindrical functions, each of which is labelled by a piecewise analytic, closed, finite graph. Hence any element of the Hilbert space is associated with at most a countable infinite set of closed graphs. From the properties of $`d\mu _H`$ it is easy to see that given any such state $`|\psi >L^2(\overline{𝒜/𝒢},d\mu _H)`$, and any loop $`\alpha `$ which does not belong to the countable set of graphs associated with $`|\psi >`$, $`<\psi |\widehat{T}_\alpha ^0|\psi >=0`$ and $`\mathrm{\Delta }\widehat{T}_\alpha ^0=\frac{1}{2}`$. Clearly, there are uncountably many loops of the type $`\alpha `$. It follows that there is no state for which (i) holds for all loops $`\gamma `$ in $`\mathrm{\Sigma }`$.
Hence, the most straightforward approach to an analysis of the connection degrees of freedom fails and a different approach, which relaxes (i) in some way, needs to be formulated. The remainder of this section is devoted to the construction of such an approach.
¿From our arguments above, it is clear that we must relax (i) to hold for at most a countable set of loops. It is reasonable to require that the structure of this set of loops be such that we can use them to approximate, in some way, any given loop. A lattice structure is one which has this property. Thus, we are naturally led to require that the set of loops for which (i) is imposed provides a latticization of the compact spatial manifold. (For simplicity, we shall restrict attention to lattices with a finite, though arbitrarily large, number of links).
However, an arbitrary fixed choice of such a lattice (or indeed, of any other countable set of loops) introduces a preferred structure into the description and hence breaks spatial covariance.<sup>3</sup><sup>3</sup>3Diffeomorphisms are unitarily represented on the kinematic Hilbert space. Spatial covariance implies that classical data sets differing by the action of diffeomorphisms are approximated by quantum states which differ by the action of the corresponding unitary operators. We get around this difficulty as follows.
It is essential to note that we are only interested in quantum states which approximate classical data. In particular such states approximate the data for the classical spatial metric. States which approximate only the spatial metric have been constructed in and are based on an underlying graph. If this graph does not extend into a region $`R\mathrm{\Sigma }`$ then $`R`$ has zero volume and any surface in $`R`$ has zero area. Hence a state based on such a graph does not correspond to any classical metric in the region $`R`$. It follows that the graph underlying a quasi-classical state must ‘extend into all of $`\mathrm{\Sigma }`$’ in order to approximate a classical metric on $`\mathrm{\Sigma }`$. Such graphs are called weaves . For our purposes it seems natural to require that the graph underlying any weave state which not only approximates a classical 3- metric, but also approximates a classical connection (modulo $`SU(2)`$ gauge), provides a latticization of $`\mathrm{\Sigma }`$. More precisely, we require that the graph be the 1-skeleton of some cellular complex whose topology is that of $`\mathrm{\Sigma }`$. <sup>4</sup><sup>4</sup>4Most of the weaves constructed in the literature (see and references therein) are the disjoint union of sets of loops, and do not provide a latticization of $`\mathrm{\Sigma }`$. Notable exceptions are the boundary data of spin foam models (see ). Thus, the required lattice structure is not chosen arbitrarily but is obtained from the quasiclassical state itself. It is this feature which preserves the spatial covariance of the resulting framework.
The availability of a lattice structure enables us to analyse many more functions than just the holonomies. More precisely, any function which admits a lattice approximant may be analysed using techniques from lattice gauge theory. Therefore, we shall develop our framework in such a way as to deal with any function of the classical data which admits a lattice approximant (the degree of approximation will be made quantitative shortly).
To make these ideas more precise, we define the following mathematical structures. Let $`L`$ denote a finite piecewise analytic graph which provides a latticization of the compact manifold $`\mathrm{\Sigma }`$. Note that $`L`$ belongs to an uncountably infinite label set, since the action of a diffeomorphism on $`L`$ produces a lattice $`L^{}`$ which is, in general, different from $`L`$.
We define the lattice projector $`\widehat{P}_L`$ as the projection operator which maps any state in $`L^2(\overline{𝒜/𝒢},d\mu _H)L^2(\overline{𝒜},d\mu _H)`$ to its component in the subspace spanned by spin network states which have the following properties: (a) every spin network state in the subspace is labelled by the graph $`L`$, and (b) for every such spin network state, every link of the graph $`L`$ is labelled by some non-trivial (i.e. $`j0`$) representation of $`SU(2)`$. It can be checked that
$$\widehat{P}_L\widehat{P}_L^{}=\delta _{L,L^{}}\widehat{P}_L$$
(2)
where $`\delta _{L,L^{}}=0`$ if $`LL^{}`$ and $`\delta _{L,L^{}}=1`$ if $`L=L^{}`$. Also $`\widehat{P}_L`$ is a (bounded) self adjoint operator on $`L^2(\overline{𝒜/𝒢})`$ so that
$$\widehat{P}_L=\widehat{P}_L^{}.$$
(3)
Denote the space of finite linear combinations of spin networks associated with all the graphs contained in the graph $`L`$ by $`𝒟_L`$ and its completion in $`L^2(\overline{𝒜/𝒢},d\mu _H)`$ as $`_L`$.<sup>5</sup><sup>5</sup>5Note that, since all spin network states based on $`L`$ (including those with some or all links labelled by $`j=0`$) are contained in $`_L`$ , $`_L\widehat{P}_L(L^2(\overline{𝒜/𝒢},d\mu _H))`$. Note that $`_L`$ is the Hilbert space of $`SU(2)`$ lattice gauge theory on the lattice $`L`$.
Let $`\widehat{O}_L`$ be a bounded self adjoint operator on $`_L`$ (or a densely defined symmetric operator on $`𝒟_L`$). Then define the operator $`\widehat{O}`$ as
$$\widehat{O}:=\underset{L}{}\widehat{P}_L\widehat{O}_L\widehat{P}_L.$$
(4)
Here, the sum is over all possible latticizations of $`\mathrm{\Sigma }`$. $`\widehat{O}`$ has the following well defined action on any spin network state in $`L^2(\overline{𝒜/𝒢},d\mu _H)`$. Every spin network state is associated with some unique ‘coarsest’ graph i.e. the graph which has all its edges labelled by non zero spin. Let $`\gamma _0`$ be the coarsest graph for the spin network state $`\psi _{\gamma _0}`$. Then, if $`\gamma _0`$ does not provide a latticization of $`\mathrm{\Sigma }`$, from (4), $`\widehat{O}\psi _{\gamma _0}=0`$ otherwise $`\widehat{O}\psi _{\gamma _0}=\widehat{P}_{\gamma _0}\widehat{O}_{\gamma _0}\psi _{\gamma _0}`$. This action can be extended by linearity to the dense set of finite linear combinations of spin network states in $`L^2(\overline{𝒜/𝒢},d\mu _H)`$ and thus $`\widehat{O}`$ is a densely defined operator on this dense domain.
We now use (4) to encode our ideas for the approximation of classical data $`(A_{0a}^i(x),\stackrel{~}{E}_{0i}^b(x))`$. Let the classical metric constructed from $`\stackrel{~}{E}_{0i}^a(x)`$ be $`q_{oab}`$. Let $`O_L`$ be the classical lattice approximant to the (real) classical quantity $`O`$ on the lattice $`L`$. Typically, for classical functions of interest, the lattice function $`O_L`$ is a sum over the ‘cell’ functions $`O_{I_L}`$ where $`I_L`$ labels the cells/plaquettes of the lattice $`L`$. The finer the lattice $`L`$, the closer is $`O_L`$ to the continuum function $`O`$ and the larger is the number of ‘cell’ contributions to $`O_L`$. The degree to which $`O_L`$ approximates $`O`$ can be made quantitative in terms of the length of the lattice parameters of $`L`$ as measured by $`q_{oab}`$.
Let $`\widehat{O}_L`$ be the operator corresponding to $`O_L`$. We require that $`\widehat{O}_L`$ be constructed as a self adjoint operator on $`_L`$ (or $`𝒟_L`$), from magnetic flux type operators of $`SU(2)`$ lattice gauge theory on the lattice $`L`$. Then, for calculations of expectation values in a quasi-classical state we interpret (4) as the operator corresponding to the classical quantity $`O`$. Recall that we require quasi-classical states to be associated with some lattice $`L`$. From the considerations of , it is expected that the typical link size of such a lattice as measured by $`q_{oab}`$ is of the order of the Planck length. Thus, the only term to contribute to an expectation value in a quasi-classical state in the right hand side of (4), will be one associated with a lattice with Planck size lattice parameters!
This completes the description of our proposed framework but for one last issue. Since the operator $`\widehat{O}`$ has the lattice projection operators, $`P_L`$, in its definition, it is not obvious that the usual correspondence is guaranteed between the Poisson brackets of macroscopic classical quantities $`O`$ and the commutators of the corresponding operators $`\widehat{O}`$. Thus it must be checked if this correspondence holds in expectation value in order that our candidate quasi-classical states be physically acceptable.
We can now summarize our proposed framework for analysing quasiclassicality as follows:
(1) We require that any quasi-classical state $`\psi `$, which approximates both the classical 3- metric and the conjugate connection, $`(A_0,E_0)`$, be associated with some lattice $`L_0`$, so that $`\widehat{P}_{L_0}\psi =\psi `$. <sup>6</sup><sup>6</sup>6Though we shall not do so here, it seems natural to relax this condition and only require (2) and (3) of any quasiclassical state.
(2) Given a classical function, $`O`$, we construct a corresponding operator $`\widehat{O}`$ as follows. We identify the lattice approximant $`O_L`$ to $`O`$ and construct the operator $`\widehat{O}_L`$ in the lattice gauge theory on $`L`$. Then we construct $`\widehat{O}`$ as in (4).
(3) We require that $`<O>O(A_0,E_0)`$, that $`\mathrm{\Delta }\widehat{O}`$ be small compared to typical classical values of the function $`O`$ and that the usual correspondence between commutators and Poisson brackets holds for expectation values in quasi-classical states.
We end this section with a few technical remarks. If for every $`L`$, $`\widehat{O}_L`$ is a bounded self adjoint operator on $`_L`$ then using Lemma 1, section 4.4 of , it can be verified that $`\widehat{O}`$ is an essentially self adjoint operator on the dense domain of finite linear combinations of spin networks in $`L^2(\overline{𝒜/𝒢},d\mu _H)`$.
However, typically, the operators $`\widehat{O}_L`$ of interest are (unbounded) densely defined symmetric operators on $`𝒟_L`$. Then it is straightforward to see that $`\widehat{O}`$ is a densely defined symmetric operator on the dense domain of finite linear combinations of spin network states in $`L^2(\overline{𝒜/𝒢},d\mu _H)`$.
## 3. ‘Magnetic flux’ operators.
Holonomies serve as natural candidates for the classical functions ‘$`O`$’ of section 2, in as much as the connection degrees of freedom are concerned. ¿From the considerations of section 2, we restrict attention to holonomies along loops which lie on the lattice associated with a quasiclassical state. The classical metric being approximated endows every such loop with a size. Clearly, it does not make sense to require that holonomies along Planck size loops display classical behaviour; it is only for loops of size much larger than the Planck scale, that we expect classical behaviour. Hence, we may further restrict our attention to holonomies along such “large” loops.
A different set of operators than the large loop holonomies is suggested by the lattice structure. These operators are the magnetic flux operators of lattice gauge theory which measure the non-abelian magnetic flux through the plaquettes of the lattice. They are defined in the natural way via holonomies along the plaquettes .
Since a single plaquette is typically of Planck size, we shall refer to the magnetic flux through a plaquette as the ‘microscopic’ magnetic flux. Clearly, the microscopic magnetic flux is not of direct relevance to the classical limit. It is only ‘macroscopic’ operators associate with ‘macroscopic’ length scales (i.e. length scales far above the Planck scale) that are relevant to the classical limit. The utility of the microscopic magnetic flux (or equivalently, the holonomy along a ‘microscopic’ loop) is that it serves as the lattice approximant to the curvature of the connection - the curvature is approximated on the lattice by the flux through a plaquette divided by the plaquette area. Many physically interesting functions can be constructed from the curvature (for e.g. $`D(\stackrel{}{N})=_\mathrm{\Sigma }N^a\stackrel{~}{E}_i^aF_{ab}^i`$, where $`N^a`$ is a vector field and $`F_{ab}^i`$ the curvature of the connection) and thus, admit lattice approximants built out of microscopic fluxes.
It turns out, as we show in section 3.1, that because of the differences in their algebraic properties, it is simpler to use the flux operators rather than the holonomies along macroscopic loops, to analyse the classical limit. Moreover, as discussed in section 3.2, the consideration of flux-based macroscopic operators suggests a general strategy to build states in which these operators have low relative fluctuations. For these reasons we shift focus from the holonomies of macroscopic loops to flux based macroscopic operators in our explicit constructions of section 4. As we shall see in section 5, the strategy discussed in section 3.2 is not entirely successful; nevertheless this strategy springs from an interesting idea and, among other things, this work is devoted to examining it in detail.
## 3.1. Algebraic properties of holonomies vs fluxes.
The holonomies and fluxes have very different algebraic properties. Fluxes are associated with 2d surfaces and and are additive. The flux through the union, $`S`$, of disjoint surfaces $`S_I,I=1..M`$ is the sum of the fluxes through each of the surfaces,
$$_SF_{ab}^i=\underset{I=1}{\overset{M}{}}_{S_I}F_{ab}^i.$$
(5)
Here $`F_{ab}^i`$ is the curvature of the connection pulled back to the relevant 2 surface and $`i`$ is some fixed internal $`SU(2)`$ direction. Equivalently, defining the flux $`\mathrm{\Phi }^i(S)=_SF_{ab}^i`$,
$$\mathrm{\Phi }^i(S)=\underset{I=1}{\overset{M}{}}\mathrm{\Phi }^i(S_I)$$
(6)
In contrast holonomies are associated with 1d loops and are multiplicative. Thus if $`\gamma :=\gamma _1\gamma _2\mathrm{}.\gamma _N`$ is the loop composed of the loops $`\gamma _I,I=1..N`$,
$$H_\gamma (A)=\underset{I=1}{\overset{N}{}}H_{\gamma _I}(A)$$
(7)
where the product signifies group multiplication.
Thus, (6) determines the flux through large surfaces in terms of small surfaces which combine to form the large surfaces and (7) determines the holonomy of a composite loop in terms of the holonomies of the loops which compose it.
By definition, (6) also holds for the quantum flux operators and hence for their expectation values. Thus, the expectation values of the fluxes through small surfaces determine the expectation value of the flux through the large surface via the quantum version of (6). This simplifies the construction of quasiclassical states since it suffices to restrict attention to a smaller “basis” set of surfaces from which all surfaces of interest can be composed. Similar considerations hold for gauge invariant flux based macroscopic operators.
In contrast, although (7) also holds for the holonomy operators, it does not necessarily hold for their expectation values due to quantum fluctuations. In fact, as we show below, if (7) is imposed as a relation between expectation values in a quantum state, that state cannot be quasiclassical. This complicates the construction of quasiclassical states; since we cannot restrict attention to a smaller “basis” set of loops, the holonomies have to be approximated all at once. It is in this sense that it is easier to use fluxes than holonomies.
We now prove our claim regarding the holonomy expectation values.
On $`L^2(\overline{𝒜},d\mu _H)`$ define the bounded self adjoint operators
$$\widehat{x}_\alpha ^0:=\widehat{T}_\alpha ^0,\widehat{x}_\alpha ^j:=\frac{i}{2}Tr(\widehat{H}_\alpha \sigma ^j),$$
(8)
where $`\sigma ^j`$ are the $`2\times 2`$ Pauli matrices. Since $`H_\alpha (A)SU(2)`$,
$$\underset{\mu =0}{\overset{3}{}}(\widehat{x}^\mu )^2=det\widehat{H}_\alpha =1.$$
(9)
For any state in $`L^2(\overline{𝒜},d\mu _H)`$,
$$det<\widehat{H}_\alpha >=\underset{\mu =0}{\overset{3}{}}<\widehat{x}^\mu >^2.$$
(10)
¿From (9)
$$det<\widehat{H}_\alpha >=1\underset{\mu =0}{\overset{3}{}}(\mathrm{\Delta }\widehat{x}_\alpha ^\mu )^2.$$
(11)
¿From (10) and (11),
$$0det<\widehat{H}_\alpha >1(\mathrm{\Delta }\widehat{T}_\alpha ^0)^2$$
(12)
Let $`\gamma _I,I=1..N`$ be a set of $`N`$ loops such that their composition is the loop $`\gamma `$. Thus, $`\gamma :=\gamma _1\gamma _2\mathrm{}.\gamma _N`$.
Let $`(A_{0a}^i,\stackrel{~}{E}_{0i}^a)`$ be the classical data to be approximated. Let $`ϵ>0`$ be a physically reasonable lower bound on the attainable uncertainty in the measurement of the $`\widehat{T}_{\gamma _I}^0,I=1..N`$. Thus,
$$\mathrm{\Delta }\widehat{T}_{\gamma _I}^0ϵ>0,$$
(13)
Since $`H_\gamma (A_0)=_{I=1}^NH_{\gamma _I}(A_0)`$ we impose that
$$<\widehat{H}_\gamma >\underset{I=1}{\overset{N}{}}<\widehat{H}_{\gamma _I}>.$$
(14)
$$det<\widehat{H}_\gamma >\underset{I=1}{\overset{N}{}}det<\widehat{H}_{\gamma _I}>.$$
(15)
Since $`L^2(\overline{𝒜/𝒢},d\mu _H)L^2(\overline{𝒜},d\mu _H)`$, we can use (12) to get
$$det<\widehat{H}_\gamma >\underset{I=1}{\overset{N}{}}1(\mathrm{\Delta }\widehat{T}_{\gamma _I}^0)^2.$$
(16)
From (10), (13) and (16)
$$|<\widehat{T}_\gamma ^0>|^2<det<\widehat{H}_\gamma ><(1ϵ)^N.$$
(17)
Since $`ϵ`$ is independent of $`N`$, clearly, for sufficiently large $`N`$, the above equation implies that $`|<\widehat{T}_\gamma ^0>|<<1`$. For generic $`A_{0a}^i`$ there is no reason for the classical variable $`T_\gamma ^0(A_0)`$ to be small. So if we assume that the classical connection of interest is such that
$$T_\gamma ^0(A_0)O(1),$$
(18)
then (i) is clearly violated for the loop $`\gamma `$ because $`|<\widehat{T}_\gamma ^0>T_\gamma ^0(A_0)|`$ is not much less than unity.
For loops of macroscopic size, we obtain rough estimates for $`N`$ and $`ϵ`$ as follows. Quantum gravitational fluctuations are not expected to be significant well above the Planck scale. So for the purposes of the gravitational interaction alone, energy scales of up to a few hundred $`Gev`$ (or equivalently length scales larger than $`10^{16}cm`$) can safely be considered as ‘classical’. A macroscopic size surface of the order of $`100m^2`$ contains the loops $`\gamma _I,I=1..N`$, where each $`\gamma _I`$ encloses a ‘classical’ size area of the order of $`10^{32}cm^2`$. Thus $`N`$ is of the order of $`10^{38}`$. Even if $`ϵ`$ is chosen as small as $`10^{34}`$, we obtain
$$|<\widehat{T}_\gamma ^0>|(110^{34})^{10^{38}}e^{10^4}0,$$
(19)
which clearly violates (i) for classical connections which satisfy (18)!
One way of arriving at a physically motivated choice for $`ϵ`$ is as follows. In addition to the loops $`\gamma _I`$, consider a set of surfaces $`S_J,J=1..N`$, each of classical size $`10^{32}cm^2`$ such that each $`\gamma _I`$ transversely intersects $`S_I`$ exactly once. Then, choosing the orientation of $`S_I`$ to be in the direction of $`\gamma _I`$, and denoting the area of the surface $`S_I`$ by $`A_{S_I}`$, we have
$$\{T_{\gamma _I}^0,A_{S_I}\}=\iota G_0\frac{i}{2}Tr(H_{\gamma _I}\sigma ^i)n_i,$$
(20)
where the right hand side is evaluated at the point of intersection between the loop $`\gamma _I`$ and the surface $`S_I`$. $`n_i`$ is defined as follows. Let $`E_i^a:=\frac{\stackrel{~}{E}_i^a}{\sqrt{q}}`$ where $`q`$ is the determinant of the metric constructed from the triad. Let $`n_a`$ be the unit normal to the surface $`S_I`$ defined by this metric. Then $`n_i:=n_aE_i^a`$. Thus $`n_in^i=1`$ and we expect that for a large class of connections (with less than Planck scale curvature and which also satisfy (18)) and triads, it should be true that
$$\{T_{\gamma _I}^0,A_{S_I}\}=\iota G_0\frac{i}{2}Tr(H_{\gamma _I}\sigma ^i)n_i\iota G_0O(1).$$
(21)
Note that if the above equation holds, then $`Tr(H_{\gamma _I}\sigma ^i)`$ is of order unity. This implies that the curvature of the connection, $`F_{ab}^i`$, in physically reasonable coordinates is of the order of $`10^{32}cm^2`$ which is still, for purposes of quantum gravity, classical.
For quasi-classical states we expect that the Poisson bracket to quantum commutator correspondence holds in the sense of expectation value so that
$$i\mathrm{}\{T_{\gamma _I}^0,A_{S_I}\}<[\widehat{T}_{\gamma _I}^0,\widehat{A}_{S_I}]>.$$
(22)
Combining (21) with (22) with the uncertainty principle for $`\mathrm{\Delta }\widehat{T}_{\gamma _I}^0`$, $`\mathrm{\Delta }\widehat{A}_{S_I}`$ we get
$$\mathrm{\Delta }\widehat{A}_{S_I}\mathrm{\Delta }\widehat{T}_{\gamma _I}^0\iota l_{0P}^2$$
(23)
Let us assume, to be conservative, a huge uncertainty in the measurement of area <sup>7</sup><sup>7</sup>7Note that our estimates are in the context of a thought experiment in which the only quantum effects are from the gravitational interaction. In practice, it would of course be almost impossible to directly make the appropriate measurements, due, in part, to the quantum nature of any interaction used in the measuring process. equal to $`10^{32}cm^2`$ and set $`\iota l_{0P}^2`$ to be of the order of the Planck area (the latter is consistent with the black hole entropy calculations of ). Then from (23) $`ϵ=\frac{10^{66}}{10^{32}}=10^{34}`$.
Finally, we note that (14) mirrors the relations (7) between classical holonomies. Since classical holonomies are not gauge invariant objects, it is necessary to extend our arguments to the gauge invariant context of traces of holonomies. We do this in appendix A1 by using Giles’ (re)construction of holonomies from their traces.
This completes our discussion as to why it is technically simpler to use fluxes as opposed to holonomies.
## 3.2. A general strategy for low fluctuations based on flux operators
In this section we describe a general strategy to obtain low relative fluctuations of flux based ‘macroscopic’ operators. This strategy is patterned on the mechanism for low relative fluctuations in statistical mechanics. In the statistical mechanics description of thermodynamic systems, there are ‘$`N`$’ weakly correlated degrees of freedom, $`N`$ being very large. Mean values of macroscopic quantities typically go as $`N`$ times some microscopic quantity whereas the relative fluctuations about the mean go as $`\frac{1}{\sqrt{N}}`$. It is the poor correlation between the degrees of freedom that is responsible for such low relative fluctuations.
How can we use this mechanism for low relative fluctuations in the context of our proposal? Recall from section 2 that the lattices of physical interest associated with quasiclassical states have links which are of the order of the Planck length. A ‘macroscopic’ lattice operator, $`\widehat{O}_L`$, associated with a classical function $`O`$ is typically the sum over ‘$`N`$’microscopic operators $`\widehat{O}_{I_L}`$. The index $`I_L`$ typically ranges over all the plaquettes/cells in a macroscopic volume. Since the cells are of Planck size, $`N`$ is very large. This raises the possibility of constructing states with $`\frac{1}{\sqrt{N}}`$ relative fluctuations in the measurement of $`\widehat{O}`$. We indicate how this could happen below and show that it is possible to construct such states in the next section.
¿From (2) and (4) it is easy to see that
$$\widehat{O}^2=\underset{L}{}\widehat{P}_L\widehat{O}_L\widehat{P}_L\widehat{O}_L\widehat{P}_L.$$
(24)
It can be checked that
$$<\underset{L}{}\widehat{P}_L\widehat{O}_L^2\widehat{P}_L><\widehat{O}^2>.$$
(25)
$$(\mathrm{\Delta }^{}\widehat{O})^2:=<\underset{L}{}\widehat{P}_L\widehat{O}_L^2\widehat{P}_L><O>^2(\mathrm{\Delta }\widehat{O})^2.$$
(26)
It can be verified that $`\mathrm{\Delta }^{}\widehat{O}`$ evaluated in the quasi-classical state based on the lattice $`L_0`$ is given by
$$\mathrm{\Delta }^{}\widehat{O}=\mathrm{\Delta }\widehat{O}_{L_0}.$$
(27)
But
$$\widehat{O}_{L_0}=\underset{I_{L_0}=1}{\overset{N}{}}\widehat{O}_{I_{L_0}}.$$
(28)
Then, if the $`\widehat{O}_{I_{L_0}}`$ are sufficiently uncorrelated in the state, we have for $`I_{L_0}J_{L_0}`$ that
$$<\widehat{O}_{I_{L_0}}\widehat{O}_{J_{L_0}}><\widehat{O}_{I_{L_0}}><\widehat{O}_{J_{L_0}}>.$$
(29)
Then (28) implies that
$$(\mathrm{\Delta }\widehat{O}_{L_0})^2\underset{I_{L_0}=1}{\overset{N}{}}(\mathrm{\Delta }\widehat{O}_{I_{L_0}})^2.$$
(30)
Typically, we expect $`<\widehat{O}_{I_{L_0}}>`$ and $`\mathrm{\Delta }\widehat{O}_{I_{L_0}}`$ to be of order 1 times some microscopic (in general, dimension-full) constant and $`<\widehat{O}>=_{I_{L_0}=1}^N<\widehat{O}_{I_{L_0}}>`$ to be of order $`N`$ times the same constant <sup>8</sup><sup>8</sup>8Unfortunately, as we shall see in section 5, our strategy is not entirely successful because this expectation is not quite true for the operators and the states that we examine in section 4.. Then we get
$$\frac{\mathrm{\Delta }\widehat{O}}{<\widehat{O}>}\frac{\mathrm{\Delta }^{}\widehat{O}}{<\widehat{O}>}=\frac{\mathrm{\Delta }\widehat{O}_{L_0}}{<\widehat{O}_{L_0}>}\frac{1}{\sqrt{N}}.$$
(31)
In section 4 we shall examine some classical functions and their flux-based lattice approximants, and apply the strategy of this section to construct states with low relative fluctuations of the corresponding operators.
## 4 .Kinematical $`2+1`$ gravity
In subsections 4.1-4.3, we explore our ideas in the context of 2 spatial dimensions. In 4.1 we define some macroscopic functions and construct their quantum analogs in accordance with (4). In 4.2 we construct candidate quasi-classical states. In 4.3 we show that the relative fluctuations of the macroscopic operators defined in 4.1 can go as $`\frac{1}{\sqrt{N}}`$ in accordance with the ideas of section 3.
Unfortunately, for the reason mentioned in footnote 8, it is difficult to keep the scale of the fluctuations of these operators smaller than the typical scale of the corresponding classical quantities. Hence, not all of our ideas are successfully implemented. A discussion presenting ways in which our states may be modified, or our strategy refined is also contained in section 5.
In subsection 4.4 we indicate how to generalize our constructions to three spatial dimensions. We note that the same difficulties with the scale of the fluctuations arise there, too, and hence our construction of quasi-classical states is not yet satisfactory.
### 4.1 The macroscopic observables
In two spatial dimensions the phase space variables are a densitized triad and a $`SU(2)`$ connection $`(\stackrel{~}{E}_i^a,A_a^i)`$ where $`i`$ is an $`SU(2)`$ Lie algebra index and $`a`$ is the spatial index. The metric is constructed from $`\stackrel{~}{E}_i^a`$ through $`\stackrel{~}{E}_i^a\stackrel{~}{E}^{bi}=qq^{ab}`$. In two dimensions the spatial geometry is determined if the lengths of all curves in the 2-manifold are specified. Moreover, for non-degenerate $`\stackrel{~}{E}_i^a`$ (i.e. for $`\stackrel{~}{E}_i^a`$ which define non-degenerate 2 metrics), the information in the curvature $`F_{ab}^i`$ of the connection is coded in the local expressions $`\stackrel{~}{E}_i^aF_{ab}^i`$ and $`ϵ_k^{ij}\stackrel{~}{E}_i^a\stackrel{~}{E}_j^bF_{ab}^k`$. Hence the classical functions of interest are the length of an arbitrary curve ‘$`c`$’, $`l(c)`$, the ‘vector constraint’, $`D(\stackrel{}{N})=N^a\stackrel{~}{E}_i^aF_{ab}^i`$, and the ‘scalar constraint’, $`S(N)=\text{}ϵ_k^{ij}\stackrel{~}{E}_i^a\stackrel{~}{E}_j^bF_{ab}^k`$ where $`N^a`$ is an arbitrary vector field and is a density -1 scalar.
The corresponding operators are constructed as follows. The length operator can be constructed independent of the strategy of section 3, in the same fashion as the area operator in 3d. The eigenstates of the length operator $`\widehat{l}(c)`$ are the spin network states and their eigen values have a contribution of $`\lambda _j=2l_P\sqrt{j(j+1)}`$ for every intersection of the curve $`c`$ and a link of the spin network colored by $`j`$. Here $`l_P:=\iota l_{0P}`$. Note that in the language of section 3, this operator induces length operators $`\widehat{l}_L(c)`$ in any lattice $`L`$.
The two sets of connection dependent operators can be defined first on a lattice $`L`$ and then promoted to genuine operators on $`L^2(\overline{𝒜/𝒢},d\mu _H)`$ through (4).
$$\widehat{D}_L(\stackrel{}{N})=\frac{1}{4}\underset{v,p=l_{p1}l_{p2}v}{}\widehat{F}(p)(\widehat{E}(v,l_{p1})N(v,l_{p2})\widehat{E}(v,l_{p2})N(v,l_{p1}))+\mathrm{H}.\mathrm{T}.$$
(32)
where the sum runs over all vertices and all plaquettes that contain each given vertex (at vertex $`v`$ the orientation of plaquette $`p`$ is given by an ordered pair of links $`l_{p1}l_{p2}`$), $`\widehat{F}^i(p)=\frac{i}{2}\mathrm{Tr}(H(p)\sigma _i)`$ ($`H(p)`$ is the holonomy around plaquette $`p`$) and $`\widehat{E}(v,l)`$ acts as a left invariant vector field (multiplied by a factor of $`l_P`$) on functions depending on the holonomy along the link $`l`$ oriented away from vertex $`v`$ . ‘H.T.’ refers to Hermitian transpose.
$`\widehat{E}(v,l)`$ can be interpreted as the triad operator smeared over a line transverse to the link $`l`$, but not crossing $`l`$ in the center but at $`v`$. $`\widehat{F}(p)`$ contains the information of the curvature smeared in the plaquette (plus higher order terms in the curvature that are not small in general). Thus, $`\widehat{E}`$ and $`\widehat{F}`$ are related to the triad and the curvature times factors of the lattice spacing $`a_g`$, measured by the macroscopic metric induced by the length operator in our state. Equation (32) provides a discretization of the classical vector constraints if the vector field $`N^a`$ and the collection of weights assigned to the lattice links are related by $`N(v,l_{p1})\widehat{l_{p1}}+N(v,l_{p2})\widehat{l_{p2}}=\frac{1}{a_g}\stackrel{}{N}(v)`$, with $`\widehat{l_{p1}}`$, $`\widehat{l_{p2}}`$ being unit vectors in the direction of two of the links starting at $`v`$ and forming a right-handed basis.<sup>9</sup><sup>9</sup>9We assume that the vertex is four valent and formed by the intersection of two smooth curves; in this way, the definition does not depend on the choice of links to form the basis. A definition of $`\widehat{D}_L(\stackrel{}{N})`$ which corresponds to the classical function $`D(\stackrel{}{N})`$ for states with arbitrary valence would be more cumbersome to write. Since most of the vertices in the states that we will construct are four valent the expression (32) for $`\widehat{D}_L(\stackrel{}{N})`$ is good enough for our purposes.
The other family of operators is defined by
$$\widehat{S}_L(\text{})=\frac{1}{4}\underset{v,p=l_{p1}l_{p2}v}{}\text{}_L(v)\widehat{F}^i(p)\widehat{E}^j(v,l_{p1})\widehat{E}^k(v,l_{p2})\epsilon _{ijk}+\mathrm{H}.\mathrm{T}.$$
(33)
For this family of operators, the expectation values (on states with mostly four valent vertices) will approximate the classical functions known as the scalar constraints if the scalar of density weight $`1`$ labeling the functions is related to the collection of weights assigned to the vertices by the relation $`\text{}_L(v)=\frac{1}{a_g^2}\text{}(v)`$.
### 4.2 States with $`\frac{1}{\sqrt{N}}`$ relative fluctuations
In this section we display candidate quasi-classical states which provide a realization of our idea of $`\frac{1}{\sqrt{N}}`$ relative fluctuations. As mentioned earlier and discussed in section 5, the states which we construct are not completely satisfactory quasiclassical states. Neverthless, we present the construction of the candidate quasiclassical states in detail in the hope that this may fuel future efforts towards modifying our present strategy appropriately.
We shall display candidate quasi-classical states approximating homogeneous geometries and connections. This family of states includes, for example, states that generate expectation values approximating Euclidean metrics and flat connections on a torus, as well as states which approximate round spheres with constant curvature $`SU(2)`$ connections on them<sup>10</sup><sup>10</sup>10 We remind our reader that the macroscopic observables that we are studying now are of local character and therefore two gauge inequivalent classical flat connections would appear indistinguishable to our “magnetic flux type” observables. .
To make the macroscopic geometry (locally) isotropic, the physical lattice prescribed by the state will cover space with domains with the connectivity of a regular square lattice; these domains will be separated by narrow bands. We demand the distribution of orientations of the regular domains to be isotropic. The dominant contributions to any macroscopic observable will be those coming from the interior of the regular domains, and many domains will be involved in any macroscopic measurement. Thus, macroscopic observables will lose track of the connectivity of the lattice which will only be obvious at the micro-scale.
Our lattice should be composed by regular domains of typical size $`D>>>l_P`$ and have a linear density of links $`\rho _l=\frac{k}{l_P}`$. This is the density of intersections of the lattice links with any curve which wiggles only at the macroscopic scale (technically, its radius of curvature should be macroscopic). With this linear density of links the density of plaquettes is $`\rho _p=(\frac{k\pi }{4l_P})^2`$. We will later show that $`k=\frac{2}{\sqrt{3}}`$ is the correct value of this parameter, given the form of the states described below.
The states will be constructed as products; to each regular domain we will assign a factor and a separate factor will be assigned to the region between the regular domains. The factors assigned to regular domains will also be constructed as products. Taking advantage of the regularity, the interior of the domains are divided into black and white plaquettes in an alternate fashion. In the chess-board-like geometry of the interior of the domains we will assign factors of the wave function only to the black plaquettes asking that the color $`n=2j=0`$ does not appear in the spin network decomposition of any of the factors. Due to the alternate plaquette geometry, the color assigned to the links in a spin network decomposition of the state would be exactly the one coming from its only black plaquette neighbor. In this way our quasi-classical state will provide a physical lattice. It will be important that the spin network decomposition of the state does not acquire any zero color in the region between the regular domains to make sure that the state does prescribe a physical lattice and not a collection of separate domains. A technique to fit the domains together will be described after the contributions from the interior of the domains are explained.
As we mentioned earlier, the factor of the wave function assigned to a domain is a product of factors associated to plaquettes
$$\psi _D=\underset{pBl}{}\psi _p$$
(34)
where $`Bl`$ contains alternate plaquettes. Since there are many more plaquettes in the interior of the domains than in the region between domains, to approximate any macroscopic observable we need to adjust only the factors associated to interior plaquettes. Furthermore, since we will illustrate our construction with a state approximating a homogeneous geometry and a homogeneous connection, all the factors $`\psi _p`$ from the interior plaquettes can be taken equal. We choose
$$\psi _p=\mathrm{cos}\theta \varphi _{n=1}+\mathrm{sin}\theta \varphi _{n=2}.$$
(35)
where $`\varphi _i(A)`$ is the trace of the holonomy around plaquette $`p`$ in the spin $`\frac{i}{2}`$ representation. Other choices of $`\psi _p`$ are possible; we chose the simplest states that defined a physical lattice and had small spins.
Let us now describe the assignment of factors of the wave function to the regions of the lattice that do not belong to the domains described earlier.
To simplify our work we will restrict the geometry of the lattices that we consider. First, we concentrate on the boundary of the regular domains. The boundary of the regular domains will be composed only of black plaquettes (one may construct these kind of geometries by erasing the boundary links of the white plaquettes in the boundary). In addition, we will only consider geometries where the black plaquettes in the boundary of the regular domains share at least two vertices with the black plaquettes in the domain, and if one of these black plaquettes shares only two vertices with the interior plaquettes this plaquette must be triangular (the plaquettes in the interior of the domains are all square plaquettes, but in the boundary we allow also triangular plaquettes). In these geometries it follows that all the plaquettes having a link in the boundary are black and that these boundary plaquettes have at most one vertex that is not shared by any plaquette in the regular domain. We will call these vertices black vertices. Apart from these vertices, in the boundary of the regular domains, there are vertices that are shared by interior plaquettes. We will call these vertices, white vertices.
We will take these boundary vertices as data and construct the rest of the lattice by filling the gaps in between the domains in a way that lets us assign a simple factor of the wave function to this “in between” region of the lattice.
At a bigger scale we can use the regular domains as cells of a latticization <sup>11</sup><sup>11</sup>11 The 1-skeleton of cellular complex with the topology of $`\mathrm{\Sigma }`$. of the surface $`\mathrm{\Sigma }`$. Neighboring domains are separated by bands (analog of links) and these bands meet in rotaries (analog of vertices). For convenience, in the lattices that we will consider the bands and rotaries will have no internal vertices, and the rotaries will have no internal links. In other words, the rotaries are simply cells whose links are boundary links of the bands or boundary links of the regular domains. On the other hand, the bands have interior links, but the interior links of each band are restricted to form a closed curve $`\gamma _B`$ joining black vertices (either joining black vertices from the same domain or joining black vertices of neighboring domains). See the figure.
Fig. 1 We show a region of the lattice that is in the boundary between three regular domains. The domains have square lattice connectivity in their interiors and we assign factors of the wave function to the black plaquettes in the chess-board geometry of the interiors. Separate factors are assigned to each of the bands that serves as boundary between two regular domains; these factors are spin networks of color one whose graph $`\gamma _B`$ is drawn joining the black vertices of the figure without retracing any line.
Due to the connectivity of its interior links, we can assign to a band a factor of the wave function which is simply the spin network determined by its graph and the color $`n=1`$, $`\psi _B=\psi _{\gamma _B,n=1}`$. Since the links of a band join only black vertices, when we multiply the band factors with the domain factors, the spin network decomposition of the state will not have color zero in any link. That is, our proposed wave function
$$\psi =[\underset{D}{}\psi _D][\underset{B}{}\psi _B]$$
(36)
defines a physical lattice which encodes the topology of $`\mathrm{\Sigma }`$.
It is clear that a space manifold with arbitrary topology can be covered by a lattice composed by disconnected domains with trivial topology (whose interiors have the connectivity of a regular square lattice and required link density) and joined by narrow band regions where the lattice does not need to posses any regularity. Thus, from the classical data of a Euclidean torus with a flat $`SU(2)`$ connection we can construct our candidate quasi-classical states based on the required lattice, and analogously from the classical data of a round sphere with a constant curvature connection we can construct candidate quasi-classical states.
### 4.3 Expectation values, fluctuations and correspondence
Given a macroscopic curve we want to calculate the expectation value of its length $`<\widehat{l}_L(c)>`$. The calculation is easy. In two dimensions the eigenstates of the length operator are spin network states and their eigen values have a contribution of $`\lambda _n=\frac{1}{2}l_P\sqrt{n(n+2)}`$ for every intersection with the curve. According to our conventions, the total number of intersections is $`\frac{k}{l_P}l_g(c)`$, where $`l_g(c)`$ is the length of the curve. Now we will make two approximations; we will assume that every plaquette which intersects $`c`$ intersects it twice and that the parameter $`\theta `$ is small (because, as we will see, it is linked to the contribution of one plaquette to the curvature of the connection). In this way we get
$`<\widehat{l}_L(c)>`$ $`=`$ $`2{\displaystyle \lambda _{n=1}\mathrm{cos}^2\theta }+\lambda _{n=2}\mathrm{sin}^2\theta `$ (37)
$``$ $`\lambda _{n=1}{\displaystyle \frac{k}{l_P}}l_g(c)`$
where the sum runs over all plaquettes that intersect the curve $`c`$. ¿From this formula we deduce $`k=\frac{l_P}{\lambda _{n=1}}=\frac{2}{\sqrt{3}}`$.
Now we calculate the expectation values of the connection related observables.
The connection measuring operators are a sum of terms. The dominant contributions to the expectation values come from the interior of the regular domains. Since each term in the sum only affects a few factors in the wave function, the contributions can be easily calculated. Since we deal only with homogeneous wave functions, the calculations simplify even more. If we had allowed the parameter $`\theta `$ to be a function of the plaquette the basic mechanism of our proposal would still work; we would be able to approximate many more classical configurations, but the calculations would not be as simple.
It turns out that the alternate plaquette nature of the wave function makes
$`<\widehat{D}(\stackrel{}{N})>=0`$ for all constant shifts. This would be true even if the parameter $`\theta `$ were a function of the plaquette. To see this, it is convenient to rewrite the action of the vector fields $`\widehat{E}^i(v,l_x)`$ in a way that is tailored to act on wave functions that are products of plaquette factors. When acting on fucntions of type (34), $`\widehat{E}^i(v,l_x)=\frac{il_P}{2}[\widehat{L}^i(p)\widehat{L}^i(p,\widehat{y})]`$, where $`\widehat{L}^i(p),\widehat{L}^i(p,\widehat{y})`$ are the left invariant vector fields acting on functions of the group assigned to the plaquette $`(p)`$ and the neighbor of $`(p)`$ in the $`(\widehat{y})`$ direction respectively. Due to the alternate plaquette geometry, the only terms that do not vanish in the expectation value are the ones containing the factor $`\widehat{F}(p)\widehat{L}(p)`$. The result is
$$<\widehat{D}(\stackrel{}{N})>=\underset{pBl}{}C_0\mathrm{Div}N(p)$$
(38)
where, in the plaquette $`(p)`$ defined by the vertices $`(v,v+l_{p1},v+l_{p1}+l_{p2},v+l_{p2})`$,
$`\mathrm{Div}N(p)=N(v,l_{p1})+N(v,l_{p2})N(v+l_{p_1},l_{p1})+N(v+l_{p_1},l_{p2})`$
$`N(v+l_{p_1}+l_{p2},l_{p1})N(v+l_{p_1}+l_{p2},l_{p2})+N(v+l_{p2},l_{p1})N(v+l_{p2},l_{p2})`$
and $`C_0`$ is calculated for any interior plaquette $`(p)`$ as
$$C_0=\frac{il_P}{2}<\widehat{F}(p)\widehat{L}(p)H.T.>$$
(39)
For small $`\theta `$, $`C_0\frac{5}{4}\theta l_P`$.
For similar reasons, the expectation value of $`\widehat{S}(\text{})`$ simplifies greatly and we get
$$<\widehat{S}(\text{})>=\underset{pBl}{}\text{}_L(p)C_0$$
(40)
We now discuss the fluctuations. For the length operator one can consider $`\widehat{l}_L(c)=_p\widehat{l}_L(c_p)`$. Here, the sum is over black plaquettes which intersect the curve ‘$`c`$’ and $`c_p`$ refers to the segment of $`c`$ in the black plaquette $`p`$. Then, one can easily check that
$$\mathrm{\Delta }\widehat{l}_L(c).6l_P\mathrm{cos}(\theta )\mathrm{sin}(\theta )\sqrt{N}$$
(41)
where $`N`$ is the number of black plaquettes which intersect the curve. For small $`\theta `$, $`\mathrm{\Delta }\widehat{l}_L(c).6l_P\theta \sqrt{N}`$. This is consistent with the fact that for $`\theta =0`$ our states are eigen states of the length operator.
In the case of the vector constraint operator and the scalar constraint operator the calculations are not as simple and multiple contributions appear. Nonetheless, due to the nature of our state, in the regular domains a plaquette is significantly correlated with only a few nearby plaquettes. The boundaries of the regular domains contribute only a small amount to $`(\mathrm{\Delta }\widehat{D}(\stackrel{}{N}))^2`$ and hence it is easy to verify that $`(\mathrm{\Delta }\widehat{D}(\stackrel{}{N}))^2`$ is proportional to the number of plaquettes ‘$`N`$’ in the regular domains. Similar considerations apply to $`(\mathrm{\Delta }\widehat{S}(\text{}))^2`$. Thus our idea of $`\sqrt{N}`$ fluctuations is successfully implemented in the states we have displayed.
One difference (in detail) from the calculation of length fluctuations is that $`\mathrm{\Delta }\widehat{D}(\stackrel{}{N})`$ and $`\mathrm{\Delta }\widehat{S}(\text{})`$ do not vanish when $`\theta =0`$. For example, an important contribution to $`(\mathrm{\Delta }\widehat{D}(\stackrel{}{N}))^2`$ comes from terms of the form $`F(p)L(p)l_P`$; we get (to second order in $`\theta `$)
$$(\mathrm{\Delta }F(p)L(p)l_P)[\frac{7}{8}\frac{7}{16}\theta ^2]l_P$$
(42)
This is when $`(p)`$ is a black plaquette; for white plaquettes we get
$$(\mathrm{\Delta }F(p)L(p)l_P)O(l_P)$$
(43)
regardless of $`\theta `$.
The correspondence between commutator expectation values and Poisson brackets is a more involved calculation and we have not investigated this in any detail. Therefore, we restrict ourselves to the following remarks.
In the case of the length operators it is easy to see that
$$<[l_L(c),l_L(c^{})]>=0.$$
(44)
In three spatial dimensions there are general reasons to expect that, in quasi-classical states, the expectation value of the commutator of the area operators and its fluctuations are small ; the argument also applies to the length operators in the two dimensional case. With regard to the calculations for $`\widehat{D}(\stackrel{}{N})`$ and $`\widehat{S}(\text{})`$, the results of applied to the present context support the correspondence between Poisson brackets and commutator expectation values in quasiclassical states.
### 4.4 Extension to 3+1 dimensions
There is a natural analog of the set of observables that determine our quasi-classicality criterion. For the geometry, the area operators and for the connection we could consider the induced connection on surfaces with arbitrary embedding and measure the connection with the same type of “magnetic flux type” operators (that would have the smearing surface as an extra label). This set of operators seems to be large enough and would be very close to the 2D case studied here. However, we have not done any serious study of its properties. Other families may prove to be better.
Our family of candidate quasi-classical states is tightly tied to a two-dimensional space. However, the main idea is easily extendible to other dimensions. Now we describe it briefly.
The three-dimensional chess-board geometry inside the regular domains is such that black cubes meet only at their vertices. At a vertex ($`v_0`$) two opposite octants are black and the rest are white; one can color the whole lattice translating the painted cubes meeting at $`v_0`$ in the three cartesian directions by an even number of steps. To each black cell we assign the factor $`\psi =\mathrm{cos}\theta \psi _2+\mathrm{sin}\theta \psi _4`$ with $`\psi _{2n}`$ being the spin network state with color $`2n`$ in the edges of the black cube. (Other choices with smaller colors are also possible.)
The factors assigned to the bands in the two dimensional case were found using a procedure that can be adapted to the three-dimensional case. We require that at the boundary of the regular domains the black cells share at least three vertices with the other black cells in the domain. Then we change the shape of the boundary black cells to have only one free vertex. These free vertices are the black vertices needed to construct the lattice in the band region and assign a factor of the wave function to each band. We use these factors that tie neighboring domains with a single spin network of color one per band.
In this way we construct a family of states each of which defines a physical lattice. By adjusting the multitude of free parameters (density of intersections of the lattice links with surfaces that look flat at the microscale and $`\theta `$ as a function of the cells of the regular domains) we should be able to approximate any given classical data. Also, we can restrict to homogeneous states that we would only be able to approximate homogeneous classical data.
## 5. Discussion
It is important to clarify that our intent is not to provide an alternative quantization to that of loop quantum gravity. Loop quantum gravity is a theory still under construction and thus, a yet incomplete enterprise. Even at the kinematic level, as we have argued in section 2, the theory is incomplete in that its most straightforward interpretation does not lead us to the classical limit. We view this work as an attempt to remedy this particular instance of incompleteness by providing a framework to discuss quasiclassicality. As we have stressed before, this framework is applicable only to the calculation of expectation values and fluctuations in quasiclassical states. Thus, we do not yet understand the transition from the fully quantum regime to the regime in which our framework is proposed to apply, namely the semiclassical regime. However, it is clearly necessary to establish some framework which defines a satisfactory notion of quasiclassicality, before the study of this transition can be undertaken and therein, we believe, lies the virtue of this work.
After these preliminary remarks, we discuss the following issues which arise in the context of our proposal.
(i)Superselection sectors: Given a quasiclassical state associated with a lattice ‘$`L`$’, it is clear that no operator of the form (4) maps the state out of the space of spin networks based on $`L`$. Thus, if operators of the form (4) were the only operators in loop quantum gravity, we would be faced with the existence of uncountably many superselected sectors, one for every choice of lattice. However, as stressed in our preliminary remarks above, operators of the form (4) have been constructed solely to probe the classical limit in terms of their expectation values and fluctuations in quasiclassical states. There exist for example, in addition to such operators, microscopic (Planck scale) operators in loop quantum gravity which need not be of the form (4). Such operators can map quasiclassical states out of the putative superselected sectors.
As mentioned earlier in this section, we do not yet understand the relation between the fully quantum regime and the semiclassical regime as defined through our proposal. Hence we do not know the role of these Planck scale operators in the semiclassical regime. Thus, even if there are superselected sectors at the kinematic level, these sectors may disappear when the dynamical aspects of quantum gravity at the Planck scale are incorporated.
(ii) Ambiguities in the construction of the operator $`\widehat{O}`$ corresponding to the classical quantity $`O`$: On a fixed lattice ‘$`L`$’ , there are (infinitely) many microscopically distinct lattice approximations to the same continuum quantity. Thus, there are infinitely many, distinct ways to construct $`\widehat{O}`$ through (4). It is not clear if we should demand that our state be quasi-classical with respect to all possible choices of $`\widehat{O}`$, and if so, whether there exist any such states.
(iii)The algebra of operators of the type $`\widehat{O}`$: A qualitatively different ambiguity results from an examination of the algebra of operators of the type $`\widehat{O}`$. Let the quasi-classical state of interest be associated with the lattice $`L`$. Consider the operators $`\widehat{A}`$ and $`\widehat{B}`$ constructed from $`A_L`$ and $`B_L`$ through (4). For simplicity, assume $`[\widehat{A},\widehat{B}]=0=[\widehat{A}_L,\widehat{B}_L]`$. Then the operator corresponding to the quantity $`AB`$ can be constructed either as
$$\widehat{AB}=\underset{L}{}P_L\widehat{A}_L\widehat{B}_LP_L$$
(45)
or as
$$\widehat{A}\widehat{B}=\underset{L}{}P_L\widehat{A}_LP_L\widehat{B}_LP_L.$$
(46)
Since $`\widehat{AB}\widehat{A}\widehat{B}`$, there is an ambiguity in the definition of the operator corresponding to $`AB`$.
There is a special case in which this ambiguity is irrelevant. If the quasiclassical state based on a lattice is such that $`\mathrm{\Delta }^{}A`$, $`\mathrm{\Delta }^{}B`$ (see equation (26)) are small compared to $`<A>,<B>`$ then it can be shown that the uncertainty principle implies
$$\frac{<\widehat{AB}>}{|<\widehat{A}><\widehat{B}>|}=\frac{<\widehat{A}\widehat{B}>}{|<\widehat{A}><\widehat{B}>|}+ϵ,$$
(47)
where
$$|ϵ|\frac{\mathrm{\Delta }^{}\widehat{A}}{|<\widehat{A}>|}\left(1+(\frac{\mathrm{\Delta }^{}\widehat{B}}{|<\widehat{B}>|})^2\right)^{\frac{1}{2}}.$$
(48)
Thus, in this special case, this particular ambiguity in the definition of the operator corresponding to $`AB`$ is of no consequence.
(iv) How small is the microscopic ‘magnetic flux’?: Our construction of states with small fluctuations is based on the premise that every macroscopic quantity is $`N`$ times some microscopic quantity with $`N`$ very large. Therefore, it is essential that the characteristic scale of the microscopic quantity be much smaller than that of the macroscopic quantity. In this regard, the ‘magnetic’ flux presents the following dilemma. <sup>12</sup><sup>12</sup>12Although our arguments involve quantities which are not $`SU(2)`$ gauge invariant, it is easy to see that our conclusions apply to any gauge invariant quantities constructed from the magnetic field such as $`D`$ and $`S`$ of the previous section
The classical ‘magnetic’ field is related to the spatial and extrinsic curvatures through
$$F_{ab}^i=R_{ab}^i+2\iota D_{[a}K_{b]}^i+\iota ^2ϵ_{jk}^iK_a^iK_b^j,$$
(49)
where $`D_a`$ is the operator compatible with the triad and $`R_{ab}^i`$ is its curvature. $`K_a^i`$ is closely connected with the extrinsic curvature when all the constraints of general relativity are imposed. If the Immirzi parameter, $`\iota `$, is of order unity, then in any physically reasonable coordinates, it is clear that the classical scale for $`F_{ab}^i`$ is much smaller than an inverse Planck area. Hence, the magnetic flux through a plaquette of Planck size should be much less than unity. The microscopic flux operator for a Planck size plaquette ‘$`p`$’ of the lattice associated with a quasiclassical state is $`Tr\widehat{H}_p\tau ^i`$. This operator ‘lives’ on the copy of $`SU(2)`$ associated with ‘$`p`$’ and clearly its fluctuations are of order unity for the type of state contemplated in section 4. This translates to huge fluctuations of order inverse Planck area in the associated microscopic magnetic field. Hence the microscopic field fluctuation is much larger than the macroscopic scale and our ideas do not apply. It can be seen that such large fluctuations in the microscopic field, for our states, result in large fluctuations in the macroscopic field. For the macroscopic field averaged over a surface of macroscopic area $`A`$, the fluctuations turn out to be of the order $`\frac{1}{\sqrt{A}l_P}`$ where $`l_P`$ is the Planck length (see equation (50)). Thus, the macroscopic fluctuations swamp out typical classical values!
Nevertheless, let us see how far we can push our ideas. We need to somehow magnify the typical macroscopic scale. Notice that this is possible (from (49)) if we choose a large value of $`\iota `$. Then small fluctuations of the extrinsic curvature magnify to large fluctuations of the magnetic field/flux. Thus, it is possible to salvage our ideas by appealing to a large $`\iota `$. However, in such a case,it is not clear that the curvature can be identified with the plaquette holonomy since this identification assumes that the plaquette flux is small. Neverthless, if we ignore this objection and choose $`\iota `$ to be large and if we still identify $`\iota l_{0P}^2`$ (in 3d) with the Planck area, $`l_P^2`$, then it is clear that $`G_0`$ cannot take the value of Newton’s constant but must be interpreted as a bare constant.
Another way to improve matters, say in 3d, is to decrease the effective ‘magnetic field’ by increasing the plaquette size. This is possible if the quasiclassical state is defined by high spins so that the area of a plaquette is of the order of $`l_{typical}^2=j_{typical}l_P^2`$. Here $`j_{typical}`$ characterizes the scale of the (high) spins. Note that fluctuations in area will then be characterized by $`l_{typical}^2`$ rather than $`l_P^2`$; hence $`l_{typical}`$ must be much less than the macroscopic scale.
For example, in 3d, this idea applied to the (non gauge invariant) magnetic flux, $`\mathrm{\Phi }^i(S)`$ (see section 3), through a surface $`S`$ of area $`A`$ results in the following estimates. Let $`N`$ be the number of plaquettes tiling $`S`$. Then we have $`A=Nl_{typical}^2`$ and $`\mathrm{\Delta }\mathrm{\Phi }^i(S)\sqrt{N}`$. Then the fluctuation in the average magnetic field $`B^i=\frac{\mathrm{\Phi }^i(S)}{A}`$ is
$$\mathrm{\Delta }B^i:=\frac{\mathrm{\Delta }\mathrm{\Phi }^i(S)}{A}\frac{1}{\sqrt{Al_{typical}^2}}$$
(50)
instead of $`\frac{1}{\sqrt{Al_P^2}}`$.
The emergence of a scale defined by the quasiclassical state between the macroscopic scale and the Planck scale can be argued, independently of our specific “$`\frac{1}{\sqrt{N}}`$” inspired constructions. The area operators (length in 2d) are the fundamental metric dependent operators. The uncertainities in the measurement of connection dependent operators are constrained through the uncertainity principle by the size of the fluctuations in the area (length) and the commutator between the area and the connection dependent operators. The larger the permissible uncertainity in the area, the smaller is the achievable uncertainity in the connection operators. The scale of area fluctuations defines $`l_{typical}`$ and a characteristic ‘spin’, $`j_{typical}:=(\frac{l_{typical}}{l_P})^2`$. Clearly, $`j_{typical}`$ must characterize the scale of spins occurring in a spin network decomposition of the quasiclassical state. As mentioned earlier, for smoothness of the macroscopic geometry, $`l_{typical}`$ must be much less than the macroscopic scale.
A similar picture of the classical limit arises in quantum Regge calculus. The relation between the Ponzano-Regge-Turaev-Viro partition function and the Regge action for three dimensional Euclidean spacetimes holds in the large $`j`$ limit . This means that classical smooth spacetimes have origin in states whose quantum geometry defines a scale $`j_{typical}l_P`$ which is macroscopically small (to approximate a smooth geometry at macroscopic scales) and at the same time is much bigger than the Planck scale.
(v) The possibility of incorporating spatial diffeomorphism invariance into our proposal: Since our constructions do not use any external fixed structures, they are covariant with respect to spatial diffeomorphisms. Hence they ought to generalise to a spatially diffeomorphism invariant setting. Such a setting is provided by the Rovelli model which combines the Hussain-Kuchař model with a matter reference system. In the context of our constructions, the lattice associated with a quasiclassical state for the classical data $`(\stackrel{~}{E}_0^a,A_{0a})`$ can be specified through the choice of a particular eigenstate of the fermion fields in the Rovelli model. The fermion fields define surfaces and the cells of the lattice can be located through the intersections of these surfaces. Let us refer to the eigenstate of the fermion fields which specifies a lattice $`L`$ as $`|L_F>`$. In the Rovelli model, it is possible to construct classical diffeomorphism invariant ‘gravitational’ quantities by involving the reference matter fields in their definition. Our proposal would indicate that an analysis of the classical limit for such diffeomorphism invariant configurations of the ‘gravitational’ field and the matter reference system, can be done in terms of diffeomorphism invariant operators of the form
$$\widehat{O}:=\underset{L}{}P_{L_F}P_L\widehat{O}_LP_LP_{L_F}$$
(51)
Here $`O_L`$ is the lattice approximant of the diffeomorphism invariant classical quantity $`O`$ and $`P_{L_F}=|L_F><L_F|`$ is the projector onto the ‘reference system lattice’. The subsequent considerations of section 3 can be also be suitably generalised to the Rovelli model. The quasiclassical state thus constructed will be one for the ‘gravitational’ variables only- the matter variables are still very quantum because the ‘matter part of the state’, $`|L_F>`$, is an eigenstate of the matter fields.
## 6. Conclusions
In this work we have shown that there are no quantum states in the kinematical Hilbert space of loop quantum gravity which approximate, in expectation value, classical holonomies along all possible loops and that, at best, it may be possible to approximate only a countable number of classical holonomies. Since the holonomy variables are the primary variables of the loop approach, a new framework to analyse the classical limit of kinematic loop quantum gravity is needed which takes into account the above result.
We have proposed a framework in which the choice of a countable number of loops is made without breaking spatial covariance by identifying the loops with those which are contained in the graph underlying the quasiclassical state itself. Since the graph is required to be a lattice we are able to import techniques from lattice gauge theory to examine various operators of interest (see section 2). This part of our work is quite robust.
Next, inspired by the mechanism for low relative fluctuations in statistical mechanics, we explicitly constructed candidate quasiclassical states in 2 spatial dimensions. Although we could successfully implement this mechanism for low relative fluctuations, the states were not completely satisfactory because the fluctuations (as opposed to the relative fluctuations) were very large. More precisely, we were able to construct states which had fluctuations of order $`\sqrt{N}`$ times some naturally occurring microscopic unit, with $`N`$ large. Under the assumption that typical classical values were of order $`N`$ times this unit, these states had $`\frac{1}{\sqrt{N}}`$ relative fluctuations. However, on closer examination we found that this assumption was unwarranted and that the microscopic unit was not small enough. As a result the fluctuations swamped out typical classical values. Nevertheless, the fluctuations were reduced drastically in size as compared to the fluctuations at the Planck scale. For example, at the Planck scale, curvature fluctuations are expected to be of the order of the inverse Planck area; the mechanism of low relative fluctuations reduced the fluctuations in the macroscopic curvature by a factor of $`\frac{l_P}{\sqrt{A}}`$ where $`A`$ is the macroscopic area (see (iv), section 5).
Even if our particular explicit construction of candidate quasiclassical states is irrelevant, it is still true that our proposal establishes a connection with lattice gauge theory and reinforces the ‘weave’ based picture of discrete space. The very fact that we have made a connection to lattice gauge theory techniques raises the issue of ‘bareness’ of the gravitational coupling and the possible existence of several phases and length scales in our quantum theory. In lattice gauge theory, the coupling is renormalized, and phases appear due to dynamical considerations. The considerations of the previous section point towards the need of considering scenarios for different phases and renormalization of coupling constants at the kinematic level. Certainly not much more can be inferred in the absence of dynamics, i.e., the construction of a projector into the space of physical states (in Hamiltonian language, the imposition of the diffeomorphism and, especially, the Hamiltonian constraint).
Since our proposal is new and unconventional, it is essential to confront our constructions with physically reasonable criteria and modify our proposal accordingly. We have attempted to do this to some extent in the previous section, but the consequences of our formalism need to be explored thoroughly before accepting it as a viable approach towards an analysis of the classical limit.
Nevertheless, given that a new framework is needed which identifies a countable set of loops, it seems inevitable that projectors onto this set of loops (such as we have defined) play a crucial role. This fact, along with the need to preserve spatial covariance and the requirement of hermiticity of the operator versions of real classical functions, naturally point towards our specific proposal.
Loop quantum gravity is a very conservative approach to the problem of quantum gravity in that it is an attempt to combine the principles of quantum mechanics with that of general relativity in accordance with tried and tested rules. We believe that the real virtue of the loop quantum gravity approach is that it captures, in a clear way, the points of tension between quantum mechanics and general relativity and hence suggests new ideas beyond the scope of its own framework, which may relax this tension.
In this respect, our work seems to emphasize structures intrinsic to the quantum states as important and hence points away from the embedded spin networks of Rovelli and Smolin to the intrinsically defined spin networks of Penrose . In closing we note that the considerations of this work, the qualitative similarity of the resulting description of classical space with the quantum statistical mechanics description of a classical solid and considerations such as that of Jacobson , reinforce the idea that the dynamics of general relativity (and particularly the Hamiltonian constraint) may arise as a coarse grained/statistical description of fundamental degrees of freedom at the Planck scale.
Acknowledgements: We are indebted to the anonymous referee for invaluable comments and suggestions. JAZ was partially supported by CONACyT-990443.
## Appendix
### A1
Let the space of loops with base point $`x_0`$ be $`_{x_0}`$. Denote the trivial loop by $`e`$. As in , consider the free vector space $`_{x_0}`$ generated by loops in $`_{x_0}`$. On $`_{x_0}`$, define the product law
$$(\underset{i=1}{\overset{n}{}}a_i\alpha _i)(\underset{j=1}{\overset{m}{}}b_j\beta _j):=(\underset{i=1}{\overset{n}{}}\underset{j=1}{\overset{m}{}}a_ib_j\alpha _i\beta _j)$$
(52)
and the involution
$$(\underset{i=1}{\overset{n}{}}a_i\alpha _i)^{}:=\underset{i=1}{\overset{n}{}}a_i^{}\alpha _i^1.$$
(53)
Here, $`a_i,b_j`$ are complex numbers and $`\alpha _i,\beta _j_{x_0}`$.
Fix a connection $`A_{0a}^i`$ and extend the definition of holonomy trace to $`_{x_0}`$ by linearity so that
$$T_{_{i=1}^na_i\alpha _i}^0(A_0)=\underset{i=1}{\overset{n}{}}a_iT_{\alpha _i}^0(A_0).$$
(54)
Next, define
$$I_{A_0}:=\{\underset{i=1}{\overset{n}{}}a_i\alpha _i_{x_0}|\underset{i=1}{\overset{n}{}}a_iT_{\alpha _i\beta }^0(A_0)=0\mathrm{for}\mathrm{every}\beta _{\mathrm{x}_0}\}.$$
(55)
It can be checked that $`I_{A_0}`$ is a two sided ideal in $`_{x_0}`$.
Note that, since $`T^0`$ is an $`SU(2)`$ trace, under involution
$$T_{(_{i=1}^na_i\alpha _i)^{}}^0(A_0)=\underset{i=1}{\overset{n}{}}a_i^{}T_{\alpha _i}^0(A_0)=(\underset{i=1}{\overset{n}{}}a_iT_{\alpha _i}^0(A_0))^{}.$$
(56)
We choose $`A_0`$ such that there exists some loop $`\tau _{x_0}`$ for which
$$|T_\tau ^0(A_0)|1.$$
(57)
Define the complex numbers $`l_1(\tau ),l_2(\tau )`$ as
$$l_1(\tau ):=T_\tau ^0(A_0)+i(1(T_\tau ^0(A_0))^2)^{\frac{1}{2}},$$
(58)
$$l_2(\tau ):=T_\tau ^0(A_0)i(1(T_\tau ^0(A_0))^2)^{\frac{1}{2}}$$
(59)
and $`\rho _1(\tau ),\rho _2(\tau )_{x_0}`$ as
$$\rho _1(\tau ):=(l_1(\tau )l_2(\tau ))^1(\tau l_2e),$$
(60)
$$\rho _2(\tau ):=(l_2(\tau )l_1(\tau ))^1(\tau l_1e).$$
(61)
It can be checked that mod $`I_{A_0}`$,
$$\rho _i(\tau )\rho _j(\tau )=\delta _{ij}\rho _i(\tau ),\rho _1(\tau )+\rho _2(\tau )=e$$
(62)
and that for any $`\alpha _{x_0}`$
$`T_{\rho _1(\tau ^1)\alpha }^0(A_0)`$ $`=`$ $`T_{\rho _2(\tau )\alpha }^0(A_0).`$ (63)
$`T_{(\rho _1(\tau ))^{}\alpha }^0(A_0)`$ $`=`$ $`T_{\rho _1(\tau )\alpha }^0(A_0).`$ (64)
We shall further restrict attention to $`A_{0a}^i`$ such that there exists some $`a_{x_0}`$ for which
$$C:=T_{\rho _1(\tau )a\rho _2(\tau )a^{}}^0(A_0)0.$$
(65)
Using the algebraic properties of the $`T^0`$ variables and (56), (62), (63) and (64) it can be verified that
$$U_\alpha (A_0):=\left(\begin{array}{cc}2T_{\alpha \rho _1(\tau )}^0(A_0)\hfill & \frac{2T_{\rho _1(\tau )\alpha \rho _2(\tau )a^{}}^0(A_0)}{\sqrt{2C}}\hfill \\ \frac{2T_{\rho _2(\tau )\alpha \rho _1(\tau )a}^0(A_0)}{\sqrt{2C}}\hfill & 2T_{\alpha \rho _2(\tau )}^0(A_0)\hfill \end{array}\right)$$
(66)
is an $`SU(2)`$ matrix such that $`U_\alpha (A_0)U_\beta (A_0)=U_{\alpha \beta }(A_0)`$ with $`\frac{1}{2}TrU_\alpha (A_0)=T_\alpha ^0(A_0)`$. Details of this construction maybe found in .
We note that the proof of the above properties of $`U_\alpha (A_0)`$ depend solely on the algebraic properties of the $`T^0`$ (and their extensions to $`_{x_0}`$) and the property (56) of the $`T^0`$ under involution; and is independent of the particular connection $`A_0`$. <sup>13</sup><sup>13</sup>13 Provided, of course, that the various expressions in (66) are well defined. That they are indeed well-defined is guaranteed by the requirements (57) and (65). These algebraic properties are shared by the $`\widehat{T}^0`$ operators and the property (56) translates to adjointness properties of the $`\widehat{T}^0`$ operators. Moreover, since these operators form a commutative algebra it can be verified that substituting $`\widehat{T}^0`$ for all occurrences of $`T^0(A_0)`$ in the above construction, yields an $`SU(2)`$ valued operator $`\widehat{U}_\alpha `$ such that $`\widehat{U}_\alpha \widehat{U}_\beta =\widehat{U}_{\alpha \beta }`$ and $`\frac{1}{2}Tr\widehat{U}_\alpha =\widehat{T}_\alpha ^0`$. <sup>14</sup><sup>14</sup>14 The counterpart of the restrictions (57) and (65) is the fact that some of the operators encountered ( namely, $`(\widehat{l}_1(\tau )\widehat{l}_2(\tau ))^1`$ and $`\widehat{C}^{\frac{1}{2}}`$) are unbounded. We assume that mathematical subtleties related to domain issues of unbounded operators can be taken care of in a more careful treatment.
Now we can substitute $`\widehat{H}`$ by $`\widehat{U}`$ in the arguments of section 2 and obtain (17), this time, in a gauge invariant context with the (weak) restrictions (57) and (65) on the classical connection $`A_{0a}^i`$. |
warning/0001/astro-ph0001399.html | ar5iv | text | # Locally Optimally-Emitting Clouds and the Variable Broad Emission Line Spectrum of NGC 5548
## 1 INTRODUCTION
An important goal of quasar research is to understand the origin and physics of the gas which reprocesses a substantial fraction of the energy generated by the quasar central engine. Since these broad line-emitting regions (BLRs) cannot (yet) be imaged directly, we must infer their properties from the IR – X-ray spectra. Spectral synthesis codes are one of the necessary tools used in the interpretation of these clues, and sophisticated numerical simulations of the broad emission lines (BELs) of quasars and active galactic nuclei (AGN) sprang into existence some 20 years ago (Davidson 1977; Davidson & Netzer 1979; Kwan & Krolik 1981). These were limited to single-slab photoionization calculations, more because of a lack of computer power than because of a lack of observational constraints, though future observations were to better define the breadth of the BLR.
The multi-wavelength monitoring campaigns of the past decade were launched to take advantage of the variable nature of the ionizing continuum and its reverberation signatures in the BELs (Blandford & McKee 1982). The results of these campaigns demonstrated the existence of compact yet geometrically extended and ionization stratified broad emission line regions whose characteristic sizes scale roughly as $`R_{BLR}0.1\mathrm{pc}(L_{46})^{1/2}`$ (Peterson 1993; Netzer & Peterson 1997), where $`L_{46}`$ is the quasar’s mean ionizing luminosity in units of $`10^{46}`$ ergs s<sup>-1</sup>. These and other observations and the increase in computer power spawned a new generation of photoionization models (Rees, Netzer, & Ferland 1989; Goad, O’Brien, & Gondhalekar 1993, hereafter GOG93), in which the gas density, column density, and covering fraction were allowed to vary systematically with distance from the continuum source in a geometrically thick BLR. In this scenario each of these parameters are power law functions in radius, meant to mimic a single pressure law governing the conditions of the line emitting gas through the BLR. Most recently Kaspi & Netzer (1999) applied the pressure law model to their photoionization calculations and took advantage of the constraints provided by the observed integrated flux light curves of five emission lines in the well-studied Seyfert 1, NGC 5548 (Clavel et al. 1991, hereafter C91). Figure 1 shows a plot of its mean UV spectrum from the 1993 HST monitoring campaign. Kaspi & Netzer concluded that the total hydrogen column densities must be at least $`10^{23}`$ cm<sup>-2</sup> at a distance of 1 light-day from the continuum source and that the hydrogen gas densities here must lie between $`10^{11}n_H(\mathrm{cm}^3)10^{12.5}`$.
If the properties of the line emitting gas are controlled by a single pressure law, what is it that sets and normalizes it to account for the surprising similarity of quasar/AGN emission line spectra through several orders of magnitude in luminosity? Can more than one “pressure law” exist? Recently, Baldwin et al. (1995) proposed that in a geometrically extended BLR, a range in line-emitting gas properties (e.g., density, column density) might exist at each radius, and showed that under these conditions the BEL spectrum of quasars and AGN may be dominated by selection effects introduced by the atomic physics and general radiative transfer within the large pool of line emitting entities. This model was dubbed “locally optimally-emitting clouds,” or LOC. They showed that a typical quasar spectrum results from the summation of this amalgam of clouds, and that ionization stratification and the luminosity-radius relationship that produces the similar spectra are natural outcomes. In the next section we confront the LOC model with the time-averaged and time-variable spectra of one of the most intensively studied AGN, NGC 5548, to gain further insights into the model’s strengths and weaknesses and into the physical characteristics of this object’s BLR. In $`\mathrm{\S }`$ 3 we discuss the results and mention a new and potentially powerful technique in deriving the physical parameters of the BLR (Horne, Korista, & Goad 1999) — one that takes the most general approach to the LOC picture. The conclusions are given in $`\mathrm{\S }`$ 4.
## 2 PUTTING A SIMPLE LOC MODEL TO THE TEST
Here we will confront the predictions of the LOC model with the observed time-averaged and time-variable BEL spectra of one of the best-studied Seyfert 1 galaxies, NGC 5548 (C91; Korista et al. 1995, hereafter K95; Peterson et al. 1999; references therein). It is not our intention here to derive the line emitting geometry and dynamics of NGC 5548, but rather to test the viability of the LOC model under simple assumptions by comparing the predicted spectrum and emission line light curves with that of a well-studied AGN.
### 2.1 Photoionization Grid Computations & Assumptions
Using Ferland’s spectral synthesis code, Cloudy (v90.04) (Ferland 1997; Ferland et al. 1998) we generated a grid of photoionization models of BEL emitting entities, here assumed to be simple slabs (hereafter “clouds”) each of which we assumed has constant gas density and a clear view to the source of ionizing photons. The continuum incident upon the clouds does not include the diffuse emission from other BEL clouds, nor do we consider the effects of cloud-cloud shadowing. The grid dimensions spanned 7 orders of magnitude in total hydrogen gas number density, $`7\mathrm{log}n_H(\mathrm{cm}^3)14`$, and hydrogen-ionizing photon flux, $`17\mathrm{log}\mathrm{\Phi }_H(\mathrm{cm}^2\mathrm{s}^1)24`$ (see Korista et al. 1997, hereafter K97), and stepped in 0.125 decade intervals in each dimension (3,249 separate Cloudy models). We will call the plane defined by these two parameters the density – flux plane. For the present simulations we assumed all clouds have a single total hydrogen column density, $`N_H=10^{23}\mathrm{cm}^2`$, though in practice the computations of those clouds with very low ionization parameter, $`U_H\mathrm{\Phi }_H/(n_Hc)10^5`$, stopped when the electron temperature fell below 4000 K. Below this temperature the gas is mainly neutral and there is very little contribution to the optical/UV emission lines. Each individual spectral simulation was iterated until the hydrogen and helium line optical depths converged to 20% or better on successive iterations. The emitted spectrum is not all that sensitive to the cloud column density over the range $`10^{22}N_H(\mathrm{cm}^2)10^{24}`$, since the emitting volumes of most collisional excitation metal lines are fully formed within clouds of column densities $`10^{22}10^{23}\mathrm{cm}^2`$, given a significant range of gas density and ionizing flux (K97; Goad & Koratkar 1998). However, it should be kept in mind that the variations in an emission line spectrum driven by a variable ionizing continuum can differ significantly for column densities spanning $`10^{22}10^{24}\mathrm{cm}^2`$ (GOG93; Shields, Ferland, & Peterson 1995; Kaspi & Netzer 1999). We also assumed that the diffuse emission forms in gas with only thermal motions — this may not be the case and local extra-thermal gas motions could have a significant impact on the diffuse emission spectrum through de-saturation of optically thick lines, altering the radiative transfer, and increasing the contribution of photon pumping to the line emission (Shields, Ferland, & Peterson 1995). Next, we initially assumed solar gas abundances (Grevesse & Anders 1989; Grevesse & Noels 1993); however, for reasons discussed below we altered the gas abundances slightly based upon comparison of the models with the observed time-averaged spectrum of NGC 5548. Finally, we assumed an incident continuum spectral energy distribution (SED) that closely resembles that inferred for NGC 5548 by Walter et al. (1994; model A) using simultaneous IUE and ROSAT/PSPC observations (see also Gondhalekar, Goad, & O’Brien 1996), and that this continuum is emitted isotropically and does not change shape substantially as the luminosity varies. With an average ionizing photon energy of about 84 eV, this continuum is significantly harder than that inferred by Mathews & Ferland (1987) for typical quasars. This is necessary in order to reproduce the heating per photoionization reflected in the observed mean C IV/Ly$`\alpha `$ flux ratio ($``$ 1), large even for Seyfert 1 spectra.
The equivalent width (EW) contour maps in the density – flux plane for six of the seven UV emission line/line blends considered here are shown in Figure 2a and for Mg II $`\lambda `$2800 in Figure 2b. The EW is proportional to the total energy emitted by the line, and so is a measure of the continuum reprocessing efficiency for that emission line. The value of the EW of each point lying in the grid assumes full geometric coverage of the continuum source by that particular cloud. For example, the classical BEL cloud parameters of density and ionizing flux lie at the location of the “star” in the EW contour grids in Figure 2 (e.g., Davidson & Netzer 1979). The value of the Ly$`\alpha `$ EW at this location within the density – flux plane is approximately 800 Å. It was this predicted EW from a single cloud coupled with the observed EWs of Ly$`\alpha `$ that led early researchers to deduce the value of the cloud covering fraction in quasars ($``$ 10%) and Seyfert galaxies ($``$ 20%). The reader may consult K97 for a brief discussion of the distribution of EW contours for the various emission lines. All EWs in Figure 2 are measured with respect to the incident continuum at 1215 Å, and so a ratio of the EW contours of two emission lines yields their flux ratio. Finally, lines of constant $`\mathrm{log}U_H`$ run at 45° angles, from lower left to upper right, in each of density – flux planes in Figure 2.
### 2.2 The UV Broad Emission Line Spectrum of NGC 5548
In this section we derive from the data a time-averaged, rest frame, dereddened, velocity-integrated BEL spectrum of NGC 5548, and then simulate it with a simple LOC model.
#### 2.2.1 The Time-Averaged Emission Line Spectrum
Here we establish the velocity-integrated time-averaged ultraviolet BEL spectrum of NGC 5548. Because of its high quality we chose the unweighted mean spectrum from the 1993 HST observing campaign (Figure 1; see also K95). We use an unweighted mean spectrum because we want the average emission line fluxes without regard to differences in S/N in the individual spectra. We list the total measured (observed frame, reddened) line fluxes in column 2 of Table 1. Since Mg II $`\lambda `$2800 flux was not measured during the 1993 HST campaign, we used a slightly smaller value than its mean flux from the 1989 IUE campaign to reflect the lower mean continuum flux and its expected small response to continuum variations. The measured Mg II flux is also problematic because of its blending with surrounding Fe II emission; see Goad et al. (1999) for a recent deblending analysis for another Seyfert 1. Note that these values are essentially those that appear in column 5 of Table 24 in K95, with the exceptions of C IV and He II $`+`$ O III\]. The direct integrated fluxes of these two sets of emission lines were reported in that table, and those measurements did not include the region of overlap lying between them. Here we use their mean fitted fluxes. Goad & Koratkar (1998) isolated the UV narrow line spectrum from a single HST spectrum (1992 July; Crenshaw, Boggess, & Wu 1993) when the continuum and broad emission lines were in a near historic low state (in 1992 April). We list the measured (observed frame, reddened) narrow emission line fluxes in column 3 of Table 1; most are taken from Goad & Koratkar. The N V value was derived from the analyses in Korista et al. (1995) and that of the Si IV $`+`$ O IV\] blend is our recent estimate. The identification of the narrow line emission contributions of this latter septuplet emission line blend whose individual narrow line widths are expected to be $``$ 5 Å (FWHM) is difficult since their positions in wavelength are spread over $``$ 10 Å. Whatever its value, the observations would indicate that the narrow line contribution at $`\lambda `$1400 is likely to be small.
Correcting the broad line fluxes for reddening was not straight forward. Galactic H I measurements place $`E(\mathrm{B V})`$ near 0.03 (Murphey et al. 1996). However, Kraemer et al. (1998) found an observed narrow emission line ratio of He II 1640/4686 that indicated $`E(\mathrm{B V})0.07`$, placing about $`E(\mathrm{B V})0.04`$ somewhere within NGC 5548. This line ratio is expected to remain near its simple Case B value under most conditions, lying near 7 for conditions in the narrow emission line regions, and thus should be a robust reddening indicator (Seaton 1978; MacAlpine 1981; Ferland et al. 2000). But does this extra reddening lie within the narrow line emitting gas, or in a screen covering the narrow and broad emission line regions plus the continuum, or some combination? For conditions present within the broad emission line region, this ratio should lie between 7 and 9 (Ferland et al. 2000), and the observed broad-line ratio might point to the amplitude of the reddening through the sight-line to the broad emission line region. Unfortunately, the isolation of the BEL components of the He II lines is made difficult due to their breadth and blending with other broad lines (Wamsteker et al. 1990). The broad emission line of $`\lambda `$1640 is blended with the extreme red wing of C IV and emission from O III\] $`\lambda `$$`\lambda `$1661,1666. An unreported analysis of K95 attempted to isolate the broad He II emission using the rms spectrum as a guide, and found that approximately 2/3 of the broad emission from the He II $`+`$ O III\] blend belonged to He II, though the significance of this finding is difficult to quantify. We have adopted this estimated He II/O III\] ratio for the present analysis. On the other hand, we know of no attempt to isolate the BEL component of He II $`\lambda `$4686, blended with moderately strong emission from both Fe II and the extreme wing of H$`\beta `$. Fortunately, whatever the case may be, the reddening correction is small, and other uncertainties are at least as large. Here we adopt the Galactic reddening value, $`E(\mathrm{B V})=0.03`$ ($`R(V)=3.1`$; extinction curve: Cardelli, Clayton, & Mathis 1989), to correct the BEL and continuum fluxes, and assume the remaining reddening occurs within the narrow emission line gas of NGC 5548.
The UV BEL fluxes corrected for narrow line fluxes and Galactic reddening are given in column 4 of Table 1. Finally, column 5 lists the derived time-averaged UV BEL luminosities ($`H_\mathrm{o}=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$; $`q_\mathrm{o}=0.5`$; $`z=0.0172`$) and their adopted uncertainties (in linear luminosity values, not logarithm values). The latter are rather coarse and meant to merely reflect estimates of the uncertainties associated with the measurements of their narrow and broad line fluxes, but do not reflect the uncertainties associated with the adopted cosmological parameters or reddening/extinction correction.
Finally, for our choice of SED and measured time-averaged value of $`\mathrm{log}\lambda L_{\lambda 1350}43.54`$ ergs s<sup>-1</sup>, the hydrogen-ionizing luminosity is $`\mathrm{log}L_{ion}44.26`$ ergs s<sup>-1</sup>. At this luminosity a $`\mathrm{log}\mathrm{\Phi }_H=20`$ $`\mathrm{cm}^2\mathrm{s}^1`$ in Figure 2 corresponds to a distance from the continuum source of $`R12.6(75/H_\mathrm{o})`$ light-days. It should be kept in mind, however, that because of reverberation effects, the measured energy in the UV continuum is not precisely related to the energy that is measured in the lines, even if the form of the ionizing SED is known. A monitoring campaign should be of sufficient duration such that a given line emitting region has been sampled at least once by the range of Fourier frequency components of the variable incident continuum. Whether such “steady-state” conditions are achievable before non-reverberation (e.g., dynamical) effects alter the line emitting regions is uncertain (see Perry, van Groningen, & Wanders 1994; Wanders & Peterson 1996), though the observations do indicate that the reverberation time scales are generally much shorter than the dynamical time scales (Peterson 1993).
#### 2.2.2 Simulating the Time-Averaged UV BEL Spectrum
The first set of assumptions concerning the integration of emission from the clouds in our grid involved the simplification of the geometry: a spherically symmetric distribution of BEL clouds, and we did not consider either the effects of line beaming (Ferland et al. 1992; O’Brien, Goad, & Gondhalekar 1994) or continuum beaming (Wanders & Goad 1996).
To derive an integrated emission line spectrum, the spectrum of each cloud lying within the density – flux plane was assigned a weight in two dimensions: gas density and distance from the ionizing source assuming $`\mathrm{\Phi }_HL/R^2`$. Without specifying the particular shape of the emitting entities, this is equivalent to a two dimensional ($`n_H`$, $`R`$), spherically symmetric, function in effective “cloud” covering fraction. As in Baldwin et al. (1995) and Ferguson et al. (1997) we made the simplifying assumptions that this function is analytic, separable and a power law in each of the two variables (i.e., $`f(R)R^\mathrm{\Gamma }`$ and $`g(n_H)n_H^\beta `$; see equations 1 & 2 in Ferguson et al.). These assumptions are not central to the LOC model, but were chosen merely for their simplicity given the void of observational constraints. Baldwin (1997) found that composite quasar spectra were best matched if the power law indices for the two weighting functions lay near $`1`$. This is equivalent to a cloud covering fraction distribution, $`C_f(R,n_H)`$, with equal weighting per decade in the density – flux plane. In this case, the emission line EW (i.e., continuum reprocessing efficiency) contours in Figure 2 are also proportional to the emission lines’ relative luminosity distributions. Steeper radial and/or flatter gas density cloud distribution functions concentrate the emission at high continuum fluxes and gas densities where the emission is mainly thermalized and inefficient. Flatter radial and/or steeper gas density cloud distributions concentrate the emission at large radii and low gas densities. With minimal line thermalization, the line emission from these latter types of clouds is efficient. However, broad emission line reverberation and line profile studies indicate that significant line luminosity must arise from smaller radii as well (Peterson & Wandel 1999; Wandel, Peterson, & Malkan 1999). In this analysis we adopted an index of $`1`$ for the weighting along gas density, but allowed for a range in possible radial covering fraction power law index. The latter was a parameter in the optimization process, explained below.
Next, using the adopted gas density distribution function we summed the emission along the density axis for each radius, producing a radial surface emissivity function for each of the lines and line blends considered (see Figure 3). We considered densities in the range $`8\mathrm{log}n_H(\mathrm{cm}^3)12`$ at each radius. We did not include in our sum the contributions from transparent clouds with very large $`U_H`$. While low density clouds lying very near to the continuum source may have dimensions that rival their distances from the continuum source, they are also virtually transparent (when Thomson thin) so their existence is irrelevant for the purposes of this study. Clouds with gas densities $`n_H<10^8`$ cm<sup>-3</sup> emit unobserved forbidden lines, and our simple constant density, $`10^{23}`$ cm<sup>-2</sup> column density model clouds become geometrically large relative to the characteristic size of the BLR (C91; Peterson et al. 1991). It is also true that at the distances from the continuum source at which these low density clouds are efficient emitters ($`\mathrm{log}\mathrm{\Phi }_H10^{18}`$ $`\mathrm{cm}^2\mathrm{s}^1`$), the temperatures of grains, if present, lie below their sublimation points. Netzer & Laor (1993) suggested grain survival as a natural mechanism to cut off the broad emission at large radii; this would serve to demarcate the boundary between the BLR and NLR in AGN. Above gas densities of $`10^{12}`$ cm<sup>-3</sup>, the majority of clouds are mainly continuum emitters, and most of the lines are thermalized, assuming thermal local line widths (Rees, Netzer, & Ferland 1989; K97). The notable exceptions to this rule are the excited-state recombination lines of H, He I, and He II, which continue to emit efficiently at these high densities (see K97). The radiative transfer of the Balmer lines is probably the least understood and most uncertain of all the prominent AGN emission lines. In addition the general methods employed in codes like Cloudy to determine ionization and thermal equilibria begin breaking down above densities of $`10^{12}`$ cm<sup>-3</sup>. So while gas densities of $`10^{14}`$ cm<sup>-3</sup> may be present within the BLR, and may solve the long standing Ly$`\alpha `$/H$`\beta `$ problem (recently discussed in Netzer et al. 1995 and Baldwin 1997), we chose not to include this gas in our simulations, nor did we use the Balmer lines to constrain our simulations<sup>1</sup><sup>1</sup>1With the gas density distribution function and upper limit to the gas density of $`10^{12}`$ cm<sup>-3</sup> adopted here, the integrated emission from the models presented here find Ly$`\alpha `$/H$`\beta `$ ratios of about 25. This is smaller than the classical high-density Case B ratio of 34, but still significantly larger than observed, $``$ 11. Such integrated ratios are attainable with the gas density distribution adopted here and an upper limit to the gas density approaching $`10^{14}`$ cm<sup>-3</sup> (Baldwin 1997), but we have excluded the extremely dense gas from consideration here.. The EWs of all emission lines considered here peak at or well below gas densities of $`10^{12}`$ cm<sup>-3</sup>. In summary, the gas density distribution function was fixed and not part of the optimization process.
Finally, we fixed the inner radius of the BLR to 1 light-day. This is not a feature of the LOC picture, and it was done only to accommodate our chosen analytic cloud covering fraction distribution function of physical radii, as appropriate for the luminosity of NGC 5548 ($`R(L/\mathrm{\Phi }_H)^{1/2}`$). As long as this choice of inner radius is small its impact on the results is minor, since for the adopted hydrogen-ionizing luminosity the opt-UV emission line gas at smaller radii must be very high density ($`10^{12}`$ cm<sup>-3</sup>) and/or very high column density. Otherwise the gas will not emit opt-UV emission lines. In order to account for the response of the emission lines to a variable ionizing continuum, the emission line surface emissivity curves were tabulated down to a radius of about ⅓ light-day and out to a radius of about ⅓ pc (Figure 3).
To accommodate the fact that several of the measured emission lines are actually blends, we summed the simulated emission from blended species. This obviated the problem of relying heavily upon the results of uncertain deblending analyses. Thus henceforth, Ly$`\alpha `$ is the sum of Ly$`\alpha `$ $`\lambda `$1216, He II $`\lambda `$1216, and O V\] $`\lambda `$1218. Si IV is the sum of Si IV $`\lambda `$1397, O IV\] $`\lambda `$1402, and S IV\] $`\lambda `$1405. He II is the sum of He II $`\lambda `$1640, O III\] $`\lambda `$1663, and Al II $`\lambda `$1670. C III\] is a sum of C III\] $`\lambda `$1909, Si III\] $`\lambda `$1892, and Al III $`\lambda `$1860. N V $`\lambda `$1240, C IV $`\lambda `$1549, and Mg II $`\lambda `$2800 were treated as unblended emission lines. While N V is certainly blended with the red wing of Ly$`\alpha `$, we assume here that N V dominates the measurement.
Using the simulated annealing scheme described in Goad & Koratkar (1998) to minimize $`\chi ^2`$ between the predicted and observed emission line luminosities, we varied the outer radius ($`R_{out}`$), the power law index on the radial cloud covering fraction ($`\mathrm{\Gamma }`$; see GOG93), and the normalization to the integrated cloud covering fraction ($`C_f`$) to fit the time-averaged BEL spectrum in Table 1. The first two parameters adjust the relative spectrum and the last normalizes the spectrum to the correct luminosity. For a wide range of combinations of these parameters the emission lines belonging to the $`\alpha `$-production elements O, Si, and Mg (O VI $`\lambda `$1034<sup>2</sup><sup>2</sup>2The strength of this line in NGC 5548 has not been reported in the refereed literature. Based upon reports from other Seyfert 1 spectra, we have adopted the ratio O VI/C IV = 0.5 as an upper limit to its strength. While we did not include this upper limit in the optimization process, we did confirm that this upper limit was realized in all acceptable models for $`\mathrm{\Gamma }>1.4`$., O III\] $`\lambda `$1663, Si IV $`+`$ O IV\] $`\lambda `$1400, Si III\] $`\lambda `$1892, Mg II $`\lambda `$2800) were all too strong by factors 1.5–2 compared to their best estimated observed intensities relative to Ly$`\alpha `$, He II, and C IV. Note that oxygen and silicon are each represented by a resonance line and a lower ionization intercombination line; each pair of lines together span large regions in the density – flux plane (Figure 2a and K97). As approximate measures of the total heating and photoionization rates, respectively, the intensities of C IV and Ly$`\alpha `$–He II are to first order independent of the gas abundances. For illustrations of various emission line sensitivities to gas abundances in the density – flux plane, see K97 and Korista, Baldwin, & Ferland (1998). Given the results of this preliminary analyses, we considered the possibility that the gas metal abundances could lie below solar, and we simply scaled the metal/hydrogen abundance ratios to ½ their solar values. However, we left carbon and nitrogen at their solar abundance values, since the few spectral constraints, notably C III\] and N V, did not indicate subsolar abundances for these two elements. Carbon, nitrogen, the $`\alpha `$-elements, and the iron peak elements all have somewhat different stellar population progenitors and need not scale together (e.g., Pagel 1997). The He/H abundance ratio was also left at its solar value. Because the overall metallicity of the simulated gas is slightly sub-solar, the equilibrium electron temperatures within the clouds are slightly elevated over their solar abundance counterparts, and the intensity of the major coolant of the BEL gas, C IV $`\lambda `$1549, is enhanced accordingly. This resulted in a generally closer match to the observed Ly$`\alpha `$/C IV ratio for our choice of continuum SED. While we ascribe no great significance to these adopted gas abundances, they are less arbitrary than the assumption of solar abundances. A much more complete analysis of parameter space (cloud distribution functions, SEDs) will be required to acquire more accurate measures of the gas abundances.
Figure 4a shows the envelopes in minimum $`\chi ^2`$ (solid curves) as functions of each of the three parameters, as determined by the simulated annealing process. The lower dashed lines show the $`1\sigma `$ confidence level ($`\chi ^2=\chi _{min}^2+1.00=2.01`$) for one interesting parameter. The upper dashed lines show the $`1\sigma `$ confidence contour ($`\chi ^2=\chi _{min}^2+4.72=5.73`$) for $`NM=73=4`$ interesting parameters. We consider all satisfactory models to lie below the upper dashed lines and above the solid curves. Figure 4b shows the confidence contours of $`\mathrm{log}\chi ^2`$ as a function of $`C_f`$ and $`\mathrm{\Gamma }`$ for fixed values of $`\mathrm{log}R_{out}`$ incremented at 0.2 dex. The contours are in steps of 0.25 dex with the outer value equal to 1.50 dex in every case. Satisfactory models lie within the bold dashed contour representing $`\chi ^2=\chi _{min}^2+4.72=5.73`$. These appear for outer radii $`\mathrm{log}R_{out}1.75`$. The relative emission line spectrum is driven by $`\mathrm{\Gamma }`$ and $`R_{out}`$, and it is apparent from the location of the bold contour in Figure 4b that a broad inverse relationship exists between these two parameters. Radial cloud covering fractions which fall off more steeply with radius generally require larger outer radii. This is because some of the emission lines are emissive almost exclusively at larger radii (e.g., C III\] and Mg II). Ly$`\alpha `$ and C IV are emissive at intermediate and large radii, while N V, He II, and Si IV are also emissive at small radii (see Figures 2 and 3). Since the line emission from clouds at small radii is inefficient for the two strongest and best measured emission lines (Ly$`\alpha `$ and C IV), larger integrated cloud covering fractions must result from steeper radial covering fraction distributions. The luminosities of Ly$`\alpha `$ and C IV and their ratio provide the tightest constraints on the models. Together, Figures 4a and 4b show that satisfactory fits to the time-averaged emission line spectrum are possible for $`R_{out}60`$ light-days, $`1.6\mathrm{\Gamma }0.5`$, and $`0.33C_f0.80`$, although the condition that $`C_f50\%`$ (as well as O VI/C IV $`0.5`$) constrains $`\mathrm{\Gamma }1.4`$. These results are not surprising considering the observed intensity of Mg II and its theoretical EW contours in the density – flux plane, the analysis of Baldwin (1997), and the observed large EW of broad Ly$`\alpha `$ ($`160`$ Å), respectively. The models of Ferland et al. (1992), Shields & Ferland (1993), and Goad & Koratkar (1998) used a single cloud in their attempts to reproduce the observed properties of Ly$`\alpha `$ and C IV in NGC 5548, and their required covering fractions exceeded 30% – 40%. Thus any model that includes additional emission contributions from other clouds for other emission lines must necessarily have a larger covering fraction. Figure 2 also shows that the lowest ionization parameter clouds included in our models ($`5.5\mathrm{log}U(H)4`$) emit little else but Mg II $`\lambda `$2800 (plus optical H, He I lines not modeled here; K97).
In choosing one particular fit to the time-averaged spectrum in order to illustrate that model’s emission line variability properties, we considered the steepest radial covering fraction index for which an integrated cloud covering fraction $`C_f50\%`$ resulted: $`\mathrm{\Gamma }1.2`$. A steep radial cloud covering fraction distribution results in broader distributions in the emission line lags (GOG93), and a broad distribution in lags is observed for NGC 5548. This choice of $`\mathrm{\Gamma }`$ is also in general agreement with that found by Kaspi & Netzer (1999), $`1.33`$. Figure 4b shows that good solutions with $`\mathrm{\Gamma }<1.2`$ exist, but at the cost of increasingly larger outer radii and integrated cloud covering fractions. Since we do not know the origin of the emitting gas, the outer radial boundary of the BLR is only loosely constrained from any time-averaged, profile-integrated emission line spectrum. However, a significant contribution of emission from very large radii will dampen the emission line variability amplitudes and in most dynamical models will produce narrow emission lines. Models with integrated cloud covering fractions $`C_f>50\%`$ are surely affected significantly by cloud-cloud shadowing and diffuse emission from other clouds, and while one could argue that an integrated cloud covering fraction of even 50% is significant in this sense, we adopted $`C_f50\%`$ as a reasonable validity limit of these integrated cloud models. These were the only “filters” we applied as we considered our choice of successful model fit to the time-averaged spectrum.
The first three columns of Table 2 lists the emission lines or line blends, their simulated time-averaged luminosities, and their simulated luminosity-weighted radii in light days for the model parameters of $`R_{out}140`$ light-days, $`\mathrm{\Gamma }1.2`$, and an integrated cloud covering fraction $`C_f50\%`$ ($`\chi ^23.33`$). If grains were present, their temperatures would lie near or above their sublimation values near this choice of outer radius. All predicted time-averaged line luminosities are within 0.1 dex of their target values (Table 1), based upon the mean spectrum of the 1993 HST campaign. This and similar models predict the following line ratios within the blends: (He II $`\lambda `$1216 $`+`$ O V\] $`\lambda `$1218)/Ly$`\alpha `$ $`\lambda `$1216 $`0.06`$, O IV\] $`\lambda `$1402/Si IV $`\lambda `$1397 $`0.5`$, O III\] $`\lambda `$1663/He II $`\lambda `$1640 $`0.5`$, and Si III\] $`\lambda `$1892/C III\] $`\lambda `$1909 $`0.3`$. The emissivity (or luminosity) weighted radius ($`R_L`$) is proportional to the emission line response function centroid that in turn is equal to the continuum – emission line cross-correlation function centroid for linear line responses to continuum variations (Koratkar & Gaskell 1991; GOG93). For this and all satisfactory solutions to the time-averaged spectrum we find the following sequence in increasing $`R_L`$: N V, Si IV, He II, C IV, Ly$`\alpha `$, C III\], Mg II. This order very nearly corresponds to the one of increasing emission line lags observed in NGC 5548. Thus a simple LOC model can reproduce the observed spectrum and the general observed trends in ionization stratification within the BLR of this and other objects.
### 2.3 Emission Line Light Curves and Lags from an LOC Model
Using the emission line emissivities, the above adopted model parameters ($`R_{out}`$, $`\mathrm{\Gamma }`$, $`C_f`$), and the simple geometrical assumptions ($`\mathrm{\S }`$ 2.2.2) we computed one dimensional emission line response functions, $`\mathrm{\Psi }(\tau )\eta (R)R^2F(R)`$, where $`F(R)`$ is the emission line surface flux at radius $`R`$, $`\eta (R)`$ is the responsivity of the cloud at radius $`R`$, and the lag $`\tau =\frac{R}{c}(1+\mathrm{cos}\theta `$), with $`\theta `$ measuring the azimuthal angle. We then convolved these emission line response functions with observed UV $`\lambda `$1350 continuum light curves (C91; K95) to generate the emission line light curves and emission line – continuum cross correlation function (CCF) peak and centroid (at 50% peak) lags (Peterson & White 1994). We present these light curves and lags here, and discuss their comparisons with the observations in the next section.
Figure 5a shows the comparison between the simulated (emissivity and responsivity-weighted) and observed light curves for Ly$`\alpha `$, N V, Si IV, C IV, He II, C III\], and Mg II (recall that several of these lines are blends) from the 1988–1989 AGN Watch IUE campaign for NGC 5548 (C91). Figure 5b shows a similar comparison for the same sets of lines minus Mg II from the 1993 HST campaign (K95). In the latter we also utilized a smoothed version of the measured noisy UV continuum measured by IUE just prior to the HST campaign. The error bars on the observed data points do not reflect the systematic errors present to varying degrees in these data (C91; K95). The model emissivity and responsivity-weighted CCF peak and centroid lags are given in columns 4–7 of Table 2; the measured lags are reported in C91 and K95 and we discuss these further in the next section. The responsivity $`\eta (R)`$ of an emission line is proportional to the slope $`dF_{line}/dF_{cont}`$ (GOG93), and is a function of radius in our simple spherical BLR. We used the “local” responsivity approximation, in that at every radius each emission line was assigned the local value of the response to a small variation in the continuum flux, given the luminosity/flux normalization for the time-averaged spectrum. This is appropriate as long as either the responsivity does not change dramatically with radius or the continuum variations are not too large. While the first assumption does break down at small radii (see Figure 3), these clouds generally do not contribute substantially to the integrated emission line luminosities. Note that we have approximated what should be $`\eta (R,n_H,N_H)`$ as $`\eta (R)`$. Finally, when generating the simulated emission line light curves, we did not alter the shape of the continuum. The opt-UV continuum in NGC 5548 has been demonstrated to harden with increasing luminosity of the continuum source (e.g., Romano & Peterson 2000). Marshall et al. (1997) showed that over short time intervals at least the EUV continuum is correlated with but varies with a larger amplitude than does the UV continuum. However, the detailed nature of continuum variability across the energy bands remains a mystery (Nandra et al. 1998). Kaspi & Netzer tried a variety of different schemes to alter the SED with the UV luminosity in order to produce a better match to the observed light curves. They found that if their adopted SED’s EUV break-point energy shifted from 3 Rydbergs to 5 Rydbergs with increasing UV luminosity, their models could better reproduce the He II light curve by increasing this line’s variability amplitude. The variability amplitude of Ly$`\alpha `$ increased as well, which was an improvement, although its mean flux value became too large, and the overall quality of their fits to Ly$`\alpha `$, C IV, C III, and Mg II diminished.
We emphasize that at no point did we attempt to fit the observed light-curves or lags — we made an ad hoc choice of parameters ($`R_{out}`$, $`\mathrm{\Gamma }`$, $`C_f`$) from a range of solutions which produced a match to the integrated mean emission line spectrum within the uncertainties, and which would result in a reasonable spread in the emission line lags. This was done purposely to test whether or not broad but simple cloud distribution functions which lead to matches to a time-averaged spectrum might also predict the continuum – line reverberation. Additionally, given our present simplistic approach to the LOC picture, we saw no reason to over-fit the data.
## 3 DISCUSSION OF RESULTS AND THE FUTURE
### 3.1 Lags
The emission lines’ luminosity-weighted radii in Table 2 can be roughly eye-balled in Figure 2 by mentally centroiding the EW contours (allowing for the integrand limits in gas density and radius), since for $`\mathrm{\Gamma }=\beta =1`$ the EW contours are directly proportional to those of luminosity. This was pointed out by Baldwin et al. (1995). However, the emission line lags will be biased toward the response of emission line gas from smaller radii which can respond more rapidly and more coherently to the continuum flux variations than gas at larger radii. This explains in part why the predicted emissivity-weighted lags are 3 to 5 times smaller than the corresponding values of $`R_L/c`$ (Pérez, Robinson, & de la Fuente 1992). The model responsivity-weighted lags will be generally longer than the emissivity-weighted lags because the responsivity $`\eta (R)`$ is proportional to the slope $`dF_{line}/dF_{cont}`$ which generally flattens at small radii for most lines due to effects of ionization and thermalization.
Since the CCF, used to measure the emission line lags, is equivalent to the convolution of the emission line response function with the autocorrelation function of the continuum, the measured lags for an emission line will depend upon the variability nature of the driving continuum, even if the parameters which govern the distribution of emitting gas in phase space are time steady. A continuum variation with a characteristic time scale $`\tau _{cont}`$ will most effectively probe line-emitting regions at distances $`Rc\tau _{cont}/(1+\mathrm{cos}\theta )`$. Differences in the emission line lags are observed between the two campaigns (C91; K95) and predicted in Table 2. These differences may also occur due to a line’s luminosity-weighted radius that migrates in and out with the mean ionizing luminosity of the continuum source (O’Brien, Goad, & Gondhalekar 1995). This is a consequence of non-linearity in the emission line response, and is accommodated to some degree in our locally-linear response approximation. Two other possible reasons for the observed changes in the emission line lags are: (1) the BLR is non-stationary on a time scale of 4 years that separates the campaigns (Wanders & Peterson 1996), and (2) the finite nature of the monitoring campaigns coupled with the dominance of long time-scale trends in the continuum flux variations (Welsh 1999).
The predicted lags of Ly$`\alpha `$ and C IV in Table 2 lie fairly near their reported values for the two monitoring campaigns (12 days and 8 – 16 days respectively \[C91\]; 7.5 – 6.9 days and 4.6 – 7.0 days, respectively \[K95\]). Those of the subordinate lines N V and He II are long compared to their observed values (4 days and 4 – 10 days, respectively \[C91\]; 1.4 – 2.4 days and 1.7 – 1.8 days, respectively \[K95\]), while those of the $`\lambda `$1900 blend and Mg II are too short compared to their observed values (26 – 32 days and $`34`$ days, respectively \[C91\]). The predicted lags of the $`\lambda `$1400 blend are too short compared to the measured values from 1989 campaign spectra ($`12`$ days \[C91\]), and perhaps a bit too long compared to the values measured from the 1993 campaign spectra (3.5 – 4.8 days \[K95\]). We compare the model and observed light curves in Figures 5a,b, and we will discuss these in more detail in $`\mathrm{\S }`$ 3.2. Assuming these differences to be significant, they suggest some clues as to how the actual gas distribution may differ from the one we have derived from the mean spectrum, and we speculate here.
That the observations suggest longer Mg II lags relative to that of Ly$`\alpha `$ than produced in this model may imply the presence of high density gas at radii larger than our model’s outer radius (refer to Figure 2b). Or perhaps the high density gas at larger radii intercepts a larger fraction of the incident continuum than our monotonic radial covering fraction function would predict. This component may be denser than $`10^{12}`$ cm<sup>-3</sup> and may also be in part the same gas which emits the Balmer emission lines (see K97) not modeled here. We also suggest a possible reason that Kaspi & Netzer’s models generally far underproduced Mg II emission: at the larger radii where this line is emissive, their pressure law resulted in clouds with gas densities that emit Mg II less than optimally. Their imposed outer radius of 100 light-days also did not help in this respect. On the other hand, it is our models’ inclusion of this gas that emits Mg II $`\lambda `$2800 and little else other than Balmer emission that helps drive the predicted integrated cloud covering fractions to $`40\%`$. It must also be kept in mind that the conditions under which Mg II $`\lambda `$2800 is emitted are affected by the radiative transfer of the Balmer lines and continuum, and so is probably the least accurate of the seven line and line blends simulated here.
In most of the models computed here, the lag of the $`\lambda `$1900 blend is predicted to be just a bit longer than that of Ly$`\alpha `$, whereas the observations show significantly longer lags for the $`\lambda `$1900 blend, albeit with considerable uncertainties. A glance at Figure 2a shows that the $`\lambda `$1900 blend has a secondary peak in optimal emission near the coordinates ($`\mathrm{log}n_H9.25,\mathrm{log}\mathrm{\Phi }_H19`$). This emission is almost entirely that of C III\] $`\lambda `$1909 (K97) and lies at a far higher ionization parameter than the main diagonal ridge of optimal emission ($`\mathrm{log}U_H2.5`$) for C III\] and the blend. Lying near the $`10^{23}`$ cm<sup>-2</sup> column density-imposed ionization “cliff,” this emission’s strength depends much more sensitively to column density — it forms near the back boundary of the cloud for column densities $`N_H10^{22}`$ cm<sup>-2</sup>. Tests show that an integration over a $`10^{22}`$ cm<sup>-2</sup> column density grid with identical boundary conditions to those here result in a 10% increase in the luminosity-weighted radius of $`\lambda `$1900 relative to that of Ly$`\alpha `$. The presence of a range of column densities, such that the lower gas density clouds (that are emissive in opt-UV lines only at larger radii) have predominantly lower column densities, would further separate the model lags of the $`\lambda `$1900 and Ly$`\alpha `$ lines. All else being equal this would also reduce somewhat the predicted $`\lambda `$1900 intensity as well as the C III\]/Si III emission line ratio.
Finally, it is seen in Figure 2a that the He II blend emission peaks near the $`10^{23}`$ cm<sup>-2</sup> column density-imposed ionization “cliff” for gas densities $`9.5\mathrm{log}n_H11.5`$. Clouds at smaller radii than this blend’s peak in Figure 2a would emit more efficiently in both He II $`\lambda `$1640 and O III\] $`\lambda `$1663 if these clouds had higher column densities. Tests show that an integration over a $`10^{24}`$ cm<sup>-2</sup> column density grid with identical boundary conditions to those here result in a 15% decrease in the luminosity-weighted radius of the He II blend relative to that of Ly$`\alpha `$. The presence of a range of column densities, such that higher column densities were more prevalent for gas densities $`10^{11}`$ cm<sup>-3</sup>, would further separate the model lags of the He II blend and Ly$`\alpha `$. To a smaller extent this would also apply to the N V line as well. These adjustments may not be enough, however, to overcome the significant disparities between the predicted and observed He II blend and N V emission line lags with respect to the continuum. The various problems concerning the measured intensities and lags of the He II spectrum in Seyfert 1 galaxies are to be discussed in greater detail in a work in progress (Ferland et al. 2000).
### 3.2 Light Curves
The predicted emission line lags tell only part of the story. As pointed out by Kaspi & Netzer, the light curves of many emission lines contain far more constraints than either just a mean spectrum or the mean spectrum plus emission line lags. The offsets between the model and observed light curves are for the most part explained by the differences between the model and observed mean broad emission line luminosities (Tables 1 and 2). That the model light curves match as well as they do those from a campaign that occurred 4 years prior to the HST campaign is remarkable (Figure 5a). This may mean that the “cloud” parameters are reasonably steady on this time scale (but see Wanders & Peterson 1996). The responsivity-weighted emission line light curves generally came closer to matching the observed light curves. We invite the reader to compare the results in Figure 5a to those of the more restrictive single pressure law models of Kaspi & Netzer (1999; Figures 7 and 9) for the five emission lines in common. It should be kept in mind that in the latter work, the uni-dimensional radial power-law cloud parameters were optimized to fit the 5 emission line light curves explicitly.
Disregarding the simple offsets between the model and observed mean luminosities, the greatest mismatches in Figure 5 occur in the variation amplitudes of the two recombination lines, Ly$`\alpha `$ and He II (blended with O III\]), for both campaigns. The model light curves of Kaspi & Netzer (1999) for these two lines suffered similarly until they allowed for a continuum luminosity dependent SED, as mentioned above. Because of this deficiency the observed inverse correlation between the continuum flux and the C IV/Ly$`\alpha `$ ratio (Pogge & Peterson 1992) is not predicted by the model presented here. A variable SED and/or the inclusion of a range of cloud column densities (Shields, Ferland, & Peterson 1995), as discussed above, may resolve this shortcoming.
### 3.3 Future Directions
The most important feature of the LOC model is its assumption of a large pool of “line emitting entities” from which to draw the emission, and the strong influence of the natural selection effects introduced by the atomic physics and general radiative transfer. While reasonably successful in result, the analysis presented here is actually quite restrictive in its approach to the LOC model. Here and in Baldwin et al. (1995) the cloud parameter distributions were limited to simple power law functions, and the gas density distribution function was fixed over the full breadth of the broad line region. We have also considered only a single column density. Nature need not choose this rigid distribution, and indeed we have discussed how the loosening of some of these assumptions might improve the model’s match to the observations. We do not know yet the origin of the “line emitting entities” and the LOC model does not directly address this origin, except that it does not need to be one of finely tuned parameters. The LOC model requires only that there be a wide range of “cloud” properties throughout the BEL geometry. We point out that while the LOC model was developed within the cloud paradigm, other origins for the line emitting matter may fall within its general philosophy. For example, an illuminated wind from an accretion disk (e.g., Murray & Chiang 1998), while differing in detail, has an illuminating continuum cutting through gradients in gas density and column density within some differential spherical radius. Such gradients will depend upon the angle of the emitting gas above the disk midplane.
The strength in the approach adopted here is in its simplicity in reproducing general emission line properties of both composite quasar spectra (Baldwin et al. 1995) and those of a well observed AGN. However, the derivation of the detailed properties of the broad line emitting regions and thus understanding their true nature will require loosening the simplifying restrictions imposed here. In fact one could hope to derive a more general multi-dimensional cloud distribution function whose form need not be analytic, e.g., $`f(n_H,N_H,R,\theta ,v_R)`$ using time-variable spectra of sufficient quality, where $`\theta `$ is the azimuthal angle and $`v_R`$ is the radial velocity. Using fake AGN integrated emission line flux and continuum flux light curves and a two-dimensional distribution function $`f(n_H,R)`$, Horne, Korista, & Goad (1999) outlined a means to unify the methods of maximum entropy echo-mapping (Horne, Welsh, & Peterson 1991; Krolik et al. 1991) with photoionization modeling. In this approach, a general cloud distribution function is constrained mainly by the data and the multi-dimensional spectral simulation grid, rather than partly by the simplifying analytic assumptions made here. Horne et al. also take into account the effects of anisotropic line emission for simple cloud geometries, and consider a range of symmetric BLR geometries. This new technique combines the direct and indirect methods for solving this complex problem. In later work we will apply this method to the light curves of the 1989 IUE monitoring campaign of NGC 5548, with the hope of learning something concrete about the distribution of the line emitting entities within its broad line region. Eventually we hope to incorporate a more general cloud distribution function, such as $`f(n_H,N_H,R,\theta ,v_R)`$, in our analysis. However, to take full advantage of this we may need a data set of higher quality (mainly longer duration) than even the 1993 HST campaign. These analyses should impose boundary conditions upon the distribution functions describing the BEL gas, and therefore constrain scenarios for the physical origins and dynamics of this gas.
## 4 CONCLUSIONS
Spanning over 10 billion years of cosmic history and 5 orders of magnitude in energy, the general similarity of quasar/AGN spectra is astounding. Baldwin et al. (1995) suggested the locally optimally-emitting clouds (LOC) picture as a path to advancing our understanding the broad emission lines of AGN: that selection effects (atomic physics and general radiative transfer) operating in a large “pool” of line-emitting environments govern the spectra of quasars. In this picture, the physical characteristics of the line-emitting entities (e.g. gas density, column density) are not unique but are broadly distributed along the radial dimensions of the BLR. The continuum SED and gas chemical abundances are the primary drivers of this natural selection process. In the case of NGC 5548 we find that the ionizing continuum SED must be significantly harder than that of Mathews & Ferland (1987) to reproduce the observed C IV/Ly$`\alpha `$ emission line ratio — in concurrence with interpolation of contemporaneous multi-wavelength observations, and the gas abundances are roughly solar with some hint that some elements have sub-solar abundances. More accurate derivations of each of these will require analyses beyond the scope of this paper.
Using very simple (power law) “cloud” distribution functions in both radius from the central source and gas density for a fixed cloud column density, we tested the LOC picture against the spectroscopic observations of NGC 5548. For a fixed but broad cloud distribution function in gas density, the outer radius ($`R_{out}`$), power law index of the radial cloud covering fraction function ($`\mathrm{\Gamma }`$), and the integrated cloud covering fraction ($`C_f`$) were optimized to predict the time-averaged UV spectrum from the 1993 HST campaign. Satisfactory fits to the time-averaged emission line spectrum are possible for a broad range of parameters: $`R_{out}60`$ light-days, $`1.6\mathrm{\Gamma }0.5`$, and $`0.33C_f0.80`$, although the condition that $`C_f50\%`$ (as well as O VI/C IV $`0.5`$) constrains $`\mathrm{\Gamma }1.4`$. The maximum outer radius of the BLR is only loosely constrained by a time and velocity-averaged emission line spectrum, but considerations of grain survival at low incident continuum photon fluxes, the breadths of the emission line profiles, and the observed responses of the emission lines probably conspire to limit $`R_{out}200`$ light-days. Consideration of the individual emission line light curves would have placed stronger constraints upon $`R_{out}`$, $`\mathrm{\Gamma }`$ and $`C_f`$, had we chosen to do so.
Given that we did not optimize our models to fit directly the light curves of the emission lines, but merely that of the time-averaged spectrum from the 1993 HST campaign, the similarities of the predicted emission line light curves and their lags compared to the observed ones are remarkable. We believe this is a demonstration of the natural predictive power of the LOC picture. The differences between the observed emission line light curves and lags and those predicted by our model suggest to us the following possible general improvements to the LOC model presented here. A range of column densities is present such that higher column density clouds are predominant for higher gas density clouds, while lower column density clouds are predominant for the lower gas density clouds. Because taken together most opt-UV emission lines are visible only for a range in $`\overline{U}_H\overline{L}/(R^2n_H)`$ up to some maximum value of $`\overline{U}_H`$, a broadly characteristic radial dependence of column density may be involved. While this is broadly consistent with the findings of Kaspi & Netzer, there is no reason to believe that such column densities should be uniquely defined at every radius. The gas density distribution function need not remain constant over all radii as assumed here, and in particular there may be a greater concentration of very high density gas at the larger radii than our simple cloud distribution functions would allow.
The next step involves loosening the simplifying constraints imposed here and deriving a general distribution function of cloud properties, e.g., $`f(n_H,N_H,R,\theta ,v_R)`$, constrained mainly by the observed spectra and the spectral simulations. This can be realized through the marriage of echo-mapping techniques with spectral simulation grids, using the constraints provided by a high quality temporal spectroscopic data set.
Finally, we wish to respond to a statement that appears in Kaspi & Netzer (1999). They wrote “Finally, we must comment that present day LOC models are too general and do not contain full treatment of shielding and mixing of the various coexisting components.” We are not certain what was meant by “are too general.” Being “general” is the main point of the LOC picture. We have shown here and elsewhere the LOC picture has been applied, that there is a fairly wide range in the way that BLR gas can be distributed in ($`R,n_H,N_H`$) space and still produce spectra that match typical quasar/AGN spectra. Here we fit the mean integrated fluxes of a series of lines in a specific AGN. Many LOC/pressure-law models can do that, but to match the variability requires that the mean formation radius of the lines show a large degree of variation and in such a way that high ionization lines form characteristically closer in to the central ionizing continuum source. As first pointed out by Baldwin et al., this appears to be a natural consequence of the LOC picture. By contrast, in the single pressure-law model, a given pressure-law constrains the run of density and column density with radius to vary in a very specific manner, and their starting values must be normalized to a given ionization parameter/incident flux at a pre-defined radius. Amongst other things the LOC picture’s generality frees one to investigate some important global parameters over the population of quasars, such as the continuum SED and gas abundances (e.g., Korista, Baldwin, and Ferland 1998). In regards to the “shielding and mixing” comment, the models of broad emission line spectra of AGN presented here, in Baldwin et al. (1995), and elsewhere are and have been internally self-consistent. We discussed a more mature approach to the LOC picture in $`\mathrm{\S }`$ 3.3.
We thank an anonymous referee for his/her constructive comments. This work benefited substantially from the support of a PPARC grant of Keith Horne’s and we would like to thank Keith and the University of St. Andrews for their hospitality. MRG acknowledges support through a PPARC fellowship during the completion of this work. We are also grateful to Gary Ferland for maintaining his freely distributed code, Cloudy, and to Jack Baldwin for his inspiration. We also thank Jack for his careful reading and suggestions which improved the manuscript. |
warning/0001/gr-qc0001080.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The gravitomagnetic clock effect (GCE) consists in the loss of synchrony of identical clocks carried around a massive spinning body, in opposite directions. This effect is a consequence of general relativity and its presence has been foreseen in connection with the so called gravitomagnetism, i.e. that part of the gravitational field which, in weak field approximation behaves as the magnetic part of the electromagnetic interaction . The GCE has been considered as an interesting an promising means to test the general relativistic influence of the angular momentum of a mass on the structure of space time nearby, and in particular on the pace of clocks orbiting around the body .
Actually the GCE is strictly akin to the Sagnac effect, which is a special relativistic effect induced by pure rotations, first considered as a purely classical effect by G. Sagnac , further on recognized in its real nature and studied by several author (see for instance ).
A relevant aspect both of GCE and of Sagnac effect is their genuine geometrical character. The present paper moves from the remark that the world line of an object steadily moving along a spacially circular trajectory around a symmetry axis of the gravitational field is a helix, to derive a general method for describing the GCE from different viewpoints in general terms, without being a priori limited to the Kerr metric or to its weak field limit as in , and . Finally the general method, when applied to the mentioned special cases, will reproduce the known results, but will allow also for an easier recognition of the situations more viable for experimentally detecting time delays and synchrony defects.
## 2 General geometric features of the gravitomagnetic clock effect
As already said in the introduction the space time of a massive body steadily rotating about an axis at rest with respect to distant galaxies has a simple and interesting property: the world line of any object orbiting the central mass at a constant distance from the axis is a helix.
In other words, for any such object a cylinder (bidimensional surface) exists on which the world line is drawn; this surface is flat and when opened in a plane the world line becomes a simple straight line. This is true irrespective of the global curvature of space time along the fourdimensional trajectory.
Viewing things this way allows one to catch at a glance both the Sagnac and the GCE. Consider the $`1+1`$ cylinder opened in a plane, shown on figure 1.
Oblique lines are the developed helices; on the horizontal axis angular coordinates are reported, the vertical one shows a variable proportional to the inertial time of an observer at rest with the symmetry axis; points at $`\varphi =\pi `$ are identified with points at $`\varphi =\pi `$ and the same $`t`$. Objects steadily rotating with different speeds are represented by different slope lines; it is evident that such lines cannot cross each other on the vertical of the origin (vertical dotted line on O) unless they are symmetric with respect to it, i. e. unless the speeds are the same in magnitude but oppositely directed. Similarly it is evident that also the intervals between two successive crosspoints (two conjunctions) are different if measured along the two lines, unless the slopes are symmetric: the proper times between conjunctions are different.
These simple facts are indeed the graphic explanation of the Sagnac effect: a rotating observer is in his turn represented by a helix in space time (straight line on fig. 1), whose interceptions with the world lines of the test bodies determine proper intervals and proper time differences typical of the effect. The situation is not different when a gravitational field is present, provided it possesses an axial symmetry and no angular momentum as it is the case for a Schwarzschild metric. Simply the gravitational field will affect the conversion between the proper time of the rotating probes and the coordinate time, but all the typical Sagnac phase effects will be there.
When an angular momentum must be accounted for, the scheme remains in principle the same, but a careful discussion of the viewpoint of different observers needs be made.
### 2.1 Graphic representation of the non zero angular momentum case
When the gravity source possesses an angular momentum two objects rotating in opposite directions with the same coordinate speed have again a conjunction at the same coordinate time and the same proper time. However two freely orbiting objects on the same circular trajectory and opposite directions do so in general at different angular coordinate speeds, have conjunctions at different proper times and cross the line of sight of a distant inertial observer at different coordinate times. The image given in fig. 1 may be retained but for the fact that now the axes of $`t`$’s and $`\varphi `$’s are no longer orthogonal: the surface on which the world lines of circularly rotating objects are drawn, are still flat and the plane representation of the situation is shown on fig. 2.
Everything may be described by the use of simple methods of bidimensional Minkowskian geometry. The reference frame is the one drawn in the figure. The world line of a rotating object may still be written as
$$t=\frac{\varphi }{\omega }$$
(1)
provided it passes through the origin event. We assume that when $`\omega >0`$ the object is corotating with the source of the gravitational field; the reverse when $`\omega <0`$.
Considering the rotation symmetry, we must complement (1) with the condition that, when the running event reaches the borders of fig. 1 or fig. 2 the representative angle bounces back ($`\omega >0`$) or ahead ($`\omega <0`$) by a $`2\pi `$ term. In other words the world line becomes
$$t=\frac{\varphi \pm 2\pi }{\omega }$$
(2)
where the $`+`$ sign corresponds to corotation and the $``$ one to counter-rotation. After one more turn an additional $`2\pi `$ is introduced, and so on.
Equipped with these simple definitions and rules we immediately see that the coordinate time for a complete revolution (distant observer view point) which is of course $`T=\frac{2\pi }{\left|\omega \right|}`$ corresponds to the proper time interval of the revolving object
$`\tau `$ $`=`$ $`{\displaystyle \frac{1}{c}}\sqrt{g_{tt}T^2+2g_{t\varphi }2\pi T+g_{\varphi \varphi }4\pi ^2}`$
$`=`$ $`{\displaystyle \frac{2\pi }{c\left|\omega \right|}}\sqrt{g_{tt}+2g_{t\varphi }\omega +g_{\varphi \varphi }\omega ^2}`$
In flat space time this would be $`\tau =T\sqrt{1\beta ^2}`$; in the general case the presence of a term containing $`g_{t\varphi }`$ accounts for the non orthogonality of the reference axes (polar coordinates are understood).
It is useful to work out the position of the first intersection event (conjunction) of two objects endowed with different angular velocities $`\omega _1`$and $`\omega _2`$; let us assume by default that $`\omega _1>0`$ and $`\omega _1>\omega _2`$. If it is also $`\omega _2>0`$ the first conjunction is found when the equation
$$\frac{\varphi +2\pi }{\omega _1}=\frac{\varphi }{\omega _2}$$
is satisfied, i.e. when
$$\{\begin{array}{c}\varphi =\varphi _a=2\pi \frac{\omega _2}{\omega _1\omega _2}\hfill \\ t=t_a=\frac{2\pi }{\omega _1\omega _2}\hfill \end{array}$$
The proper times of the two objects at the conjunction are
$$\tau _{1,2}=\frac{2\pi }{c\left(\omega _1\omega _2\right)}\sqrt{g_{tt}+2g_{t\varphi }\omega _{1,2}+g_{\varphi \varphi }\omega _{1,2}^2}$$
corresponding to a synchrony defect
$`\delta \tau _{12a}`$ $`=`$ $`\left(\tau _1\tau _2\right)_a`$ (4)
$`=`$ $`{\displaystyle \frac{2\pi }{c\left(\omega _1\omega _2\right)}}\left(\sqrt{g_{tt}+2g_{t\varphi }\omega _1+g_{\varphi \varphi }\omega _1^2}\sqrt{g_{tt}+2g_{t\varphi }\omega _2+g_{\varphi \varphi }\omega _2^2}\right)`$
If it is $`\omega _2<0`$ and $`\omega _2>\omega _1`$the values are
$$\{\begin{array}{c}\varphi _b=2\pi \frac{\omega _1}{\omega _1\omega _2}\hfill \\ t_b=t_a\hfill \\ \delta \tau _{12b}=\delta \tau _{12a}\hfill \end{array}$$
(5)
Finally, when $`\omega _2<0`$ and it is $`\omega _2\omega _1`$
$$\{\begin{array}{c}\varphi _c=2\pi \frac{\omega _1+\omega _2}{\omega _1\omega _2}=\varphi _a+\varphi _b\hfill \\ t_c=\frac{4\pi }{\omega _1\omega _2}=2t_a\hfill \\ \delta \tau _{12c}=2\delta \tau _{12a}\hfill \end{array}$$
A relevant situation is that of circular geodetic motion at a constant coordinate radius $`r`$. To study this case let us start from a metric whose non zero elements are $`g_{tt}`$, $`g_{t\varphi }`$, $`g_{rr}`$, $`g_{r\theta }`$, $`g_{\theta \theta }`$ and $`g_{\varphi \varphi }`$; all of these elements depend on $`r`$ and $`\theta `$ only. Imposing the conditions $`r=`$ constant, $`\theta =`$ constant $`=\pi /2`$ with a symmetry such that all the metric elements are extremal for the chosen $`\theta `$ value, the equations of geodesic motion lead to the expression
$$g_{\varphi \varphi ,r}\omega ^2+2g_{t\varphi ,r}\omega +g_{tt,r}=0$$
(6)
Commas mean partial differentiation with respect to the variable after them.
Angular velocities of (spacely) circular geodesic motion are then
$$\omega _\pm =\frac{g_{t\varphi ,r}\pm \sqrt{g_{t\varphi ,r}^2g_{\varphi \varphi ,r}g_{tt,r}}}{g_{\varphi \varphi ,r}}$$
(7)
This can be written
$$\omega _{1,2}=\omega _0\pm \omega _{}$$
(8)
where $`\omega _0=g_{t\varphi ,r}/g_{\varphi \varphi ,r}`$ and $`\omega _{}=\sqrt{\omega _0^2g_{tt,r}/g_{\varphi \varphi ,r}}`$.
Using (8) and arranging summations so that $`\omega _1>\omega _2`$, we can calculate the synchrony defect at conjunction for two freely counter-orbiting objects.
To end this section let us still consider the situation as viewed by an observer who rotates with an angular velocity $`\mathrm{\Omega }`$ of his own. In the proper time of this observer the revolution period of a prograde orbiting object is deduced from (4) with $`\omega _1=\omega _0+\omega _{}`$ and $`\omega _2=\mathrm{\Omega }`$ ($`\omega _1>\mathrm{\Omega }>0`$), obtaining
$`\tau _+`$ $`=`$ $`{\displaystyle \frac{2\pi }{c\left(\omega _1\mathrm{\Omega }\right)}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
$`=`$ $`{\displaystyle \frac{2\pi }{c\left(\omega _0+\omega _{}\mathrm{\Omega }\right)}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
The retrograde case (now $`\omega _1=\mathrm{\Omega }`$ and $`\omega _2=\omega _0\omega _{})`$ corresponds to a revolution time
$`\tau _{}`$ $`=`$ $`{\displaystyle \frac{2\pi }{c\left(\mathrm{\Omega }\omega _2\right)}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
$`=`$ $`{\displaystyle \frac{2\pi }{c\left(\mathrm{\Omega }+\omega _{}\omega _0\right)}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
Then the proper time difference between the conjunctions with the observer will be
$`\delta \tau `$ $`=`$ $`{\displaystyle \frac{2\pi }{c}}{\displaystyle \frac{\left(2\mathrm{\Omega }\omega _1\omega _2\right)}{\left(\omega _1\mathrm{\Omega }\right)\left(\mathrm{\Omega }\omega _2\right)}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
$`=`$ $`{\displaystyle \frac{4\pi }{c}}{\displaystyle \frac{\mathrm{\Omega }\omega _0}{\omega _{}^2\left(\mathrm{\Omega }\omega _0\right)^2}}\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}`$
The same quantity, expressed in terms of coordinate times, is
$$\delta t=2\pi \left(\frac{1}{\omega _0+\omega _{}\mathrm{\Omega }}\frac{1}{\mathrm{\Omega }+\omega _{}\omega _0}\right)=\frac{c\delta \tau }{\sqrt{g_{tt}+2g_{t\varphi }\mathrm{\Omega }+g_{\varphi \varphi }\mathrm{\Omega }^2}}$$
Finally it must be remarked that an inertial distant observer finds also a difference in revolution times between pairs of freely counter-rotating objects. One obtains
$$\delta T=2\pi \left(\frac{1}{\omega _0+\omega _{}}\frac{1}{\omega _0\omega _{}}\right)=4\pi \frac{\omega _{}}{\omega _{}^2\omega _0^2}$$
(10)
An interesting category of observers are the so called locally non rotating observers (LNRO) or Bardeen observers . An LNRO is an observer who does not rotate with respect to matter radially falling towards the central mass; when the latter is spinning it drags, in a sense, the space time around it (though the image of a ”drag” is not really appropriate as pointed out in ), so that a locally ”non rotating” observer is actually seen as rotating from another inertial far away observer (distant stars); its angular velocity is $`\mathrm{\Omega }_{LNRO}=g_{t\varphi }/g_{\varphi \varphi }`$ and its motion is in general non geodesic. In fig. 2 $`\mathrm{\Omega }_{LNRO}`$ is a measure of the slope of the cylinder.
Another situation that could be of importance for experimentation is the one of a rotating observer (angular speed $`\mathrm{\Omega }`$) who sends with opposite but locally equal velocities (in the tangent space) two objects along his own path. If $`𝐮_o`$, $`𝐮_1`$ and $`𝐮_2`$ are the fourvelocities of the observer and the two objects the condition for the equality of the velocities with respect to the observer is
$$𝐮_o𝐮_1=𝐮_o𝐮_2$$
(11)
Using (1) and the normalization condition for the fourvelocity of the observer (11) transforms into
$$\frac{g_{tt}+2g_{t\varphi }\omega _2+g_{\varphi \varphi }\omega _2^2}{g_{tt}+2g_{t\varphi }\omega _1+g_{\varphi \varphi }\omega _1^2}=\left(\frac{g_{tt}+\left(g_{t\varphi }+g_{\varphi \varphi \mathrm{\Omega }}\right)\omega _2+g_{t\varphi }\mathrm{\Omega }}{g_{tt}+\left(g_{t\varphi }+g_{\varphi \varphi \mathrm{\Omega }}\right)\omega _1+g_{t\varphi }\mathrm{\Omega }}\right)^2$$
Solving for $`\omega _2`$ one finds, besides the trivial solution $`\omega _2=\omega _1`$, the relevant result
$$\omega _2=\frac{g_{tt}\omega _12g_{tt}\mathrm{\Omega }g_{\varphi \varphi }\mathrm{\Omega }^2\omega _12g_{t\varphi }\mathrm{\Omega }^2}{g_{tt}+2g_{\varphi \varphi }\mathrm{\Omega }\omega _1g_{\varphi \varphi }\mathrm{\Omega }^2+2g_{t\varphi }\omega _1}$$
(12)
Combining (12) with (2.1) it is possible to find the time delay registered by the observer at the passing by him of the two apparently equal velocity objects. The result is
$$\delta \tau =\frac{4\pi }{c}\frac{g_{\varphi \varphi }\mathrm{\Omega }+g_{t\varphi }}{\sqrt{g_{\varphi \varphi }\mathrm{\Omega }^2+2\mathrm{\Omega }g_{t\varphi }+g_{tt}}}$$
(13)
As it can be seen $`\delta \tau `$ does not depend on $`\omega _1`$ i.e. it is independent from the actual velocity of the objects with respect to the observer. In Minkowski space time (13) reproduces the formula of the Sagnac effect.
## 3 Special cases
The time delays and synchrony defects determined in the preceding section may be specialized to various different metric tensors. Two cases are particularly of interest either in principle or for practical reasons: the Kerr metric and the weak field approximation of the metric of a spinning object. The relevant metric elements ($`\theta =\pi /2`$) are:
$$\begin{array}{ccc}\text{Metric elements}\hfill & \text{Kerr}\hfill & \text{Weak field}\hfill \\ g_{tt}\hfill & c^2\left(12G\frac{M}{c^2r}\right)\hfill & c^2\left(12G\frac{M}{c^2r}\right)\hfill \\ g_{\varphi \varphi }\hfill & 2a^2G\frac{M}{c^2r}r^2a^2\hfill & r^2\hfill \\ g_{t\varphi }\hfill & 2aG\frac{M}{cr}\hfill & 2aG\frac{M}{cr}\hfill \end{array}$$
Introducing these expressions into the formulas of the preceding section we obtain
$`\omega _{0K}`$ $`=`$ $`{\displaystyle \frac{aGMc}{a^2GMc^2r^3}}`$
$`\omega _{0wf}`$ $`=`$ $`aG{\displaystyle \frac{M}{cr^3}}`$ (14)
and
$$\{\begin{array}{c}\omega _K=c\frac{\sqrt{GMc^2r^3}}{a^2GMc^2r^3}\hfill \\ \omega _{wf}=\sqrt{\frac{GM}{r^3}\left(1+\frac{a^2GM}{c^2r^3}\right)}\hfill \end{array}$$
(15)
Combining these formulas we obtain
$$\{\begin{array}{c}\omega _{1,2K}=\frac{c}{a\pm c\sqrt{\frac{r^3}{GM}}}\hfill \\ \omega _{1,2wf}=\sqrt{\frac{GM}{r^3}}\left(\pm 1\frac{a}{c}\sqrt{\frac{GM}{r^3}}\right)\hfill \end{array}$$
(16)
It is clearly $`\omega _2>\omega _1`$ consequently the synchrony defect between two counter-orbiting objects is obtained from (5) and (4). Let us directly calculate the result in weak field approximation, keeping the first order (in $`a`$ and $`GM/c^2`$) terms only:
$$\delta \tau _{12}6\pi \frac{GM}{c^2r}\frac{a}{c}$$
(17)
From the view point of a distant inertial observer the difference in revolution times for the two counter-orbiting objects is obtained from (10) and (16):
$$\delta T=4\pi \frac{a}{c}$$
(18)
This exact result (in Kerr geometry) is remarkably independent both from the $`r`$ parameter of the orbit and from the gravity constant $`G`$. Actually in weak field approximation and for a spherical homogeneous mass it turns also to be independent from the very mass of the central object since then it is: $`a=2R^2\mathrm{\Omega }_0/5c`$ ($`R`$ is the radius of the body and $`\mathrm{\Omega }_0`$ is its rotation speed).
From (2.1) we see that the readings of clocks attached to our two objects after what is seen as a complete revolution by the distant observer are
$$\tau _{1,2}=\frac{2\pi }{c\omega _{1,2}}\sqrt{g_{tt}+2g_{t\varphi }\omega _{1,2}+g_{\varphi \varphi }\omega _{1,2}^2}$$
The difference between these readings corresponds, in weak field approximation, to (18): $`\tau _1\tau _24\pi \frac{a}{c}`$. It is also the value that would be found, with the same approximation, by a LNRO.
In WFA formula (13) becomes
$`\delta \tau `$ $``$ $`{\displaystyle \frac{4\pi }{c^2}}{\displaystyle \frac{r^2\mathrm{\Omega }2aG\frac{M}{cr}}{\sqrt{12G\frac{M}{c^2r}\frac{r^2\mathrm{\Omega }^2}{c^2}+4\mathrm{\Omega }\frac{a}{c}G\frac{M}{c^2r}}}}`$ (19)
$``$ $`{\displaystyle \frac{4\pi }{c^2}}r^2\mathrm{\Omega }\left(1+G{\displaystyle \frac{M}{c^2r}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{r^2\mathrm{\Omega }^2}{c^2}}\right)8\pi {\displaystyle \frac{a}{c}}{\displaystyle \frac{GM}{c^2r}}`$
There is a contribution to $`\delta \tau `$ depending on $`a`$ but not on $`\mathrm{\Omega }`$ and reproducing the result obtained as first order relativistic correction to the Sagnac effect .
## 4 Conclusion
We have shown how a simple geometric vision of the world lines of steadily rotating objects in axisymmetric metrics endowed with angular momentum allows for a description and explanation of the GCE. The method evidences and recovers some interesting results that could lead to experimental verifications. Using satellites on circular trajectories one would find a synchrony defect between counter-orbiting identical clocks given (in WFA) by (17), wich in the case of Earth is $`10^{16}`$s. A much bigger effect is seen considering revolution times with respect to a fixed direction in space (with respect to fixed stars); in that case the two rotation directions correspond to differences in period length given exactly by (18) which for the Earth is $`10^7`$s.
Another interesting possibility would be to work with a satellite (the observer) sending light signals in opposite directions along (non geodesic) closed paths; the relevant formula in this case would be (19) with $`\mathrm{\Omega }`$ given by (16), which produces
$$\delta \tau \pm 4\frac{\pi }{c^2}r^2\sqrt{G\frac{M}{r^3}}12\pi \frac{a}{c}\frac{GM}{c^2r}$$
(20)
The upper (lower) sign corresponds to a prograde (retrograde) orbiting observer. The first term in (20) is the Sagnac effect, the second one is the correction induced by the angular momentum of the source of gravity. Again, in the case of the Earth, the correction is in the order of $`10^{16}`$s (a hundredth of a period for visible light).
Finally, in case of experiments on the surface of the Earth (non geodesic equatorial observer, equal speed objects/signals in opposite directions) the formula is again (19) with $`\mathrm{\Omega }`$ coinciding with the angular speed of the Earth $`\mathrm{\Omega }_0`$. Using the expression of $`a`$ appropriate for this case, the formula reads
$$\delta \tau \frac{4\pi }{c^2}R^2\mathrm{\Omega }_0\left(1+\frac{1}{5}G\frac{M}{c^2R}+\frac{1}{2}\frac{R^2\mathrm{\Omega }_0^2}{c^2}\right)$$
The ”correction” originated by the angular momentum of the planet is still in the order of $`10^{16}`$s. |
warning/0001/hep-ph0001273.html | ar5iv | text | # 1 Introduction
## 1 Introduction
With the Giga–$`Z`$ option of the Tesla project one may expect the production of about $`10^9`$ $`Z`$ bosons at resonance . This huge rate, about a factor 100 higher than rates at LEP 1, allows one to study a number of problems with unprecedented precision. Among them is the search for lepton-flavour changes in $`Z`$ decays:
$`Ze\mu ,\mu \tau ,e\tau .`$ (1.1)
Non-zero rates are expected if neutrinos are massive and mix . Often one considers the branching ratio for the production of the following sum of charged states:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )={\displaystyle \frac{\mathrm{\Gamma }(Z\overline{l}_1l_2+l_1\overline{l}_2)}{\mathrm{\Gamma }_Z}}.`$ (1.2)
First predictions for flavour-changing $`Z`$ decays in the framework of the Standard Model , using techniques developed in , were given in . The best direct limits are obtained by searches at LEP 1 (95% c.l.) :
$`\mathrm{BR}(Ze^{}\mu ^\pm )`$ $`<`$ $`1.7\times 10^6\text{[14]},`$ (1.3)
$`\mathrm{BR}(Ze^{}\tau ^\pm )`$ $`<`$ $`9.8\times 10^6\text{[15, 14]},`$ (1.4)
$`\mathrm{BR}(Z\mu ^{}\tau ^\pm )`$ $`<`$ $`1.2\times 10^5\text{[14, 16]}.`$ (1.5)
A careful analysis shows, taking into account realistic conditions at future experiments, that the sensitivities for the branching ratios could be improved considerably at the Giga–$`Z`$ , namely down to:
$`\mathrm{BR}(Ze^{}\mu ^\pm )`$ $`<`$ $`2\times 10^9,`$ (1.6)
$`\mathrm{BR}(Ze^{}\tau ^\pm )`$ $`<`$ $`f\times 6.5\times 10^8,`$ (1.7)
$`\mathrm{BR}(Z\mu ^{}\tau ^\pm )`$ $`<`$ $`f\times 2.2\times 10^8,`$ (1.8)
with $`f=0.2÷1.0`$.
These numbers may be confronted with expectations derived from the signals for $`\nu _\mu \nu _\tau `$ oscillations in atmospheric neutrino experiments . They are at the 90% c.l. compatible with the following parameter set :
$`\mathrm{\Delta }m_{\nu _\mu \nu _\tau }^2`$ $``$ $`(2÷8)\times 10^3\mathrm{eV}^2,`$ (1.9)
$`\mathrm{sin}^2(2\vartheta _{\mu \tau })`$ $``$ $`0.8÷1.`$ (1.10)
There is also evidence for $`\nu _e\nu _\mu `$ oscillations from solar neutrino experiments , being compatible with:
$`\mathrm{\Delta }m_{\nu _e\nu _\mu }^2`$ $``$ $`10^{10}÷10^5\mathrm{eV}^2.`$ (1.11)
From reactor searches, there are no hints of $`\nu _e\nu _\tau `$ oscillations . For more details see e.g. the review and references therein.
The good news from the evidences for neutrino oscillations is that they suggest non-vanishing rates for reaction (1.1). The bad news is, that these rates are, if derived with (1.9)–(1.11) in the minimally extended Standard Model ($`\nu `$SM), extremely small<sup>1</sup><sup>1</sup>1Our estimate is in clear distinction to Eqn. (6) of , where from the data a limit was derived which corresponds to $`\mathrm{BR}(Z\mu ^{}\tau ^\pm )𝒪(10^8÷10^5)`$. :
$`\mathrm{BR}(Z\mu ^{}\tau ^\pm )`$ $``$ $`10^{54},`$ (1.12)
$`\mathrm{BR}(Ze^{}\mu ^\pm )\mathrm{BR}(Ze^{}\tau ^\pm )`$ $`<`$ \[-0.07cm\] $``$ $`4\times 10^{60}.`$ (1.13)
This is derived in Section 2 and is also in accordance with older calculations .
How do these small numbers arise? In Born approximation, the lepton-flavour changing $`Z`$ decay into two charged leptons is forbidden in the $`\nu `$SM due to the GIM mechanism . However, it may take place if $`n`$ types of neutrinos have masses $`m_i`$ and mix with each other (with mixing matrix $`𝐕_{ij}`$), i.e. if symmetry eigenstates and mass eigenstates are different in the lepton sector. Then, the virtual exchange of these neutrinos produces an effective lepton-flavour changing vertex and the corresponding branching ratio has the following structure:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )`$ $`=`$ $`{\displaystyle \frac{\alpha ^3}{192\pi ^2s_W^6c_W^2}}{\displaystyle \frac{M_Z}{\mathrm{\Gamma }_Z}}|𝒱(M_Z^2)|^210^6|𝒱(M_Z^2)|^2.`$ (1.14)
The form factor $`𝒱(Q^2)`$ depends on the details of the interaction:
$`𝒱(Q^2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}𝐕_{l_1i}𝐕_{l_2i}^{}V(m_i^2/M_W^2).`$ (1.15)
The vertex function $`V`$ was calculated in 1982/83 independently by three groups for sequential Dirac particles, and in a more general context later .
The function $`V`$ depends quadratically on the mass both in the small neutrino mass limit and in the large neutrino mass limit :
$`V(\lambda _i1)V(0)`$ $``$ $`(2.562.30\times i)\lambda _i+𝒪\left(\lambda _i^2\mathrm{ln}\lambda _i\right),`$ (1.16)
$`V(\lambda _i1)V(0)`$ $``$ $`{\displaystyle \frac{1}{2}}\lambda _i+𝒪(\mathrm{ln}\lambda _i),`$ (1.17)
with
$`\lambda _i`$ $`=`$ $`{\displaystyle \frac{m_i^2}{M_W^2}}.`$ (1.18)
The constant terms are not shown in (1.16), (1.17) since they drop out in (1.15) due to the unitarity of the mixing matrix, and the branching ratio becomes proportional to the fourth power of the neutrino massses. It is this behaviour which makes the expected rates so extremely small for the experimentally evidenced tiny neutrino masses. For values of $`m_i`$ of the order of $`M_W`$, the vertex $`V`$ is of the order one, and could become large if the $`m_i`$ would be much bigger than $`M_W`$.
The evidence of tiny neutrino masses may also be indicative for a mechanism which produces at the same time very large masses. Heavy neutrinos are expected by some GUTs and string-inspired models , and are suggested by the seesaw mechanism . Therefore, the above observations motivate us to have a closer look at the prospects of observing lepton-flavour changes with the Tesla Linear Collider. For the Giga–$`Z`$ option, it might not be sufficient to apply the well-known and simple approximations for large masses $`m_i`$, but also the medium- or even small-mass cases may be of experimental interest.
To be concrete, we will explore the following scenarios:
* The $`\nu `$SM. We treat the known light neutrinos ($`\nu _e,\nu _\mu ,\nu _\tau `$) as massive Dirac particles. Individual lepton numbers $`L_e,L_\mu ,L_\tau `$ are not conserved any more. The lepton sector is then in exact analogy to the quark sector. As a by-product, the $`Z`$ decay amplitude into two quarks of different flavours can be read off from our general expressions.
* The $`\nu `$SM sequentially extended with one heavy ordinary Dirac neutrino. This case implies the existence of a heavy charged lepton as well<sup>2</sup><sup>2</sup>2A fourth generation of quarks is also needed to keep the theory anomaly free.. It is not a very favoured scenario but we consider it as a simple application of the expressions of case (i) for heavier neutrinos. Again, total lepton number $`L`$ is conserved.
* The $`\nu `$SM extended with two heavy right-handed singlet Majorana neutrinos. Not only individual, but also total $`L`$ is, in general, not conserved since the presence of Majorana mass terms involves mixing of neutrinos and their charge-conjugate partners (antineutrinos), with opposite fermion-number. For two equal and heavy masses this case reduces to the addition of one heavy singlet Dirac neutrino . In this latter case $`L`$ is recovered .
With our numerical estimates we will be as model-independent as possible and will assume no constraints on neutrino masses or mixings, except for the ones imposed by the unitarity of the leptonic and neutrino mixing matrices, by the present bounds on lepton universality, CKM unitarity, and the measured $`Z`$ boson invisible width , and by oscillation experiments (see for a review).
In the following sections, we will discuss the predictions of the three scenarios for the lepton-number changing $`Z`$ decay. In Appendices, the generalization of Lagrangian and Feynman rules of the Standard Model to the general case of Dirac and Majorana masses is explained, experimental limits on neutrino mixings and masses are quoted, and the calculation of the vertex function is sketched.
## 2 Predictions from the $`\nu `$SM
The amplitude for the decay of a $`Z`$ boson into two charged leptons with different flavour, $`l_1`$ and $`l_2`$, is given in a self-explanatory notation by:
$`={\displaystyle \frac{ig\alpha _W}{16\pi c_W}}𝒱(Q^2)\epsilon _Z^\mu \overline{u}_{l_2}(p_2)\gamma _\mu (1\gamma _5)u_{l_1}(p_1),`$ (2.1)
where
$`\alpha _W{\displaystyle \frac{\alpha }{s_W^2}},`$ (2.2)
and the form factor $`𝒱`$ depends on $`Q^2=(p_2p_1)^2`$ and can be written as:
$`𝒱(Q^2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}𝐕_{l_1i}𝐕_{l_2i}^{^{}}V(\lambda _i),`$ (2.3)
$`V(\lambda _i)`$ $`=`$ $`\left[v_W(i)+v_{WW}(i)+v_\varphi (i)+v_{\varphi \varphi }(i)+v_{W\varphi }(i)+v_\mathrm{\Sigma }(i)\right],`$ (2.4)
with $`𝐕_{ij}`$ being the leptonic CKM mixing matrix. In general, there are besides the vector and axial-vector couplings $`f_V`$ and $`f_A`$ in (2.1) also contributions of the $`f_S,f_P,f_M,f_E`$ types, but for the production of on-shell fermions (with their masses being neglected) they vanish here. Further, it is $`f_V=f_A=𝒱(Q^2)`$ due to the presence of $`W`$ bosons coupling only to left-handed fermions. The contributions from the one-loop diagrams of Figure 1 depend on $`\lambda _i`$ and additionally on
$`\lambda _Q`$ $`=`$ $`{\displaystyle \frac{Q^2}{M_W^2}},`$ (2.5)
which on the $`Z`$ boson mass shell becomes
$`\lambda _Z`$ $`=`$ $`{\displaystyle \frac{M_Z^2}{M_W^2}}={\displaystyle \frac{1}{c_W^2}}1.286.`$ (2.6)
In terms of the usual vector and axial-vector couplings,
$`v_i`$ $`=`$ $`I_3^{i_L}2Q_is_W^2=I_3^{i_L}(14s_W^2|Q_i|),`$ (2.7)
$`a_i`$ $`=`$ $`I_3^{i_L},`$ (2.8)
the individual contributions are:
* From vertex diagrams:
$`\mathrm{D1}:v_W(i)`$ $`=`$ $`(v_i+a_i)\left[\lambda _Q(C_0+C_{11}+C_{12}+C_{23})2C_{24}+1\right]`$ (2.9)
$`(v_ia_i)\lambda _iC_0,`$
$`\mathrm{D2}:v_{WW}(i)`$ $`=`$ $`2c_W^2(2I_3^{i_L})\left[\lambda _Q(\overline{C}_{11}+\overline{C}_{12}+\overline{C}_{23})6\overline{C}_{24}+1\right],`$ (2.10)
$`\mathrm{D3}:v_\varphi (i)`$ $`=`$ $`(v_i+a_i){\displaystyle \frac{\lambda _i^2}{2}}C_0`$ (2.11)
$`(v_ia_i){\displaystyle \frac{\lambda _i}{2}}\left[\lambda _QC_{23}2C_{24}+{\displaystyle \frac{1}{2}}\right],`$
$`\mathrm{D4}:v_{\varphi \varphi }(i)`$ $`=`$ $`(12s_W^2)(2I_3^{i_L})\lambda _i\overline{C}_{24},`$ (2.12)
$`\mathrm{D5}:v_{W\varphi }(i)`$ $`=`$ $`2s_W^2(2I_3^{i_L})\lambda _i\overline{C}_0;`$ (2.13)
* From self-energy corrections to the external fermion lines:
$`\mathrm{D}\mathrm{\Sigma }:v_\mathrm{\Sigma }(i)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(v_i+a_i4c_W^2a_i)\left[(2+\lambda _i)B_1+1\right].`$ (2.14)
The one-loop tensor integrals $`C_0`$, $`\overline{C}_0`$, $`C_{ij}`$, $`\overline{C}_{ij}`$, and $`B_1`$ are defined in Appendix C. With our numerical results for the Dirac case, we rely on two calculations, an old one and also this new, completely independent one. In the latter, the numerical evaluation of the tensor integrals is performed with the help of the computer program package LoopTools .
Numerical results are shown in Figure 2. The quantity presented is related to a branching ratio definition often used in the literature:
$`B_Z`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }(Zf_1^{}f_2^\pm )}{2\times \mathrm{\Gamma }(Z\overline{\nu }_l\nu _l)}}=\left({\displaystyle \frac{\alpha }{\pi }}\right)^2{\displaystyle \frac{N_c}{16s_W^4}}|𝒱(M_Z^2)|^2,`$ (2.15)
with $`N_c`$ as colour factor. Using
$`\mathrm{\Gamma }(Z\overline{\nu }_l\nu _l)`$ $`=`$ $`{\displaystyle \frac{\alpha _W}{24c_W^2}}M_Z,`$ (2.16)
we get a useful relation to the branching ratio introduced in (1.2) and (1.14):
$`\mathrm{BR}(Zf_1^{}f_2^\pm )`$ $`=`$ $`{\displaystyle \frac{2\times \mathrm{\Gamma }(Z\overline{\nu }_l\nu _l)}{\mathrm{\Gamma }_Z}}B_Z=0.1333B_Z.`$ (2.17)
Figure 2 shows the contribution from one neutrino generation to the branching ratio as a function of the neutrino mass without the influence of the mixing matrix elements. We choose two interesting mass regions, namely that corresponding to the findings of neutrino oscillation searches (Figure 2(a)) and also that with potential predictions in the reach of the Giga–$`Z`$ (Figure 2(b)). The latter figure nicely agrees with the earlier calculations , and the small mass limit with . The dotted lines in Figure 2(b) correspond to the approximation $`\lambda _Z=0`$: the large mass limit is reproduced quite well, while the small mass limit differs from the correct result. This is discussed in Appendix D.3.
We will now estimate the branching ratio under the assumption that there are three generations of light neutrino flavours with a unitary mixing matrix $`𝐕`$ as evidenced by experiment. The general form of this matrix may be chosen to be :
$`𝐕_{l\nu }`$ $`=`$ $`\left(\begin{array}{ccc}𝐕_{e1}& 𝐕_{e2}& 𝐕_{e3}\\ 𝐕_{\mu 1}& 𝐕_{\mu 2}& 𝐕_{\mu 3}\\ 𝐕_{\tau 1}& 𝐕_{\tau 2}& 𝐕_{\tau 3}\end{array}\right)`$
$`=`$ $`\left(\begin{array}{ccc}c_{12}c_{13}& c_{13}s_{12}& s_{13}\\ c_{23}s_{12}e^{i\delta }c_{12}s_{13}s_{23}& c_{12}c_{23}e^{i\delta }s_{12}s_{13}s_{23}& c_{13}s_{23}\\ s_{23}s_{12}e^{i\delta }c_{12}c_{23}s_{13}& c_{12}s_{23}e^{i\delta }c_{23}s_{12}s_{13}& c_{13}c_{23}\end{array}\right)\left(\begin{array}{ccc}e^{i\alpha }& 0& 0\\ 0& e^{i\beta }& 0\\ 0& 0& 1\end{array}\right)`$
Here, we have three mixing angles and one CP-violating phase as in the quark CKM case, plus two CP-violating phases $`\alpha ,\beta `$ if neutrinos are Majorana particles (they are strictly neutral so that less phase factors may be ‘eaten’ by redefining complex fermion fields). Current data suggest the following form of this matrix:
$`𝐕_{l\nu }=\left(\begin{array}{ccc}c_{12}& s_{12}& 0\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{2}}s_{12}& \frac{1}{\sqrt{2}}c_{12}& \frac{1}{\sqrt{2}}\end{array}\right),`$ (2.19)
where we have assumed $`\alpha =\beta =\delta =0`$, extracted from the $`s_{13}=0`$, and further assumed $`s_{23}=1/\sqrt{2}`$ (corresponding to maximal mixing) and left the $`\theta _{12}`$ free.
With
$`M_W`$ $`=`$ $`80.41\mathrm{GeV},`$ (2.20)
$`M_Z`$ $`=`$ $`91.187\mathrm{GeV},`$ (2.21)
$`\mathrm{\Gamma }_Z=\mathrm{\Gamma }(Z\mathrm{all})`$ $`=`$ $`2.49\mathrm{GeV},`$ (2.22)
we get after trivial calculations:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )`$ $`=`$ $`{\displaystyle \frac{\alpha _W^3M_Z}{192\pi ^2c_W^2\mathrm{\Gamma }_Z}}|𝐕_{l_11}𝐕_{l_21}^{}[V(\lambda _1)V(0)]`$
$`+𝐕_{l_12}𝐕_{l_22}^{}[V(\lambda _2)V(0)]+𝐕_{l_13}𝐕_{l_23}^{}[V(\lambda _3)V(0)]|^2,`$
with
$`{\displaystyle \frac{\alpha _W^3M_Z}{192\pi ^2c_W^2\mathrm{\Gamma }_Z}}=1.127\times 10^6.`$ (2.24)
For small neutrino masses, we show in Section D.2:
$`V(\lambda _i)V(0)=a_1(\lambda _Z)\lambda _i.`$ (2.25)
The resulting branchings are, without approximations yet, but using the information from three-generation unitarity:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )`$ $`=`$ $`{\displaystyle \frac{\alpha _W^3M_Z}{192\pi ^2c_W^2\mathrm{\Gamma }_Z}}|a_1|^2\left|𝐕_{l_11}𝐕_{l_21}^{}\lambda _{12}𝐕_{l_13}𝐕_{l_23}^{}\lambda _{23}\right|^2,`$ (2.26)
with
$`\lambda _{ij}`$ $`=`$ $`|\lambda _i\lambda _j|`$ (2.27)
and the $`a_1`$ is given in (D.67):
$`|a_1|^2=11.832.`$ (2.28)
There are two different cases to be considered when using now the specific mixing matrix (2.19):
$`\mathrm{BR}(Ze^{}\mu ^\pm )\mathrm{BR}(Ze^{}\tau ^\pm )`$ $``$ $`1.333\times 10^5{\displaystyle \frac{c_{12}^2s_{12}^2}{2}}\lambda _{12}^2,`$ (2.29)
and
$`\mathrm{BR}(Z\mu ^{}\tau ^\pm )`$ $`=`$ $`1.333\times 10^5{\displaystyle \frac{1}{4}}\left|s_{12}^2\lambda _{12}\lambda _{23}\right|^2.`$ (2.30)
From the mass estimate (1.9) we get as additional input from atmospheric neutrino studies:
$`\lambda _{23}{\displaystyle \frac{(2÷8)\times 10^3\mathrm{eV}^2}{M_W^2}}=(3÷12)\times 10^{25},`$ (2.31)
and from solar neutrino searches, (1.11):
$`\lambda _{12}1.5\times \left(10^{32}÷10^{27}\right).`$ (2.32)
It is easy now to see that the expected rates for the lepton-flavour changing $`Z`$ decays are limited to:
$`\mathrm{BR}(Ze^{}\mu ^\pm )\mathrm{BR}(Ze^{}\tau ^\pm )`$ $`<`$ $`3.75\times 10^{60},`$ (2.33)
$`\mathrm{BR}(Z\mu ^{}\tau ^\pm )`$ $``$ $`(3÷48)\times 10^{55}.`$ (2.34)
In (2.33), we assumed arbitrarily a maximal mixing $`s_{12}=1/\sqrt{2}`$. These rates are extremely small. In fact, we will neglect the effects of the light neutrino sector in the next sections, where we extend the $`\nu `$SM to accommodate heavy neutrinos, taking massless the known ones.
## 3 Predictions from the $`\nu `$SM plus One Heavy Dirac Neutrino
Here, the only effective difference to the case before is the existence of a fourth generation with a sequential Dirac neutrino of mass $`m_N`$. In this case, the branching ratio gets:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )`$ $`=`$ $`{\displaystyle \frac{\alpha _W^3M_Z}{192\pi ^2c_W^2\mathrm{\Gamma }_Z}}\left|𝐕_{l_1N}𝐕_{l_2N}^{}\right|^2\left|V(\lambda _N)V(0)\right|^2,`$ (3.1)
with $`V(\lambda _N)`$ given in (2.4). The numerical results depend crucially on the mixing between the light and heavy leptons. An optimistic assumption would be maximal mixing:
$`\left|𝐕_{l_1N}𝐕_{l_2N}^{}\right|^2=\left({\displaystyle \frac{1}{\sqrt{2}}}\right)^4={\displaystyle \frac{1}{4}}.`$ (3.2)
This is of course unrealistic. Stringent (though indirect) limits may be derived from the analysis of flavour-diagonal reactions as advocated e.g. in as well as from the lepton flavour-changing process $`\mu e\gamma `$ . A short summary may be found in Appendix B. There the matrix $`𝐁`$ is, for this particular case of heavy Dirac neutrinos, $`𝐁=𝐕`$.
In order to be definite, we show in Figure 3a, solid line, the predictions for $`\mathrm{BR}(Z\mu ^\pm \tau ^{})`$ and assume the upper limit of the mixings allowed from (B.9)–(B.11) and (B.12):
$`\left|𝐕_{\mu N}𝐕_{\tau N}^{}\right|^2<1.5\times 10^4.`$ (3.3)
For the other two lepton flavour-changing $`Z`$ decay channels, the corresponding graphs scale simply in accordance with the ratios of the mixing matrix elements:
$`\left|𝐕_{eN}𝐕_{\tau N}^{}\right|^2`$ $`<`$ $`1.9\times 10^4,`$ (3.4)
$`\left|𝐕_{eN}𝐕_{\mu N}^{}\right|^2`$ $`<`$ $`1.2\times 10^4.`$ (3.5)
These coupling factors are to be compared to (3.2): we observe a suppression of the expected branching ratio (for a given mass $`m_N`$) by more than three orders of magnitude.
As mentioned in Appendix B, these bounds on light-heavy mixings from flavour-diagonal processes are improved by flavour-changing processes involving the first two lepton families. In fact, from
$`\mathrm{BR}(\mu e\gamma )<1.2\times 10^{11},`$ (3.6)
one obtains, for heavy neutrinos, the following (nearly) mass-independent limit :
$`\left|𝐕_{eN}𝐕_{\mu N}^{}\right|^2`$ $`<`$ $`1.4\times 10^8,`$ (3.7)
much more stringent than (3.5).
The Giga–$`Z`$ discovery range is indicated in Figure 3a, using the maximum values of the mixings allowed by (3.3). The range will be limited to the large neutrino mass limit if we believe in the relevance of the above mixing bounds. Then, the approximations of (D.43) apply. Numerically, this means for $`\lambda _Z=1.286`$:
$`\left|V(\lambda _N)V(0)\right|^2`$ $`=`$ $`{\displaystyle \frac{1}{4}}\lambda _N^2+1.44\lambda _N\mathrm{ln}\lambda _N3.49\lambda _N+2.07\mathrm{ln}^2\lambda _N+𝒪(\mathrm{ln}\lambda _N).`$
If neutrinos with a mass of several hundred GeV or more would exist, there is a good chance to observe some effect of them from the $`Z`$ decays under study. Indeed one cannot constrain heavy neutrino masses from these processes since they also depend on the mixings, and vice versa. The fact that in Figure 3a the discovery reach of LEP or Giga–$`Z`$ cuts the curves means: if no event is observed, mixings must be smaller than the ones employed (we took present upper bounds) for neutrino masses above the intersection point; or if some effect was observed then the plot would provide a lower bound for the heavy neutrino mass. In fact, it is roughly only the product $`m_N^4\left|𝐕_{l_1N}𝐕_{l_2N}^{}\right|^2`$ which can be constrained, assuming only one heavy Dirac neutrino.
## 4 Predictions from the $`\nu `$SM plus Two Heavy <br>Right-Handed Singlet Majorana Neutrinos
Some basic features of Lagrangians with Majorana mass terms and their relations to the simpler Dirac case are summarized in Appendix A. For a derivation of Feynman rules with Majorana particles we refer to . The couplings of the virtual neutrinos to the $`W`$ bosons are left-handed, $`v_i+a_i=1`$, $`v_ia_i=0`$, and $`2I_3^{i_L}=1`$, while those to the $`Z`$ boson and to the Higgs particles are non-diagonal and contain right-handed admixtures. This may be seen from (A.68)–(A.70). The amplitude for the decay (1.1) is again given by Eqn. (2.1), but this time the form factor $`𝒱_M`$ is non-diagonal not only in the external charged leptons, but also in the virtual neutrinos due to the $`Z\nu _i\nu _j`$ coupling, see (A.69):
$`𝒱_M(Q^2)`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{n_G+n_R}{}}}𝐁_{l_1i}𝐁_{l_2j}^{}V(i,j),`$ (4.1)
$`V(i,j)`$ $``$ $`V(\lambda _i,\lambda _j,𝐂_{ij})`$
$`=`$ $`\left[v_W(i,j)+\delta _{ij}v_{WW}(i)+v_\varphi (i,j)+\delta _{ij}v_{\varphi \varphi }(i)+\delta _{ij}v_{W\varphi }(i)+\delta _{ij}v_\mathrm{\Sigma }(i)\right].`$
The new non-diagonal terms arise in diagrams D1 and D3 of Figure 1:
$`\mathrm{D1}:v_W(i,j)`$ $`=`$ $`𝐂_{ij}\left[\lambda _Q(C_0+C_{11}+C_{12}+C_{23})2C_{24}+1\right]`$ (4.3)
$`+𝐂_{ij}^{}\sqrt{\lambda _i\lambda _j}C_0,`$
$`\mathrm{D3}:v_\varphi (i,j)`$ $`=`$ $`𝐂_{ij}{\displaystyle \frac{\lambda _i\lambda _j}{2}}C_0+𝐂_{ij}^{}{\displaystyle \frac{\sqrt{\lambda _i\lambda _j}}{2}}\left[\lambda _QC_{23}2C_{24}+{\displaystyle \frac{1}{2}}\right].`$ (4.4)
The matrices $`𝐁`$ and $`𝐂`$ are introduced in Appendix A.3. Both contributions are quite similar to (2.9) and (2.11) when there the right-handed couplings are retained. The other contributions are the same as in the ordinary Dirac case. The vertex reads:
$`𝒱_M(Q^2)`$ $`=`$ $`{\displaystyle \underset{i,j=1}{\overset{n_G+n_R}{}}}𝐁_{l_1i}𝐁_{l_2j}^{}\left[\delta _{ij}F(\lambda _i)+𝐂_{ij}G(\lambda _i,\lambda _j)+𝐂_{ij}^{}\sqrt{\lambda _i\lambda _j}H(\lambda _i,\lambda _j)\right],`$ (4.5)
to be compared to the case of $`n_H`$ heavy sequential Dirac neutrinos ($`𝐁=𝐕`$):
$`𝒱(Q^2)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n_G+n_H}{}}}𝐁_{l_1i}𝐁_{l_2i}^{}V(\lambda _i).`$ (4.6)
We now have to go into a specific model in order to make definite predictions. For the case chosen, namely the Standard Model extended by two heavy Majorana singlets, the sums over virtual neutrinos involve $`n_G`$=3 light and $`n_R`$=2 heavy neutrinos. The model is described in Appendix A.3. For some important properties of the mixing matrices B and C see (A.73)–(A.76). The light neutrino sector is practically massless. Using the properties (A.73)–(A.76), one can then write the form factor $`𝒱_M`$ in terms of the heavy-neutrino sector only. Actually, it is:
$`V(\lambda _i)=F(\lambda _i)+G(\lambda _i,\lambda _i).`$ (4.7)
Further, the functions $`F,G,H`$ are uniquely defined and from (A.73)–(A.76), and taking $`\lambda _i=0`$ for $`i=1,\mathrm{},n_G`$, it is straightforward to prove that:
$`\mathrm{BR}(Zl_1^{}l_2^\pm )`$ $`=`$ $`{\displaystyle \frac{\alpha _W^3M_Z}{192\pi ^2c_W^2\mathrm{\Gamma }_Z}}|{\displaystyle \underset{i,j=1}{\overset{n_R}{}}}𝐁_{l_1N_i}𝐁_{l_2N_j}^{}`$ (4.8)
$`\times \{\delta _{N_iN_j}[F(\lambda _{N_i})F(0)+G(\lambda _{N_i},0)+G(0,\lambda _{N_i})2G(0,0)]`$
$`+𝐂_{N_iN_j}[G(\lambda _{N_i},\lambda _{N_j})G(\lambda _{N_i},0)G(0,\lambda _{N_j})+G(0,0)]`$
$`+𝐂_{N_iN_j}^{}\sqrt{\lambda _{N_i}\lambda _{N_j}}H(\lambda _{N_i},\lambda _{N_j})\left\}\right|^2.`$
It is possible to express the couplings $`𝐁,𝐂`$ on the right-hand side by the mass ratio $`r=m_{N_2}^2/m_{N_1}^2`$ plus the three light-heavy mixings $`s_{\nu _e},s_{\nu _\mu },s_{\nu _\tau }`$. The relations are explicitly given in (A.77)–(A.81). The upper limits for these mixings (given in Appendix B.1) have been used to obtain the graphs in Figure 3(b). The mass ratio $`r`$ has been taken as a free parameter. Perturbative unitarity constrains the masses of the neutrinos so that they cannot be arbitrarily heavy. This is discussed in Appendix B.2. We see again from the figure that Giga–$`Z`$ has a discovery potential, preferentially in the large neutrino mass region. Similar curves can be obtained for $`Ze\mu `$ and $`Ze\tau `$.
The distinguished case of two Majorana singlet neutrinos with equal masses, forming effectively a singlet Dirac particle , results for $`r=1`$, see the solid line. This line has been taken over from Figure 3(a), there as the dashed line, in order to show that due to the different coupling structure the simple sequential Dirac neutrino case does not constitute a limiting case for large masses. The deviations are due to the terms from the non-diagonal elementary $`Z`$ couplings proportional to $`𝐂`$ and $`𝐂^{}`$ in (4.8). The $`𝐂^{}`$ terms drop out for $`r=1`$.
In contrast, predictions for Majorana and Dirac neutrinos approach each other in the limit of small masses. This may be seen from (4.8) using the unitarity relations (A.73)–(A.76) and the Taylor series expansion of the vertex function in powers of the neutrino mass (Appendix D.2). This phenomenon is just another example of what is called in the literature the “practical Dirac-Majorana confusion theorem” (see also the recent discussion in and references therein).
At the end of this section, we would like to comment on the literature for the reaction (1.1). The early papers did not include the more “realistic” and physically most interesting models with Majorana masses. These cases were studied in detail by and, in the context of left-right symmetric models, by . While we reproduce the large mass limit given there, we obtain slight deviations in the medium mass case as given in Eqn. (B1) of . Further, the Figures 7 and 8 in are not exactly reproducible in the intermediate mass range. However, the large mass limit seems to be the only potentially relevant case.
## 5 Summary
From our study of the decays $`Ze\mu ,e\tau ,\mu \tau `$, in the context of the Giga–$`Z`$ option of the Tesla linear collider project, we conclude:
* Neglecting the influence of the mixing angles, the expected branching ratios depend on the neutrino masses $`m_i`$ and are of the order $`\mathrm{BR}10^5(m_i/M_W)^4`$ both in the small and large mass limits.
* If there exist only the known three generations of ultra-light neutrinos, $`m_\nu M_W`$, there is absolutely no hope to see any effect.
* For heavier neutrinos with masses of the order of the weak scale, $`m_\nu M_W`$, one has to calculate the form factors describing the vertex without any approximation. The necessary exact formulas have also been given.
* The light-heavy mixing angles are not so strongly restricted as with the minimal (one-family) seesaw mechanism when interfamily seesaw type models are considered. But, unfortunately, these mixings have already been constrained to be very small, so that the lepton flavour-violating reactions under study are beyond the discovery reach of the Giga–$`Z`$ for $`m_\nu M_W`$ .
* Since, ignoring the light-heavy mixings, the branching ratios are proportional to the fourth power of the neutrino masses, there is a discovery potential in the large mass case, $`m_\nu M_W`$. Nevertheless, they are limited in practice by potentially small mixing factors (bound independently of the heavy neutrino masses by other experiments) and upper neutrino mass bounds from unitarity considerations.
* In fact, there is an interplay between heavy masses and light-heavy mixings: the mixings must be small for very large neutrino masses, since otherwise the scattering matrix elements would grow above the unitarity limit.
Summarizing, the Giga–$`Z`$ offers nearly three orders of magnitude gain of sensitivity compared to LEP 1. This opens quite interesting opportunities to search for lepton-number changing processes, if there exist heavy neutrinos sufficiently mixing with the light sector, within a quite broad allowed region according to present limits.
## Acknowledgements
We would like to thank J. Gluza, A. Pilaftsis, R. Rückl, and G. Wilson for helpful discussions. We also wish to acknowledge J. Gluza and G. Wilson for carefully reading the manuscript. The work of J.I. has been partially supported by the Spanish CICYT and Junta de Andalucía, under contracts AEN96-1672 and FQM101, respectively.
## Appendix A Lagrangians and Feynman Rules
We make extensive use of the notion of Majorana particles, since in general neutrinos may be of this type, and in fact in GUTs often exactly this happens. Majorana particles are neutral fermions $`\psi `$, fulfilling
$`\psi ^c=\psi ,`$ (A.1)
where $`\psi ^c`$ is the charge-conjugate of $`\psi `$.
Some introduction on notations are given in Appendix A.1. We observe experimental evidences for ultra-light neutrinos in the known fermion families on the one hand, and on the other there are unifying theories with potentially ultra-heavy Majorana neutrinos. That both phenomena might be related will be made plausible by a toy example in Appendix A.2. A more general ansatz for Majorana mass terms is finally given in Appendix A.3.
### A.1 Dirac neutrinos rewritten
The mass-term Lagrangian for Dirac neutrinos is:
$`_D`$ $`=`$ $`m_D\left(\overline{\nu _L}\nu _R+\overline{\nu _R}\nu _L\right){\displaystyle \frac{1}{2}}\left(\overline{\chi ^0}\right)\left(\begin{array}{cc}0& m_D\\ m_D& 0\end{array}\right)\left(\chi ^0\right),`$ (A.2)
where we introduce (self-conjugate) Majorana fields $`(\chi ^0)`$
$`\left(\chi ^0\right)`$ $`=`$ $`\left(\begin{array}{c}\nu _L+\nu _L^c\\ \nu _R+\nu _R^c\end{array}\right),`$ (A.3)
with
$`\nu ^cC\overline{\nu }^T`$ (A.4)
and $`C`$ being the charge-conjugation matrix. The mass matrix can be brought into a diagonal form in the basis of $`(\chi )`$:
$`_D`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\overline{\chi }\right)\left(\begin{array}{cc}m_D& 0\\ 0& m_D\end{array}\right)\left(\chi \right)={\displaystyle \frac{1}{2}}m_D\left(\overline{\chi }_1\chi _1+\overline{\chi }_2\chi _2\right)`$ (A.5)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(\overline{\xi }\right)\left(\begin{array}{cc}m_D& 0\\ 0& m_D\end{array}\right)\left(\xi \right)={\displaystyle \frac{1}{2}}m_D\left(\overline{\xi }_1\xi _1+\overline{\xi }_2\xi _2\right).`$
In the last step, the field component $`\xi _1`$ is introduced as a chiral transform of $`\chi _1`$ in order to make the positive mass eigenvalue explicit. The fields $`(\chi )`$ and $`(\chi ^0)`$ are related by the unitary matrix $`𝐔`$, and $`\xi _1`$ additionally by $`\gamma _5`$:
$`(\chi ^0)=𝐔(\chi );𝐔`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cc}\hfill 1& \hfill 1\\ \hfill 1& \hfill 1\end{array}\right),`$ (A.8)
$`\xi _1=\gamma _5\chi _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\nu _L+\nu _L^c\nu _R+\nu _R^c\right),`$ (A.9)
$`\xi _2=\chi _2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(\nu _L+\nu _L^c+\nu _R+\nu _R^c\right).`$ (A.10)
With the above chain of equations we have shown that one ordinary Dirac neutrino is equivalent to two Majorana neutrinos of equal mass and opposite CP parities. Evidently, if the mass matrix is not of type (A.2), then true Majorana particles are realized.
### A.2 The seesaw mechanism
The seesaw mechanism allows to understand the lightness of the known neutrinos by the introduction of heavier ones. However, it is not the only possible solution to this puzzle; see e.g. for a nice discussion.
In Appendix A.1 it was shown that a Dirac-neutrino mass term may be written as symmetric 2$`\times `$2 matrix with only off-diagonal elements. Consider now the case where only one right-handed singlet neutrino is added to the ordinary left-handed doublets of the SM. Let’s take one generation of left-handed light neutrinos for simplicity and with no loss of generality. The corresponding mass matrix is in general
$`𝐌=\left(\begin{array}{cc}m_L& m_D\\ m_D& m_R\end{array}\right)`$ (A.13)
and can be diagonalized by
$`𝐔`$ $`=`$ $`\left(\begin{array}{cc}\hfill \mathrm{cos}\theta _\nu & \hfill \mathrm{sin}\theta _\nu \\ \hfill \mathrm{sin}\theta _\nu & \hfill \mathrm{cos}\theta _\nu \end{array}\right),`$ (A.16)
with
$`\mathrm{tan}2\theta _\nu `$ $`=`$ $`{\displaystyle \frac{2m_D}{m_Rm_L}},\mathrm{cos}2\theta ={\displaystyle \frac{m_Rm_L}{\sqrt{(m_Rm_L)^2+4m_D^2}}}`$ (A.17)
yielding two eigenstates with different masses
$`m_\nu ,m_N`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{m_R+m_L\left[\left(m_Rm_L\right)^2+4m_D^2\right]^{1/2}\right\}.`$ (A.18)
That is, the most general mass term of one four-component self-conjugate field describes two Majorana particles with different masses. There is a particular configuration that corresponds to one Dirac neutrino, as shown in Appendix A.1.
In the SM, $`m_L=0`$ since the SM Higgs sector consists of a Higgs doublet (see Appendix A.3). Take now $`m_Rm_D`$. Then, the two physical states<sup>3</sup><sup>3</sup>3 Actually, one of the mass eigenvalues in (A.18) is negative. This is not a problem since then the true mass eigenstate is a chiral transform of the original field which has a mass term of opposite sign as in (A.10) . are a light and a heavy neutrino with masses
$`m_\nu `$ $``$ $`m_D^2/m_R,`$ (A.19)
$`m_N`$ $``$ $`m_R,`$ (A.20)
and the light-heavy mixing angle is:
$`s_\nu \mathrm{sin}\theta _\nu m_D/m_R\sqrt{m_\nu /m_N}.`$ (A.21)
This is the minimal seesaw mechanism: the larger the heavy scale is, the smaller the light neutrino masses are, and, at the same time, the smaller the mixing becomes.
Actually, by taking mN
>
[-0.07cm]
MZsubscript𝑚𝑁
>
[-0.07cm]
subscript𝑀𝑍m_{N}\raisebox{-3.69899pt}{~{}\shortstack{$>$ \\
[-0.07cm] $\sim$}}~{}M_{Z} and the present bounds on the light neutrino masses ,
$`m_{\nu _e}`$ $`<`$ \[-0.07cm\] $``$ $`2.5\text{eV},`$ (A.22)
$`m_{\nu _\mu }`$ $`<`$ \[-0.07cm\] $``$ $`160\text{keV},`$ (A.23)
$`m_{\nu _\tau }`$ $`<`$ \[-0.07cm\] $``$ $`15\text{MeV},`$ (A.24)
we get from (A.17) the light-heavy mixing angles:
$`s_{\nu _e}^2`$ $`<`$ \[-0.07cm\] $``$ $`3\times 10^{11},`$ (A.25)
$`s_{\nu _\mu }^2`$ $`<`$ \[-0.07cm\] $``$ $`2\times 10^6,`$ (A.26)
$`s_{\nu _\tau }^2`$ $`<`$ \[-0.07cm\] $``$ $`2\times 10^4.`$ (A.27)
These mixing angles are too small to be constrained by the experimental LEP and low-energy limits (B.6)–(B.8) .
Nevertheless, one might have $`m_L0`$, by introducing a Higgs triplet in the SM, <sup>4</sup><sup>4</sup>4The introduction of a Higgs triplet affects the value of the parameter $`\rho =M_W^2/(M_Z^2c_W^2)`$. and hope for a conspiracy between $`m_L`$ and $`m_R`$ to get light-heavy mixings of order one. This does not seem very reasonable. Thus, the minimal seesaw apparently lacks of phenomenological interest.
### A.3 Majorana neutrinos
In the interfamily seesaw-type models, several right-handed neutrinos are introduced and large light-heavy mixings can be obtained even for $`m_L=0`$ .
Therefore following , we now introduce in our study two heavy right-handed singlet neutrinos and treat the light-heavy mixings $`s_{\nu _l}`$ and the heavy masses $`m_{N_1},m_{N_2}`$ as free phenomenological parameters.
Consider an extension of the SM that incorporates $`n_R`$ right-handed neutrinos $`\nu _R`$ (singlets under SU(2)$``$U(1)) to the already present set of $`n_G=3`$ generations of left-handed neutrino doublets. We keep the Higgs sector untouched.
It is convenient to arrange all the independent neutrino degrees of freedom into two vectors of left-handed fields:
$`(\nu ^0)`$ $``$ $`(\nu _e,\nu _\mu ,\nu _\tau ),`$ (A.28)
$`(N^0)`$ $``$ $`(\nu _R^c)_ii=1,\mathrm{},n_R.`$ (A.29)
The charge-conjugates of right-handed neutrinos have been introduced:
$`\nu _R^cC\overline{\nu _R}^T.`$ (A.30)
The conjugate field has chirality and lepton/fermion numbers opposite to the original field (e.g. the $`\nu _R^c`$’s are left-handed).
In the basis
$`(n^0)=(\nu ^0,N^0)`$ (A.31)
the most general mass-term Lagrangian reads:
$`_M`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\overline{n_{}^{0}{}_{}{}^{c}})𝐌(n^0)+\text{h.c.}={\displaystyle \frac{1}{2}}(n^0)^TC𝐌(n^0)+\text{h.c.},`$ (A.32)
where the Majorana fields
$`(\chi ^0)(n^0)+(n_{}^{0}{}_{}{}^{c}),`$ (A.33)
containing both chiralities, have been introduced.
The mass matrix $`𝐌`$ is symmetric and can be written in a block form:
$`𝐌=\left(\begin{array}{cc}𝐦_L& 𝐦_D^T\\ 𝐦_D& 𝐦_R\end{array}\right).`$ (A.36)
It has dimension $`n_G+n_R`$ and can be diagonalized by a unitary matrix $`𝐔`$ of the same dimension,
$`\widehat{𝐌}=𝐔^T\mathrm{𝐌𝐔}=(𝐔^{})^{}\mathrm{𝐌𝐔}`$ $`=`$ $`\mathrm{𝐝𝐢𝐚𝐠}(m_{\nu _e},m_{\nu _\mu },m_{\nu _\tau },m_{N_1},m_{N_2})`$ (A.37)
$``$ $`\mathrm{𝐝𝐢𝐚𝐠}(0,0,0,m_{N_1},m_{N_2}),`$
since we are mostly interested in a heavy neutrino sector consisting of two right-handed neutrinos. The interaction eigenstates, in terms of the left- and right-handed components of the mass eigenstates, are given by:
$`\text{Left-handed: }(n^0)`$ $``$ $`\left(\begin{array}{c}\nu ^0\\ N^0\end{array}\right)=\left(\begin{array}{c}\nu _L\\ \nu _R^c\end{array}\right)=𝐔P_L(\chi )`$ (A.42)
$`\text{Right-handed: }(n_{}^{0}{}_{}{}^{c})`$ $``$ $`\left(\begin{array}{c}\nu _{}^{0}{}_{}{}^{c}\\ N_{}^{0}{}_{}{}^{c}\end{array}\right)=\left(\begin{array}{c}\nu _L^c\\ \nu _R\end{array}\right)=𝐔^{}P_R(\chi ),`$ (A.47)
with
$`P_{R,L}{\displaystyle \frac{1}{2}}(1\pm \gamma _5).`$ (A.48)
The diagonal blocks in $`𝐌`$, connecting states with opposite fermion number $`(𝐦_L,𝐦_R)`$ are called “Majorana” mass terms, while the off-diagonal ones $`(𝐦_D)`$ are “Dirac” mass terms that conserve fermion number.<sup>5</sup><sup>5</sup>5 Non-diagonal Dirac mass terms produce transitions between states with different individual lepton numbers but the total lepton number, i.e. the fermion number, is conserved. The mass terms arise, after the spontaneous symmetry breaking (SSB) of SU(2)$`{}_{L}{}^{}`$U(1)$`{}_{Y}{}^{}`$U(1)<sub>Q</sub>, from the Yukawa coupling of the fermion fields to the neutral Higgs field. The SU(2)$`{}_{L}{}^{}`$U(1)<sub>Y</sub> quantum numbers of a Higgs doublet and the fermion bilinears that can be constructed from the doublet and singlet neutrinos are
$`\mathrm{\Phi }=\left(\begin{array}{c}\varphi ^+\\ \varphi ^0\end{array}\right)`$ $``$ $`(\mathrm{𝟐},1)`$ (A.51)
$`\mathrm{\Delta }L=0:\overline{\nu _R}\nu _L`$ $``$ $`(\mathrm{𝟏},0)(\mathrm{𝟐},1)=(\mathrm{𝟐},1)`$ (A.52)
$`\mathrm{\Delta }L=\pm 2:\overline{\nu _L^c}\nu _L`$ $``$ $`(\mathrm{𝟐},1)(\mathrm{𝟐},1)=(\mathrm{𝟏},2)(\mathrm{𝟑},2)`$ (A.53)
$`\mathrm{\Delta }L=\pm 2:\overline{\nu _R^c}\nu _R`$ $``$ $`(\mathrm{𝟏},0)(\mathrm{𝟏},0)=(\mathrm{𝟏},0)(\text{neutral}).`$ (A.54)
Therefore in a model with only Higgs doublets, the entry $`𝐦_L=\mathrm{𝟎}`$.
Our convention for the covariant derivative acting on a fermion doublet is
$`D_\mu =_\mu +ig{\displaystyle \frac{\stackrel{}{\tau }}{2}}\stackrel{}{W}_\mu +ig^{}{\displaystyle \frac{Y}{2}}B_\mu `$ (A.55)
with
$`Q_f=I_3^{f_L}+Y/2.`$ (A.56)
After SSB one gets
$`e`$ $`=`$ $`gs_W=g^{}c_W,`$ (A.57)
$`\left(\begin{array}{c}Z_\mu \\ A_\mu \end{array}\right)`$ $`=`$ $`\left(\begin{array}{cc}\hfill c_W& \hfill s_W\\ \hfill s_W& \hfill c_W\end{array}\right)\left(\begin{array}{c}W_\mu ^3\\ B_\mu \end{array}\right),`$ (A.64)
$`W_\mu ^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(W_\mu ^1\pm iW_\mu ^2\right).`$ (A.65)
The interaction Lagrangians of neutrinos with $`W`$, $`Z`$ and charged Goldstone bosons $`\varphi ^\pm `$, in the weak basis<sup>6</sup><sup>6</sup>6 The interaction of the $`Z`$ to two (right-handed) charge-conjugate neutrinos (non-existent in the ordinary SM) has been included in (A.67), making use of the bilinear transformation properties under charge-conjugation: $`\overline{\psi }\gamma ^\mu P_{L,R}\psi \overline{\psi ^c}\gamma ^\mu P_{R,L}\psi ^c`$.
$`_W`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}W_\mu ^{}{\displaystyle \underset{i,j=1}{\overset{n_G}{}}}\overline{l}_i\gamma ^\mu P_L\nu _{L_j}+\text{h.c.},`$ (A.66)
$`_Z`$ $`=`$ $`{\displaystyle \frac{g}{2c_W}}Z_\mu {\displaystyle \underset{k=1}{\overset{n_G}{}}}\left[\overline{\nu _{L_k}}\gamma ^\mu P_L\nu _{L_k}\overline{\nu _{L_k}^c}\gamma ^\mu P_R\nu _{L_k}^c\right],`$ (A.67)
become, in terms of physical neutrinos $`\chi _j`$ with masses $`m_j`$, $`j=1,\mathrm{},n_G+n_R`$:
$`_W`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}W_\mu ^{}{\displaystyle \underset{i=1}{\overset{n_G}{}}}{\displaystyle \underset{j=1}{\overset{n_G+n_R}{}}}𝐁_{l_ij}\overline{l}_i\gamma ^\mu P_L\chi _j+\text{h.c.},`$ (A.68)
$`_Z`$ $`=`$ $`{\displaystyle \frac{g}{2c_W}}Z_\mu {\displaystyle \underset{i,j=1}{\overset{n_G+n_R}{}}}\overline{\chi }_i\gamma ^\mu (P_L𝐂_{ij}P_R𝐂_{ij}^{})\chi _j,`$ (A.69)
$`_{\varphi ^\pm }`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}\varphi ^{}{\displaystyle \underset{i=1}{\overset{n_G}{}}}{\displaystyle \underset{j=1}{\overset{n_G+n_R}{}}}𝐁_{l_ij}\overline{l}_i\left({\displaystyle \frac{m_{l_i}}{M_W}}P_L{\displaystyle \frac{m_j}{M_W}}P_R\right)\chi _j+\text{h.c.},`$ (A.70)
with
$`𝐁_{l_ij}`$ $``$ $`{\displaystyle \underset{k=1}{\overset{n_G}{}}}𝐕_{l_ik}^{}𝐔_{kj},`$ (A.71)
$`𝐂_{ij}`$ $``$ $`{\displaystyle \underset{k=1}{\overset{n_G}{}}}𝐔_{ki}^{}𝐔_{kj}.`$ (A.72)
$`𝐕_{l_ij}`$ and $`𝐁_{l_ij}`$ are the leptonic CKM mixing matrices and its generalized version, respectively . For Dirac particles $`𝐁=𝐕`$, $`𝐂_{ij}=\delta _{ij}`$ and $`𝐂_{ij}^{}=0`$. For Majorana particles, in contrast, there are NC couplings of different flavours (FCNC) and with both left- and right-handed components.
The matrix $`𝐕_{l_ij}`$ is quadratic of dimension $`n_G`$ and $`𝐁_{l_ij}`$ is rectangular $`n_G\times (n_G+n_R)`$ and incorporates lepton-flavour changing mixings. The matrix $`𝐂_{ij}`$ is quadratic, has dimension $`(n_G+n_R)`$, and causes flavour non-diagonal $`Z\chi _i\chi _j`$ interactions.
These interaction Lagrangians involve Majorana fermions that complicate the evaluation of $`S`$ matrix elements, since extra Wick contractions survive in comparison to the case with only Dirac fermions. Following , one can write Feynman rules resembling the ones of the Dirac fermions, based on the well defined fermion flow, rather than on the fermion number flow which is not preserved in the vertices with Majorana fermions: after fixing an arbitrary orientation (fermion flow) for a given diagram, the vertices can be read off from the Lagrangian as usual, but for every vertex $`\overline{f}_1\mathrm{\Gamma }f_2`$ one has to add the reversed one, $`\overline{f_1^c}\mathrm{\Gamma }^{}f_2^c`$, with $`\mathrm{\Gamma }^{}=C\mathrm{\Gamma }^TC^1`$. The fermion propagators are the usual ones. This effectively yields the same result for the vertex $`Wl_i\chi _j`$ but a factor two larger for the vertex $`Z\chi _i\chi _j`$ in comparison to the case of Dirac neutrinos, since for two Majorana fermions ($`\chi =\chi ^c`$) $`\mathrm{\Gamma }=\mathrm{\Gamma }^{}`$.
Important note: elsewhere in the text we refer to $`\nu _i`$ and $`N_i`$ as the neutrino physical states, rather than $`\chi _i`$, to simplify the presentation, but with no possible confusion, since we always work in the physical basis.
The matrices $`𝐁`$ and $`𝐂`$ obey a number of useful relations :
$`{\displaystyle \underset{j=1}{\overset{n_G+n_R}{}}}𝐁_{l_1j}𝐁_{l_2j}^{}`$ $`=`$ $`\delta _{l_1l_2},`$ (A.73)
$`{\displaystyle \underset{k=1}{\overset{n_G+n_R}{}}}𝐂_{ik}𝐂_{jk}^{}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n_G}{}}}𝐁_{l_ki}𝐁_{l_kj}^{}=𝐂_{ij},`$ (A.74)
$`{\displaystyle \underset{k=1}{\overset{n_G+n_R}{}}}𝐁_{lk}𝐂_{kj}`$ $`=`$ $`𝐁_{lj},`$ (A.75)
$`{\displaystyle \underset{k=1}{\overset{n_G+n_R}{}}}m_k𝐂_{ik}𝐂_{jk}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n_G+n_R}{}}}m_k𝐁_{lk}𝐂_{ki}^{}={\displaystyle \underset{k=1}{\overset{n_G+n_R}{}}}m_k𝐁_{l_1k}𝐁_{l_2k}=0.`$ (A.76)
In the case of $`n_R=2`$ the matrix elements involving heavy neutrinos can be obtained from (A.74), (A.76) in terms of the light-heavy mixing angles (B.3) and the ratio of the two heavy masses squared $`rm_{N_2}^2/m_{N_1}^2`$, assuming that the light sector consists of massless neutrinos :
$`𝐁_{l_kN_1}`$ $`=`$ $`{\displaystyle \frac{r^{1/4}}{\sqrt{1+r^{1/2}}}}s_{\nu _k},`$ (A.77)
$`𝐁_{l_kN_2}`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{1+r^{1/2}}}}s_{\nu _k},`$ (A.78)
and
$`𝐂_{N_1N_1}`$ $`=`$ $`{\displaystyle \frac{r^{1/2}}{1+r^{1/2}}}{\displaystyle \underset{k=1}{\overset{n_G}{}}}s_{\nu _k}^2,`$ (A.79)
$`𝐂_{N_2N_2}`$ $`=`$ $`{\displaystyle \frac{1}{1+r^{1/2}}}{\displaystyle \underset{k=1}{\overset{n_G}{}}}s_{\nu _k}^2,`$ (A.80)
$`𝐂_{N_1N_2}`$ $`=`$ $`𝐂_{N_2N_1}={\displaystyle \frac{ir^{1/4}}{1+r^{1/2}}}{\displaystyle \underset{k=1}{\overset{n_G}{}}}s_{\nu _k}^2.`$ (A.81)
## Appendix B Constraints on Heavy Neutrinos
### B.1 Experimental bounds
Several kinds of constraints on heavy-neutrino masses and light-heavy mixings can be obtained from experiment.
Direct production searches establish the following limits on the neutrino masses at 95% c.l. :
$`\text{Stable neutrinos: }m_N`$ $`>`$ $`45.0(39.5)\text{GeV}`$ (B.1)
$`\text{Unstable neutrinos: }m_N`$ $`>`$ $`69.0(58.2)\text{GeV},`$ (B.2)
for Dirac (Majorana) particles, respectively.
A general formalism to describe light-heavy mixings was developed in . The mixing angles in our notation correspond to
$`s_{\nu _k}^2{\displaystyle \underset{i}{}}|𝐁_{l_kN_i}|^2.`$ (B.3)
Indirect constraints on the masses and bounds on the mixings are provided by two categories of LEP and low energy experiments:
(i) Flavour-diagonal processes. They include mass-independent and model-independent light-heavy mixing constraints from tests of lepton universality and CKM unitarity and measurements of the $`Z`$ boson invisible width, as well as other less sensitive studies like $`W`$ mass measurements and low energy experiments like neutrino scattering, atomic parity violation, etc. Flavour-conserving leptonic decays $`Zl^{}l^+`$ depend on masses and mixings through loop contributions and provide alternative constraints .
(ii) Flavour-changing processes. They include rare processes like $`\mu e\gamma `$, $`\mu ee^+e^{}`$, $`\mu e`$ conversion in nuclei, $`\tau l_al_b^+l_c^{}`$ and $`Zl_a^{}l_b^+`$ . They are mass-dependent (except for $`\mu e\gamma `$ to a good approximation). One can get also less stringent light-heavy mixing constraints from oscillation experiments .
The most stringent present bounds on the light-heavy mixings are provided by the flavour-diagonal processes. Exceptions are the ones involving the first two lepton families such as $`\mu e\gamma `$, $`Ze\mu `$ and $`\mu ee^+e^{}`$ .
For illustration we show an example of the flavour-diagonal constraints. The effective muon decay constant $`G_\mu `$ is related to the coupling $`G_F`$ of the standard model by :
$`G_\mu =c_{\nu _e}c_{\nu _\mu }G_F.`$ (B.4)
The unitarity constraint for the first row of the CKM quark-mixing matrix implies
$`{\displaystyle \underset{i=1}{\overset{3}{}}}|V_{ui}|^2`$ $`=`$ $`\left({\displaystyle \frac{c_{\nu _e}G_F}{G_\mu }}\right)^2={\displaystyle \frac{1}{c_{\nu _\mu }^2}}=0.9992\pm 0.0014.`$ (B.5)
As a summary, from (90% c.l.):
$`s_{\nu _e}^2`$ $`<`$ $`0.0071(0.005),`$ (B.6)
$`s_{\nu _\mu }^2`$ $`<`$ $`0.0014,`$ (B.7)
$`s_{\nu _\tau }^2`$ $`<`$ $`0.033(0.01),`$ (B.8)
where the most conservative bounds are obtained assuming any kind of heavy neutrinos and the ones in brackets correspond to the case of SU(2) singlets.
From the most recent update<sup>7</sup><sup>7</sup>7 These latest bounds are more conservative than the earlier ones in due to the fact that present determinations of the elements of the first row of the CKM matrix are not compatible with unitarity, and hence this constraint is eliminated from the analysis. by , assuming only heavy singlets, one gets (90% c.l.):
$`s_{\nu _e}^2`$ $`<`$ $`0.012,`$ (B.9)
$`s_{\nu _\mu }^2`$ $`<`$ $`0.0096,`$ (B.10)
$`s_{\nu _\tau }^2`$ $`<`$ $`0.016.`$ (B.11)
A final remark is in order here. Using the Schwartz inequalities ,
$`|{\displaystyle \underset{i}{}}𝐁_{l_aN_i}𝐁_{l_bN_i}^{}|^2<s_{\nu _a}^2s_{\nu _b}^2,`$ (B.12)
one can infer indirect upper limits on the off-diagonal mixings (relevant for the flavour-changing processes) from the previous flavour-diagonal constraints. Nevertheless, for our scenario (iii), as already mentioned, the mixings can be obtained exactly from the properties of $`𝐁`$ and $`𝐂`$ and such inequalities are not needed.
### B.2 Decoupling and neutrino-mass upper limits from perturbative unitarity
The heavy-neutrino masses are restricted by the perturbative unitarity condition on the decay width of heavy neutrinos ,
$`\mathrm{\Gamma }_{N_i}{\displaystyle \frac{1}{2}}m_{N_i}.`$ (B.13)
The total decay width of a heavy Dirac neutrino (four d.o.f.) with mass $`m_{N_i}M_W,M_Z,M_H`$ is :
$`\mathrm{\Gamma }_{N_i}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{3}{}}}\mathrm{\Gamma }(N_il_k^{}W^+)+{\displaystyle \underset{k=1}{\overset{3}{}}}\left[\mathrm{\Gamma }(N_i\nu _kZ)+\mathrm{\Gamma }(N_i\nu _kH)\right]`$ (B.14)
$``$ $`{\displaystyle \frac{\alpha _W}{8M_W^2}}m_{N_i}^3{\displaystyle \underset{k=1}{\overset{3}{}}}|𝐁_{l_kN_i}|^2,`$
and a factor two larger for a heavy Majorana neutrino (two d.o.f.). Therefore, the perturbative unitarity bound expressed in (B.13) reads
$`m_{N_i}^2{\displaystyle \underset{k=1}{\overset{3}{}}}|𝐁_{l_kN_i}|^2=m_{N_i}^2𝐂_{N_iN_i}\{\begin{array}{c}4M_W^2/\alpha _W\text{ for a Dirac neutrino}\hfill \\ \\ 2M_W^2/\alpha _W\text{ for a Majorana neutrino},\hfill \end{array}`$ (B.18)
(no summation over repeated indices is understood) which shows implicitly that heavy neutrinos decouple , in accordance with the Appelquist-Carazzone theorem : the unacceptable large-mass behaviour of the amplitudes ($`m_N^2`$) is actually cured when the light-heavy mixing ($`m_N^2`$) is taken into account .
Taking the values of $`𝐁_{l_kN_i}`$ in terms of the light-heavy mixing angles (B.3) for one Dirac or for two heavy Majorana neutrinos (A.77), (A.78) one can get the following upper limits:
$`m_N^2`$ $``$ $`{\displaystyle \frac{4M_W^2}{\alpha _W}}\left[{\displaystyle \underset{k=1}{\overset{3}{}}}s_{\nu _k}^2\right]^1`$ (B.19)
for a heavy Dirac neutrino, and
$`m_{N_1}^2{\displaystyle \frac{1}{r}}m_{N_2}^2`$ $``$ $`{\displaystyle \frac{2M_W^2}{\alpha _W}}{\displaystyle \frac{1+r^{1/2}}{r}}\left[{\displaystyle \underset{k=1}{\overset{3}{}}}s_{\nu _k}^2\right]^1`$ (B.20)
for two heavy Majorana singlets. The latter bound is very stringent when $`m_{N_1}`$ and $`m_{N_2}`$ are very different. The upper mass limits on heavy neutrinos, from (B.19), (B.20) are then:
$`\text{scenario (ii): }m_N^2`$ $`<`$ \[-0.07cm\] $``$ $`(4.2,4.4\text{TeV})^2`$ (B.21)
$`\text{scenario (iii): }m_{N_1}^2{\displaystyle \frac{1}{r}}m_{N_2}^2`$ $`<`$ \[-0.07cm\] $``$ $`{\displaystyle \frac{1+r^{1/2}}{r}}\times (3.0,3.1\text{TeV})^2,`$ (B.22)
using the bounds (B.6)–(B.8) and (B.9)–(B.11), respectively.
## Appendix C The Vertex Function
We use the notations of . The vertex function contains the following one-loop integrals in $`D`$ dimensions :
$`{\displaystyle \frac{i}{16\pi ^2}}A(m_0^2)`$ $`=`$ $`\mu ^{4D}{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{1}{𝒟_0}},`$ (C.1)
$`{\displaystyle \frac{i}{16\pi ^2}}\{B_0,B^\mu \}(p_1^2;m_0^2,m_1^2)`$ $`=`$ $`\mu ^{4D}{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{\{1,q^\mu \}}{𝒟_0𝒟_1}},`$ (C.2)
$`{\displaystyle \frac{i}{16\pi ^2}}\{C_0,C^\mu ,C^{\mu \nu }\}(p_1^2,Q^2,p_2^2;m_0^2,m_1^2,m_2^2)`$ $`=`$ $`\mu ^{4D}{\displaystyle \frac{d^Dq}{(2\pi )^D}\frac{\{1,q^\mu ,q^\mu q^\nu \}}{𝒟_0𝒟_1𝒟_2}},`$ (C.3)
with
$`Q^2=(p_2p_1)^2`$ (C.4)
and
$`𝒟_0`$ $`=`$ $`q^2m_0^2+iϵ,`$ (C.5)
$`𝒟_1`$ $`=`$ $`(q+p_1)^2m_1^2+iϵ,`$ (C.6)
$`𝒟_2`$ $`=`$ $`(q+p_2)^2m_2^2+iϵ.`$ (C.7)
They are decomposed into tensor integrals according to their Lorentz structure:
$`B^\mu `$ $`=`$ $`p^\mu B_1,`$ (C.8)
$`C^\mu `$ $`=`$ $`p_1^\mu C_{11}+p_2^\mu C_{12},`$ (C.9)
$`C^{\mu \nu }`$ $`=`$ $`p_1^\mu p_1^\nu C_{21}+p_2^\mu p_2^\nu C_{22}+(p_1^\mu p_2^\nu +p_2^\mu p_1^\nu )C_{23}+g^{\mu \nu }C_{24}.`$ (C.10)
We employ the following dimensionless tensor integrals and their abbreviations:
$`B_1`$ $``$ $`B_1(0;\lambda _i,1)`$ (C.11)
$`=`$ $`B_1(0;m_i^2,M_W^2),`$
$`C_{\{0,11,12,21,22,23\}}`$ $``$ $`C_{\{\mathrm{}\}}(0,\lambda _Q,0;1,\lambda _i,\lambda _j)`$ (C.12)
$`=`$ $`M_W^2C_{\{\mathrm{}\}}(0,Q^2,0;M_W^2,m_i^2,m_j^2),`$
$`C_{24}`$ $``$ $`C_{24}(0,\lambda _Q,0;1,\lambda _i,\lambda _j)`$ (C.13)
$`=`$ $`C_{24}(0,Q^2,0;M_W^2,m_i^2,m_j^2),`$
$`\overline{C}_{\{0,11,12,21,22,23\}}`$ $``$ $`C_{\{\mathrm{}\}}(0,\lambda _Q,0;\lambda _i,1,1)`$ (C.14)
$`=`$ $`M_W^2C_{\{\mathrm{}\}}(0,Q^2,0;m_i^2,M_W^2,M_W^2),`$
$`\overline{C}_{24}`$ $``$ $`\overline{C}_{24}(0,\lambda _Q,0;\lambda _i,1,1)`$ (C.15)
$`=`$ $`C_{24}(0,Q^2,0;m_i^2,M_W^2,M_W^2).`$
On the $`Z`$ mass shell it is $`\lambda _Q=\lambda _Z`$. The $`\lambda _Q,\lambda _Z,\lambda _i`$ are introduced in (1.18), (2.5), and (2.6).
The explicit expressions for the loop functions are:
$`C_{\{0,11,23\}}`$ $`=`$ $`{\displaystyle _0^1}𝑑x{\displaystyle _0^x}{\displaystyle \frac{dy}{𝒟}}\{1,y,(1x)y\},`$ (C.16)
$`C_{12}(\lambda _i,\lambda _j)`$ $`=`$ $`C_{11}(\lambda _j,\lambda _i),`$ (C.17)
$`C_{24}`$ $`=`$ $`{\displaystyle \frac{1}{2(D4)}}{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑x{\displaystyle _0^x}𝑑y\mathrm{ln}𝒟,`$ (C.18)
with
$`𝒟`$ $`=`$ $`D_{ijk}`$ (C.19)
$`=`$ $`\lambda _Qxy+(\lambda _k\lambda _j)x+(\lambda _Q+\lambda _i\lambda _k)y+\lambda _jiϵ,`$
and the correspondences are: for the non-abelian diagrams (with elementary $`ZWW`$, $`ZW\varphi `$, $`Z\varphi \varphi `$ vertices), ($`i,j`$) are virtual $`W,\varphi `$ bosons and $`k`$ a neutrino, while for the abelian diagrams $`k`$ is the $`W,\varphi `$ boson and ($`i,j`$) are neutrinos. On the $`Z`$ boson mass shell:
$`𝒟`$ $``$ $`D_{ijW}=\lambda _Zxy+(1\lambda _j)x+[\lambda _Z+(\lambda _i1)]y+\lambda _jiϵ,`$ (C.20)
$`\overline{𝒟}`$ $``$ $`D_{WWi}=\lambda _Zxy(1\lambda _i)x+[\lambda _Z(\lambda _i1)]y+1iϵ.`$ (C.21)
In the Dirac case, it is $`\lambda _i=\lambda _j`$. Further,
$`B_1`$ $`=`$ $`{\displaystyle \frac{1}{D4}}+{\displaystyle _0^1}x𝑑x\mathrm{ln}[x+\lambda _i(1x)iϵ]`$ (C.22)
$`=`$ $`{\displaystyle \frac{1}{D4}}+{\displaystyle \frac{\lambda _i}{2(1\lambda _i)}}\left(1+{\displaystyle \frac{\lambda _i\mathrm{ln}\lambda _i}{1\lambda _i}}\right){\displaystyle \frac{1}{4}}.`$
As may be seen, the tensor integrals $`B_1`$, $`C_{24}`$ and $`\overline{C}_{24}`$ are ultraviolet–divergent in $`D`$ dimensions. We mention that the functions are defined here with Minkowskian metric, so that all the functions introduced, except for $`B_0,B_1,C_{24},\overline{C}_{24}`$, have different sign from those used in e.g. (Euclidean metric). In principle, the divergent parts depend on the regularization scheme and could differ by a finite, universal term yet. Recalling the UV-behaviour of the divergent integrals $`B_1`$, $`C_{24}`$ and $`\overline{C}_{24}`$ (see (C.22) and (C.18)), we get divergent, mass-dependent contributions from individual diagrams to the vertex $`V`$.
For the Dirac case:
$`v_\varphi (i)`$ $``$ $`(v_i+a_i)\times \left[\text{finite}\right](v_ia_i)\times \left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _i}{D4}}+\text{finite}\right],`$ (C.23)
$`v_{\varphi \varphi }(i)`$ $``$ $`(2I_3^{i_L})(12s_W^2)\times \left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _i}{D4}}+\text{finite}\right],`$ (C.24)
$`v_\mathrm{\Sigma }(i)`$ $``$ $`{\displaystyle \frac{1}{2}}(v_i+a_i4c_W^2a_i)\left[{\displaystyle \frac{2}{D4}}+{\displaystyle \frac{\lambda _i}{D4}}+\text{finite}\right].`$ (C.25)
The sum of the terms proportional to $`\lambda _i/(D4)`$ vanishes. The constant divergent terms (not shown here) sum up for individual vertices $`V(\lambda _i)`$, but vanish due to the unitarity of the mixing matrix for the complete vertex function (1.15):
$`𝒱`$ $``$ $`(2I_3^{i_L})4c_W^2{\displaystyle \frac{1}{D4}}\delta _{l_1l_2}=0.`$ (C.26)
For the Majorana case ($`2I_3^{i_L}=1`$, $`v_i=a_i=1/2`$):
$`v_\varphi (i,j)`$ $``$ $`𝐂_{ij}\times \left[\text{finite}\right]+𝐂_{ij}^{}\times \left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\sqrt{\lambda _i\lambda _j}}{D4}}+\text{finite}\right],`$ (C.27)
$`v_{\varphi \varphi }(i,j)`$ $``$ $`\delta _{ij}(12s_W^2)\times \left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _i}{D4}}+\text{finite}\right],`$ (C.28)
$`v_\mathrm{\Sigma }(i,j)`$ $``$ $`\delta _{ij}{\displaystyle \frac{1}{2}}(1+2s_W^2)\left[{\displaystyle \frac{2}{D4}}+{\displaystyle \frac{\lambda _i}{D4}}+\text{finite}\right].`$ (C.29)
The mass-dependent divergent terms in (C.28) and (C.29) cancel each other, and the one proportional to $`𝐂_{ij}^{}`$ in (C.27) drops out due to the unitarity condition (A.76) when the sum over virtual neutrinos is performed in (4.1). The constant divergent terms vanish again due to the unitarity relations (A.73)–(A.76):
$`𝒱_M`$ $``$ $`4c_W^2{\displaystyle \frac{1}{D4}}\delta _{l_1l_2}=0`$ (C.30)
## Appendix D The Vertex Function for Large and Small Neutrino Masses
### D.1 The vertex for large neutrino mass, $`\lambda _i1`$
We now consider the vertex function in the limit of large Dirac-neutrino masses $`(\lambda _i=\lambda _j)`$. The leading terms of the one-loop functions for $`\lambda _Q=\lambda _Z`$ are:<sup>8</sup><sup>8</sup>8 The functions $`\overline{C}_{11}`$, $`\overline{C}_{12}`$, $`\overline{C}_{23}`$ do not contribute to the large mass limit of the vertex, and we reproduce only their first leading terms, of order $`\mathrm{ln}\lambda /\lambda `$. If these functions are of relevance for an application, one should determine also the terms of order $`1/\lambda `$.
$`C_0`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}+{\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i^2}}\left(12+\lambda _Z\right){\displaystyle \frac{1}{12\lambda _i^2}}+\mathrm{},`$ (D.31)
$`C_{11}`$ $`=`$ $`C_{12}={\displaystyle \frac{1}{4\lambda _i}}+\mathrm{},`$ (D.32)
$`C_{23}`$ $`=`$ $`{\displaystyle \frac{1}{18\lambda _i}}+\mathrm{},`$ (D.33)
$`C_{24}`$ $`=`$ $`{\displaystyle \frac{1}{2(D4)}}{\displaystyle \frac{1}{4}}\mathrm{ln}\lambda _i+{\displaystyle \frac{1}{8}}+\left(9+\lambda _Z\right){\displaystyle \frac{1}{36\lambda _i}}+\mathrm{},`$ (D.34)
$`B_1`$ $`=`$ $`{\displaystyle \frac{1}{(D4)}}+{\displaystyle \frac{1}{2}}\mathrm{ln}\lambda _i{\displaystyle \frac{3}{4}}+{\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i}}{\displaystyle \frac{1}{2\lambda _i}}+\mathrm{},`$ (D.35)
$`\overline{C}_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i}}\left[14a(y)\right]{\displaystyle \frac{1}{\lambda _i}}+\mathrm{},`$ (D.36)
$`\overline{C}_{11}`$ $`=`$ $`\overline{C}_{12}={\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i}}+\mathrm{},`$ (D.37)
$`\overline{C}_{23}`$ $`=`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i}}+\mathrm{},`$ (D.38)
$`\overline{C}_{24}`$ $`=`$ $`{\displaystyle \frac{1}{2(D4)}}{\displaystyle \frac{1}{4}}\mathrm{ln}\lambda _i+{\displaystyle \frac{3}{8}}+\left(6+\lambda _Z\right){\displaystyle \frac{\mathrm{ln}\lambda _i}{12\lambda _i}}`$ (D.39)
$`+\left[30+5\lambda _Z+24(4\lambda _Z)a(y)\right]{\displaystyle \frac{1}{72\lambda _i}}+\mathrm{},`$
with
$`y`$ $`=`$ $`\sqrt{1/\lambda _Z1/4},`$ (D.40)
$`a(y)`$ $`=`$ $`y\mathrm{arctan}(1/2y).`$ (D.41)
The large Dirac-neutrino mass limit of the vertex function is:
$`V(\lambda _i)`$ $`=`$ $`I_3^{\nu _L}[\lambda _i+(3{\displaystyle \frac{\lambda _Z}{6}}(12s_W^2))\mathrm{ln}\lambda _i`$
$`+{\displaystyle \frac{1}{18}}\left(66\lambda _Z+96s_W^2+5s_W^2\lambda _Z\right)`$
$`+{\displaystyle \frac{1}{3}}(8+2\lambda _Z32s_W^24s_W^2\lambda _Z)a(y)+{\displaystyle \frac{8c_W^2}{D4}}]+𝒪({\displaystyle \frac{\mathrm{ln}\lambda _i}{\lambda _i}}).`$
Numerically, this means for the value $`\lambda _Z=1.286`$:
$`V(\lambda _N)V(0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\lambda _4+2.88\mathrm{ln}\lambda _44.47(2.52+2.11\times i)\right].`$ (D.43)
Here, we subtracted from the large mass limit of the vertex function its value at zero mass (D.66) in order to obtain (3).
#### D.1.1 A relation to the $`Zb\overline{b}`$ vertex
The loop contributions to the vertex function $`V(\lambda )`$ are gauge-invariant. We calculate them in the ’t Hooft-Feynman gauge. They describe not only the flavour-changing $`Zf\overline{f}^{}`$ vertex, but also the mass dependent terms of the flavour-diagonal vertex. There is one case where the effect is quite visible, namely the $`Zb\overline{b}`$ vertex with virtual $`t`$ quark exchanges. The exact one-loop expression for the $`t`$ quark mass dependent part $`𝒲`$ of the $`Zb\overline{b}`$ vertex correction was first given in Eqn. (22) of (form factor $`\delta \kappa `$, calculated in the unitary gauge). The large $`t`$ quark mass limit is (Eqn. (2.4.30) of ):
$`(Z\overline{b}b)`$ $``$ $`ϵ^\mu \overline{u}\left[\gamma _\mu v_b\gamma _\mu \gamma _5a_b+\gamma _\mu (1\gamma _5)𝒲(\lambda _t,\lambda _Z)\right]u,`$ (D.44)
$`{\displaystyle \frac{𝒲(\lambda _t,\lambda _Z)}{a_b}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{\pi }}{\displaystyle \frac{1}{4s_W^2}}|P_{tb}|^2V_t,`$ (D.45)
$`V_t`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\lambda _t+\left({\displaystyle \frac{8}{3}}+{\displaystyle \frac{\lambda _Z}{6}}\right)\mathrm{ln}\lambda _t+𝒪(1)\right].`$ (D.46)
This agrees with the leading terms in (D.1).
### D.2 The vertex for small neutrino mass, $`\lambda _i1`$
The limits of the $`C`$ functions for $`\lambda _i=\lambda _j`$ and $`\lambda _Q=\lambda _Z\lambda _i`$ are:
$`C_0`$ $`=`$ $`c_0+{\displaystyle \frac{2}{\lambda _Z}}\lambda _i\mathrm{ln}\lambda _i{\displaystyle \frac{2}{\lambda _Z(1+\lambda _Z)}}(1+\lambda _Z+\mathrm{ln}\lambda _Zi\pi )\lambda _i+\mathrm{},`$ (D.47)
$`\lambda _ZC_{11}`$ $`=`$ $`\lambda _ZC_{12}`$ (D.48)
$`=`$ $`c_0+1\mathrm{ln}\lambda _Z+i\pi \lambda _i\mathrm{ln}\lambda _i+\left[c_0+{\displaystyle \frac{2}{1+\lambda _Z}}(\mathrm{ln}\lambda _Zi\pi )\right]\lambda _i`$
$`+\mathrm{},`$
$`\lambda _Z^2C_{23}`$ $`=`$ $`(\lambda _Z+2)c_0+{\displaystyle \frac{\lambda _Z}{2}}+22(\mathrm{ln}\lambda _Zi\pi )`$ (D.49)
$`\left[4(c_01)2{\displaystyle \frac{2+\lambda _Z}{1+\lambda _Z}}(\mathrm{ln}\lambda _Zi\pi )\right]\lambda _i+\mathrm{},`$
$`4\lambda _ZC_{24}`$ $`=`$ $`{\displaystyle \frac{2\lambda _Z}{D4}}2(1+\lambda _Z)c_0+3\lambda _Z+2(\lambda _Z+2)(\mathrm{ln}\lambda _Zi\pi )`$ (D.50)
$`+4\left[c_01+\mathrm{ln}\lambda _Zi\pi \right]\lambda _i+\mathrm{},`$
$`B_1`$ $`=`$ $`{\displaystyle \frac{1}{D4}}{\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}\lambda _i+\mathrm{},`$ (D.51)
$`\overline{C}_0`$ $`=`$ $`\overline{c}_0\lambda _i\mathrm{ln}\lambda _i(B1)\lambda _i+\mathrm{},`$ (D.52)
$`\lambda _Z\overline{C}_{11}`$ $`=`$ $`\lambda _Z\overline{C}_{12}=(\overline{c}_0B+1)(\lambda _i1)+\mathrm{},`$ (D.53)
$`\lambda _Z^2\overline{C}_{23}`$ $`=`$ $`2(\overline{c}_0B+1)+{\displaystyle \frac{\lambda _Z}{2}}\left[\lambda _Z\overline{c}_04(\overline{c}_0B+1)\right]\lambda _i+\mathrm{},`$ (D.54)
$`2\lambda _Z\overline{C}_{24}`$ $`=`$ $`{\displaystyle \frac{\lambda _Z}{D4}}(\overline{c}_0B+1)+{\displaystyle \frac{3\lambda _Z}{2}}\pi \lambda _Zy+2\lambda _Zy\mathrm{arctan}(2y)`$ (D.55)
$`+\left[2(\overline{c}_0B+1)\lambda _Z\overline{c}_0\right]\lambda _i+\mathrm{},`$
with
$`B`$ $`=`$ $`2y\left[\mathrm{arctan}(2y)+\mathrm{arctan}\left({\displaystyle \frac{\lambda _Z1}{3\lambda _Z}}2y\right)\right]1.75,`$ (D.56)
for $`\lambda _Z=1.286`$, $`y0.73`$, and the values of $`C_0`$ and $`\overline{C}_0`$ at $`\lambda _i=0`$ :
$`\lambda _Zc_0`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{6}}\text{Li}_2\left({\displaystyle \frac{1}{1+\lambda _Z}}\right){\displaystyle \frac{1}{2}}\mathrm{ln}^2(1+\lambda _Z)+\pi \mathrm{ln}(1+\lambda _Z)\times i,`$ (D.57)
$`\lambda _Z\overline{c}_0`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{6}}\text{Li}_2(1\lambda _Z)+2\mathrm{}e\text{Li}_2\left[(\lambda _Z1)\left({\displaystyle \frac{\lambda _Z}{2}}1+\lambda _Zy\times i\right)\right]`$ (D.58)
$`2\mathrm{}e\text{Li}_2\left(1{\displaystyle \frac{\lambda _Z}{2}}\lambda _Zy\times i\right)`$
Only the functions $`\overline{C}`$ develop imaginary parts, and only for
$`\lambda _Z>4\lambda _i.`$ (D.59)
The functions $`C_0,C_{11},\overline{C}_0`$ contain terms of the order $`\lambda _Z\mathrm{ln}\lambda _Z`$ at $`\lambda _i\lambda _Z`$, but these terms cancel in the form factor $`V`$. They survive in the case $`\lambda _Z0`$ at constant $`\lambda _i`$.
A further note: the Euler dilogarithm is badly converging on the unit circle. In fact, for one of the $`\text{Li}_2`$ above it occurs $`|1\lambda _Z/2\pm i\times \lambda _Zy|=1`$. In that case, one may use:
$`\mathrm{}e\text{Li}_2\left(e^{i\varphi }\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Li}_2(1)+{\displaystyle \frac{1}{4}}(\varphi \pm \pi )^2,`$ (D.60)
taking the value of $`\varphi `$ fulfilling the condition:
$`(\varphi \pm \pi )(\pi ,\pi ).`$ (D.61)
The resulting small mass limit of the vertex is given in for the case of sequential Dirac neutrinos:
$`V(\lambda _i1,\lambda _Z)`$ $`=`$ $`I_3^{i_L}{\displaystyle \frac{8c_W^2}{D4}}+a_0+a_L(\lambda _i\mathrm{ln}\lambda _i)+a_1\lambda _i+\mathrm{}`$ (D.62)
The divergent constant was left out in the introduction, Eqn. (1.16).
We want to stress that in this limit,
$`a_L=0.`$ (D.63)
The ansatz in Eqn. (1) of allowing for $`a_L0`$ is too general for the case of small neutrino masses in this respect. Such a term occurs for $`\lambda _Q=0`$, see section D.3.
A numerical value of $`a_0`$ was given in . The analytical expression is:
$`a_0`$ $`=`$ $`{\displaystyle \frac{(1+\lambda _Z)^2}{\lambda _Z}}c_0+{\displaystyle \frac{2}{\lambda _Z^2}}(1+2\lambda _Z)(\overline{c}_0B)+{\displaystyle \frac{6}{\lambda _Z}}[\pi y2y\mathrm{arctan}(2y)]`$ (D.64)
$`+{\displaystyle \frac{2+3\lambda _Z}{2\lambda _Z}}(\mathrm{ln}\lambda _Z\pi \times i){\displaystyle \frac{1}{4\lambda _Z^2}}(7\lambda _Z^2+14\lambda _Z8).`$
The general expression for $`a_1`$ is :
$`a_1(\lambda _Z)`$ $`=`$ $`{\displaystyle \frac{2}{\lambda _Z}}(1+\lambda _Z)c_0+{\displaystyle \frac{1}{2\lambda _Z^2}}(4\lambda _Z^25\lambda _Z6)\overline{c}_0{\displaystyle \frac{2}{\lambda _Z}}(\mathrm{ln}\lambda _Z\pi \times i)`$ (D.65)
$`+{\displaystyle \frac{1}{8\lambda _Z^2}}(25\lambda _Z^238\lambda _Z24)+{\displaystyle \frac{1}{2\lambda _Z}}(2\lambda _Z)\pi y`$
$`+{\displaystyle \frac{1}{\lambda _Z^2}}(\lambda _Z^2+7\lambda _Z+6)y\mathrm{arctan}(2y)`$
$`+{\displaystyle \frac{3}{\lambda _Z^2}}(3\lambda _Z+2)y\mathrm{arctan}\left({\displaystyle \frac{\lambda _Z1}{3\lambda _Z}}2y\right),`$
with $`y`$ and $`B`$ from (D.40) and (D.56). Further, we used $`c_0`$ and $`\overline{c}_0`$ of (D.57) and (D.58). With the inputs $`M_W=80.410`$ GeV and $`M_Z=91.187`$ GeV, these formulae yield:
$`a_0`$ $`=`$ $`1.2584+1.0524\times i,`$ (D.66)
$`a_1`$ $`=`$ $`2.56232.2950\times i.`$ (D.67)
### D.3 The vertex for $`\lambda _Q=0`$
The first calculation of the non-diagonal $`Zf_1f_2`$ vertex seems to be in , where the approach was simplified considerably by the approximation $`Q^2=0`$. We should mention that this limit is of physical relevance only in the large neutrino mass limit, since only then it is $`\lambda _i\lambda _Z\lambda _Q=0`$. If we are interested in applications where $`\lambda _Q>4\lambda _i`$, the amplitude gets essentially complex valued. The limit of , however, implies automatically the relation $`\lambda _i>\lambda _Q`$, even for $`\lambda _i1`$, and the amplitude is essentially real. So one cannot expect a continuous behaviour for the small or medium neutrino mass range. The virtual fermions are considered Dirac particles $`(\lambda _i=\lambda _j)`$.
The vertex may be easily derived from the formulae given in :
$`V(\lambda _i,\lambda _Q=0)`$ $`=`$ $`I_3^{i_L}\left[\lambda _i\left(\lambda _i10\right)_1+8+6+{\displaystyle \frac{8c_W^2}{D4}}\right],`$ (D.68)
with
$`_1`$ $`=`$ $`{\displaystyle _0^1}𝑑x{\displaystyle \frac{x}{(1\lambda _i)x+\lambda _i}}={\displaystyle \frac{\lambda _i\mathrm{ln}\lambda _i}{(1\lambda _i)^2}}+{\displaystyle \frac{1}{1\lambda _i}},`$ (D.69)
$``$ $`=`$ $`{\displaystyle _0^1}𝑑xx\mathrm{ln}\left[(1\lambda _i)x+\lambda _i\right]={\displaystyle \frac{\lambda _i^2\mathrm{ln}\lambda _i}{2(1\lambda _i)^2}}+{\displaystyle \frac{1}{2(1\lambda _i)}}{\displaystyle \frac{1}{4}}.`$ (D.70)
The divergent part is independent of $`\lambda _i`$ and of $`\lambda _Q`$.
Explicitly:
$`V(\lambda _i,\lambda _Q=0)`$ $`=`$ $`I_3^{i_L}\left(3{\displaystyle \frac{\lambda _i^2\mathrm{ln}\lambda _i}{(1\lambda _i)^2}}+2{\displaystyle \frac{\lambda _i\mathrm{ln}\lambda _i}{(1\lambda _i)^2}}{\displaystyle \frac{\lambda _i^2}{1\lambda _i}}+6{\displaystyle \frac{\lambda _i}{1\lambda _i}}+{\displaystyle \frac{8c_W^2}{D4}}\right),`$
with $`I_3^{i_L}`$ being the weak isospin of the virtual neutrinos. The constant finite terms of the vertex vanish for small $`\lambda _i`$:
$`V(\lambda _i1,\lambda _Q=0)`$ $`=`$ $`I_3^{i_L}\left(2\lambda _i\mathrm{ln}\lambda _i+6\lambda _i+{\displaystyle \frac{8c_W^2}{D4}}\right).`$ (D.72)
This has to be compared to (D.62) and (D.63).
The large mass limit is:
$`V(\lambda _i1,\lambda _Q=0)`$ $`=`$ $`I_3^{i_L}\left(\lambda _i+3\mathrm{ln}\lambda _i5+{\displaystyle \frac{8c_W^2}{D4}}\right).`$ (D.73)
This has to be compared to (D.1) and (D.43).
Finally, for the sake of completeness, the value of the vertex at the weak scale:
$`V(\lambda _i=1,\lambda _Q=0)`$ $`=`$ $`I_3^{i_L}\left({\displaystyle \frac{3}{2}}+{\displaystyle \frac{8c_W^2}{D4}}\right).`$ (D.74) |
warning/0001/hep-ph0001136.html | ar5iv | text | # ϕ→𝜋⁰𝜋⁰𝛾 DECAY WITHIN A 𝑈(3)×𝑈(3) LINEAR SIGMA MODEL
## I Introduction
Over the past years experimental evidence has accumulated for the existence of light scalar mesons , and different proposals exist for the $`\overline{q}q`$ lowest lying scalar meson nonet. The Particle Data Group (PDG) candidates for the ground state $`\overline{q}q`$ scalar nonet are : the $`f_0(980)`$, $`f_0(1370)`$ and the $`f_0(4001250)`$ ( or $`\sigma `$ meson) for two states in the $`I=0`$ sector; the $`a_0(980)`$ and $`a_0(1450)`$ for the isovector scalar meson, and the $`K_0^{}(1430)`$ for the isospinor scalar meson.
The small decay rate into two photons is among the most important drawbacks for the identification of the $`a_0(980)`$ and $`f_0(980)`$ as the $`\overline{q}q`$ scalar isovector and isosinglet respectively. These decay rates have been calculated using a variety of approaches , in particular, in different versions of the quark model . The generally accepted conclusion, is that the $`a_0(980),f_0(980)\gamma \gamma `$ decay widths are not consistent with a $`q\overline{q}`$ structure.
Moreover, the nearby mass degeneracy of these mesons suggest they are the scalar analogous of the $`\omega `$ and $`\rho `$ system, i.e. the $`a_0(980),f_0(980`$ are expected to be $`\overline{q}q`$ ( $`q=u,d`$ ) which however contradicts the strong coupling of the $`f_0(980)`$ meson to the $`\overline{K}K`$ system.
Other possibilities such as a molecule picture and a $`\overline{q}q\overline{q}q`$ structure have been explored. Recently, it has been argued that the four-quark picture for these mesons is consistent not only with the two photon decays of these states but also with the $`\mathrm{\Phi }f_0\gamma \pi \pi \gamma `$ decay .
An alternative approach to the hadron physics in this energy region is provided by the $`U(3)\times U(3)`$ chiral model which incorporates a nonet of scalar as well as a nonet of pseudoscalar particles . In fact, the $`U_A(1)`$ component of the $`U(3)\times U(3)`$ symmetry exhibited by the light sector of QCD in the massless quark limit, is broken at the quantum level which in the model amounts to the possibility of incorporating otherwise forbidden terms in the interaction lagrangian. In this model, the $`f_0(980)`$ turns out to be a mostly $`\overline{s}s`$ meson whereas the $`a_0(980)`$ meson is the chiral partner of the pion ; the reason for the nearby degeneracy of the $`a_0(980)`$ and the $`f_0980)`$ being that the $`U_A(1)`$ anomaly pushes up the $`a_0(980)`$ mass while leaving untouched the $`f_0(980)`$ meson.
The model has shown to be phenomenologically succesful . In particular the $`a_0(980)\gamma \gamma `$ and $`f_0(980)\gamma \gamma `$ decays are consistenly accounted for in this framework , providing thus an explanation to the failure of the quark model calculations which do not take into account $`U_A(1)`$-breaking induced interactions.
This year, DA$`\varphi `$NE, the high luminosity $`\varphi `$ factory, will perform precise measuraments of $`\varphi `$ radiative decays. The Novosibirsk CMD and SND collaborations already reported, among others, measurements of $`\varphi \pi ^0\pi ^0(\eta )\gamma `$ and $`\pi ^+\pi ^{}\gamma `$ . On the theoretical side the $`\varphi \pi \pi \gamma `$ has been considered by a number of authors . In particular Bramon, Grau and Pancheri (B.G.P.) considered vector meson and chiral loop contributions. By itself the vector meson contribution turns out to be small whereas the chiral loops lead to a broad pion invariant mass spectrum which could easily be distinguished by experiments.
In this contribution we report calculations for the $`\varphi \pi ^0\pi ^0\gamma `$ decay within the $`U(3)\times U(3)`$ model where intermediate scalar resonances naturally appear. From the theoretical point of view this is a clean process, since no final state radiation exists (as compared to the decay in charged pions) and the pseudoscalar mixing angle is not involved (as in the $`\pi ^0\eta \gamma `$ case).
## II Scalar meson contributions to $`\varphi \pi ^0\pi ^0\gamma `$
The process under consideration is generated at one loop level. The diagrams contributing to this process are depicted in Fig1.
The amplitude arising from the scalar contributions is given by:
$$M(\varphi \pi ^0\pi ^0\gamma )=eG(m_{\pi \pi }^2)T_{\mu \nu }\eta ^\mu ϵ^\nu $$
(1)
where
$$T_{\mu \nu }=Q.kg_{\mu \nu }k_\mu Q_\nu $$
(2)
and
$$G(m_{\pi \pi }^2)=\frac{g_{\varphi K^+K^{}}}{2\pi ^2M_\varphi ^2}FL(m_{\pi \pi }^2)$$
(3)
with the loop function
$$L(m_{\pi \pi }^2)=\frac{1}{2(ab)}\frac{2}{(ab)^2}[f(\frac{1}{b})f(\frac{1}{a})]+\frac{a}{(ab)^2}[g(\frac{1}{b})g(\frac{1}{a})].$$
(4)
where
$`f(z)=\{\begin{array}{cc}\left(arcsin\left[\frac{1}{2\sqrt{z}}\right]\right)^2\hfill & \text{ }z>\frac{1}{4}\hfill \\ \frac{1}{4}\left(ln\frac{\eta _+}{\eta _{}}i\pi \right)^2\hfill & z<\frac{1}{4}\hfill \end{array}`$
$`g(z)=\{\begin{array}{cc}(4z1)^{\frac{1}{2}}\mathrm{arcsin}\left[\frac{1}{2\sqrt{z}}\right]\hfill & \text{ }z>\frac{1}{4}\hfill \\ \frac{1}{2}(14z)^{\frac{1}{2}}(\mathrm{ln}\frac{\eta _+}{\eta _{}}i\pi )\hfill & z<\frac{1}{4}\hfill \end{array}`$
with
$$\eta _\pm =\frac{1}{2}[1\pm (14z)^{\frac{1}{2}}],a=\frac{M_\varphi ^2}{m_{k^+}^2},b=\frac{m_{\pi \pi }^2}{m_{k^+}^2}.$$
(5)
The F factor appearing in Eq.(3) contain the information on the coupling constants.
$$F=2(g_{KK\pi \pi }\frac{g_{\sigma \pi \pi }g_{\sigma KK}}{m_{\pi \pi }^2m_\sigma ^2+i\mathrm{\Gamma }_\sigma m_\sigma }\frac{g_{f\pi \pi }g_{fKK}}{m_{\pi \pi }^2m_f^2+i\mathrm{\Gamma }_fm_f})$$
(6)
The three and four-meson couplings are given by the model as
$`g_{\sigma KK}`$ $`=`$ $`{\displaystyle \frac{m_\sigma ^2m_K^2}{2f_K}}(cos\varphi \sqrt{2}sin\varphi );g_{\sigma \pi \pi }={\displaystyle \frac{m_\sigma ^2m_\pi ^2}{2f_\pi }}cos\varphi `$ (7)
$`g_{fKK}`$ $`=`$ $`{\displaystyle \frac{m_f^2m_K^2}{2f_K}}(sin\varphi +\sqrt{2}cos\varphi );g_{f\pi \pi }={\displaystyle \frac{m_\sigma ^2m_\pi ^2}{2f_\pi }}sin\varphi `$ (8)
$`g_{KK\pi \pi }`$ $`=`$ $`{\displaystyle \frac{m_\sigma ^2m_K^2}{4f_Kf_\pi }}.`$ (9)
The pion invariant mass spectrum is obtained
$$\frac{d\mathrm{\Gamma }}{dm_{\pi \pi }}=\frac{\alpha _{em}}{4\pi }\frac{m_{\pi \pi }}{M_\varphi }(\frac{g_{\varphi KK}}{4\pi })^2(\frac{1}{4\pi })^2(\frac{M_\varphi }{m_K})^4|L(m_{\pi \pi }^2)|^2|F|^2(1\frac{m_{\pi \pi }^2}{M_\varphi ^2})^3\sqrt{1\frac{4m_\pi ^2}{m_{\pi \pi }^2}}$$
(11)
We observe from Eqs(7,8) that the energy spectrum depend on the scalar mixing angle (in the $`\{|S>,|NS>\}`$ basis) $`\varphi `$ and the scalar meson masses. In particular is highly sensitive to the chosen value for the mixing angle. It is worth to remark that for a sigma meson mass above 600 MeV the theoretical predictions for the energy spectrum yield a desastrous result as compared with the experimental results. Our results reduces to those of B.G.P. -upto an overall normalization factor- in the very heavy (and non-mixed) scalars limit (as compared to the typical 1GeV scale). We have been unable to trace back the difference in normalization of the two approaches (a factor $`\sqrt{(}3)`$ in the amplitud).
The LSM results for the energy spectrum in Eq.(11) are shown in Fig. 2. We use $`\varphi =9^0`$, $`m_f=980`$ MeV, $`\mathrm{\Gamma }_f=70`$MeV and $`m_\sigma =560`$ MeV in the numerical evaluations. The sigma width is dictated by the model as
$$\mathrm{\Gamma }_\sigma =\frac{3m_\sigma ^3}{32\pi f_\pi ^2}((1\frac{m_\pi ^2}{m_\sigma ^2})cos(\varphi ))^2\sqrt{14\frac{m_\pi ^2}{m_\sigma ^2}}.$$
(12)
In table 1 we also show the theoretical predictions arising from Eq( 11) for the total and partial (i.e. integrated over a limited region of the energy spectrum) Branching Ratios. For comparison we also included the experimental results reported by the Novosibirsk groups. Within experimental errors, agreement is satisfactory.
Table 1
| $`m_{\pi ^0\pi ^0}(MeV)`$ | $`BR(CMD2)(\times 10^4)`$ | $`BR(SND)(\times 10^4)`$ | $`BR_{TH}(\times 10^4)`$ |
| --- | --- | --- | --- |
| $`>550`$ | $`1.06\pm 0.09\pm 0.06`$ | | $`0.99`$ |
| $`>700`$ | $`0.92\pm 0.08\pm 0.06`$ | $`1.00\pm 0.07\pm 0.12`$ | $`0.90`$ |
| $`>900`$ | $`0.57\pm 0.06\pm 0.04`$ | $`0.50\pm 0.06\pm 0.06`$ | $`0.53`$ |
| $`total(>2m_\pi )`$ | $`1.08\pm 0.17\pm 0.09`$ | $`1.14\pm 0.10\pm 0.12`$ | $`1.08`$ |
So far we have used a Breit-Wigner to describe the sigma (and $`f_0(980)`$) propagator. It has been argued that the inclusion of the sigma width in this way strongly breaks chiral symmetry . If we modify the sigma vertices in such a way that the Goldstone Boson nature of the pions is preserved as discussed in , the curve in Fig. 2 is modified. In this case agreement with experimental results is obtained for a lower sigma mass $`m_\sigma =430MeV`$ and the same mixing angle $`\varphi =9^{}`$.
Summarizing, the VEPP-2M SND and CMD2 experimental results for the $`\varphi \pi ^0\pi ^0\gamma `$ results are consistent with a mostly $`\overline{s}s`$ $`f_0(980)`$ meson provided we take into account the effects of the $`U_A(1)`$ breaking in the scalar sector. This process gives also support to the existence of a scalar meson resonance ($`\sigma `$) in the 400-600 MeV. The process under consideration is highly sensitive to the scalar mixing angle and experimental results for this process require $`\varphi 9^0`$ which is consistent with other estimates .
## III Acknowledgements
This work was supported by Conacyt Mexico under project I27604-E. One of us (J.L.L.) whish to acknowledge financial support from Conacyt, Mexico and Concyteg-Guanajuato, Mexico.
## IV References |
warning/0001/astro-ph0001258.html | ar5iv | text | # Probing The Structure of Space-Time with Cosmic Rays
## I Introduction
The study of interactions of high energy Cosmic Rays (nucleons and gamma rays) can provide severe tests of the validity of Special Relativity (SR). Indeed, SR is at the very base of our theories for the description of the Universe, and its validity is generally not questioned. However, several attempts have been made to put under scrutiny the postulates that SR relies upon, like the constancy of the velocity of light, and these studies have contributed accurate limits on possible violations (see e.g. and references therein).
Although it is very well proven that SR and the underlying Lorentz invariance (LI) provide a suitable framework for our low energy effective theories, it is not clear whether more ambitious theories, aiming to a global description of our World, including gravity, are or need to be Lorentz invariant. Theoretically, the need for a unified theory of gravity has led to several models, some of which automatically imply the breaking of LI. This is qualitatively understandable since the very concept of continuous space-time is likely to be profoundly changed by quantum gravitational effects . In this prospective we think it is worth keeping an open mind with respect to the validity of what we define as fundamental theories, and always put them under experimental scrutiny.
This attitude appears even more justified at the present time: observationally, high energy astrophysics is providing a range of opportunities to probe energies (i.e. Lorentz factors, or speeds) much larger than the ones ever obtained, or obtainable in the future, in accelerator experiments. In particular we concentrate our attention on the interactions of cosmic gamma rays and nucleons with some type of universal photon background, i.e. the cosmic microwave (CMB), the far infrared (FIRB), and the radio background. These are the interactions responsible for gamma ray absorption from distant sources (through production of $`e^+e^{}`$) , that should appear as a cutoff in the gamma ray spectrum, and of the well known Greisen-Zatsepin-Kuzmin (GZK) cutoff due to the photopion production in collisions of ultra high energy cosmic rays (UHECRs) off the CMB photons. All these processes have the same general structure: in a Lorentz invariant picture, these reactions would be examples of very low energy processes in the center of momentum, that appear boosted to very large Lorentz factors in the laboratory frame. Testing the presence of the thresholds for these processes is therefore a test of SR up to the large Lorentz factors in the boost.
Our approach is entirely phenomenological, and reasonably model independent. We do not propose, if not as examples, any specific model. We show that when a very general form of LI violation <sup>*</sup><sup>*</sup>*Clearly, breaking LI can have several implications, as for instance the existence of a preferred reference frame, which is the one in which all the calculations need to be carried out, since breaking LI also invalidates the transformations that allow us to change reference frame. This reference frame could be identified with the one comoving with the expansion of the universe (no peculiar motion), in which the microwave background is completely isotropic. In fact, there is only one frame with this property, being all other frames experiencing the dipole anisotropy, and therefore distinguishable. is explicitly allowed in the dispersion relation between energy and momentum of particles, the threshold momenta for some of these reactions may drastically change or even become unphysical (i.e. the process does not occur), unless the LI violation is introduced at a length scale much smaller than the Planck scale. A similar approach, although more specific, was proposed in (see Sect. 3).
The aim of the present paper is to discuss the possibility that future CR experiments could test Special Relativity and/or the structure of space-time . In our opinion, the (lack of) knowledge of the sources of the CR’s under consideration does not allow us to support the idea that the present experimental situation gives evidence for violation of LI or of continuity of space-time. Rather we stress that a verification of the existence of the quoted thresholds would entail a lower limit on the mass scale of effective LI violations.
The paper is planned as follows: in section 2 we discuss our parametrization of the LI violations, and discuss briefly some models in which these parametrizations hold; in section 3 we apply our calculations to the case of pair production and photopion production in high energy cosmic ray interactions. We conclude in section 4.
## II Breaking of Lorentz Invariance
Lorentz invariance implies that the modulus of any four-vector is unchanged when changing reference frame; for instance for the four-momentum of a particle we have (we always put $`c=\mathrm{}=1`$)
$$P_\mu P^\mu =E^2p^2=const=m^2$$
$`m`$ being the invariant particle mass and $`p=|\stackrel{}{p}|`$.
Violations of Lorentz invariance will therefore affect in general the dispersion relation above. Without referring to any specific model (later we will discuss explicit implementations), we write a modified dispersion relation obeying the following postulates:
1. Violations are a high energy effect, i.e. they vanish at small momenta.
2. Violations are universal, i.e. do not depend on the particle type, if not (possibly) through the particle mass.
3. Rotation invariance remains exact.
Clearly requirements 2 and 3 may be relaxed and in fact there are examples in this sense . We write then the modified dispersion relation as follows:
$$E^2p^2m^2=p^2f(\frac{p}{M})+m^2g(\frac{p}{M})$$
(1)
where the mass $`M`$ parametrizes the violation of Lorentz invariance (or an essential discreteness of space-time, as for instance suggested by some models of quantum gravity (e.g. and )). Even in the framework defined above this is not the most general violation term. However, in the regime we are interested in, $`mpM`$, the left hand side of eq. (1) is small compared to $`p^2`$ and $`E^2`$ so that the other possible terms we can write (containing for instance $`E`$) differ from those in eq. (1) by higher order corrections Also, terms proportional to $`mp`$ can be added. They however do not modify the general framework so we will not include them for sake of clarity..
Since $`p/M1`$ the functions $`f`$ and $`g`$ can be Taylor-expanded to give:
$$E^2p^2m^2=p^2(f(0)+f^{}(0)\frac{p}{M}+f^{\prime \prime }(0)\frac{p^2}{M^2}+\mathrm{})+$$
$$m^2(g(0)+g^{}(0)\frac{p}{M}+g^{\prime \prime }(0)\frac{p^2}{M^2}+\mathrm{})$$
(2)
In the limit $`M\mathrm{}`$ one must recover the Lorentz invariant dispersion relation, so $`f(0)=g(0)=0`$; moreover the linear term might be absent, as we will see later; if it is present, the quadratic term is negligible at the momenta we consider. The coefficients in front of the first and second derivatives are left unconstrained. However, if $`|f^{}(0)|,|g^{}(0)|1`$ or $`|f^{\prime \prime }(0)|,|g^{\prime \prime }(0)|1`$, then the functions $`f`$ and $`g`$ would be strongly varying and possibly oscillating, which would be unphysical. Moreover, in cases where the calculations of the LI breaking can be carried out explicitly (see section 2A), these coefficients turn out to be of order unity (the same motivation justifies the assumption that these coefficients are not much smaller than unity). Therefore we decided to adopt here a phenomenological approach and assume that possible small deviations from unity are embedded in the mass scale $`M`$. The validity of this approach will be checked a posteriori.
We are thus led to the following classification of Lorentz non-invariant dispersion relations:
$$I_\pm :E^2p^2m^2\pm \frac{p^3}{M}(\pm \frac{m^2p}{M})$$
(3)
$$II_\pm :E^2p^2m^2\pm \frac{p^4}{M^2}(\pm \frac{m^2p^2}{M^2})$$
(4)
where $`I`$ ($`II`$) stands for first (second) order modification and the terms in parenthesis come from the expansion of $`g`$.
We can define a critical momentum $`p_c`$ where the correction, for massive particles, equals $`m^2`$, which is the momentum for which we expect that deviations from normal relativistic kinematics become relevant. We have $`p_c2\times 10^{15}`$ eV ($`10^{13}`$ eV) for protons (electrons) in the case $`I_\pm `$ and $`p_c3\times 10^{18}`$ eV ($`10^{17}`$ eV) for protons (electrons) in the case $`II_\pm `$. In the cases in parenthesis in eqs. (3) and (4) we have always $`p_cM`$, so these modifications do not lead to observable consequences at the energies we are interested in; we will therefore put $`g(\frac{p}{M})=0`$ in the following. We will come back to this point when we will describe specific examples of Lorentz violating theories . It is clear that the values of $`p_c`$ given above are calculated in a specific frame; in fact, generally, violations of LI imply the existence of a privileged frame. We choose this frame to be the one comoving with the expansion of the universe (we name it “universal frame”), and we argue that this choice is in fact not arbitrary: this is the only possible frame where the microwave background is isotropic (the same holds for the other backgrounds, provided the sources are homogeneously and isotropically distributed). Moreover, neglecting the proper motion of the Earth, this is the reference frame in which we live and measure the thresholds for physical processes. On a more practical ground, it is worth noticing that giving up LI, the Lorentz transformations do not correctly give the transformations laws of energy and momentum between different frames although in principle it is possible to write modified transformations of energy and momentum (see for instance ) which reproduce, at least at the perturbative level chosen, the dispersion relations above. In this case, however, there is much more arbitrariness than in modifying the dispersion relation and we do not pursue this approach in this paper. It is worth noting however that the Lorentz Transformations are derived in a LI theory by the requirement of the invariance of a fundamental interval, so in a sense a modification of the dispersion relation is more fundamental.
There is a further, important point to discuss before computing particle production thresholds. In fact, when one gives up relativistic invariance, energy-momentum conservation is not anymore guaranteed, and relativistic quantum field theory may fail, so there is no guidance in deriving cross sections. However, our point of view here is to derive the consequences of experimental verification of the existence of the particle thresholds, so we will assume in the following exact energy-momentum conservation and relativistic dynamics in the preferential frame.
### A Models
In the following we will briefly discuss some models in which modified dispersion relations are actually obtained. It has been argued on very general grounds that quantum gravity effects do modify particle propagation at scales close to the Planck scale. In particular, modified uncertainty relations and existence of a $`minimum`$ proper length (see e.g. and references therein) (implying modifications of Lorentz transformations), and light-cone fluctuations (see e.g. and references therein) are fairly general implications of quantum gravity and modify the dispersion relations.
An approach where a fundamental mass/momentum scale $`M`$ is introduced to parametrize the deviations from LI can be pursued by writing the commutators between (space) boosts and translation generators of the Poincare’ algebra in a modified form , leading to
* $`II_+m^2=M^2\mathrm{sin}^2(\frac{E}{M})p^2E^2p^2\frac{E^4}{12M^2}`$
* $`II_{}m^2=M^2\mathrm{sinh}^2(\frac{E}{M})p^2E^2p^2+\frac{E^4}{3M^2}`$
A theory with both exact conservation of energy and momentum and LI violations can be obtained if it is possible to construct a local theory which is symmetric under the modified Poincaré group
Models giving the $`I_\pm `$ dispersion relations can be derived from some quantum gravity approach. The basic idea is that quantum fluctuations of gravity cause, at a scale of Plank mass, the vacuum to behave like a stochastic medium and this introduces non-zero, energy dependent non-diagonal terms in the metric ,
* $`I_+`$: $`g_{\mu \nu }g_{\mu \nu }+h_{0i}h_{0i}=U_i\stackrel{}{U}=\frac{E}{M}\widehat{p}`$
* $`I_{}`$: $`g_{\mu \nu }g_{\mu \nu }+h_{0i}h_{0i}=U_i\stackrel{}{U}=\frac{E}{M}\widehat{p}`$ ,
with consequent modifications of the dispersion relations. Similar modifications may hold in brane models with large (or even infinite) extra dimensions and $`TeV`$ scale quantum gravity , where in some cases Poincare’ invariance is broken explicitely .
In the models presented above, quantum gravity effects are described by fluctuations around a flat background metric. An entirely different approach is followed in the loop approach to quantum gravity (see e.g and references therein) in which the geometry itself emerges non perturbatively. It has been shown that this leads to an essential discretization of space, and to modifications of the dispersion relations which fall in the class $`I_\pm `$ if parity is broken; the violation might be milder (or absent) in parity conserving models.
We finish quoting a pioneering approach by Kirshnitz and Chechin , stimulated by the appearance of the papers by Kuzmin and Zatsepin , and Greisen . It is a classical approach, in which the free lagrangian of a (classical) point particle is modified. The theory is defined replacing the pseudo-Euclidean space-time of SR with a Finslerian space . In this approach the dispersion relation becomes
$$m^2=[1+\overline{g}(\alpha \frac{p^2}{E^2})](E^2p^2),$$
(5)
where $`m`$ is the particle mass, and $`\overline{g}`$ is a homogeneous function of the dimensionless parameter $`\alpha `$ that parametrizes the LI violations, $`\alpha =\frac{m^2}{M^2}`$ in terms of the scale M. This gives rise to a mild violation, disappearing for massless particles. With an appropriate choice of the function $`\overline{g}`$, this gives rise to the terms in parenthesis in eqs. (3) and (4).
## III Threshold calculations with modified dispersion relations
In this section we describe the calculation of the kinematic thresholds for some processes in the framework of the modified dispersion relations between energy and momentum introduced in the previous section. We choose two processes which are of astrophysical relevance and that will be accessible to next generation cosmic ray experiments: pair production in photon-photon scattering and photopion production in nucleon-gamma scattering. The first process is responsible for the absorption of high energy gamma rays from distant sources, while the second is responsible for the well known (and currently unobserved) GZK cutoff.
### A $`e^+e^{}`$ production in photon-photon interactions
The process under investigation is $`\gamma \gamma e^+e^{}`$. The energy of the background photons is taken as $`\omega =|k|`$ (where $`k`$ is the photon momentum) since at the typical background momenta (FIRB: $`\omega 0.01`$ eV, CMB: $`6\times 10^4`$ eV, and radio: $`4\times 10^9`$ eV for the peak of the corresponding radiance distributions) the corrections are entirely negligible. We compute the threshold, assuming that the CR particle and the background photon collide head-on, and the final particles are collinear. This is indeed not an arbitrary configuration, but the one that provides the minimum energy for which the process can occur in the universal frame.
Rotational invariance then implies that in this reaction the momenta of final particles are equal, and that particles move in the direction of the primary, so that the problem is one-dimensional; the energy of the background photon is assumed to be equal to the average for the photon background in consideration.
Let then ($`E,p`$) be the energy and (modulus of) the momentum of the incident photon. Writing the relations of conservation of energy and momentum in the laboratory frame, and using the modified dispersion relation between energy and momentum, after some trivial algebra, and neglecting sub-leading terms in the range of momenta we are considering, we get the following general equations for the threshold:
$$I_\pm :\pm \alpha _Ix^3+x1=0,$$
(6)
$$II_\pm :\pm \alpha _{II}x^4+x1=0,$$
(7)
where $`x=p_{th}/p_0`$, $`p_0=m^2/\omega `$ is the threshold for $`M\mathrm{}`$
(i.e. the usual threshold, $`3\times 10^{13}`$ eV for the FIRB as background, $`5\times 10^{14}`$ eV for the CMB, and finally $`6\times 10^{19}`$ for the radio background), $`m`$ the electron mass, and:
$$\alpha _I=\frac{p_0^3}{8m^2M};\alpha _{II}=\frac{3p_0^4}{16m^2M^2}.$$
(8)
The modified thresholds are the positive real solutions of the equations above, and in particular, if more than one positive solution is present, the one which goes to $`x=1`$ as $`M\mathrm{}`$. In table I we report the solutions (if any) of the eqs. (6) and (7) in the case $`M=M_P`$, reasonable in quantum gravity models, for the infrared, microwave, and radio background.
The general feature of these solutions is that a positive modification of the dispersion relation tends to move the thresholds towards lower momentum values; negative modifications tend to lead to complex solutions, meaning that the kinematics of the process becomes forbidden, i.e. the threshold disappears. Case II is less predictive, as it is obvious being the modification of second order in the (small) quantity $`p/M`$. Clearly a more detailed calculation, accounting for the integration over different scattering angles and over the whole spectrum of the photon background, could give slightly different results, not changing however our basic findings.
To assess the role of CR experiments to detect violations of LI we now assume that in present or future experiments it will be possible to isolate the effect of pair production on different photon backgrounds, by measuring the respective thresholds, with an experimental uncertainty on the measurement of the threshold energy of $`100\%`$. This is a conservative (and probably pessimistic) approach, unless possibly for the radio case. We therefore derive a lower limit on $`M(M_P)`$, which we report in table II.
Note that all these limits are more stringent than the few limits on $`M`$ already obtained (see for instance ). Exceptions to this statement come from the cases $`II_\pm `$ for the IR and microwave backgrounds, where the limits on $`M`$ are appreciably smaller that $`M_P`$, leaving more room for non Lorentz invariant theories.
We now discuss the present situation of VHE/UHE $`\gamma `$-ray astronomy in view of the possibility of testing LI. In doing this it is important to keep in mind the quite different observational situation existing for the CMB versus the IR and Radio backgrounds: while the former is determined with very high accuracy, the IR and Radio backgrounds are very poorly known and difficult to access through observational investigation, mainly due to the emission and/or absorption processes in our Galaxy.
Observationally, the infrared gamma-ray cutoff seems to be the one more easily approachable. Indeed signs of a cutoff in the TeV spectrum of a few blazars seem to be already present () and helped to impose some constraints on the extragalactic far-infrared background (FIRB). Unfortunately, as it is clear from table I, the case of interaction with the FIRB is not very predictive, in the sense that the threshold is not appreciably changed with respect to the case of a Lorentz-invariant theory, for most of the possible parametrizations of LI breaking (using $`M=M_P`$). The case $`I_{}`$ is an exception: in this case there is no acceptable solution of the equations that define the threshold, which means that the process is not possible at all. In other words, in this case no cutoff should be seen in astrophysical observations because photons can come unattenuated from all distances. Viceversa, if a cutoff is observed, then a lower limit on $`M`$ can be imposed.
As far as the photon absorption on the CMB radiation is concerned, the experimental situation is more uncertain, due to the higher energy of the $`\gamma `$-rays involved (although there might be suggestion of attenuation due to CMB, see e.g. ). This is theoretically more restrictive: in the scenario $`I_+`$ the threshold is moved to smaller values. For $`M=M_P`$, photons with energy as small as $`25`$ TeV from a distant source should be absorbed. In order for the threshold to be unchanged, the scale $`M`$ must exceed $`800M_P`$. In the scenario $`I_{}`$ the process becomes not allowed, while in the other cases the threshold remains unchanged.
The case of scattering off the radio background is the most interesting. We predict, for $`M=M_P`$, either no threshold (cases $`I_{}`$ and $`II_{}`$) or thresholds which are much smaller that the canonical ones. This could be of great importance for top-down models of UHECRs, where gamma rays are supposed to have an important role in the composition (e.g. see ). In particular the proton/gamma ratio is determined by the interaction of the gamma rays with the radio background at frequencies smaller than a few MHz. Unfortunately, as mentioned above, the radio background at these frequencies is extremely uncertain, and not accessible to any direct measurement, due to the strong free-free absorption in the disc and halo of our own Galaxy.
### B Pion photoproduction in UHECR interactions: the GZK cutoff
This case is slightly more involved since the masses of the final particles are different, so that even for a rotation invariant modification, final momenta are not in the same ratio as the masses. However we checked that assuming the ratio of final momenta as in the Lorentz invariant theory, we introduce only higher order corrections. Using this prescription we obtain the following two equations for the threshold, with the same symbols used in the previous section:
$$I_\pm :\pm \alpha _Ix^3+x1=0$$
(9)
$$II_\pm :\pm \alpha _{II}x^4+x1=0$$
(10)
where $`x=\frac{p_{th}}{p_0}`$ and $`p_0=\frac{m_\pi ^2+2m_\pi m_p}{4\omega }`$ is the conventional threshold (for $`M\mathrm{}`$), $`m_p`$ is the proton (neutron) mass, and $`m_\pi `$ is the pion mass. The coefficients $`\alpha _{I,II}`$ are defined as follows:
$$\alpha _I=\frac{2p_0^3}{(m_\pi ^2+2m_\pi m_p)M}\frac{m_\pi m_p}{(m_\pi +m_p)^2};$$
(11)
$$\alpha _{II}=\frac{3p_0^4}{(m_\pi ^2+2m_\pi m_p)M^2}\frac{m_\pi m_p}{(m_\pi +m_p)^2}.$$
Fixing $`M=M_P`$ we have the solutions reported in table III.
The general trend of the solutions is the same as in the case of pair production. Here however we can easily see that the consequences are even more evident: in all cases either the threshold disappears and the process becomes forbidden (cases $`I_{}`$ and $`II_{}`$), or the threshold is appreciably lowered (cases $`I_+`$ and $`II_+`$). For $`I_{}`$ or $`II_{}`$ parametrizations, the GZK cutoff is completely washed out and particles (nucleons) should be able to reach us from any distance. In the other cases the threshold falls in a region that is easy to discard even on the basis of the present data, unless the scale of breaking of LI is $`3\times 10^{13}`$ ($`500`$) times larger than the Planck mass for the case $`I_+`$ ($`II_+`$).
As in the previous section, if one assumes that the apparent absence of the GZK cut-off is an incidental fact, possibly due to our ignorance of the sources of CR’s at these energies, and that CR experiments may find some remnant of the cut-off, we can derive limits on the parameters of LI violations. These limits are reported in table III, again assuming a $`100\%`$ error on the location of the cut-off. From the values reported it is clear that UHECR experiments soon to be available will provide a powerful tool to explore the very-small-scale structure of space-time.
## IV Conclusions
Special Relativity is at the base of our understanding of the physical world. However any physical concept should always be put under stringent experimental verification. This is particularly true in the case of SR, specially in connection with the quest for a fundamental theory of Nature, which includes gravity. Any such theory will imply a full knowledge of what the vacuum really is at the level where quantum fluctuations build the space-time, possibly modelling it with a non-trivial geometry. Some first attempts to reach this goal seem to suggest that Lorentz invariance, one of the building blocks of our current low energy theories, could be broken at extremely high energy.
Lacking a true theory of quantum gravity, the most we can do at present is to adopt a phenomenological approach and ask ourselves whether there is any probe that can be used to check the validity of SR at very high energy. This is precisely the approach adopted in this paper, where high energy cosmic rays (nucleons and gamma rays) have been used as probes. We found that the thresholds for pair production and photopion production off some universal photon backgrounds, which are experimentally accessible or will be in the next generation of cosmic ray experiments, are often profoundly affected by the possibility of breaking LI at supra-Planck scales. As a consequence, any quantum gravity theory in which LI is a casualty must face the cosmic ray bound in order to be viable.
Our calculations were based on a perturbative but quite general modification of the dispersion relation between energy and momentum of a particle. This modification affects in a fundamental way the calculation of the thresholds for pair production and photopion production, making these processes forbidden in some cases, or lowering the thresholds to questionable values in others. Building on this approach we proposed a procedure to obtain very strong constraints on the energy (or length) scale at which a possible LI violation could occur.
This procedure relies upon the possibility that astrophysical observations will be able in the near future to find evidence for these processes, through cutoffs in the TeV spectra of distant blazars or/and through the discovery of possible consequences of the GZK cutoff in the spectrum of UHECRs or/and through studies of the composition of UHECRs. With a few exceptions, these limits are all higher than the Planck scale, which is an indication that breaking LI is not necessarily a safe ingredient of unified theories that require it.
We stress that our approach is purely kinematical and no dynamical effect is considered. Moreover, we assume perfect energy-momentum conservation in order to compute the thresholds. While these are clearly important issues for any calculation aiming to predict values for specific particle production thresholds, and consequently absorption cutoffs in CR spectra, our approach is rather the opposite: we want to discuss the consequences of a possible experimental verification of the presence of particle production thresholds. If a given absorption threshold is experimentally detected, discarding the possibility of miraculous compensations between relativity violations, non conservation of energy-momentum and non-relativistically invariant dynamics, we are forced to conclude that possible violations of LI are smaller than the sensitivity of the experiment, and, to make this statement quantitative, we choose to use a purely kinematical, energy-momentum conserving parametrization.
Although a very tempting possibility, we do not believe that the current experimental situation allows us to draw definitive statements : recent observations of Markarian 501 might suggest the presence of a cutoff in the TeV region, but the unaffected spectrum is not known well enough to exclude that the observed effect is the artifact of a cutoff in the production spectrum. Hopefully the situation will improve with the next generation gamma ray detectors, complemented by neutrino and X-ray detectors that could clarify the origin of the gamma ray emission. Concerning $`\gamma `$rays absorption on CMB photons, experimental data are more scarce, and the situation is not likely to improve in the near future.
Pair production on the universal radio background becomes relevant for ultra-high-energy gamma rays, usually produced in top-down models of UHECRs . The presence or absence of a threshold for this process strongly affects our predictions of the fluxes at the Earth: although very uncertain, the radio background should allow typical gamma ray pathlengths (in a Lorentz invariant world) of the order of $`220`$ Mpc . If a violation of LI made the pair production kinematically forbidden, gamma rays could reach us from any distance and contribute an enormous flux of particles above the GZK cutoff.
As for UHECRs, although the (small) number of events above $`5\times 10^{19}`$ eV seems already incompatible with the presence of the GZK cutoff, the lack of knowledge of the sources does not allow any firm statement on the propagation of primaries in the Universe. The situation will improve dramatically in the near future, as new experiments like HiRes and Auger will collect data and reliable measurements on observables different from the energy spectrum (anisotropy, clustering) will be available with reasonable statistics.
Our conclusions can be summarized in the following points: 1) it is a 40-years old idea that the incompleteness of the present theories could rely upon our ignorance of the vacuum at very high energies. When we will have that knowledge we will probably have a quantum theory of gravity; 2) several attempts to quantize gravity have naturally led to the requirement of violations of the LI; 3) experiments on cosmic rays can represent the most easily approachable tool to probe the structure of space-time on the very small scales; 4) LI violations affect the thresholds for elementary processes relevant for cosmic ray astrophysics and can be observationally tested.
After the submission of this paper we have noticed some papers which reach similar conclusions.
## V Acknowledgements
We are grateful to V. Berezinsky and G. Di Carlo for very useful and stimulating discussions. The work of P.B. was funded by the DOE and the NASA grant NAG 5-7092 at Fermilab. |
warning/0001/hep-th0001030.html | ar5iv | text | # A geometric discretisation scheme applied to the Abelian Chern-Simons theory
## 1 INTRODUCTION
A very useful way to regularize a quantum field theory is provided by the lattice formulation introduced by Wilson . (See for example, for a detailed treatment). However, this formulation has difficulty in capturing topological features of a field theory, for example, the topological theta-term in QCD. It is therefore of interest to investigate alternative discretisation schemes. In this paper we describe an alternative scheme which is applicable to antisymmetric tensor field theories including Abelian gauge theories and fermion field theory in the Kähler–Dirac framework , and which is well-suited for capturing the topological features of such theories. It is based on developing analogies between the different types of fields and the way they appear in a quantum field theory with a corresponding list of discrete variables and operators. The method is valid for any arbitrary compact 3-manifold without boundary. We illustrate the scheme by applying it to the pure Abelian Chern-Simons gauge theory in 3 dimensions, and to a doubled version, the so-called Abelian BF gauge theory. In the latter case the topological features of the theory are completely reproduced by our discretisation scheme, even before taking the continuum limit. We illustrate this by explicit numerical calculations of the discrete partition function when the spacetime is $`S^3`$ or is a lens space, L(p,1).
The discretisation scheme involves several ingredients which are not used in the standard Wilson lattice formulation. These include a triangulation of the spacetime (i.e. a decomposition into hyper-tetrahedra rather then hyper-cubes) and a mathematical tool called the Whitney map . These have previously been used to discretise field theories in where various convergence results (e.g. convergence of the discrete action to the continuum action) were established. Triangulations of spacetime are also used in discretisations of other quantum field theories, for example in quantum gravity (see for a review). The thrust of the present work is quite different. We set up the discretisation in such a way that the geometric structures of the continuum field theory are mirrored by analogous structures in the discrete formulation. As we will see, this requires a certain doubling of fields. With this doubling the topological features of Abelian Chern-Simons theory are completely captured by the formulation. An alternate scheme of discretisation proposed previously in which does not involve the doubling of fields fails to capture the topological features of the Abelian Chern-Simons theory as we will show numerically in Section 5. A mathematical treatment of this discretisation scheme has been given in . A short version of some of the results has been published in .
Our aims in this paper are, firstly, to make the techniques and results of accessible to a wider audience, and secondly, to demonstrate a practical numerical implementation of the discretisation scheme. Numerical implementations are developed for Abelian Chern-Simons theory defined on the three dimensional sphere, $`S^3`$, and on the lens spaces $`L(p,1)`$ for $`p=1,2`$ and $`3`$. Our discretisation scheme is the only one which numerically reproduces the exact topological results for Abelian Chern-Simons theory .
This paper is organised as follows. In Section 2 the discretisation scheme is described together with a summary of the topological results used to set it up. In Section 3 some features of the Abelian Chern-Simons theory on general 3-manifolds are reviewed. In Section 4 the discretisation scheme is applied to this theory. An expression for the partition function of the resulting discrete theory is derived in terms of the data specifying the triangulation of the spacetime. In Section 5 the numerical evaluations of the partition functions corresponding to these triangulations are presented. In Section 6 we summarise our conclusions.
## 2 THE DISCRETISATION SCHEME
Let us start by considering a general quantum field theory defined on an arbitrary manifold M of dimension D. Suppose the theory has fields $`\varphi ^p(\stackrel{}{x})`$ where $`\stackrel{}{x}M`$, and where $`\varphi ^p`$ is a $`p`$-form (antisymmetric tensor field of degree $`p`$ defined on $`M`$). The Lagrangian for the system involves the fields, the Laplacian operator, and possibly (as for the Chern-Simon theory) an antisymmetrized first-order differential operator.
In differential geometric terms, the theory is constructed using the following objects which are defined on the manifold $`M`$: $`p`$-forms $`\varphi ^p`$ which are generalised antisymmetric tensor fields, the exterior derivative $`d:\varphi ^p\varphi ^{p+1}`$, the Hodge star operator $`:\varphi ^p\varphi ^{Dp}`$, which is required to define scalar products, and the wedge operator $`\varphi ^p\varphi ^q=\varphi ^{p+q}`$.
We want to construct discrete analogues of these objects. We begin by summarizing the basic properties of our operators of interest . On a manifold M of dimension D, the operations ($`,,d`$) on p-forms, $`\varphi ^p(p=0,\mathrm{},D)`$, satisfy the following:
1. $`\varphi ^p\varphi ^q=(1)^{pq}\varphi ^q\varphi ^p`$.
2. $`d(\varphi ^p\varphi ^q)=d\varphi ^p\varphi ^q+(1)^p\varphi ^pd\varphi ^q`$.
3. $`\varphi ^p=\varphi ^{Dp}`$.
4. $`=(1)^{Dp+1}`$.
5. $`d^2=0`$, $`(d^{})^2=0`$.
6. $`d^{}=(1)^{D(p+1)+1}d`$, ($`d^{}`$ is the adjoint of $`d`$).
The following definitions will also be required
* The Laplacian on p-forms $`\mathrm{\Delta }_p=d_{p1}d_p^{}+d_{p+1}^{}d_p`$.
* The inner product $`\varphi _p,\varphi _{p}^{}{}_{}{}^{}=_M\varphi _p\varphi _{p}^{}{}_{}{}^{}`$.
A few examples might now be helpful. Firstly, consider QED in 4 dimensions. The gauge field $`A=\varphi ^1`$ is a 1-form. The electromagnetic field is a 2-form given by $`F=dA`$. The action for the gauge field in QED is given by
$$S(A)=(F,F)=dA,dA=_MdAdA.$$
(1)
Thus S involves the operators \*, d and the wedge product. Similarly for Abelian Chern-Simons theory the gauge field $`A`$ and electromagnetic field $`F=dA`$ are 1-forms and 2-forms respectively, as for QED. The action for the theory is given by
$$S(A)=_MAdA=A,dA$$
(2)
where the spacetime M is 3-dimensional. Note that in both cases we can think of the action $`S(A)`$ as a quadratic functional of the gauge field $`A`$.
We would like to discretise the fields of $`\varphi ^p`$, the inner product $`,`$, and the operators $`(,,d)`$ such that discrete analogues of their continuum interrelationships hold. To do this it is necessary to first introduce a few basic ideas of discretisation. We start by discretising the manifold $`M`$. This involves replacing $`M`$ by a collection of discrete objects, known as simplices, glued together. We need a few definitions .
Firstly, for $`p0`$, a $`p`$-simplex $`\sigma ^{(p)}=[v_0,\mathrm{},v_p]`$ is defined to be the convex hull in some Euclidean space $`^D`$ of a set of $`p+1`$ points $`v_0,v_1,\mathrm{},v_p^m`$. Here the vertices $`v_i`$ are required to span a $`p`$-dimensional space. This requirement will hold so long as the equations, $`_{i=0}^p\lambda _iv_i=0`$ and $`_{i=0}^p\lambda _i=0`$ admit only the trivial solution $`\lambda _i=0`$ for $`i=0,\mathrm{},p`$ for $`\lambda _i`$ real.
A few examples might clarify the geometry. Consider $`\sigma ^{(0)}=[v_0]`$. This is a point or 0-simplex. Next $`\sigma ^{(1)}=[v_0,v_1]`$ is a line segment or 1-simplex. An orientation can be assigned by the ordering of the vertices, in which case $`\sigma ^{(1)}=[v_1,v_0]`$ for example. The faces of a 1-simplex are its vertices $`[v_0]`$ and $`[v_1]`$ which are 0-simplices. $`\sigma ^{(2)}=[v_0,v_1,v_2]`$ is a triangle or 2-simplex. We note that an even permutation of the vertices has the same orientation as $`\sigma ^{(2)}`$ while an odd permutation reverses it and will be written as $`\sigma ^{(2)}`$. The faces of a 2-simplex are its edges $`[v_0,v_1]`$, $`[v_1,v_2]`$, and $`[v_2,v_0]`$. Finally, $`\sigma ^{(3)}=[v_0,v_1,v_2,v_3]`$ is a tetrahedron or 3-simplex. Its faces are the four triangles $`[v_0,v_1,v_2],[v_0,v_2,v_3],[v_0,v_1,v_3]`$ and $`[v_1,v_2,v_3]`$ which bound it.
An important feature of our discretisation scheme is that the original simplices are subdivided by using simplex barycentres. Geometrically the barycentre of a $`p`$-simplex, $`\sigma ^{(p)}`$, is the point which represents its “center of mass”. We denote the barycentre of $`\sigma ^{(p)}=[v_0,\mathrm{},v_p]`$ as the point
$$\widehat{\sigma }^{(p)}=\frac{1}{p+1}\underset{i=0}{\overset{p}{}}v_i.$$
(3)
As an example, the barycentre of $`\sigma ^{(1)}=[v_0,v_1]`$ is the midpoint of the line segment which joins the vertices $`v_0`$ and $`v_1`$.
We can now describe a particular way that a given manifold $`M`$ can be discretised. Let $`S`$ be a collection of simplices $`\{\sigma _i^{(n)}\}`$, $`n=0,1,\mathrm{},D`$, with the property that the faces of the simplices which belong to $`S`$ also belong to it. The elements of $`S`$ glued together in the following way is known as a simplicial complex:
1. $`\sigma _i^{(n)}\sigma _j^{(k)}=0`$ if $`\sigma _i^{(n)}`$, $`\sigma _j^{(k)}`$ have no common face.
2. $`\sigma _i^{(n)}\sigma _j^{(k)}0`$ if $`\sigma _i^{(n)}`$, $`\sigma _j^{(k)}`$ have precisely one face in common, along which they are glued together.
In many cases of interest(including all 3-manifolds and all differentiable manifolds ), $`M`$ can be replaced by a complex $`K`$ which it is topologically equivalent to. $`K`$ is then said to be a triangulation of $`M`$(Note this triangulation is not unique). In this way of discretising M, the building blocks are zero, one, $`\mathrm{}`$, D-dimensional objects, all of which are simplices e.g.generalised oriented tetrahedra.
We now observe that the same manifold can be discretised in many different ways. In the discretisation described, we used simplices. We could just as well have used generalised oriented cubes.
There is another method of discretising a manifold which is the dual of the simplicial discretisation just described. It associates with a simplicial complex $`K`$ a dual complex $`\widehat{K}`$. We proceed to describe this construction. We will see that the basic objects of the dual complex $`\widehat{K}`$ are again zero, one, two, $`\mathrm{}`$, D dimensional objects, but this time they are no longer simplices. We illustrate the method by considering a manifold, $`M`$, which is a disc. This is a manifold with a boundary. We triangulate this by the simplicial complex, $`K`$, shown in Figure 2.
Now consider the barycentres of the building blocks of the simplicial complex $`K`$. We have the following list shown in Table 1.
Pictorially the simplicial complex $`K`$ with its barycentres is shown in Figure 3.
We can now construct the dual description of a triangulation K. Geometrically this utilises a discrete analogue of the Hodge $``$ operator. Recall the * operator maps a p-form to a $`(Dp)`$-form , where D is the dimension of the manifold M on which the p-form is defined. In the dual geometrical decomposition of the manifold, we want to set up a correspondence between a p-dimensional object and a $`(Dp)`$ dimensional object.
This is done as follows. We first construct $`(Dp)`$ dimensional objects whose vertices are barycentres of a sequence of successively higher dimensional simplices, where each simplex is a face of the following one. In other words $`(Dp)`$ dimensional objects of the form $`\{\widehat{\sigma }_p,\widehat{\sigma }_{p+1},\mathrm{},\widehat{\sigma }_D\}`$, where $`\sigma _n`$ is a face of $`\sigma _{n+1}`$. The orientation of these are set so as to be compatible with the manifold. Joining these objects together gives us the dual of $`\sigma _p`$.
Thus for instance the map * acts on $`[v_0]`$ as follows.
$$_K:[v_0]ϵ_{01}[\widehat{v_0},\widehat{\sigma _1},\widehat{\sigma }][\widehat{v_0},\widehat{\sigma _3},\widehat{\sigma }]ϵ_{03}.$$
Where the orientation of each of the small triangles has to be coherent with the orientation of the original triangle. This is shown in fig. 3 and leads to mapping $`[v_0]`$ to the shaded two dimensional region. The orientations of the simplices are specified by arrows in the figure. Coherence of orientation means, for example, that the arrow of an edge agrees with the arrow of the triangle to which it belongs. Next we consider $`[v_0,v_1]`$. This is a 1-simplex and is to be mapped to a (2-1)=1 dimensional object. The map is defined as
$$_K:[v_0,v_1][\widehat{\sigma _1},\widehat{\sigma }].$$
Again the orientation of $`[\widehat{\sigma _1},\widehat{\sigma }]`$ has to be coherent with the orientation of the triangles already introduced when the map for $`[v_0]`$ was considered. Similarly
$`_K:[v_1,v_2]`$ $``$ $`[\widehat{\sigma _2},\widehat{\sigma }],`$
$`{}_{K}{}^{}:[v_2,v_0]`$ $``$ $`[\widehat{\sigma _3},\widehat{\sigma }],`$
and finally
$`_K:[v_0,v_1,v_2]`$ $`[\widehat{\sigma }].`$
We then have the alternate discretisation $`\widehat{K}`$ for M shown in fig. 4.
Note when two edges are glued together they must have opposite orientations. We can now give the general rule for mapping an n-simplex $`\sigma _n=[v_0,\mathrm{},v_n]`$ to a ($`Dn)`$ dimensional object ( (D-n) cell) as follows: We think of $`\sigma _n`$ as an element of a simplicial complex K. We have
$$_K:[v_0,\mathrm{},v_n][\widehat{\sigma _n},\widehat{\sigma _{n+1}},\mathrm{},\widehat{\sigma _D}],$$
where $`\widehat{\sigma }_{n+1}`$ is the barycentre of an (n+1)-simplex which has $`\sigma _n`$ as a face. $`\widehat{\sigma }_{n+2}`$ is the barycentre of an (n+2)-simplex which has $`\sigma _{n+1}`$ as a face and so on. These objects have to be coherently oriented with respect to $`[v_0,\mathrm{},v_n]`$. The set of these cells constitutes the dual space $`\widehat{K}`$ of K.
By this procedure we claim, a discrete version of the Hodge star operation * has been constructed. Let us explain. The Hodge * operator involves forms. It maps p-forms in D dimensions to a (D-p)-form. The $`_K`$ map involves not forms but geometrical objects. There is a simple correspondence relation between these two cases. Given a p-form, $`\varphi _p`$, and a p-dimensional geometrical space, $`\mathrm{\Sigma }_p`$, the p-form can be integrated over $`\mathrm{\Sigma }_p`$ to give a number. Thus $`\mathrm{\Sigma }_p`$ and $`\varphi _p`$ are objects that can be paired. We can write this as a pairing
$$(\varphi _p,\mathrm{\Sigma }_p)=_{\mathrm{\Sigma }_p}\varphi _p.$$
In order to proceed, we need to introduce some more structure. We start by associating with a simplicial complex K, containing $`\{\sigma _p^i\}`$ $`(i=1,\mathrm{},K_p;p=0,\mathrm{},D)`$ a vector space consisting of finite linear combinations over the reals of the p-simplices it contains. This vector space is known as the space of p-chains,$`C_p(K)`$. For two elements $`\sigma _p^i,\sigma _p^jC_p(K)`$, a scalar product $`(\sigma _p^i,\sigma _p^j)=\delta _j^i`$ can be introduced. An oriented p-simplex changes sign under a change of orientation i.e. if $`\sigma _p=[v_0,\mathrm{},v_p]`$ and $`\tau `$ is a permutation of the indices $`[0,\mathrm{},p]`$, then $`[v_{\tau (0)},\mathrm{},v_{\tau (p)}]=(1)^\tau [v_0,\mathrm{},v_p]`$, with $`\tau `$ denoting the number of transpositions needed to bring $`[v_{\tau (0)},\mathrm{},v_{\tau (p)}]`$ to the order $`[v_0,\mathrm{},v_p]`$.
Given the vector space $`C_p(K)`$, the boundary operator $`^K`$ can be defined as
$$^K:C_p(K)C_{p1}(K).$$
It is the linear operator which maps an oriented p simplex $`\sigma ^{(p)}`$ to the sum of its (p-1) faces with orientation induced by the orientation of $`\sigma ^p`$. If $`\sigma ^p=[v_0,\mathrm{},v_p]`$, then
$$\sigma ^p=\underset{i=0}{\overset{p}{}}(1)^i[v_o,\mathrm{},\widehat{v_i},\mathrm{}v_p],$$
where $`[v_o,\mathrm{},\widehat{v_i},\mathrm{},v_p]`$ means that the vertex $`v_i`$ has been omitted from $`\sigma ^p`$ to produce the face “opposite” to it.
Given that $`C_p(K)`$ is a vector space, it is possible to define a dual vector space $`C^p(K)`$, consisting of dual objects known as cochains; that is we can take an element of $`C_p(K)`$ and an element $`C^p(K)`$ to form a real number. Since the space $`C_p(K)`$ has a scalar product, namely if $`\sigma _p^i,\sigma _p^jC_p(K)`$ then $`(\sigma _p^i,\sigma _p^j)=\delta _{ij}`$. We can use the scalar product to identify $`C_p(K)C^p(K)`$, so that we can consider oriented p-simplices as elements of $`C^p(K)`$ as well as $`C_p(K)`$. We can write our boundary operation as
$$([v_0,\mathrm{},\widehat{v_i},\mathrm{},v_p],^K[v_0,\mathrm{},v_p])=(1)^i.$$
This suggests introducing the adjoint operation $`d_K`$ defined as
$$(d^K[v_0,\mathrm{},\widehat{v_i},\mathrm{},v_p],[v_0,\mathrm{},v_p])=([v_0,\mathrm{},\widehat{v_i},\mathrm{},v_p],^K[v_0,\mathrm{},v_p]).$$
This is the coboundary operator which maps $`C_p(K)C_{p+1}(K)`$.
Indeed we have
$$d^K[v_0,\mathrm{},v_p]=\underset{v}{}[v,v_0,\mathrm{},v_p],$$
where the sum is over all vertices $`v`$ such that $`[v,v_0,\mathrm{},v_p]`$ is a (p+1) simplex.
The boundary operators $`_K`$ and the coboundary operator $`d_K`$ have the property $`_K_K=d_Kd_K=0`$. Furthermore,
$`d_K:C_p`$ $``$ $`C_{p+1},`$
$`_K:C_p`$ $``$ $`C_{p1}.`$
These operators are the discrete analogues of the operators $`d`$ and $`(1)^{D(p+1)+1}d=d^{}`$ which act on forms.
These operators could be defined only when a scalar product(“metric”) was introduced in the vector space $`C_p`$’s. At this stage we have a discrete geometrical analogue of $`d`$, $`d^{}`$ and $``$. We have also commented on the fact that the operation $``$ maps simplices into dual cells i.e. not simplices. If the original simplicial system is described in terms of the union of the vector spaces of all p-chains then the space into which elements of the vector space are mapped by * is not contained within this space, unlike the situation for the Hodge star operation on forms. We will see that this difference leads inevitably to a doubling of the fields when discretisation, preserving topological structures, is attempted.
We now need a way to relate a p-chain to a p-form. This together with a construction which linearly maps p-forms to p-simplices will allow us to translate expressions in continuum QFT to a corresponding discrete geometrical objects. We start with the construction of the linear maps from p-chains to p-forms due to Whitney .
In order to define this map, we need to introduce barycentric coordinates associated with a given p-simplex $`\sigma ^p`$. Regarding $`\sigma ^p`$ as an element of some $`^N`$, we introduce a set of real numbers $`(\mu _0,\mathrm{},\mu _p)`$ with the property
$`\mu _i`$ $``$ $`0,`$
$`{\displaystyle \underset{i}{}}\mu _i`$ $`=`$ $`1.`$
A point $`x\sigma ^p`$ can be written in terms of the vertices of $`\sigma ^p`$ and these real numbers as
$$x=\underset{i=0}{\overset{p}{}}\mu _iv_i.$$
Note if any set of $`\mu _i=0`$ then the vector $`x`$ lies on a face of $`\sigma ^p`$. One can think of $`x`$ as the position of the center of mass of a collection of masses $`(\mu _0,\mathrm{},\mu _p)`$ located on the vertices $`(v_0.\mathrm{},v_p)`$ respectively. Setting $`\mu _i=0`$ for instance means the center of mass will be in the face opposite the vertex $`v_i`$. The Whitney map can now be defined. We have
$$W^K:C^p(K)\mathrm{\Phi }^p(K),$$
where $`\mathrm{\Phi }^p(K)`$ is a p-form. If $`\sigma ^pC^p(K)`$ then
$$W[\sigma ^p]=p!\underset{i=0}{\overset{p}{}}(1)^i\mu _id\mu _0\mathrm{}\widehat{d\mu _i}\mathrm{}d\mu _p,$$
where $`\widehat{d\mu _i}`$ means this term is missing, and $`(\mu _0,\mathrm{},\mu _p)`$ are the barycentric coordinate functions of $`\sigma ^p`$.
We next construct the linear map from p-forms to p-chains. This is known as the de Rham map. We have
$$A^K:\mathrm{\Phi }^p(K)C^p(K),$$
defined by
$$<A^K(\mathrm{\Phi }^p),\sigma ^p>=_{\sigma ^p}\mathrm{\Phi }^p,$$
for each oriented p-simplex $`K`$.
A discrete version of the wedge product can also be defined using the Whitney and de Rham maps such that $`^K:C^p(K)\times C^q(K)C^{p+q}(K)`$ as follows:
$$x^Ky=A^K(W^K(x)W^K(y)).$$
It has many of the properties of the continuous wedge product in that it is skewsymmetric and obeys the Leibniz rule but it is nonassociative.
At this stage we have introduced all the building blocks necessary to discretise a system preserving geometrical structures. We summarize the properties of the maps introduced in the form of a theorem :
1. $`A^KW^K`$ = Identity.
2. $`dW^K=W^Kd^K`$, where $`d:\varphi ^p\varphi ^{p+1}`$.
3. $`_\beta W^K(\alpha )=<\alpha ,\beta >`$, $`\alpha ,\beta K`$
4. $`d^KA^K=A^Kd`$.
This theorem shows how $`d^K`$ can be considered as the discrete analogue of $`d`$. We now show how $`^K`$ can be considered as discrete analogue of $``$. For this we need barycentric subdivision.
We recall that given a simplicial complex $`\{\sigma _i^p\}`$, i=1,$`\mathrm{},K_p;p=0,\mathrm{},D`$. A set of points(vertices) could be assigned to each simplex, namely $`\widehat{\sigma _i^p}`$. These are the barycentres. These vertices, regarded as vertices of a simplex, subdivide the original simplices to give a finer triangulation of the original manifold. This is a barycentric subdivision map $`BK`$. Clearly the procedure can be repeated to give finer and finer subdivisions in which the simplices become “smaller”. The procedure is illustrated for a triangle in Fig. 5.
Note that all the barycentres are present as vertices of the barycentric subdivision and that $`^K`$ acting on simplices belonging to the simplicial complex, $`K`$, associated with $`[v_0,v_1,v_2]`$ leads to objects which are not, in general, simplices of $`[v_0,v_1,v_2]`$ but belong to a different space $`\widehat{K}`$. However both $`K`$ and $`\widehat{K}`$ are contained in the barycentric subdivision $`B[v_0,v_1,v_2]`$. This is a crucial observation. In order to construct the star map, two geometrically distinct spaces were introduced. The original simplicial decomposition $`K`$ with its associated set of p-chains $`C^p(K)`$ and the dual cell decomposition $`\widehat{K}`$ with its associated set of p-chains $`C^p(\widehat{K})`$. These spaces are distinct. However both belong to the first barycentric subdivision of $`K`$. This allows the use the $`^K`$ operation if we think of $`K`$ and $`\widehat{K}`$ as elements of $`BK`$.
We proceed as follows. Let BK and $`\widehat{K}`$ denote the barycentric subdivision and dual triangulation, and
$$^K:C^p(K)C^{np}(\widehat{K}).$$
However $`C^p(K)`$ and $`C^p(\widehat{K})`$ are both contained in $`C^p(BK)`$ as we have seen. Let
$$W^{BK}:C^p(BK)\varphi ^p(M),$$
denote the Whitney map. Then we have for $`xC^p(K),yC^{np1}(\widehat{K})`$:
1. $`<^Kx,y>=\frac{(n+1)!}{p!(np)!}_MW^{BK}(Bx)W^{BK}(By).<^{\widehat{K}}y,x>=\frac{(n+1)!}{p!(np)!}_MW^{BK}(By)W^{BK}(Bx).`$
2. $`^K=(1)^{np+1}^{\widehat{K}}d^{\widehat{K}}^K`$ on $`C^p(K)`$.
$`^{\widehat{K}}=(1)^{nq+1}^Kd^K^{\widehat{K}}`$ on $`C^q(\widehat{K})`$.
These are the discrete analogues of the interrelationships between $`d,d^{},`$ and $`<,>`$ in the continuum.
Note $`K\widehat{K}`$ and that properties of $`_K,d_K`$ analogous to those for differential forms only hold if $`K`$ , $`\widehat{K}`$ are both regarded as elements of BK. This feature of the discretisation method is, as we shall see, crucial if we want to preserve topological properties of the original system. If a discretisation method is introduced without the $``$ operation in it then as we shall see in section 5 the topological properties of the partition function for the Abelian Chern-Simons gauge theory do not hold.
We proceed to apply these ideas to the Abelian Chern-Simons gauge theory on a compact three manifold $`M`$. First we summarize properties of the continuum field theory.
## 3 Schwarz’s topological field theory and the Ray-Singer torsion
We begin our treatment of continuum field theory by describing Schwarz’s method for evaluating the partition function of the Chern-Simons gauge theory on a three-manifold $`M`$ . We assume that the first real homology group ( to be defined shortly) of the manifold vanishes; this is done for the sake of simplicity. Schwarz’s method is applicable for arbitrary compact 3-manifolds without boundary. Such manifolds are called “homology 3-spheres”. The main example we have in mind are the 3-sphere $`S^3`$ and the lens spaces $`L(p,1)`$, p=1, 2, $`\mathrm{}`$( for a definition and the basic properties of lens spaces see ). The fields of the theory are the 1-forms on M i.e. $`\omega \mathrm{\Omega }^1(M)`$. (In terms of a local coordinate system $`(X^\mu )`$ on M we have $`\omega (x)=\omega _\mu (x)dx^\mu `$.) The action is
$`S(\omega )={\displaystyle _M}\omega d_1\omega ={\displaystyle _M}𝑑x^1𝑑x^2𝑑x^3ϵ^{\mu \nu \rho }\omega _\mu _\nu \omega _\rho .`$ (4)
Here and in the following $`\mathrm{\Omega }^q(M)`$ denotes the space of $`q`$-forms on M (i.e. the antisymmetric tensor fields of degree q) and $`d_q:\mathrm{\Omega }^q(M)\mathrm{\Omega }^{q+1}(M)`$ i.e. the exterior derivative. It has the property $`d_qd_{q1}=0`$ so $`\mathrm{Im}(d_{q1})\mathrm{Ker}(d_q)`$, where $`Im(d_{q1})`$ is the image of the operator $`d_{q1}`$ while $`Ker(d_q)`$ is the null space of the operator $`d_q`$ and the cohomology spaces $`H^q(M)`$ are defined by
$$H^q(M)=\mathrm{Ker}(d_q)/\mathrm{Im}(d_{q1}).$$
The $`H^q(M)`$ are Abelian groups which contain topological information about the manifold. The vanishing of $`H^1(M)`$, for instance, holds if the manifold is simply connected that is any loop in $`M`$ can be smoothly deformed to any other loop in $`M`$ . Note that $`\mathrm{\Omega }^0(M)`$ is the space of functions on M and since $`d_0`$ is the derivative $`\mathrm{Ker}(d_0)`$ consists of the constant functions i.e.
$$H^0(M)=\mathrm{Ker}(d_0)/0=\mathrm{Ker}(d_0)=.$$
Our requirement on M that $`H^1(M)=0`$ implies that
$$\mathrm{Im}(d_0)=\mathrm{Ker}(d_1).$$
A choice of metric on M determines an inner product in the spaces $`\mathrm{\Omega }^q(M)`$ and allows the action (4) to be written as
$`S(\omega )=\lambda <\omega ,(d_1)\omega >,`$ (5)
where * is the Hodge star operator. ( See for background on this and other differential-geometric constructions.) In order to evaluate the partition function of this action by Schwarz’s method requires the introduction of the resolvent for $`S(\omega )`$. The partition function is defined as
$`Z(\lambda )`$ $`=`$ $`N{\displaystyle 𝑑\omega e^{iS(\omega )}}.`$
The main problem in evaluating $`Z(\lambda )`$ is to properly deal with the zeroes of $`S(\omega )`$. These zero modes contain topological information regarding the manifold as the space of zero modes is given by $`\mathrm{Ker}(d_1)`$ and hence should not be discarded. Schwarz introduced an algebraic method (the resolvent method) for dealing with this problem. Although it is only valid for $`S(\omega )`$’s which are quadratic in $`\omega `$, it can be used to analyse $`S(\omega )`$’s constructed on arbitrary compact manifolds without boundary. For systems of this type Schwarz’s method is an algebraic analogue of the problem of gauge fixing. The advantage of the resolvent method is that it can be easily extended to deal with the process of discretisation as we will show.
The resolvent is defined to be the following chain of maps
$`0^{\varphi _0}\mathrm{\Omega }^0(M)^{d_0}\mathrm{Im}(d_0)=\mathrm{Ker}(d_1)\mathrm{Ker}(S)0.`$ (6)
These chain of maps form an exact sequence, that is, the image of a map is the kernal of the map which follows. With the help of the resolvent, Schwarz was able to show that the partition function for the theory was given by
$$Z(\lambda )=e^{\frac{i\pi }{4}\iota }(\frac{\lambda }{\pi })^{\frac{\zeta }{2}}\stackrel{}{det}((d_1)^2)^{\frac{1}{4}}\stackrel{}{det}(d_0^{}d_0)^{\frac{1}{2}}det(\varphi _0^{}\varphi _0)^{\frac{1}{2}},$$
(7)
where $`\iota `$ is a non-topological geometry dependent function and $`\zeta `$ as shown in is given by
$$\zeta =\mathrm{dim}H^0(M)\mathrm{dim}H^1(M).$$
$`\iota `$ is part of a phase factor and thus the absolute value of the partition function $`Z(\lambda )`$ is a topological quantity. This will be our main concern .
For completeness we give a quick proof of this result, ignoring phase factors and constants.
Introducing a metric in the space of $`\omega `$’s allows us to write
$`Z(\lambda )`$ $`=`$ $`N{\displaystyle _{\mathrm{Ker}d_1(\mathrm{Ker}d_1)^{}}}𝑑\omega e^{iS(\omega )},`$
$`=`$ $`\text{Vol}(\mathrm{Ker}d_1)(detd_1^{}d_1)^{\frac{1}{4}}N.`$
We proceed to rewrite Vol($`\mathrm{Ker}d_1`$) using the exact sequence associated with $`\mathrm{Ker}d_1`$ and the manifold $`M`$. This procedure gives an expression for the partition function $`Z`$ containing information about the spaces $`\mathrm{Ker}d_1`$ and $`(\mathrm{Ker}d_1)^{}`$. Simply dropping Vol(Ker$`d_1`$) leads to a loss of information. We have Vol($`\mathrm{Ker}S`$)=Vol($`\mathrm{Ker}d_1`$) = Vol($`\mathrm{Im}d_0`$) by assumption (if $`H_1(M)`$ is non-trivial, this equation has to be modified ). Also,
$`d_0_{(\mathrm{Ker}d_0)^{}}:(\mathrm{Ker}d_0)^{}`$ $``$ $`(\mathrm{Im}d_0)`$
$`\text{Vol}(\mathrm{Im}d_0)`$ $`=`$ $`\stackrel{}{det}d_0\text{Vol}(\mathrm{Ker}d_0)^{}.`$
Note
$`\mathrm{\Omega }^0`$ $`=`$ $`\mathrm{Ker}d_0(\mathrm{Ker}d_0)^{}`$
and
$`\varphi _0:_0`$ $``$ $`H_0`$
$`\text{Vol}(H_0)`$ $`=`$ $`det\varphi _0\text{Vol}(_0).`$
where $`_0`$ represents the space of harmonic $`0`$-forms. Note harmonic $`p`$-forms are solutions of $`(d^{}d+dd^{})\varphi _p=0`$. In this space the Hodge star operator is present and hence a scalar product and volume can be defined. The map $`\varphi _0`$ introduced relates the space of harmonic $`0`$-forms to the space of de Rham cohomology $`H_0`$. By a theorem of Hodge this space of harmonic $`p`$-forms is isomorphic to the space of $`H_p`$. The space of de Rham cohomology does not have a metric and hence we define the volume in this space with the help of the map $`\varphi _0`$. Therefore,
$`\text{Vol}(\mathrm{Ker}d_0)^{}`$ $`=`$ $`\text{Vol}(\mathrm{\Omega }_0)[\text{Vol}(\mathrm{Ker}d_0)]^1,`$
$`=`$ $`\text{Vol}(\mathrm{\Omega }_0)(\text{Vol}(H_0))^1,`$
$`=`$ $`(\text{Vol}(\mathrm{\Omega }_0))(det\varphi _0)^1(\text{Vol}(_0))^1.`$
So that finally we get
$$\text{Vol}(\mathrm{Ker}S)=\stackrel{}{det}d_odet\varphi _0^1\text{Vol}(\mathrm{\Omega }_0)(\text{Vol}_\mathcal{0})^1.$$
Choosing
$$N\text{Vol}(\mathrm{\Omega }_0)(\text{Vol}_\mathcal{0})^1=1,$$
we get
$$Z=\text{Vol}(\mathrm{Ker}S)(detd_1^{}d_1)^{\frac{1}{4}}=\stackrel{}{det}d_0detd_1d^{}^{\frac{1}{4}}det\varphi _0^1.$$
A more careful calculation gives the determinant in (7). The quantities $`\iota `$ and $`\zeta `$ in (7) also need to be zeta-regularised - this was done in , where it was shown that the regularised $`\zeta `$ is given by
$$\zeta =\mathrm{dim}H^0(M)\mathrm{dim}H^1(M).$$
So in the present case, where $`H^0(M)`$ and $`H^1(M)=0`$, we have
$`\zeta =10=1.`$ (8)
Using the formulae $`d_1^{}=d_1`$ and $`=1`$ (modulo a possible sign), we get $`d_1=d_1^{}`$ and therefore $`(d_1)^2=d_1d_1=d_1^{}d_1`$, which gives
$`\stackrel{}{det}((d_1)^2)=\stackrel{}{det}(d_1^{}d_1).`$ (9)
Substituting (8) and (9) in (6) we get
$`Z(\lambda )=e^{\frac{i\pi }{4}\iota }({\displaystyle \frac{\lambda }{\pi }})^{\frac{\zeta }{2}}\stackrel{}{det}(d_1^{}d_1)^{\frac{1}{4}}\stackrel{}{det}(d_0^{}d_0)^{\frac{1}{2}}det(\varphi _0^{}\varphi _0)^{\frac{1}{2}}.`$ (10)
We now rewrite the product of determinants in (10) in terms of the Ray-Singer torsion of M. Since the Hodge star operator * is unitary with $`=1`$ and $`d_0^{}=d_2`$ (modulo a possible sign) we have
$`\stackrel{}{det}(d_0^{}d_0)`$ $`=`$ $`\stackrel{}{det}(d_2d_0)=\stackrel{}{det}((d_2d_0)),`$
$`=`$ $`\stackrel{}{det}(d_2d_0)=\stackrel{}{det}(d_2d_2^{}),`$
$`=`$ $`\stackrel{}{det}(d_2^{}d_2).`$
It follows that
$`\stackrel{}{det}(d_0^{}d_0)^{\frac{1}{2}}\stackrel{}{det}(d_1^{}d_1)^{\frac{1}{4}}=(\stackrel{}{det}(d_0^{}d_0)^{\frac{1}{2}}det(d_1^{}d_1)^{\frac{1}{2}}\stackrel{}{det}(d_2^{}d_2)^{\frac{1}{2}})^{\frac{1}{2}}.`$ (11)
It is possible to rewrite $`det(\varphi _0^{}\varphi _0)`$ using a standard result of manifold theory in a different form(see ). We start by noting that the integration map
$$_M:H^3(M)$$
is an isomorphism, i.e. for each $`r`$ there is a unique class $`[\alpha ]H^3(M)`$ such that $`\alpha =r`$. (Note that the integration map is well defined on $`H^3(M)`$ since $`_M\alpha +d\beta =_M\alpha `$ i.e. $`_M𝑑\beta =0`$ by Stokes theorem.) Also from the definition
$$H^3(M)=\mathrm{\Omega }^3(M)/\mathrm{Im}(d_2)=(\mathrm{Im}(d_2)\mathrm{Im}(d_2)^{})/\mathrm{Im}(d_2)=\mathrm{Im}(d_2)^{},$$
it follows that the map given by
$`\mathrm{Im}(d_2)^{}`$ (12)
is also an isomorphism. Now define the map
$$\varphi _3:\mathrm{Im}(d_2)^{}$$
to be the inverse of (12). Then using the properties of the Hodge star operator it can be shown that
$$det(\varphi _3^{}\varphi _3)=det(\varphi _0^{}\varphi _0)^1.$$
It follows that
$`det(\varphi _0^{}\varphi _0)^{\frac{1}{2}}=(det(\varphi _0^{}\varphi _0)^{\frac{1}{2}}det(\varphi _3^{}\varphi _3)^{\frac{1}{2}})^{\frac{1}{2}}.`$ (13)
Substituting (13) and (11) in (10) we get
$$Z(\lambda )=(\frac{\lambda }{\pi })^{\frac{1}{2}}\tau _{RS}(M)^{\frac{1}{2}},$$
(14)
where
$$\tau _{RS}(M)=det(\varphi _0^{}\varphi _0)^{\frac{1}{2}}det(\varphi _3^{}\varphi _3)^{\frac{1}{2}}\stackrel{}{det}(d_0^{}d_0)^{\frac{1}{2}}\stackrel{}{det}(d_1^{}d_1)^{\frac{1}{2}}\stackrel{}{det}(d_2^{}d_2).$$
(15)
This quantity $`\tau _{RS}(M)`$ is the Ray-Singer torsion of M . It is a topological invariant of M i.e. is independent of the metric of M. Thus the modulus $`Z(\lambda )`$ of the partition function, given by (13)-(14), is a topological invariant.
We are now ready to construct a discrete version of the preceding topological field theory which reproduces the continuum expression for the partition function where subdivision invariance is the discrete property corresponding to topological invariance. We will see that in order to do this it is crucial that there is an analogue of the Hodge star operator in the discrete theory. As we will see in the next section, this requires a field doubling. Therefore we consider a doubled version of the preceding theory, with the fields $`\omega _1`$ and $`\omega _2`$ in $`\mathrm{\Omega }^1(M)`$ and with the action functional (5) changed by
$$S(\omega )=\lambda <w,(d_1)w>\stackrel{~}{S}(\omega _1,\omega _2)=\lambda <({}_{\omega _2}{}^{\omega _1}),({}_{d_1\mathrm{\hspace{0.17em}\hspace{0.17em}0}}{}^{\mathrm{\hspace{0.17em}\hspace{0.17em}0}d_1})({}_{\omega _2}{}^{\omega _1})>.$$
(16)
The reason for this specific choice of action $`\stackrel{~}{S}(\omega _1,\omega _2)`$ for the doubled theory will become clear in the next section. An obvious generalisation of the proceeding, with $`T\stackrel{~}{T}=({}_{d_1\mathrm{\hspace{0.17em}\hspace{0.17em}0}}{}^{\mathrm{\hspace{0.17em}\hspace{0.17em}0}d_1})`$ shows that the partition function of the doubled theory
$$\stackrel{~}{Z}(\lambda )=_{\mathrm{\Omega }^1(M)\times \mathrm{\Omega }^1(M)}𝒟\omega _1𝒟\omega _2e^{\lambda \stackrel{~}{S}(\omega _1,\omega _2)}$$
can be evaluated to obtain the square of (15)
$$\stackrel{~}{Z}(\lambda )=Z(\lambda )^2=(\frac{\lambda }{\pi })^1\tau _{RS}(M).$$
(17)
Note that there is no phase factor here. This is because the quantity $`\iota =d_++d_{}`$ for the action $`\stackrel{~}{S}`$ in (16) vanishes since $`\stackrel{~}{T}=({}_{d_1\mathrm{\hspace{0.17em}\hspace{0.17em}0}}{}^{\mathrm{\hspace{0.17em}\hspace{0.17em}0}d_1})`$ has a symmetric spectrum.
## 4 The discrete version of the topological field theory
We proceed to construct a discrete version of Abelian Chern-Simons gauge theory. The Whitney map enables the Abelian Chern-Simons theory to be discretised by replacing the gauge field (1-form) $`A\mathrm{\Omega }^1(M)`$ by the discrete analogue, a 1-cochain $`xC^1(K)`$.
The most immediate way to do this is to construct the action $`S_K`$ of the discrete theory by
$$\lambda S_K(x)=\lambda S(W^K(x))=\lambda _M𝑑W^K(x)W^K(x).$$
This can be shown to coincide with the discrete action for the Abelian Chern-Simons theory introduced in . This prescription fails however, in the sense that the resulting partition function $`Z_K(\lambda )`$ is not a topological invariant i.e. is not independent of $`K`$, and does not reproduce the continuum expression for the partition function. We demonstrate this by considering the resolvent for $`S_K`$ obtained in an analogous way to the resolvent of the continuum action $`S`$ described in the previous section. Let $`T_K:C^1(K)C^1(K)`$ denote the self-adjoint operator on $`C^1(K)`$ determined by
$$S_K(x)=_MdW^K(x)W^K(x)=<T_Kx,x>.$$
Then
$$\mathrm{Ker}(T_K)\mathrm{Ker}(d_1^K).$$
Since for $`x\mathrm{Ker}(d_1^K)`$ we have
$`<T_Kx,x>`$ $`=`$ $`{\displaystyle _M}𝑑W^K(x)W^K(x),`$ (18)
$`=`$ $`{\displaystyle _M}W^K(d_1^Kx)W^K(x).`$ (19)
Thus the discrete analogue of the resolvent (6) is a resolvent for $`S_K`$:
$$0^{\varphi _0}\mathrm{\Omega }^0(M)^{d_0^K}\mathrm{Ker}(d_1^K)\mathrm{Ker}(T_K)=\mathrm{Ker}(S_K)0.$$
The resulting partition function is the discrete analogue of the partition function $`Z(\lambda )`$:
$$Z_K(\lambda )=\stackrel{}{det}((\varphi _0^K)^{}\varphi _0^K)^{\frac{1}{2}}\stackrel{}{det}((d_0^K)^{}d_0^K)^{\frac{1}{2}}\stackrel{}{det}(\frac{i\lambda }{\pi }T_K)^{\frac{1}{2}}.$$
In the following formula for $`T_K`$ was obtained:
$$T_K[v_0,v_1]=\frac{1}{6}[v_2,v_3]$$
where the sum is over all $`1`$-simplices $`[v_2,v_3]`$ such that $`[v_0,v_1,v_2,v_3]`$ is a $`3`$-simplex with orientation compatible with the orientation of $`M`$.
It is possible to show that $`det((\varphi _0^K)^{}\varphi _0^K)=N_0^K=\mathrm{dim}C^{(K)}=`$ the number of vertices of $`K`$. Then the failure of the discretisation prescription can be demonstrated by showing that the quantity
$$Z(\lambda )^2=\frac{1}{N_0^K}\stackrel{}{det}(_1^Kd_0^K)\stackrel{}{det}(T_K^2)^{\frac{1}{2}}$$
is not independent of $`K`$.
A discrete version of the doubled topological field theory with action (16) has been constructed in in such a way that the expression (17) for the continuum partition function is reproduced. We briefly describe this in the following.
The discretisation prescription is
$`(\omega _1,\omega _2)\mathrm{\Omega }^1(M)\times \mathrm{\Omega }^1(M)(x,y)C^1(K)\times C^1(\widehat{K}),`$ (20)
$`S(\omega )=\lambda <({}_{\omega _2}{}^{\omega _1}),({}_{d_1\mathrm{\hspace{0.17em}\hspace{0.17em}0}}{}^{\mathrm{\hspace{0.17em}\hspace{0.17em}0}d_1})({}_{\omega _2}{}^{\omega _1})>\stackrel{~}{S}_K(x,y)=\lambda <({}_{y}{}^{x}),({}_{^{\widehat{K}}d^{\widehat{K}}\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}}{}^{\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}^Kd^K})({}_{y}{}^{x})>,`$ (21)
where $`K`$ is the simplicial complex triangulating M, $`\widehat{K}`$ is its dual, $`C^q(K)`$, $`C^p(\widehat{K})`$, $`d^K`$ and $`d^{\widehat{K}}`$ are as described in the previous section. The analogue of the Hodge star operator is the duality operator $`^K`$. This is a map $`^K:C^q(K)C^{zq}(\widehat{K})`$ ( and $`^{\widehat{K}}:C^p(\widehat{K})C^{zp}(K)`$ which explains the need for field doubling and the expression (21) for the discrete action $`\stackrel{~}{S}_K(x,y)`$. There is a natural choice of resolvent for $`\stackrel{~}{S}_K(x,y)`$, analogous to the resolvent (6) in the continuum case. It is
The partition function is
$$\stackrel{~}{Z}_K(\lambda )=_{C^1(K)\times C^1(\widehat{K})}𝒟x𝒟ye^{\stackrel{~}{S}_K(x.y)}.$$
(22)
Evaluating this by Schwarz’s method with the resolvent above leads to:
$`\stackrel{~}{Z}_K(\lambda )=`$ $`({\displaystyle \frac{\lambda }{\pi }})^{1+N_0^KN_1^K}det((\varphi _0^K)^{}\varphi _0^K)^{\frac{1}{2}}\stackrel{}{det}((d_0^K)^{}d_0^K)^{\frac{1}{2}}\stackrel{}{det}((d_1^K)^{}d_1^K)^{\frac{1}{4}},`$
$`det((\varphi _0^{\widehat{K}})^{}\varphi _0^{\widehat{K}})^{\frac{1}{2}}\stackrel{}{det}((d_0^{\widehat{K}})^{}d_0^{\widehat{K}})^{\frac{1}{2}}\stackrel{}{det}((d_1^{\widehat{K}})^{}d_1^{\widehat{K}})^{\frac{1}{4}}.`$
There is no phase factor in (22) since $`\zeta `$ vanishes just like in (17). We have also used the fact that $`\zeta =1N_0^K+N_1^K`$ which is shown in .
Now rewrite the determinant involving $`\widehat{K}`$ \- objects in terms of determinants of $`K`$-objects. Modulo a possible sign $`\pm `$ we have the formulae
$`(^{\widehat{K}})^1`$ $`=(^{\widehat{K}})^{}=^K,`$ (23)
$`(^K)^1`$ $`=(^K)^{}=^{\widehat{K}},`$ (24)
$`(d_q^K)^{}`$ $`=^{\widehat{K}}d_{nq1}^{\widehat{K}}^K,`$ (25)
$`(d_p^{\widehat{K}})^{}`$ $`=^Kd_{np1}^K^{\widehat{K}}.`$ (26)
( The $`\pm `$ signs are omitted because they will all cancel out in the following calculation.) Now:
$`\stackrel{}{det}((d_0^{\widehat{K}})^{}d_0^{\widehat{K}})`$ $`=\stackrel{}{det}(^Kd_2^K^{\widehat{K}}d_0^{\widehat{K}}),`$ (27)
$`=\stackrel{}{det}(^{\widehat{K}}(^Kd_2^K^{\widehat{K}}d_0^{\widehat{K}})^K)`$ (28)
$`=\stackrel{}{det}(d_2^K^{\widehat{K}}^K),`$ (29)
$`=\stackrel{}{det}(d_2^K(d_2^K)^{})`$ (30)
$`=det((d_2^K)^{}d_2^K),`$ (31)
and
$`\stackrel{}{det}((d_1^{\widehat{K}})^{}d_1^{\widehat{K}})`$ $`=\stackrel{}{det}(^Kd_1^K^{\widehat{K}}d_1^{\widehat{K}}),`$ (32)
$`=\stackrel{}{det}(^{\widehat{K}}(^Kd_1^K^{\widehat{K}}d_1^{\widehat{K}})^K)`$ (33)
$`=\stackrel{}{det}(d_1^K^{\widehat{K}}^K),`$ (34)
$`=\stackrel{}{det}(d_1^K(d_1^K)^{})`$ (35)
$`=det((d_1^K)^{}d_1^K).`$ (36)
The integration map (12) has a discrete analogue
$`\mathrm{Ker}(d_2^K)^{}`$ $``$ $`,`$
$`a`$ $``$ $`<a,[M]>,`$ (37)
where $`[M]C_3(K)`$, the orientation cycle of $`M`$, i.e. the sum of all 3-simplices of $`K`$, oriented so that their orientations are compatible with the orientation of $`M`$. ( Note that $`a\mathrm{Ker}(d_2^K)^{}C^3(K)`$ can be evaluated on any element $`\sigma C_3(K)`$ to get a real number $`<a,\sigma >`$.) Define the map
$$\varphi _3^K:\mathrm{Ker}(d_2^K)^{}$$
(38)
to be the inverse of (37). Then using the properties of $`^K`$ and $`^{\widehat{K}}`$, it can be shown that
$$det((\varphi _3^K)^{}\varphi _3^K)=det((\varphi _0^{\widehat{K}})^{}\varphi _0^{\widehat{K}})^1.$$
(39)
Now using (39), (31) and (36) we can rewrite (22) as
$$\stackrel{~}{Z}_K(\lambda )=(\frac{\lambda }{\pi })^{1+N_0^KN_1^K}\tau _K(M),$$
(40)
where
$$\tau _K(M)=det((\varphi _0^K)^{}\varphi _0^K)^{\frac{1}{2}}det((\varphi _3^K)^{}\varphi _3^K)^{\frac{1}{2}}\underset{q=0}{\overset{2}{}}\stackrel{}{det}((d_q^K)^{}d_q^K)^{\frac{1}{2}(1)^q}.$$
(41)
This quantity $`\tau _K(M)`$ is the R-torsion of the triangulation $`K`$ of $`M`$. It is a combinatorial invariant of $`M`$ i.e. is independent of the choice of triangulation $`K`$ .
This is the untwisted torsion of $`M`$, more generally the torsion can be “twisted” by a representation of $`\pi _1(M)`$. The factors involving the determinants $`det^{}((d_q^K)^{}d_q^K)`$ constitute the usual Reidemeister torsion of $`M`$ . When these are put together with the factors involving $`det((\varphi _i^K)\varphi _i^K)`$, $`i=0,3`$, as in (41) we get the R-torsion “as a function of the cohomology” introduced and shown to be triangulation-independent in .
The expression (41) for $`\tau _K(M)`$ is analogous to the expression (15) for the R-torsion $`\tau _{RS}(M)`$, and in fact it has been shown that these torsions are equal
$$\tau _K(M)=\tau _{RS}(M).$$
It follows that the partition function (40) of the discrete theory coincides with the partition function (15) of the continuum theory, except for the $`K`$-dependent quantities $`N_0^K`$ and $`N_1^K`$ appearing in (40). These quantities can be removed by a suitable $`K`$-dependent renormalisation of the coupling parameter $`\lambda `$.
It is possible to show that
$`det((\varphi _0^K)^{}\varphi _0^K)`$ $`=N_0^K`$
$`det((\varphi _3^K)^{}\varphi _3^K)`$ $`={\displaystyle \frac{1}{N_3^K}}.`$
We will use this result in our numerical work.
## 5 Numerical Results
We are now in a position to proceed to numerically evaluate the discrete expressions for the torsion obtained. This allows us to check the underlying theoretical ideas by numerically verifying that the discrete expressions agree with expected analytic results. It also allows us to check that the results obtained are subdivision invariant. The subdivision invariance of torsion is demonstrated by showing that if any simplex of the triangulation is subdivided, the value of the torsion does not change. This is what is meant by topological invariance in the discrete setting. The expected analytic result for torsion for a lens space $`L(p,1)`$ is (See )
$$T(L(p,1))=\frac{1}{p}.$$
Thus $`T(S^3)=T(L(1,1))=1`$.
We also show, numerically, that the discrete expression for the Chern-Simons partition function obtained without using the * operator is not a topological invariant. This shows very clearly the importance of the doubling construction method used in the discretisation method, for capturing topological information.
In order to proceed, we need to efficiently triangulate the spaces $`S^3`$ and L(p,q). First we triangulate $`S^3`$. We do this by considering a four dimensional simplex $`[v_0,v_1,v_2,v_3,v_4]`$ and observing that the boundary of this object is precisely the triangulation $`K`$ of $`S^3`$ which we require. Next we turn to spaces L(p,q) which we need to triangulate in order to proceed. An efficient triangulation of this space has been constructed by Brehm and Swiatkowski . We use this procedure for our computations .
We can now summarise our numerical results. The R-torsion for a simplicial complex $`K`$, with dual cell complex $`\widehat{K}`$, for either $`S^3`$ or L(p,1) involves evaluating
$$T=\sqrt{\frac{1}{N_0N_3}det(_1d_0)det(_2d_1)^1det(_3d_2)},$$
where $`N_i=`$nos. of $`i`$-simplices in $`K`$, $`_1,_2,_3`$ are boundary operators on $`K`$ and $`d_0,d_1,d_2`$ are coboundary operators on $`K`$. Note that, in the discrete setting, these operators can be expressed as matrices. We do this by using the vertices, edges, faces and tetrahedra of our complex as a basis list. Any $`p`$-simplex in the complex can then be expressed in terms of this. Since we know what $`d`$ maps the various basis list elements to, we can set the coefficients of its matrix representation. When we say $`det(_1d_0)`$, for example, we simply mean the determinant of the matrix which results from multiplying the matrices corresponding to the operators $`_1`$ and $`d_0`$.
If our complex consisted of just one triangle $`[0,1,2]`$, the basis list would be:
$$\begin{array}{cc}0\hfill & [0],\\ 1\hfill & [1],\\ 2\hfill & [2],\\ 3\hfill & [0,1],\\ 4\hfill & [0,2],\\ 5\hfill & [1,2],\\ 6\hfill & [0,1,2].\end{array}$$
We know that $`d[0]=[1,0]+[2,0]`$. This can be expressed in terms of our basis list as $`d`$ acting on basis element $`0`$ going to $`34`$. So the matrix $`d`$ for this complex is:
$$\left(\begin{array}{ccccccc}0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0\\ 1& 0& 1& 0& 0& 0& 0\\ 0& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 1& 1& 1& 0\end{array}\right).$$
Note that the last column is zero since $`d`$ has nothing to map a $`2`$-simplex to and that the first three rows are zero since nothing is mapped to $`0`$-simplices. If we act on $`0`$ with this matrix we get $`34`$, as expected, and $`d^2=0`$.
If we do not use $`\widehat{K}`$ then from section 4,
$$\widehat{T}=Z(1)^2=\frac{1}{N_0^K}(\stackrel{}{det}_1^Kd_0^K)(detT_K^2)^{\frac{1}{2}}.$$
We evaluated the quantity $`\widehat{T}`$ numerically for various triangulations $`K`$ of $`S^3`$ and found the following results shown in Table 2 for the change of $`\widehat{T}`$ under subdivision where $`mn`$ corresponds to a triangulation of $`S^3`$ with $`n`$ vertices and $`m`$ tetrahedra and where $`X_i=det_id_{i1}`$.
It is clear that $`\widehat{T}`$ for $`S^3`$ is not subdivision invariant. We next evaluate $`T`$ for $`S^3`$ and for L(p,1) and check that it is indeed subdivision invariant and agrees with the analytic calculations for L(p,1), with p=2,3,4 and 5. These results are shown in Tables 4-6.
As a check on the numerical method we also count the number of zero modes of the Laplacian operator on the different p-chain spaces. These numbers gives the dimension of the Homology groups and are shown in Table 3.
### 5.1 $`S^3`$
As a further check, a systematic way of carrying out subdivision known as the Alexander moves was used to study the subdivision invariance of the torsion. There are four such moves in 3 dimensions. They are best explained by example. We have
1. $`[0,1,2,3,4]>[x,1,2,3,4]+[0,x,2,3,4]`$.
2. $`[0,1,2,3,4]>[x,1,2,3,4]+[0,x,2,3,4]+[0,1,x,3,4]`$.
3. $`[0,1,2,3,4]>[x,1,2,3,4]+[0,x,2,3,4]+[0,1,x,3,4]+[0,1,2,x,4]`$.
4. $`[0,1,2,3,4]>[x,1,2,3,4]+[0,x,2,3,4]+[0,1,x,3,4]+[0,1,2,x,4]+[0,1,2,3,x]`$.
These are all natural operations in that the first (nth) move corresponds to adding a vertex splitting the 1-simplex (n-simplex) $`[0,1]`$ ($`[0,1,\mathrm{},n]`$) and connecting it to all the vertices resulting in two (n+1) tetrahedra.
The torsion T, in terms of its component determinants, and the way they change under the Alexander moves is exhibited in the table below,
where $`X_i=det_id_{i1}`$ and $`a2s3`$ means the complex which resulted after the type 2 Alexander moves were performed on the $`S^3`$. It’s clear from this that
$$T=\frac{1}{N_3N_0}det(_1d_0)(det_2d_1)^1det(_3d_2)$$
is subdivision invariant and thus a topological invariant of the manifold.
As a final check we tried several other subdivisions. We took a given triangulation and barycentric subdivided one or more ($`n`$) of its faces to get $`1s3,2s3,\mathrm{}`$ ($`ns3`$). The results are in Tables 4 and 5.
### 5.2 Lens spaces
We conclude with the results for the lens spaces. The results are shown in Table 6 . As can be seen, these agree extremely well with the known analytic result T(L(p,1))=$`1/p`$.
## 6 Conclusions
The method of discretisation introduced works extremely well. The main point of the method is to construct discrete analogues for the set $`(\mathrm{\Omega }^p,d,,)`$. Previous work in the direction has neglected the Hodge star operator, * . We have thus demonstrated that the Hodge star operator plays a vital role in the construction of topological invariant objects from field theory. We were able to construct an expression for the partition function which is correct even as far as overall normalisation is concerned. Mathematically the equivalence between the Ray-Singer torsion and the combinatorial torsion of Reidemeister was proved independently in 1976 by Cheeger and Müller . It is nice to see the result emerge in a direct manner by a formal process of discretisation. On the way we also had to double the original system so that K, the triangulation , and $`\widehat{K}`$, its dual, are both present. If this doubling and the reason for it are overlooked, then the topological information present in the discretisation is lost, as our numerical results demonstrated.
It is clear that the geometry motivated discretisation method introduced is very general and that it can be used to analyse a wide variety of physical systems. In the approach outlined we have captured topological features. In applications capturing geometrical features of a problem are also very important. We are currently investigating this aspect of our approach.
A limitation of the method is that there is no simple generalisation to deal with non-abelian theories. The discrete analogies of $`d`$, $`d^{}`$ were linear : there is no natural discrete analogue of $`d_A:=d+A`$, with $`A`$ a Lie algebra valued one form.
## Acknowledgments
The work of S.S. is part of a project supported by Enterprise Ireland. D.A. would like to acknowledge support from Enterprise Ireland and Hitachi. Samik Sen and S.S. would like to thank D. Birmingham for explaining the triangulation method of Ref. . |
warning/0001/astro-ph0001307.html | ar5iv | text | # Interpretation of the OGLE Q2237+0305 microlensing light-curve (1997-1999)
## 1 Introduction
The QSO 2237+0305, sometimes known as Huchra’s lens or the Einstein Cross (Huchra et al. 1985), is perhaps the most remarkable gravitational lens yet discovered. It comprises a foreground barred Sb galaxy (z=0.0394) whose nucleus is surrounded by four images of a radio-faint QSO (z=1.695). Ground based spectroscopic observations have verified that all four images are similar QSO’s at the same redshift (Adam et al. 1989). Broad band monitoring has shown that significant microlensing events occur (eg. Irwin et al. 1989; Corrigan et al. 1991). Since the optical depth to microlensing is of order unity at each of the image positions (eg. Kent & Falco 1988; Schneider et al. 1988; Schmidt, Webster & Lewis 1998), the magnification effects on the source can be considered as a network of caustics moving across the source plane. Strong variation in a particular image results from the source either crossing a caustic or passing close to a cusp. Q2237+0305 provides a unique opportunity to study microlensing events for two good reasons. Firstly, the relative closeness of the lensing galaxy and the nearly on-axis alignment means that the time delay between the four quasar images is less than a day. Thus it is easy to separate microlensing events from intrinsic variations of the QSO. Also, for this QSO the time-scale for microlensing events is typically 30-50 days, reduced by a factor of 10 from typical time-scales of perhaps a year.
Although Q2237+0305 has been monitored since its discovery, data collected by the OGLE collaboration over the last 3 years (Wozniak et al. (2000a,b); see also http://www.astro.princeton.edu/$``$ogle/ogle2/huchra.html) has for the first time obtained data of sufficient coverage to clearly demonstrate smooth independent flux variation between the images. In this paper we provide interpretations of several features in the image A and C light-curves, and discuss their implications for future microlensing. Secs. 2 and 3 describe the published light-curves and the models used to interpret them. In Sec. 4 we discuss the HME classes associated with different light-curve features by comparing their parameters (such as height and maximum light-curve derivative) to model distributions. We then describe the implications for future HMEs and present qualitative scenarios in Secs. 5 and 6.
## 2 Existing monitoring data for Q2237+0305
Fig. 1 shows the OGLE light-curve for Q2237+0305, containing all points taken before the end of the 1999 observing season (figure taken from the OGLE web page (see http://www.astro.princeton.edu/$``$ogle/ogle2/huchra.html). Note that this data was published in Wozniak et al. (2000a,b). The light-curves have an unprecedented sampling rate ($`1`$ week). The OGLE data adds to the approximately 10 years of previously obtained, but less densely sampled photometry presented in Schneider et al. (1988), Kent & Falco (1989), Irwin et al. (1989), Corrigan et al. (1990) and $`Ø`$stensen et al. (1996). Fig. 2 shows the complete compiled data set. Error bars are shown representing the published errors. Note that the data taken prior to OGLE is in R and r bands whereas the OGLE monitoring is in V band. The typical amplitudes of microlensing flux variation may therefore differ between data sets if the source size is colour dependent (Wambsganss & Paczynski 1991).
## 3 The Microlensing model
Our analysis involves the comparison of ensembles of microlensing models to Q2237+0305 light-curve features. Through this comparison we interpret observed HMEs and attempt to predict the features of future HMEs. This section provides a brief discussion of the models used.
Throughout the paper we use standard notation for gravitational lensing. The Einstein radius of a microlens projected into the source plane is denoted by $`\eta _0`$. The normalised shear is denoted by $`\gamma `$, and the convergence or optical depth by $`\kappa `$. The basic model for microlensing at high optical depth comprises a disc of point masses (eg. Kayser, Refsdal & Stabell 1986) with a size such that a large fraction ($`>99\%`$) of macroimage flux is recovered (Katz, Balbus & Paczynski 1986; Lewis & Irwin 1995). To construct a microlensed light-curve we use the contouring technique of Lewis et al. (1993) and Witt (1993). For the microlensing models of Q2237+0305 used in the current work we assume the macro-parameters for the lensing galaxy calculated by Schmidt, Webster & Lewis (1998). We will use the standard notation introduced by Yee (1988), to describe these images. Where required a cosmology having $`\mathrm{\Omega }=1`$ with $`H_o=75kmsec^1`$ is assumed.
We describe the microlensing rate in terms of the effective transverse velocity, which is defined as the transverse velocity that produces a microlensing rate from a static model equal to that of the observed light-curve (Wyithe, Webster & Turner 1999 (hereafter WWT99)). The effective transverse velocity therefore describes the microlensing rate due to the combination of the effects of a galactic transverse velocity and proper motion of microlenses. To calculate distributions of values for light-curve features associated with HMEs, we assume that the effective transverse velocity accurately describes not only the distribution of light-curve derivatives during an HME (Wyithe, Webster & Turner 2000a), but also the distribution of orientations between the caustic and source trajectory, and hence the event duration.
We have previously obtained the following normalised probability distributions. These were obtained under the assumption that the source size $`S\eta _o`$. Evidence in favour of this assumption was presented in Wyithe, Webster & Turner (2000c).
$`i)`$ $`p_s(S|m,v_{eff})`$, the probability that the continuum source diameter is between $`S`$ and $`S+\mathrm{d}S`$ given a mean microlens mass $`m`$ and an effective galactic transverse velocity $`v_{eff}`$ (Wyithe, Webster, Turner & Mortlock 2000).
$`ii)`$ $`p_v(v_{eff}|m)`$ the probability that the effective galactic transverse velocity is between $`v_{eff}`$ and $`v_{eff}+\mathrm{d}v_{eff}`$ given a mean microlens mass $`m`$ (WWT99).
$`iii)`$ $`p_m(m)`$, the probability that the mean microlens mass is between $`m`$ and $`m+\mathrm{d}m`$ (Wyithe, Webster & Turner 2000b).
$`p_v(v_{eff}|m)`$ and $`p_m(m)`$ were computed using flat ($`p(V_{tran})dV_{tran}`$), and logarithmic ($`p(V_{tran})\frac{dV_{tran}}{V_{tran}}`$) assumptions for the Bayesian prior for galactic transverse velocity ($`V_{tran}`$). $`p_v(v_{eff}|m)`$ was found to be insensitive to the prior assumed, however $`p_m(m)`$ showed some dependence. In the remainder of this paper we use $`p_m(m)`$ calculated using the assumption of a logarithmic prior. We note that the assumption of the flat prior raises the average light-curve derivative by a few percent.
The functions $`p_s(S|m,v_{eff})`$, $`p_v(v_{eff}|m)`$, $`p_m(m)`$ and the HME statistics presented in this paper were computed for different assumptions of smooth matter density, photometric error, and direction of the galactic transverse velocity. Since the probability functions referred to above were computed from a derivative analysis, the statistics that we compute in this paper are quantitatively similar for the different possible models. Therefore we present only results from models with no smooth matter, a transverse velocity direction along the image C-D axis and simulated photometric errors assigned according to a Gaussian distribution with half widths of $`\sigma =\mathrm{\Delta }M/2`$ in images A and B, and $`\sigma =\mathrm{\Delta }M`$ in images C and D.
Both the microlensing rate due to a transverse velocity (eg. Witt, Kaiser & Refsdal 1993), as well as the corresponding rate due to proper motions (WWT00a) are not functions of the details of the microlens mass distribution, but rather are only dependent on the mean microlens mass. We therefore limit our attention to models in which all the microlenses have the same mass since the results obtained will be applicable to other models with different forms for the mass function.
The determination of probability for the quantity $`v_{eff}\sqrt{m}`$ from the Q2237+0305 monitoring data is quite robust. However the probability for the source size is derived from a single poorly sampled HME. The small number of observations describing the 1988 peak (Irwin et al. 1989; Corrigan et al 1991) introduces the potential for a systematic error in the source size equal to the ratio of the true event length and the inferred event length of $`52`$ days (twice the estimated rise time). This can be compared to the $``$100 day separation of the two points that provide an upper-bound on the event duration. The resulting systematic error in the estimate of source size is therefore smaller than a factor of $`2`$. In addition there may also be a component of systematic error from the assumption that the 1988 peak was due to a single caustic crossing. The statistics presented in the following sections are therefore computed assuming prior probabilities for $`S`$ assuming no systematic error,
$$p_s\left(S|m,v_{eff}\right)dS,$$
(1)
and systematic errors of $`\times 2`$ and $`\times 5`$ in $`S`$:
$$p_s^{}\left(S^{}|m,v_{eff}\right)dS^{}=\frac{1}{2}p_s\left(2S|m,v_{eff}\right)dS$$
(2)
and
$$p_s^{}\left(S^{}|m,v_{eff}\right)dS^{}=\frac{1}{5}p_s\left(5S|m,v_{eff}\right)dS.$$
(3)
Our source size estimate was made from data collected in the R and r bands, while the OGLE light-curves showing the features that we wish to investigate are in the V band. This introduces the possibility for another source of systematic error if the source has significantly different sizes in the R/r and V bands.
To investigate individual HMEs we must look at light-curve statistics for single images. Therefore unlike the calculation of $`p_s`$, $`p_v`$ and $`p_m`$, which used difference light-curves, intrinsic source variation may be important. This cannot be directly measured, however in Wyithe, Webster, Turner & Agol (2000) (hereafter WWTA00) limits are placed on the intrinsic variability power-spectrum and it is shown that intrinsic variability should not be an important consideration during HMEs.
## 4 Analysis of the OGLE Light-Curve
In this section we apply simple statistics describing event heights and light-curve derivatives a postiori to specific features in the published light-curves. We also retrospectively apply the triggering function developed in WWTA00 to two observed light-curve peaks. We use the results to interpret features found in monitoring data. In particular we discuss whether the large scale variation in image C (1999) was due to the source having crossed a caustic or moved outside of a cusp.
We note that choosing light-curve features a-postiori, and comparing them to a sample of models in order to draw conclusions regarding the type of event involved has the potential to introduce a statistical bias. However in the present case we feel that our method is justified because of the well established prior knowledge that caustic crossing HMEs produce larger amplitude, shorter period and more asymmetric (in time) variation than cusp related events.
### 4.1 Statistics of the 1999 light-curve Peak in image C
The image C light-curve shows a remarkable resolved peak described by $`35`$ points on separate days spanning $`7`$ months (see Fig. 1). By December 1999 the image C light-curve had dropped to a level similar to the end of the 1998 observing season. The peak is quite symmetric having reached a height $`0.5`$ magnitudes above the 1998 level. The maximum derivative reached both before and after the event peak was $`2`$ mags/year. Images B and D remained fairly constant during this period suggesting constant intrinsic luminosity. The event properties are discussed in relation to model calculations of the probability for their values. The results are summarised in Tab. 1.
#### 4.1.1 Maximum light-curve derivative
We have calculated the cumulative probability of observing a maximum light-curve derivative $`\dot{M}`$ prior to the peak maximum:
$`P_{\dot{M}}(\dot{M}<\dot{M}_o)=`$
$`{\displaystyle }dm{\displaystyle }dS{\displaystyle }dv_{eff}(p_s(S|m,v_{eff})p_m\left(m\right)`$
$`\times p_v(v_{eff}|m)P_{\dot{M}}(\dot{M}<\dot{M}_o|S,m,v_{eff})),`$ (4)
using a sampling rate typical of the monitoring data (1 point per week). We have looked at the probability $`P`$ for cusps ($`P_{\dot{M}_C}(\dot{M}<\dot{M}_o)`$), $`ve`$ caustic crossings ($`P_{\dot{M}_{}}(\dot{M}<\dot{M}_o)`$) and $`+ve`$ caustic crossings ($`P_{\dot{M}_+}(\dot{M}<\dot{M}_o)`$). Fig. 3 shows these functions; solid, dotted and dashed lines correspond to $`P_{\dot{M}_C}(\dot{M}<\dot{M}_o)`$, $`P_{\dot{M}_{}}(\dot{M}<\dot{M}_o)`$, and $`P_{\dot{M}_+}(\dot{M}<\dot{M}_o)`$ respectively. The functions shown assume that the source size estimate $`S`$ is correct. We find that a maximum derivative of $`2`$ magnitudes per year is inconsistent at $`>99\%`$ level with the event having been a $`+ve`$ caustic crossing. In addition, we find that a $`ve`$ caustic crossing is excluded at the 95% level. A cusp cannot be excluded on the basis of the maximum derivative (though it is higher than expected). Our conclusions are barely changed (the $`ve`$ event confidence is then 90%) if we assume our source size has been underestimated by a factor of 2. If our underestimate is a factor of 5 then only the $`+ve`$ event can be excluded (at the 90% level). The results are summarised in column 5 of Tab. 1.
#### 4.1.2 Height of the peak above the previous local minimum
If the light-curve minimum preceding the 1999 image C event is assumed to have been at the level of the 1998 season then the height of the peak maximum above the previous local minimum is $`M_{peak}0.5`$ magnitudes. We calculate the cumulative probability for the difference between peak maxima and the preceding local minima ($`M_{height}`$). $`M_{height}`$ is a function of sampling rate for sharp peaks and sparse samplings. We have therefore used a sampling rate of 1 point per week and integrated the probability over $`v_{eff}`$ as well as $`S`$ and $`m`$:
$`P_M(M_{height}<M_{peak})=`$
$`{\displaystyle }dm{\displaystyle }dS{\displaystyle }dv_{eff}(p_s(S|m,v_{eff})p_m\left(m\right)`$
$`\times p_v(v_{eff}|m)P_M(M_{height}<M_{peak}|S,m,v_{eff})).`$ (5)
As before we have computed the function for cusps ($`P_{M_C}(M_{height}<M_{peak})`$), $`+ve`$ caustic crossings ($`P_{M_+}(M_{height}<M_{peak})`$), and $`ve`$ caustic crossings ($`P_M_{}(M_{height}<M_{peak})`$). Fig. 4 displays the resulting curves (assuming no systematic uncertainty in the source size estimate). $`\mathrm{\Delta }M_{peak}0.5`$ magnitudes is typical if the event is due to a cusp, and is ruled out at the 98% level if the event is $`+ve`$ caustic crossing. If the event is a $`ve`$ caustic crossing the results are inconclusive. A very similar conclusion is reached if the source size has been underestimated by a factor of 2. The results are summarised in column 6 of Tab. 1.
#### 4.1.3 A postiori application of the triggering function to event 1999C
WWTA00 describes a general function to determine how long one should wait ($`P`$) for a HME following a hypothetical observed light-curve derivative ($`T`$). If a light-curve derivative $`T_{obs}\pm \mathrm{\Delta }T_{obs}`$ is observed, predictions $`F(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$ can be made about forthcoming events specifically for the current data using the sampling rate identical to that of the observations. On the 19th of June 1999 monitoring from the OGLE collaboration (OGLE web page) showed image C rising at a rate of 1.21-1.78 mags/year. Since at this level $`P`$ is not sensitive to the sampling rate for sampling spacings smaller than two weeks, we used observations on the 10th June, 19th June and 1st July.
For the present calculation the triggering function discussed in WWTA00 is slightly modified. Rather than searching for peaks following a derivative larger than some hypothetical value, we search for peaks following light-curve derivatives in the measured range. Thus since the algorithm steps along the curve in steps equal to the sampling rate, certain events (in particular $`+ve`$ caustic crossings) are missed if the derivative jumps from below the lower bound to above the upper bound during one sampling interval. Therefore the results describe the relative likely-hood of observing the different types of event following a current light-curve derivative. As a consequence of the triggering algorithm searching for derivatives along the light-curve in one direction, then the model light-curve derivative at the point of the trigger is systematically biased toward $`T\mathrm{\Delta }T`$. $`P`$ is therefore an overestimate. This bias is minimised by reduced photometric error.
Fig. 5 shows the resulting triggering functions $`F_+`$, $`F_{}`$ and $`F_C`$ computed assuming no systematic uncertainty in the source size estimate. The solid, dotted and dashed lines correspond to $`F_C(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$, $`F_{}(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$, $`F_+(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$ respectively. We find that the observed trigger precedes a caustic crossing HME $`75\%`$ of the time. The event peak is most likely to occur $`3`$ months following such a trigger if the event was a cusp, 2 months later if it was a $`ve`$ caustic and 1 month later if it was a $`+ve`$ caustic crossing. Note that having observed a derivative in the quoted range means that a $`+ve`$ caustic crossing is very unlikely due to the rapid change in the light-curve derivative. The light-curve peaked at $`1`$ July, about 2-weeks after the observed derivative. This is surprisingly early for all types of events (though $`P`$ is an overestimate). At the $`90\%95\%`$ level, caustic crossing events are excluded. This result holds if a systematic source size error is assumed. The probabilities of both the class of event following the trigger as well as the arrival time are summarised in columns 3 and 4 of Tab. 1.
### 4.2 What class of event have we seen in image C ?
Constraints on the event type of the OGLE image C HME are placed by both the maximum derivative observed prior to the event peak, and the height of the peak above the previous minima. $`+ve`$ caustic crossings are excluded by both these measurements (even when potential systematic errors in source size are assumed). In addition, assuming no systematic uncertainty in source-size, $`ve`$ caustic crossings are excluded by the measured maximum derivative. The cusp interpretation is consistent with both measurements.
Because triggers in the 1999 image C light-curves trigger were relatively small, the fraction of events predicted for different classes of HME by the triggering function does not restrict the type of event observed. In particular, the trigger was not large enough to rule out either cusps or $`ve`$ caustic crossings. However the OGLE data shows a previous rise in the light-curve of image C occurring between $`300`$ and $`500`$ days prior to July 1999. This rise occurred between observing seasons, but the net change is $`0.8`$ magnitudes suggesting that a $`+ve`$ caustic crossing event may have been missed (see Sec. 4.4). The typical separation of double peaked events calculated by Witt, Kayser and Refsdal (1993) is 300-500 days. This fact coupled with the comparative rarity of a $`+ve`$ event following the observed trigger leads to the conclusion that if the 1999 event were a caustic crossing then it is more likely to have been the second peak in a double horned event ($`ve`$ caustic crossing) than a $`+ve`$ caustic crossing. However, based on the peak arrival time, a cusp event is more likely than a $`ve`$ caustic crossing (or a $`+ve`$ caustic crossing).
Based on this evidence we conclude that the event observed for image C was probably a cusp event rather than a caustic crossing. However, further monitoring may show that this interpretation is false, and therefore that the microlensing models are in error. Thus the event will provide a valuable test of whether our current microlensing models can be used to reliably interpret the initial stages of HME’s.
#### 4.2.1 Asymmetry of the 1999 image C event
At the time of writing the image C light-curve was still in decline, but appeared to be decelerating ($`+ve`$ second derivative). Observations of the final stages of this event will provide a further property that can be used to distinguish between the different possibilities. After the light-curve has flattened out, the difference between the light-curve minima immediately preceding and following the event can be measured:
$$\mathrm{\Delta }M_{min}=M_{min}(left)M_{min}(right).$$
(6)
We have calculated the distribution of these values as before:
$`P_{\mathrm{\Delta }M_{min}}(\mathrm{\Delta }M_{min}<\mathrm{\Delta }M_{obs})=`$
$`{\displaystyle }dm{\displaystyle }dS{\displaystyle }dv_{eff}(p_s(S|m,v_{eff})`$
$`\times p_m\left(m\right)p_v\left(v_{eff}|m\right)`$
$`\times P_{\mathrm{\Delta }M_{min}}(\mathrm{\Delta }M_{min}<\mathrm{\Delta }M_{obs}|S,m,v_{eff})).`$ (7)
The functions were calculated corresponding to cusps, $`ve`$ caustic crossings and $`+ve`$ caustic crossings ($`P_C`$, $`P_{}`$ and $`P_+`$ respectively). The resulting curves are plotted in Fig. 6 (for the case of no systematic error in the source size estimate). The results are summarised in column 7 of Tab. 1.
While the trailing minimum has not yet been observed, it may provide a discriminate between the interpretations of the 1999 image C event as a cusp and a $`ve`$ caustic. In particular, if the light-curve flattens out at a level approximately equal to that of the 1998 season ($`\mathrm{\Delta }M_{min}0`$), then the $`ve`$ caustic crossing interpretation will be ruled out at the 95% level. Assuming that the previous light-curve minimum occurred during 1998, the $`+ve`$ caustic crossing interpretation is already ruled out at $`>95\%`$.
### 4.3 Comparison with the 1998 image A event
The first unambiguous microlensing signal was the rise and fall of image A in 1988 during a period when the other images remained at a consistent level (ie. no intrinsic variability). Some insight is gained through comparison of the shape of this light-curve peak with the 1999 image C event. The two events are described by very different data sets. While the OGLE data has provided excellent coverage of the 1999 peak, allowing quantities such as the maximum derivative and peak height to be calculated, the 5 observations of the 1988 peak provide a much cruder record. In particular the values of maximum pre-peak derivative, the peak height above the minimum and the difference between the minimum cannot be computed. On the other hand, our source size estimation was determined from the assumption that the image A peak is a caustic crossing. Statistical determinations of the event type are therefore only meaningful as a check of the self-consistency of the calculations. In that vein we note that the one leading derivative that can be measured from the data is $`T5`$ mags/year. Calculations in WWTA00 suggest that the observations should have been followed by a caustic crossing in $`23`$ weeks or a cusp event in $`5`$ weeks, which lie between and after the two brightest observations respectively.
Fig. 7 shows light-curves of the two peaks, placed on the same time-axis such that the maxima approximately coincide at $`t=0`$. The image A and C peaks are shown by thick and thin lines respectively. The large initial rise of the 1988 event measures the minimum to the maximum gradient (which is surely significantly larger given the high second derivative). This sharp rise is not replicated in the 1999 event which had a maximum derivative of $`2`$ magnitudes per year. Since we have a lower bound for the maximum derivative of the 1988 image A event which is consistent with all three types of HME, the size of the observed derivative cannot be used as a discriminate (although for a cusp event there is a 90% chance of $`\dot{M}_{max}<5`$ magnitudes per year). However, the leading three points also measure a lower limit for the maximum second derivative. Fig. 8 shows probability $`P_{\ddot{M}}`$ for the maximum value of second derivative on the leading sides of cusp events (solid line) and $`+ve`$ (dashed line) and $`ve`$ (dotted line) caustic crossings computed for image A (no systematic error in source size was assumed):
$`P_{\ddot{M}}(\ddot{M}<\ddot{M}_{obs})=`$
$`{\displaystyle }dm{\displaystyle }dS{\displaystyle }dv_{eff}(p_s(S|m,v_{eff})p_m\left(m\right)`$
$`\times p_v(v_{eff}|m)P_{\ddot{M}}(\ddot{M}<\ddot{M}_{obs}|S,m,v_{eff})).`$ (8)
The derivatives were calculated using a sampling rate corresponding to the initial 3 observations of the 1988 peak. Comparing the image A second-derivative of $`\ddot{M}_{obs}0.5`$ magnitudes per year per day to Fig. 8 we infer that the 1988 peak was probably a $`+ve`$ caustic crossing.
A second feature to be noted from Fig. 7 is that the 1999 peak appears to have a much longer duration than the 1988 peak. Fig. 9 shows scatter plots of event height above full-width-at-half-maximum (fwhm) vs. fwhm for images A (light dots) and C (dark dots). The plot on the left shows the relationship for caustic crossings while the right-hand plot is for cusp events. The plots highlight the rise-time - peak-height correlation for caustic crossings, and the cloud of points due to smooth light-curve variations (both of which were pointed out by Witt & Mao (1994)). The separation of the points into two categories demonstrates the intuitive notion that cusp events have longer durations than caustic crossings. The cusp event parameters shows a peak width lower-limit corresponding to the caustic crossing correlation. Therefore, the systematic bias introduced into the source size determination by the assumption that the 1988 event was due to a caustic crossing can only result in an over estimate of source size. In addition, Fig. 9 demonstrates that a cusp event contains very little information on source size/transverse velocity due to the lack of any correlation with peak height. The inference that the 1988 peak was a caustic crossing and the 1999 peak was a cusp event is consistent with Fig. 9.
### 4.4 A 1997-1998 event in the OGLE image C light-curve ?
Fig. 2 shows a $`\mathrm{\Delta }M_{obs}0.8`$ magnitude rise between the 1997 image C minimum and the 1998 level, suggesting an event in between those observing seasons. We assume that the intrinsic source luminosity was approximately constant over this period, which is supported both by the facts that image A changed by $`0.2`$ magnitudes and that the other images show opposite trends. Comparing the image C change to the probabilities in Fig. 6 we find that $`\mathrm{\Delta }M_{obs}`$ is consistent with a $`+ve`$ caustic crossing having occurred between the 1997 and 1998 observing seasons, but rules out $`ve`$ caustic crossings ($`99\%`$) and cusp events ($`95\%`$). We therefore infer that a $`+ve`$ caustic crossing was missed between the 1997 and 1998 observing seasons.
## 5 Predictions of future HMEs
In this section we calculate probability functions for the likelihood of observing future HMEs in images A and C given current light-curves.
### 5.1 The next image C HME
In this section we assume that there was a $`+ve`$ caustic crossing between the 1997 and 1998 observing seasons, and investigate when we should next see a caustic crossing in image C. These calculations follow Witt, Kayser & Refsdal (1993) who calculated the separations in dimensionless units of the different combinations of $`+ve`$ and $`ve`$ events. However we have included both our estimates of $`m`$ and $`v_{eff}`$, and the cusp as a third class of HME.
Due to the typical diamond formation of fold caustics, the case of a $`+ve`$ followed by a $`ve`$ caustic crossing is common. Similarly, inspection of model light-curves shows that cusp events follow $`ve`$ caustic crossings as the source moves past the cusp associated with that caustic (this feature is seen in the double horned profile that is characteristic of the Chang-Refsdal lens). However we have inferred that the OGLE image C light-curve shows a $`+ve`$ caustic crossing followed by a cusp event. Such a combination is much less common and is due to the source moving past a cusp formed from two fold caustics other than the one responsible for the $`+ve`$ caustic crossing HME. It can also be seen in model light-curves. Examples of the two scenarios are shown in Fig. 10. The upper panel shows an example of a double horned event. The lower two panels show examples of a double horned event surrounding a cusp event. The source is shown passing very close to the cusp (centre panel), partly coming in contact, and passing further away, producing a lower amplitude event (lower panel). The light-curves were produced using our most likely model for the microlensing parameters ($`v_{eff}=300kmsec^1`$, $`m=0.1M_{}`$ and $`S=5\times 10^{14}cm`$). The intervals (eg. 1-2) quoted below refer to the intervals between the events labelled on these plots.
From these two scenarios, we generate the probability functions for 4 different event separation statistics (plotted in Fig. 11). The left hand and right hand panels show functions computed for transverse velocities aligned with the C-D and A-B axes respectively. Fig. 11 also contains a key corresponding to the intervals between event types shown in Fig. 10. These intervals are defined below.
Firstly, we assume that the 1999 peak was a $`ve`$ caustic crossing.
$`i)`$ $`interval`$ 1-2: The dashed lines show cumulative probabilities for the separation of two adjacent caustic crossing HMEs where the first is a $`+ve`$ caustic crossing.
Secondly, the 1999 peak is interpreted as a cusp event following a $`+ve`$ caustic crossing.
$`ii)`$ $`interval`$ 1-3: The dot-dashed line shows the distribution of caustic separations in a double horned event where a cusp lies between the two caustics. In this case the typical separation of caustics is larger due to the fact that the cusp is generally formed from independent caustics, so that the chance of a cusp lying inside a diamond or another cusp is higher if the separation (and therefore area) is larger.
$`iii)`$ $`interval`$ 1-4: The dotted lines represent the probability for the separation between a cusp and the (immediately) preceding $`+ve`$ caustic crossing.
$`iv)`$ $`interval`$ 4-3: The solid lines represent the probability for the separation between a cusp event (that has followed a $`+ve`$ caustic crossing) and the subsequent caustic crossing.
As expected $`iii)`$ and $`iv)`$ are very similar. The separation of an inferred $`+ve`$ caustic crossing between the 1997 and 1998 seasons and the 1999 image C peak is consistent with Fig. 11 whether the 1999 peak is interpreted as a $`ve`$ caustic crossing or a cusp event. However if the 1999 peak is interpreted as a cusp, we expect another caustic crossing to follow. Fig. 11 shows that in this case, the separation between the cusp peak and the subsequent caustic is likely to be $`5`$ years (although may be significantly longer).
These statistics for the separation of caustic crossings that surround a cusp does not make use of the information that the observed cusp event followed the caustic crossing by $`500`$ days. An alternative way of analysing the separation between the cusp event peak and the next caustic crossing is to calculate the ratio of the time between the cusp event peak and last caustic crossing and the time between the cusp event peak and next caustic crossing. The probability of this ratio is plotted in Fig. 12. The typical ratio is 1 and we expect it to be between $`\frac{1}{4}`$ and $`4`$, yielding a most likely arrival time for the next event of $`500`$ days, and an upper limit of $`2000`$ days ($`90\%`$). Also plotted in Fig. 12 is the probability of the ratio for a random position (dashed lines). The agreement of the curves illustrates the independence of the positions of the cusps and unassociated caustics.
### 5.2 The next image A HME
In addition to the remarkable peak observed during 1999 in the image C light-curve, the OGLE data also shows image A brightening over the entire season, with the most rapid variations occurring in the latter observations. We have applied the triggering function (as described in Sec. 4.1.3) to the image A light-curve.
On the 30th of October 1999 monitoring by OGLE (OGLE web page) showed rises in image A, of 1.41-1.88 mags/year (we used observations on the 20th of October, 30th of October and 9th of November). Fig. 13 shows the resulting triggering functions $`F_+`$, $`F_{}`$ and $`F_C`$. Functions are shown assuming that the source size estimate $`S`$ is correct. The solid, dotted and dashed lines correspond to $`F_C(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$, $`F_{}(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$, $`F_+(P|T_{obs}\pm \mathrm{\Delta }T_{obs})`$ respectively. We find that the results are similar to those obtained for the June image C trigger. The observed trigger precedes a caustic crossing HME $`80\%`$ of the time. The event peak is most likely to occur $`13`$ months following such a trigger if the event was a cusp, $`2`$ months later if it was a $`ve`$ caustic and $`1`$ month later if it was a $`+ve`$ caustic crossing. Having observed a derivative in the quoted range means that a $`+ve`$ caustic crossing is very unlikely, and a $`ve`$ caustic crossing is the most likely option. Unfortunately, these results predict that an event will occur in image A between the 1999 and 2000 observing seasons. If the impending event is assumed to be a $`ve`$ caustic crossing, with a previous minimum occurring during the 1998 season, then from Fig. 6 we predict that the image A light-curve should have a subsequent minimum at a level $`11.5`$ magnitudes fainter than the November 1999 level.
If the source is small with respect to $`\eta _o`$ and therefore the inter-caustic spacing, and the brightening of image A is due to the imminent disappearance of a pair of critical images, then the rise can be modelled using the near caustic approximation of Chang and Refsdal (1979): The flux $`f_p`$ of a point source at a small time $`\mathrm{\Delta }t=t_{caust}t`$ from a fold caustic is
$$f_p=f_o+\theta (\mathrm{\Delta }t)\frac{a_o}{\sqrt{\mathrm{\Delta }t}},$$
(9)
where $`f_o`$ is the magnification of the non-critical images, $`\theta (\mathrm{\Delta }t)`$ the Heaviside step function and $`a_o`$ a constant describing the strength of the caustic. The choice of which points to use in a fit of this type is somewhat arbitrary. We have chosen the data following JD 1400, which is after the apparent inflection in the light-curve. The fit is shown in Fig. 14, giving parameter values of $`f_o=0.40mJy`$, $`a_o=3.50mJydays^{\frac{1}{2}}`$ and $`t_{caust}=1554days`$ ($``$ 11th January 2000). The final figure is of particular value since it predicts the time of the caustic crossing. $`t_{caust}`$ agrees with the most popular value for the caustic arrival time according the triggering calculation.
## 6 Light-curve interpretations
While it is pointless to try and fit the observed light-curves directly either by a functional form (except very near a caustic crossing), or with model light-curves, it is illustrative in summary to draw by hand possible interpretations of the data. Fig. 15 shows data for images A and C from Wozniak et al. (2000a,b) and the OGLE web page. The top panel shows possible light-curves corresponding to our most probable interpretation. The solid curve shows the inferred image C event at $`4200`$ days, the 1999 light-curve peak as a cusp, and the potential second caustic crossing discussed in Sec. 5.1. The dashed line shows the potential $`ve`$ caustic crossing expected following the rise in the light-curve at the end of the 1999 observing season. The lower panel in Fig. 15 shows schematics of less favoured options: the interpretation of the 1999 image C peak as a $`ve`$ caustic crossing, and the rise in image A heralding a $`+ve`$ caustic crossing.
## 7 conclusion
We have applied simple event statistics a postiori to features in the OGLE light-curves of Q2237+0305. In the specific case of the 1999 peak we conclude that the event was due to the source passing outside of, but close to a cusp rather than to a caustic crossing. This hypothesis may be confirmed or refuted when the trailing peak minimum is observed during the 2000 observing season. In addition, we find that the image C light-curve rise between the 1997 and 1998 OGLE observing seasons was a caustic crossing that resulted in two new critical images. The hypothesis of a cusp event following the first half of a double horned event is a rare feature. However its consequence is that we expect another caustic crossing high magnification event to follow in the image C light-curve. Our models predict that this caustic crossing is most likely to arrive in $`500`$ days, and can be expected within $`2000`$ days ($`90\%`$ confidence). However it may be considerably longer before the next image C caustic crossing is observed, particularly if the transverse velocity moves perpendicular to the shear vector in image C. By applying the triggering function developed in WWTA00 to the rise in the image A OGLE light curve, we predict that a caustic crossing high magnification event will occur between the 1999 and 2000 observing seasons which will result in the loss of a pair of critical images. We therefore expect that the image image A light-curve should have a minimum at a level $`11.5`$ magnitudes fainter than November 1999.
## 8 acknowledgements
The authors would like to thank Przemyslaw Wozniak for many useful discussions. We would also like to acknowledge the OGLE collaboration for making their monitoring data publically available before publication. This work was supported by NSF grant AST98-02802. JSBW acknowledges the support of an Australian Postgraduate Award and a Melbourne University Postgraduate Overseas Research Experience Award. |
warning/0001/nucl-th0001008.html | ar5iv | text | # Aspects of thermal and chemical equilibration of hadronic matter Work supported by BMBF, GSI Darmstadt and DFG.
## 1 Introduction
Nucleus-nucleus collisions at relativistic and ultrarelativistic energies are studied experimentally and theoretically to obtain information about the properties of hadrons at high density and/or temperature as well as about the phase transition to a new state of matter, the quark-gluon plasma (QGP). In the latter deconfined partons are the essential degrees of freedom that resolve the underlying structure of hadrons . Whereas the early ’big-bang’ of the universe most likely evolved through steps of kinetic and chemical equilibrium, the laboratory ’tiny bangs’ proceed through phase-space configurations that initially are far from an equilibrium phase and then evolve by fast expansion. These ’specific initial conditions’ – on the theoretical side – have lead to a rapid development of nonequilibrium quantum field theory and nonequilibribrium kinetic theory . Presently, semiclassical transport models are widely used as approximate solutions to these theories and practically are an essential ingredient in the experimental data analysis. For recent reviews we refer the reader to Refs. .
On the other hand, many observables from strongly interacting systems are dominated by many-body phase space such that spectra and abundances look ’thermal’. It is thus tempting to characterize the experimental observables by global thermodynamical quantities like ’temperature’, chemical potentials or entropy . We note, that even the use of macroscopic models like hydrodynamics employs as basic assumption the concept of local thermal and chemical equilibrium. The crucial question, however, how and on what timescales a global thermodynamic equilibrium can be achieved, is presently a matter of debate. Thus nonequilibrium approaches have been used in the past to address the problem of timescales associated to global or local equilibration . In view of the increasing ’popularity’ of thermodynamic analyses a thorough microscopic reanalysis of this questions appears necessary especially for nucleus-nucleus collisions at ultrarelativistic energies that aim at ’detecting’ a phase transition to the QGP.
In this paper we study equilibration phenomena in ’infinite’ hadronic matter using a microscopic transport model that contains both hadron resonance and string degrees-of-freedom. With this investigation we want to provide insight into the reaction dynamics by the use of cascade-like models and also point out some of their limitations. The ’infinite’ hadronic matter is modelled by initializing the system solely by nucleonic degrees of freedom through a fixed baryon density and energy density, while confining it to a cubic box and imposing periodic boundary conditions during the propagation in time. We, furthermore, then study the expansion of the hadronic fireball after equilibration to investigate the changes in hadron spectra during the rapid explosion as well as related equilibration phenomena in realistic nucleus-nucleus collisions for light and heavy systems.
Our paper is organized as follows: In Section 2 we briefly describe the approach employed in our investigations, specify the initial conditions for a finite box with periodic boundary conditions, present our numerical results and extract various (hadronic) equilibration times as well as thermodynamical properties for different initial conditions. Section 3 is devoted to the expansion dynamics of the equlibrated fireball and a discussion of the related physical phenomena. In Section 4 we analyse reactions of colliding finite light and heavy systems and compare our result to a blast model. Section 5 concludes our study with a summary.
## 2 Equilibration and limiting temperature
To investigate the equilibration phenomena addressed above we perform microscopic calculations using the Boltzmann-Uehling-Uhlenbeck (BUU) model of Refs. . This model is based on the resonance concept of nucleon-nucleon and meson-nucleon interactions at low invariant energy $`\sqrt{s}`$ , adopting all resonance parameters from the Manley analysis . The high energy collisions – above $`\sqrt{s}`$ = 2.6 GeV for baryon-baryon collisions and $`\sqrt{s}`$ = 2.2 GeV for meson-baryon collisions – are described by the LUND string fragmentation model FRITIOF . This aspect is similar to that used in the HSD approach and the UrQMD code . For a detailed description of the underlying model at low energy we refer the reader to Ref. .
For later discussions it is essential to realize that the code respects detailed balance only for reactions of the type $`12+3`$ and approximately for $`1+23+4`$ <sup>1</sup><sup>1</sup>1In the latter case small violations of detailed balance are due to the treatment of $`t`$-channel and background contributions. where the numbers $`1,\mathrm{},4`$ are any reaction partners. This implies that in particular at high energies, where the string degrees of freedom with their decay to many ($`>2`$) final particles becomes important, detailed balance is violated. We will discuss the consequences of this violation, which is inherent in all such transport codes, at the appropriate points in the following sections.
### 2.1 A box with periodic boundary conditions
In order to study ’infinite’ hadronic matter problems we confine the particles in a cubic box with periodic boundary conditions for their propagation similar to a recent box calculation within the UrQMD model . We specify the initial conditions, i.e. baryon density $`\rho `$, strange particle density $`\rho _S`$ and energy density $`\epsilon `$ as follows: first the initial system is fixed to $`N_p=80`$ protons and $`N_n=80`$ neutrons, which are randomly distributed in a cubic box of volume $`V`$. The 3-momenta $`\stackrel{}{p}_i`$ of the nucleons in a first step are randomly distributed inside a Fermi-sphere of radius $`p_F`$ = 0.26 GeV/c (at $`\rho _0`$) and in a second step boosted by $`\pm \beta _{cm}`$ by a proper Lorentz transformation. Thus the initial baryon density $`\rho `$ is fixed as $`\rho =A/V`$, $`A=N_p+N_n`$. The strange particle density is set to zero as in related heavy-ion experiments while the energy density is defined as $`\epsilon =E/V`$, where $`E`$ is the total energy of all nucleons
$`E={\displaystyle \underset{i}{\overset{A}{}}}\sqrt{p_i^2+m_N^2}.`$ (1)
The boost velocity $`\beta _{cm}`$ is related to the initial energy density $`\epsilon `$ (excluding Fermi motion) as
$`\beta _{cm}=\sqrt{1{\displaystyle \frac{\rho ^2m_N^2}{\epsilon ^2}}}`$ (2)
using
$`\epsilon =\gamma _{cm}\rho m_N`$ (3)
with $`\gamma _{cm}=1/\sqrt{1\beta _{cm}^2}`$. Recall that $`\rho _0m_N0.15`$ GeV/fm<sup>3</sup> so that an energy density $`\epsilon 1.5`$ GeV/fm<sup>3</sup> at density $`\rho _0`$ corresponds to $`\gamma _{cm}10`$, i.e. the SPS energy $`T_{lab}185`$ A$``$GeV. We thus start with a ’true’ nonequilibrium situation in order to mimique the initial stage in a relativistic heavy-ion collision. The initial phase represents two interpenetrating, (ideally) infinitely extended fluids of cold nuclear matter.
We now propagate all particles in the box in the cascade mode (without mean-field potentials) using periodic boundary conditions, i.e. particles moving out of the box are reinserted at the opposite side with the same momentum. The phase-space distribution of particles then can change due to elastic collisions, resonance and string production and their decays to mesons and baryons again. We recall that we include all baryon resonances up to an invariant mass of 2 GeV and meson resonances up to the $`\varphi `$-meson. According to the initial conditions for $`\epsilon `$ and $`\rho `$ the factor $`\gamma _{cm}`$ in (3) determines if strings are excited in the very first collisions. This is the case for $`\gamma _{cm}>1.4`$ where the early equlibration stages are dominated by string formation and decay.
### 2.2 Chemical equilibration
Figure 1 shows the time evolution of the various particle abundances (nucleons $`N`$, $`\mathrm{\Delta }`$, $`\mathrm{\Lambda }`$, $`\pi `$, $`\eta `$, $`K^+`$ and $`K^{}`$ mesons) for density $`\rho =\rho _0`$ (left panel) at different energy densities $`\epsilon =1.1,0.52`$ and 0.22 GeV/fm<sup>3</sup> and for $`\rho =3\rho _0`$ (right panel) at $`\epsilon =3.4,1.57`$ and 0.66 GeV/fm<sup>3</sup>. These initial conditions correspond to bombarding energies $`T_{lab}`$ per nucleon of roughly 100, 20 and 2 A$``$GeV, respectively. In Fig. 1 (as well as in Figs. 2,3) we count all particles which are ’hadronized’, i.e. produced by string decay after a formation time of $`\tau _F=0.8`$ fm/c in their rest frame.
After several fm/c the number of nucleons decreases due to inelastic collisions that produce either baryon resonances or additional mesons. The number of $`\mathrm{\Delta }`$-resonances grows up to a maximum in a few fm/c, since a lot of $`\mathrm{\Delta }`$’s are produced in the first $`NN`$ collisions; their number subsequently decreases with time due to their decay and excitation of further resonances or due to reabsorption. The numbers of $`\pi `$’s and $`\eta `$’s increase very fast and reach the equilibrium value within a few fm/c whereas the strange particles ($`K^+,K^{},\mathrm{\Lambda }`$) require a much longer time for equilibration.
In Fig. 2 we present the time evolution of the particle ratios $`\pi /N`$, $`\mathrm{\Delta }/N`$, $`\mathrm{\Lambda }/N`$, $`K^+/\pi ^+`$, $`K^{}/\pi ^{}`$, $`\eta /\pi `$ for density $`\rho =\rho _0`$ at energy densities $`\epsilon =1.1,0.52`$ and 0.2 GeV/fm<sup>3</sup>, while Fig. 3 shows the same particle ratios for density $`\rho =3\rho _0`$ at energy densities $`\epsilon =3.4,1.57`$ and 0.66 GeV/fm<sup>3</sup>, respectively. The left panels in both plots correspond to the full time scale as in Fig. 1 (up to 1000 fm/c), whereas the right panels present in more detail the initial phase (up to 30 fm/c). We use the same scale for the $`y`$-axis on the right and left panels, so one can easily see that the $`\pi /N`$, $`\mathrm{\Delta }/N`$ ratios reach the equilibrium values very fast especially at low energy density since the string degrees of freedom here play a minor role and pion production basically emerges through $`\mathrm{\Delta }`$ resonance decay.
The meson-pion ratios ($`K^+/\pi ^+`$, $`K^{}/\pi ^{}`$, $`\eta /\pi `$) at high energies show a decrease in the first few fm/c and then an increase again up to the equilibrium values. This is due to the fact that the bulk of the strange mesons is produced very early (at high energy density) through string formation and decay whereas most of the pions appear later, with a delay of several fm/c, as a result of the decay of heavy vector mesons (e.g. $`\rho `$ and $`\omega `$). From the right panels of Figs. 2 and 3 one can see that in the initial stage the particle ratios containing strange to nonstrange particles – $`K^+/\pi ^+`$, $`K^{}/\pi ^{}`$, $`\mathrm{\Lambda }/N`$ – are still far off chemical equilibrium for all energies and densities and the equilibration takes up to a few hundred fm/c depending on the energy and baryon density.
For the higher energies the initial particle production proceeds via the formation and decay of string excitations. This leads in particular to a very early onset of strange particles (mainly kaons and hyperons) within the first fm/c either due to the initial strings or due to secondary or ternary baryon-baryon, meson-baryon and meson-meson induced string-like interactions (see the right panels of Figs. 2 and 3). In Ref. it was shown that these early secondary and ternary reactions can contribute up to about 50 $`\%`$ of the total strange particles obtained in a Pb + Pb reaction at CERN SPS energies and thus explain the factor of 2 in the observed relative strangeness enhancement compared to p+p reactions. This, however, does not imply that chemical equilibrium for the dominant strange particles has been achieved in this reaction, as our analysis clearly shows. In the later stages, when the system has become, more or less, isotropic in momentum space, strange particles can only be further produced by low energy hadronic reactions, which, however, have a considerable threshold and are thus strongly suppressed. This explains the long chemical equilibration times for the strange particles first demonstrated by Koch, Müller and Rafelski .
In order to define an overall chemical equilibration time we perform a fit to the particle abundances $`N(t)`$ for pions and kaons as
$`N(t)=N_{eq}\left(1\mathrm{exp}(t/\tau _{eq})\right),`$ (4)
where $`N_{eq}`$ is the equilibrium limit. The equilibration time $`\tau _{eq}`$ thus corresponds to the time $`t`$ when $`63`$% of $`N_{eq}`$ is achieved.
Figure 4 shows the equilibration time $`\tau _{eq}`$ versus energy density for $`\pi `$ and $`K^+`$ mesons at different baryon densities of $`1/3\rho _0,\rho _0,3\rho _0`$ and $`6\rho _0`$. We find that the equilibration time for pions scales as $`\tau _{eq}^\pi 1/\rho `$ or $`\mathrm{\Gamma }_\pi \rho `$, thus we present the curve only for baryon density $`\rho _0`$. Whereas $`\tau _{eq}^\pi `$ slowly grows with energy-density, $`\tau _{eq}^K`$ falls steeply with $`\epsilon `$. This marked difference is due to the fact that, on one hand, the kaon production rate increases dramatically with $`\sqrt{s}`$ whereas that of the pions, on the other hand, is more flat. With increasing energy thus more strange particles are produced through strings especially from the primary collisions with high $`\sqrt{s}`$ and the chemical equilibration is achieved faster.
In Fig. 4 we have considered an ’ideal’ situation, i.e. hadron matter at fixed energy and baryon density. In realistic heavy-ion collisions the system goes through the different stages due to interactions and expansion. However, as follows from Fig. 4, the equilibration time for strangeness is larger than 40 fm/c for all energy and baryon densities. Thus in realistic nucleus-nucleus collisions the chemical equilibration of strange particles requires also a time above 40 fm/c which is considerably larger than the actual reaction time of a few 10 fm/c or less (cf. Section 4).
The particle abundances used to extract $`\tau _{eq}`$ in Fig. 4 have been calculated without any in-medium potentials. In fact, the introduction of attractive potentials (especially for $`K^{}`$) will lower the hadronic thresholds and thus increase the scattering rate between strange and nonstrange hadrons, whereas the $`K^+`$ feels some repulsive potential and the trend goes in the opposite way. According to our calculations such in-medium modifications (in line with Ref. ) give a correction to the $`K^+`$ equilibration times by atmost 10 % and shortens the $`K^{}`$ equilibration times up to 20 % at density $`\rho _0`$.
### 2.3 Thermal equilibration and limiting temperature
In this subsection we investigate the approach to thermal equilibration. This is initially driven by the very early string phase on the momentum equilibration of the hadronic degrees of freedom, when the system is still very far from equilibrium and the energy density is sufficiently high. This one can see by looking at the quadrupole moment $`<Q_2>=<2p_z^2p_x^2p_y^2>`$ of the momentum distribution of all hadrons involved. In the left panel of Fig. 5 we present the time evolution of the quadrupole moment $`<Q_2>`$ for density $`\rho =\rho _0`$ at energy densities $`\epsilon =0.22`$, 0.3, 0.52, 0.8, 1.1 and 1.6 GeV/fm<sup>3</sup>. In order to take into account the string contributions we have counted here all particles even within the formation time. The thin solid lines indicate exponential fits of the form
$`<Q_2>(t)A_1\mathrm{exp}(t/\tau _{short})+A_2\mathrm{exp}(t/\tau _{long})`$ (5)
with two equilibration times $`\tau _{short}`$ and $`\tau _{long}`$.
The right panel of Fig. 5 shows $`\tau _{short}`$ and $`\tau _{long}`$ versus energy density $`\epsilon `$. Whereas $`\tau _{short}5`$ fm/c is roughly independent on $`\epsilon `$ the ’hadronic’ equilibration time $`\tau _{long}`$ increases with energy density. These results have to be interpreted as follows: in the initial nonequilibrium phase the string degrees of freedom are excited and decay according to many-body phase on a short time scale $`\tau _{short}`$. The string decays reduce the initial quadrupole moment (at high energy density) in time by a significant factor of about $`34`$. One can understand the result obtained for $`\tau _{short}`$ in a rather simple way. Due to our prescription of the initialization of the system the first strings on average are formed after the time $`\tau _{coll}1/((\rho /2)\sigma _{NN}v_{NN})34`$ fm/c for $`\rho =\rho _0`$. The strings then decay within their formation time $`\tau _F0.8`$ fm/c giving rise to a significant production of transversal momentum. One should point out, that according to these arguments $`\tau _{short}`$ approximately scales like $`1/\rho `$. Due to Lorentz contraction $`\tau _{short}`$ is thus considerably smaller in a real heavy-ion collision. Hence, string decays provide a very efficient source for a strong decrease in longitudinal momentum and production of transverse momentum in the very early stage of an ultrarelativistic heavy-ion collision. A decrease (increase) of the formation time $`\tau _F`$ to 0.5 fm/c (1.5 fm/c) changes $`\tau _{short}`$ on the scale of 20%, only.
After string decay, however, the emerging hadronic system still has significantly larger longitudinal than transverse momenta – the ratio increases with energy density $`\epsilon `$ – and low energy hadronic reactions are less effective in transfering longitudinal to transverse momentum or simply in production of mesons. This explains the increase of $`\tau _{long}`$ with $`\epsilon `$ in simple terms.
From the above analysis it follows that after typical relaxation times of $`\tau _{short}`$ 5 fm/c the momentum unisotropy of hot and dense matter has dropped to $`\mathrm{e}^1`$ such that one might describe the system by simple global thermodynamical variables like temperature etc. This thermal equilibrium has to be contrasted with the chemical equilibrium which – as we have shown in the preceding subsection – is reached only after much longer times ($`40`$ fm/c for strange particles, for example).
For the equilibrated system we can extract a temperature $`T`$ by fitting the particle spectra with the Bolzmann distribution
$`{\displaystyle \frac{d^3N_i}{dp^3}}\mathrm{exp}(E_i/T),`$ (6)
where $`E_i=\sqrt{p_i^2+m_i^2}`$ is the energy of particle $`i`$. We note that at the temperatures of interest here, the Bose and Fermi distributions are practically identical to a Boltzmann distribution. We find that in equilibrium the spectra of all particles can be characterized by one single temperature $`T`$. This is demonstrated in Fig. 6 where we show the spectra of nucleons ($`N`$), pions ($`\pi `$) and kaons ($`K^+`$) as a function of the kinetic energy $`Em`$ for $`\rho =\rho _0`$ at energy densities $`\epsilon =0.52,0.8`$ and 1.6 GeV/fm<sup>3</sup> (left panel) and for $`\rho =3\rho _0`$ at energy densities $`\epsilon =0.66,1.57`$ and 2.85 GeV/fm<sup>3</sup> (right panel). Here we have averaged the spectra from 950 fm/c to 1000 fm/c in order to decrease the numerical fluctuations. The spectra of $`N,\pi ,K^+`$ here can be fitted with a single temperature $`T`$ which increases with the energy density $`\epsilon `$ for both baryon densities $`\rho _0`$ and $`3\rho _0`$. We note explicitly that the slope of the equilibrium particle spectra does not depend on the formation time $`\tau _F`$.
In Fig. 7 we display the energy density $`\epsilon `$ versus temperature $`T`$ for different baryon densities $`\rho `$: $`1/3\rho _0`$ (open down triangles), $`\rho _0`$ (full squares), $`3\rho _0`$ (full dots), $`6\rho _0`$ (full up triangles). In order to compare calculations for different baryon densities we have subtracted the baryon energy density at rest, i.e. $`m_N\rho `$ (except for Fermi motion). As seen from Fig. 7 the temperature grows with energy density up to a limiting value reminiscent of a ’Hagedorn’ temperature . From our detailed investigations we obtain for the limiting temperature $`T_s150\pm 5`$ MeV which practically does not depend on baryon density. Such a singular behavior of $`\epsilon (T)`$ for $`TT_s`$ has also been found in the box calculations in Ref. for $`\rho =\rho _0`$. Our limiting temperature is slightly higher than that in Ref. ($`T_s=130\pm 10`$ MeV) due to the different number of degrees of freedom; the model contains more resonances and uses a different threshold for string excitations. Thus, there is some phenomenological sensitivity to the hadronic zoo of particles and string thresholds employed in the model.
In Fig. 8 we show the excitation function for the ratio of string energy density to the energy density of the whole system $`\epsilon _{string}/\epsilon `$ at $`\rho =\rho _0`$ when the system has equilibrated for long times. If the equilibrated system is very dense, lower energy strings are still continuously being excited and thus – because of their subsequent decay – the strings constitute a stationary portion of the total energy of the system. The relative ratio in the energy density increases with $`\epsilon `$ up to a saturation value of $`16`$% and then stays essentially constant. This reflects that the system reaches a limiting temperature, since the relative amount of string excitations compared to resonance excitations does not change any more, whereas the number of strings as well as the number of hadrons produced increases with $`\epsilon `$. This fact one might have guessed since the string production rate in equilibrium depends only on the temperature $`T`$ characterizing the Bose/Fermi distributions in the collision terms. In addition, this constant fraction, of course, also intrinsically depends on the excitation threshold and on the chosen decay (or formation) time $`\tau _F`$ of the strings.
As pointed out above, the string degrees of freedom play an essential role for particle production at high bombarding energies since they describe the continuum excitations of the system. The number of strings created is especially high at the first stages of the collision, when the energy of baryon-baryon interactions is close to the initial energy $`\sqrt{s}`$. It decreases with time to some constant value which corresponds to the equilibrium state. Because of this string-dominance one now has to worry about possible consequences of a violation of detailed balance for these degrees of freedom. As already pointed out earlier, all hadronic cascade-type approaches use the phenomenological string picture in order to describe quantitatively energetic (soft or semi-hard) inelastic reactions above some specified $`\sqrt{s}`$-threshold. In such binary hadronic reactions typically many hadronic particles and resonances are produced, the number depending on the incident energy $`\sqrt{s}`$. The ‘back reaction’ of these particles produced from decay of an excited string (or two strings in the LUND model) leading to the formation of only two energetic hadrons again is not considered as it is statistically suppressed and difficult to describe. On the other hand, in an ‘infinite’ matter calculation these back reactions have to be taken into account in order to allow for the principle of detailed balance. This is not done here as it is technically difficult to handle; it thus represents a potential ‘Achilles heal’ in a thermodynamic analysis. However, for simulating a heavy-ion collision this deficiency is not of any major importance since the excitation of strings happens in the first moment of the reaction when the phase space is still widely open and no back reaction can occur.
### 2.4 Comparison to the statistical model
In order to investigate the equilibrium behavior of hadron matter we also compare our transport (box) calculations with a simple Statistical Model (SM) for an Ideal Hadron Gas (IHG) where the system is described by a grand canonical ensemble of non-interacting fermions and bosons in equilibrium at temperature $`T`$. All baryon and meson species considered in the transport model also have been included in the statistical model. Our main objective here is to compare our results with the Hagedorn bootstrap picture of hadronic matter .
We recall that in the SM particle multiplicities $`n_i`$ and energy densities $`\epsilon _i`$ are given by
$`n_i={\displaystyle \frac{g_i}{(2\pi \mathrm{})^3}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi p^2dp}{\mathrm{exp}\left[(E_iB_i\mu _BS_i\mu _S)/T\right]\pm 1}},`$ (7)
$`\epsilon _i={\displaystyle \frac{g_i}{(2\pi \mathrm{})^3}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi E_ip^2dp}{\mathrm{exp}\left[(E_iB_i\mu _BS_i\mu _S)/T\right]\pm 1}},`$ (8)
where $`E_i=\sqrt{p^2+m_i^2}`$ is the energy of particle $`i`$, $`B_i`$ is the baryon charge, $`S_i`$ is the strangeness, and $`g_i`$ is the spin-isospin degeneracy factor. In Eqs. (7),(8) $`\mu _B`$ and $`\mu _S`$ are the baryon and strangeness chemical potentials. Here we neglect the electric chemical potential ($`\mu _n=\mu _p=\mu _B`$) since we consider an isospin symmetric system. Note, however, that in realistic collisions of heavy-ions (like Au + Au) this reduction is no longer fully appropriate. For particles with finite spectral width we include in Eqs. (7),(8) the spectral functions $`\rho _i(m)`$ with the same parametrization for the width as in the transport model,
$`n_i={\displaystyle \frac{g_i}{(2\pi \mathrm{})^3}}{\displaystyle \rho _i(m)𝑑m\underset{0}{\overset{\mathrm{}}{}}\frac{4\pi p^2dp}{\mathrm{exp}\left[(E_iB_i\mu _BS_i\mu _S)/T\right]\pm 1}},`$ (9)
$`\epsilon _i={\displaystyle \frac{g_i}{(2\pi \mathrm{})^3}}{\displaystyle \rho _i(m)𝑑m\underset{0}{\overset{\mathrm{}}{}}\frac{4\pi E_ip^2dp}{\mathrm{exp}\left[(E_iB_i\mu _BS_i\mu _S)/T\right]\pm 1}}.`$ (10)
The energy density $`\epsilon `$, baryon density $`\rho `$ and strange density of the hole system in equilibrium then given as
$`\epsilon ={\displaystyle \underset{i}{}}\epsilon _i(T,\mu _B,\mu _S)`$ (11)
$`\rho ={\displaystyle \underset{i}{}}B_in_i(T,\mu _B,\mu _S)`$ (12)
$`\rho _S={\displaystyle \underset{i}{}}S_in_i(T,\mu _B,\mu _S)\mathrm{\hspace{0.17em}0}.`$ (13)
As ’input’ for the SM we use the same $`\epsilon ,\rho `$ and $`\rho _S`$ as in the box calculations and we obtain the thermodynamical parameters – $`T,\mu _B,\mu _S`$ – by solving the system of nonlinear equations (11),(12) and (13).
Within the SM we find that the temperature increases continuously with energy density since the continuum excitations, i.e. the string degrees of freedom, are not included (full dots in Fig. 9), whereas the box calculation with strings gives the limiting temperature (full squares in Fig. 9). Both curves in Fig. 9 have been calculated for density $`\rho _0`$.
To reproduce qualitatively our box result within the SM we have to include continuum excitations in the statistical model, i.e. a Hagedorn mass spectrum for strings as defined by
$`\rho ^{str}(m)={\displaystyle \frac{\rho _0^{str}}{m^3}}\mathrm{exp}(m/T_H),`$ (14)
where $`T_H`$ denotes the ’Hagedorn’ temperature. For $`T_H`$ we use the temperature $`T_s`$ as obtained from the box calculations, i.e. $`T_H=T_s150`$ MeV. In (14) $`\rho _0^{str}`$ is a fit parameter additionally to $`T,\mu _B`$ and $`\mu _S`$ to reproduce $`\epsilon (T)`$ from the box calculations. The string multiplicities $`n_i^{str}`$ are given by
$`n_i^{str}={\displaystyle \frac{1}{(2\pi \mathrm{})^3}}{\displaystyle \underset{m_{min}}{\overset{\mathrm{}}{}}}\rho _i^{str}(m)𝑑m{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{4\pi p^2dp}{\mathrm{exp}\left[(E_iB_i\mu _BS_i\mu _S)/T\right]\pm 1}},`$ (15)
where the lowest mass in the string excitation ($`m_{min}`$) is defined by the string threshold in the transport model: $`m_{min}=2.6m_N`$ GeV for baryon strings and $`m_{min}=2.2`$ GeV for meson strings. In our transport model we include the following strings $`i`$: baryon strings $`B=1,S=0,1,2,3`$, anti-baryon strings $`B=1,S=0,1,2,3`$, meson strings $`B=0,S=0,1`$ and anti-meson strings $`B=0,S=1`$.
Before going over to the actual analysis we point out that the limiting temperature $`T_s`$ from our string model involves somewhat different physics assumptions than the Hagedorn model at temperature $`T_H`$. $`T_s`$ should not really be identified with the ’Hagedorn’ temperature $`T_H`$, though close similarities exist. In the Hagedorn picture and for temperatures close to $`T_H`$ the abundance of ‘normal’ hadrons or known resonances stays constant with increasing energy density whereas the number and energy density of the (hypothetical) bootstrap excitations diverges for $`TT_H`$. The Hagedorn model thus assumes ‘particles’ of mass $`m\mathrm{}`$ to be populated for $`TT_H`$, that dynamically can be formed in collisions of high mass hadrons for $`t\mathrm{}`$. In contrast, our string model does not include energetic string-string interactions that might produce more massive strings. (There exist some phenomenological recipes how to incorporate such interactions .) The ’high mass’ strings decay to hadrons and, because of the detailed balance problem discussed in the last subsection, are only repopulated by binary hadron-hadron or hadron-string interactions, so that their internal energy is limited and the low-energy hadronic degrees of freedom are overpopulated. This leads to the saturation of string-energy to total energy (observed in Fig. 8) to a value of $`0.16`$ in contrast to the value of 1 in the Hagedorn model.
This, however, does not imply a fundamental inconsistency for the overall properties of the system. In perfect chemical equilibrium, like in the Hagedorn model, more strings (or hypothetical resonances) would be excited which, for lower temperatures (e.g. in a nearly isentropic expansion of the system like in heavy-ion collisions), would immediately decay into a large number of hadronic particles. The violation of detailed balance in our case thus physically describes an overpopulation of hadronic particles only in stationary equilibrium. The important point, however, is the observation that in either description the system at equilibrium can not exceed the critical temperature $`T_s`$.
As seen in Fig. 9 we achieve agreement of the extended SM and our box calculations from Fig. 8 by choosing $`T_HT_s`$ in Eq. (14). In addition, from the extended SM we can also define thermodynamical parameters such as the baryon chemical potential $`\mu _B`$. In Fig. 10 we present the resulting $`T\mu _B`$ correlation, i.e. temperature $`T`$ versus baryon chemical potential $`\mu _B`$, at fixed baryon densities (in the box calculations) of $`\rho =1/3\rho _0,\rho _0`$ and various energy densities. The open triangles and squares (connected by the dashed lines) show the result of the SM without strings at densities $`1/3\rho _0`$ and $`\rho _0`$, respectively, whereas the full triangles and squares (connected by the solid lines) correspond to the thermodynamical fit of the box calculations (at $`1/3\rho _0`$ and $`\rho _0`$) including string excitations. The errorbars indicate the uncertainty in the extraction of $`\mu _B`$ in the SM; they become larger when the system is closer to $`T_H`$ due to the divergence in the energy density integral (11). The arrow at $`\mu _B=0`$ indicates the temperature $`T_s=150`$ MeV from our box calculations. The full dots in Fig. 10 correspond to chemical freeze-out points extracted in a thermodynamical model from hadron abundances ; the open dots are the thermal freeze-out points from the momentum spectra of hadrons and two-particle correlations as taken from Ref. .
Our calculations here are for nuclear matter densities $`1/3\rho _0`$ and $`\rho _0`$ whereas the freeze-out points have been extracted from heavy-ion data; the comparison thus can be only qualitative. However, one can see the general tendency: if the continuum excitations (strings) are not included in the thermodynamical analysis, one can ’extract’ much larger temperatures at high energy density simply due to the limited number of degrees of freedom involved in the model analysis. In this respect our box result is more in line with the thermal (’kinetic’) freeze-out analysis from Ref. than with the thermodynamical analysis from Ref. that is based on particle ratios and thus on chemical freeze-out. The point to make is that at higher temperatures, like e.g. the ones obtained for a ‘chemical’ freeze-out in Ref. , the consideration of continuum excitations does make a thermodynamical analysis much less certain than at lower temperatures, like e.g. at ‘thermal’ (or kinetic) freeze-out as in Ref. , where the continuum excitations do not play any significant role.
In this context we have to mention, furthermore, that a combined experimental analysis of particle spectra and HBT radii favors even lower freeze-out temperatures (below 100 MeV ). For these freeze-out conditions the pion density (for fixed charge) drops below $`10^2`$ fm<sup>3</sup>, i.e. the average distance between two pions (of different charge) becomes large than $`4.6`$ fm, which in turn is large compared to their classical interaction radius $`r_I=\sqrt{\sigma _{\pi \pi }/\pi }`$ at all relative momenta between the two pions. Since thermal freeze-out temperatures of 90-100 MeV at SPS energies can be considered as a lower bound, the ’experimental’ points in Fig. 10 have to be taken with care.
## 3 Expanding hadronic fireballs
In realistic nucleus-nucleus collisions the system rapidly expands after the possible formation of a hot hadronic fireball. The final hadronic spectra can be changed substantially during this expansion phase, i.e. the temperature extracted from the experimentally observed slopes of the spectra also contains information about the nuclear expansion dynamics.
To investigate the expansion of the hadronic fireball we initialize the system in a box with periodic boundary conditions – as described above – and propagate the system up to 500 fm/c, when equilibrium is reached. Afterwards we let the system expand without boundary conditions. Even though this is an idealized description of the expansion phase during a heavy-ion collision we hope to learn from this scenario how the expansion stage changes the picture of perfect thermal equilibrium (for an analysis of an actual collision see the discussion in the following section).
In Fig. 11 we present the time evolution of the various particle abundances (nucleons $`N`$, $`\mathrm{\Delta }`$, $`\mathrm{\Lambda }`$, $`\pi `$, $`K^+`$ and $`K^{}`$ mesons) during the expansion for density $`\rho =\rho _0`$ (left panel) at different energy densities $`\epsilon =0.22,0.3`$ and 1.1 GeV/fm<sup>3</sup> and for density $`\rho =1/3\rho _0`$ at $`\epsilon =0.84`$ GeV/fm<sup>3</sup> (upper part in the right panel), for $`\rho =\rho _0`$ at $`\epsilon =1.6`$ GeV/fm<sup>3</sup> (middle part in the right panel) and for $`\rho =3\rho _0`$ at $`\epsilon =3.4`$ GeV/fm<sup>3</sup> (lower part in the right panel). The number of stable particles ($`N,\mathrm{\Lambda },\pi ,K^+,K^{}`$) increases during the expansion up to some asymptotic value due to string and heavy resonance decay as well as inelastic interactions. One can see that the asymptotic values are reached after a few 10 fm/c from the beginning of the expansion (depending on the initial energies and baryon densities) which is comparable to the actual reaction time in heavy-ion collisions (cf. Section 4).
In Fig. 12 we show the spectra of nucleons ($`N`$), pions ($`\pi `$) and kaons ($`K^+`$) versus the kinetic energy $`Em`$ for $`\rho =\rho _0`$ at energy densities $`\epsilon =0.22,0.3`$ and 1.1 GeV/fm<sup>3</sup> before the expansion – averaged over time from 450 fm/c to 500 fm/c – (left panel) and after the expansion – averaged from 580 fm/c to 600 fm/c – (right panel). For completeness in Fig. 13 we present the result for $`\rho =1/3\rho _0`$ at $`\epsilon =0.84`$ GeV/fm<sup>3</sup> (upper part), for $`\rho =\rho _0`$ at $`\epsilon =1.6`$ GeV/fm<sup>3</sup> (middle part) and for $`\rho =3\rho _0`$ at $`\epsilon =3.4`$ GeV/fm<sup>3</sup> (lower part). In the left panels the systems are in equilibrium; the $`N,\pi ,K^+`$ spectra show a common temperature $`T`$ whereas after the expansion the slopes of the particle spectra are different; the nucleon spectra are much harder than the pion spectra, i.e. the apparent temperature of particles (after the expansion) increases with the mass $`m`$. This effect is illustrated in Fig. 14, where we show the apparent slope $`T`$ versus $`m`$ for $`\pi ,K^+,N`$ for $`\rho =\rho _0`$ at different energy densities: $`\epsilon =0.2`$ GeV/fm<sup>3</sup> (full up triangles), $`\epsilon =0.3`$ GeV/fm<sup>3</sup> (full squares), $`\epsilon =1.1`$ GeV/fm<sup>3</sup> (full dots), $`\epsilon =1.6`$ GeV/fm<sup>3</sup> (full diamonds); and for $`\rho =3\rho _0`$ at $`\epsilon =3.4`$ GeV/fm<sup>3</sup> (open down triangles). The arrow indicates the limiting temperature $`T_s=150`$ MeV before the expansion. One can see from Figs. 1214 that the ’expansion’ temperature of particles increases also with the energy density. However, if during the equilibration phase the system reaches $`T_s`$, the ’expansion’ temperatures for different particles show a universal behaviour, i.e. practically do not depend on the energy $`\epsilon `$ as well as on the baryon density $`\rho `$. This phenomenon is due to the fact that close to $`T_s`$ the initial hadron velocity distributions, reflected in the particle momentum profile, become similar for all $`\epsilon `$ and $`\rho `$ in equilibrium.
In order to investigate the origin of the enhancement in the particle slope during the expansion we have performed several illustrative calculations: at 500 fm/c – after the system has achieved equilibrium – we i) let all resonances and strings decay; in this case we find that the slopes do not change as compared to the equilibrium phase, ii) we let the system expand without interactions (allowing only decays) and find that the slopes slightly decrease in comparison to the equilibrium phase. Both examples indicate that the slope enhancement stems basically from multiple interactions of the particles in the initial stages of the expansion phase.
For analyzing the expansion flow phenomena we have performed a fit of the particle spectra (after expansion) using the blast model of Siemens and Rasmussen . In this model all particle spectra are described by a universal formula with common thermal freeze-out parameters, i.e. a temperature $`T`$ of the fireball and a radial-flow velocity $`\beta `$:
$`{\displaystyle \frac{d^3N_i}{dp^3}}=A_i\mathrm{exp}\left({\displaystyle \frac{\gamma E_i}{T}}\right)\left[{\displaystyle \frac{\mathrm{sinh}\alpha _i}{\alpha _i}}\left(\gamma +{\displaystyle \frac{T}{E_i}}\right){\displaystyle \frac{T}{E_i}}\mathrm{cosh}\alpha _i\right],`$ (16)
where $`\gamma =(1\beta ^2)^{1/2},\alpha =\gamma \beta p_i/T`$. Here $`E_i,p_i`$ are the total energy and momentum of the considered particle $`i`$ while $`A_i`$ are normalization factors.
We now try to describe the final particle spectra after the expansion by Eq. (16) with common freeze-out parameters $`T`$ and $`\beta `$. In Fig. 15 we show the result of our least-squares fit, using the MINUIT method , for the energy densities $`\epsilon =0.3`$ GeV/fm<sup>3</sup> (upper part) and 1.1 GeV/fm<sup>3</sup> (lower part) at $`\rho =\rho _0`$. The left panel shows the contour plots for the parameter errors in the $`T\beta `$ plane (for the $`\chi _{optimal}^2+1`$ level); the dot-dashed lines stand for nucleons ($`N`$), the solid lines for pions ($`\pi `$) and the dashed lines for kaons ($`K^+`$). The full symbols indicate the ’best’ values for $`T`$ and $`\beta `$ according to the $`\chi ^2`$ criteria (squares for $`N`$, dots for $`\pi `$ and triangles for $`K^+`$). The thin solid lines in the right panel demonstrate the fit of the particle spectra within the optimal parameters from MINUIT.
Since the particle spectra cover several orders of magnitude and the low energy points contribute to $`\chi ^2`$ with a larger weight than those at high energy, we use the logarithmic $`\chi ^2`$ method to give a higher weight to the tail of the spectra in the fitting procedure, i.e. we minimize $`\chi _{\mathrm{ln}}^2=\underset{i}{}(\mathrm{ln}f(x_i)\mathrm{ln}f_0(x_i))^2`$, where $`f_0`$ represent the ’experimental’ data (i.e. the results of our box calculations), $`f`$ is the value of the fit (16) at point $`x_i`$.
One can see from Fig. 15 that the ’best’ parameters $`T`$ and $`\beta `$ (as well as the contours for the parameter errors) are quite different for $`N`$, $`\pi `$ and $`K^+`$ especially for $`\epsilon =0.3`$ GeV/fm<sup>3</sup>. So we do not find (within the ’optimal’ $`\chi ^2`$) common freeze-out parameters for all spectra simultaneously. This is similar to an analysis of experimental spectra by Peitzmann et al. . For all particles we obtain different values for $`\beta `$ and much lower temperatures $`T`$ than that of the initial fireball: $`T_{in}=107`$ MeV for $`\epsilon =0.3`$ GeV/fm<sup>3</sup> and $`T_{in}=145`$ MeV for $`\epsilon =1.1`$ GeV/fm<sup>3</sup>.
Since especially the pion spectra contain large contributions from resonance decays at low $`Em`$, we have also performed fits excluding the pion spectra for $`Em0.4`$ GeV. This procedure essentially gives lower $`\beta `$ parameters and higher values for $`T`$ (open circle in the upper panel). However, the low energy cut-off is an additional parameter that allows to ’extend’ the ($`\beta ,T)`$ values to a wider range.
Thus our analysis indicates that the final particle spectra do not allow a reliable reconstruction of freeze-out parameters within the collective flow model (16). The parameters $`T`$ and $`\beta `$ obtained from the fit are very sensitive to the low energy shape of the hadron spectra or low energy cut-off applied since this region contributes with the largest weight to the $`\chi ^2`$ minimization. On the other hand, a global ’eye’ fit with the parameters given by the ’star’ for all hadrons considered gives a quite reasonable overall description of the spectra (dashed lines, r.h.s of Fig. 15), too. A very accurate deduction of one single overall fit for all hadrons by a common temperature ($`T`$) and flow velocity ($`\beta `$) parameter (see, e.g., and references therein) seems to us thus rather ambiguous. In particular, such an analysis may indicate that thermal equilibrium has been reached to a much larger extent than is actually true.
The particle flow effect due to the expansion is demonstrated in Fig. 16 where we show the velocity distributions $`dN/d\beta `$ for nucleons ($`N`$), pions ($`\pi `$) and kaons ($`K^+`$) for $`\rho =1/3\rho _0`$ at $`\epsilon =0.84`$ GeV/fm<sup>3</sup> (upper part), for $`\rho =\rho _0`$ at $`\epsilon =1.1`$ GeV/fm<sup>3</sup> (middle part) and for $`\rho =3\rho _0`$ at $`\epsilon =3.4`$ GeV/fm<sup>3</sup> (lower part). The left panel shows the $`dN/d\beta `$ distribution at equilibrium whereas the right panel corresponds to $`dN/d\beta `$ after the expansion phase. On can see that the average velocity of the particles decreases with the mass; the pions are much faster than the nucleons. They thus leave the reaction zone at the initial stage of the ongoing and rapidly evolving expansion with a higher velocity and accelerate the slower hadrons that ’feel’ the ’pion wind’ by the multiple interactions . We recall that the pion density is very high especially at high $`\epsilon `$ such that practically all other hadrons are shifted in direction of large $`\beta `$. The same effect is shown in Figs. 12,13 by the enhancement of the slope of the nucleon spectra due to the expansion.
## 4 Reactions of finite systems
In this Section we turn to realistic nucleus-nucleus collisions with the BUU transport model. We have learned from our analysis in the previous Section that even by starting from an idealized scenario of perfect thermal equilibrium, a rapid expansion stage makes the extraction of one global temperature $`T`$ and one global expansion parameter $`\beta `$ quite ambiguous. We thus expect this to become even worse for the true situation of a relativistic heavy-ion collision, where a perfect equilibrium state at some intermediate stage cannot really be assumed. Also, we note that even for a very heavy system like Pb+Pb the fraction of effective surface layer ($`4\pi R_{eff}^2\lambda `$) to total volume is still quite sizeable, resulting in a continuous emission or evaporation of particles from the outer layers before a global freeze-out of bulk particles occurs. (For Pb+Pb reactions at SPS energies – combining a hydrodynamical evolution with a nonequilibrium picture of surface emission – it has indeed been shown that at least 25$`\%`$ of all particles are continuously evaporated before a global freeze-out has occurred .)
In Fig. 17 we show the time evolution of the particle abundances (nucleons $`N`$, $`\mathrm{\Delta }`$, $`\mathrm{\Lambda }`$, $`\pi `$, $`\eta `$, $`K^+`$) for central collision of the light system <sup>12</sup>C + <sup>12</sup>C (upper part) and heavy system <sup>197</sup>Au + <sup>197</sup>Au and <sup>208</sup>Pb + <sup>208</sup>Pb (lower part) at the low energy of 1 A$``$GeV (left panel) and the high energy of 100 and 160 A$``$GeV (right panel). The number of nucleons decreases in a few fm/c due to the inelastic collisions, whereas the number of $`\mathrm{\Delta }`$-resonances increases accordingly. At low energy the $`\eta `$-mesons and strange particles (we disregard strange particles for C + C at 1 A$``$GeV due to the low statistics) appear with a delay of a few fm/c due to the fact that they are basically produced from resonance decays (the same as pions) or from secondary pion-baryon collisions, whereas at high energy they appear earlier due to the primary production mechanism through the string formation and decay.
As seen from Fig 17 the reaction time $`\tau _{reac}`$ for 1 A$``$GeV is $`2030`$ fm/c, whereas for high energies $`\tau _{reac}`$ is shorter due to a faster expansion – $`\tau _{reac}1020`$ fm/c. It has been shown in Section 2.2 that the chemical equilibration of hadronic matter under ’ideal’ conditions (box without expansion) requires a quite long time, e.g. the equilibration time $`\tau _{eq}`$ for strange particles has been found to be larger than 40 fm/c for all energies and densities (cf. Fig. 4). In realistic central nucleus-nucleus collisions, such as Au + Au, the system expands rapidly (depending on the energy) after the compression and formation of the hadronic fireball. The number of interactions, which is the dynamical origin for equilibration, decreases correspondingly very fast with time; after a few 10 fm/c the particles are moving practically freely. Thus, the reaction time even for central Au + Au collisions is much shorter than the time required for strangeness equilibration: $`\tau _{reac}\tau _{eq}`$.
The thermal equilibration time in this energy range around 0.25 GeV/fm<sup>3</sup>, as obtained from the box calculation, is about 5–7 fm/c (see Fig. 5). Notice, however, that this calculation – because of its periodic boundary conditions – probably underestimates the equilibration time. Indeed, studies of the longitudinal and transverse temperatures ($`T_L`$ and $`T_T`$, resp.) have shown that full thermal equilibrium is reached only in the very late expansion phase, when the density $`\rho `$ has dropped already below its saturation value. After a period of 10 fm/c (after first contact) one still finds $`T_L1.5T_T`$, i.e. an anisotropy of about 40 %, considerably more than indicated in Fig. 5. At the bombarding energy of 160 A$``$GeV we find a rapid decrease of the quadrupole moment in the momentum space of all hadrons by about a factor of 3 at the scale of 5 fm/c leading to longitudinally expanding matter. In view of Fig. 5 this stretched ellipsoid in momentum space becomes isotropic only on the scale of 10–20 fm/c since low energy hadronic reactions are less effective for equilibration. Since this ’hadronic’ equilibration time is larger than the reaction time for Pb + Pb at 160 A$``$GeV a substantial anisotropy remains in the hadron momentum distributions after the collision.
Contrary to the box case the $`d^3N/dp^3`$ spectra for realistic nucleus-nucleus collisions do not follow the simple exponential behaviour (6) due to the strong longitudinal expansion; especially at high bombarding energy the particle spectra show the specific ’banana’ shape (reflecting the $`pp`$ spectra at high energies). In order to exclude this simple dynamical effect related to the longitudinal expansion, we present in Fig. 18 (right panel) the transverse mass spectra $`1/m_T^2dN/dm_T`$ at mid-rapidity ($`0.5y_{cm}0.5`$) versus $`m_Tm`$ for central Au + Au collisions at 1 A$``$GeV (upper part) and for central Pb + Pb at 160 A$``$GeV (lower part) calculated at the end of the reaction. The $`m_T`$-spectra show an exponential behaviour (excluding small $`m_T`$), however, with different slopes which can not be associated directly with a temperature of a hot fireball formed at the intermediate stages of the reaction.
In line with Section 3 we have performed a fit within the blast model (16) within an interval of unit rapidity around midrapidity using MINUIT. The results of the fit are displayed in Fig. 18. The full symbols (squares for $`N`$, dots for $`\pi `$ and triangles for $`K^+`$) correspond to the ’best’ values for $`T`$ and $`\beta `$ according to the $`\chi ^2`$ criteria. The thin solid lines in the right panel demonstrate the fit of the $`m_T`$ spectra within the optimal fit parameters (we obtain a smaller $`\chi ^2`$ within the linear $`\chi ^2`$ method ($`\chi ^2=\underset{i}{}(f(x_i)f_0(x_i))^2`$), which provides a better description of the low $`m_T`$ spectra).
Similar to Section 3 (cf. Fig. 15) we obtain (within the ’optimal’ $`\chi ^2`$ criterium) quite different freeze-out parameters for $`N,\pi `$ and $`K^+`$ spectra. In order to exclude the influence of $`\mathrm{\Delta }`$\- (and other resonance) decays on the pion spectra and to investigate the sensitivity of the freeze-out parameters to the low energy cuts applied, we performed a fit of the particle spectra using the following cut-offs: $`m_Tm>0.2`$ GeV (open symbols) and $`m_Tm>0.4`$ GeV (open symbols with crosses inside) for Au + Au at 1 A$``$GeV; $`m_Tm>0.4`$ GeV (open symbols) and $`m_Tm>0.5`$ GeV (open symbols with crosses inside) for Pb + Pb at 160 A$``$GeV. As seen from the left panel of Fig. 18 the implementation of the low $`m_T`$ cut-off leads to a substantial shift of the ’optimal’ MINUIT parameters $`\beta `$ and $`T`$ especially for pions. The $`\beta ,T`$ values for the different spectra move towards to each other when discarding the low $`m_T`$ points. For Au + Au at 1 A$``$GeV our $`\beta `$ and $`T`$ parameters agree with those extracted by the TAPS collaboration using the blast model (star in the upper left plot). Here we have to mention that the cut-off $`m_Tm>0.4`$ GeV has been applied in the experimental analysis, too . For the Pb + Pb spectra at 160 A$``$GeV our freeze-out parameters are similar to those from Kämpfer ($`T=120`$ MeV, $`\beta =0.43`$; star in the lower left plot). The dashed lines in the right panel of Fig. 18 show the fit to the particle $`m_T`$ spectra for the $`\beta ,T`$ values corresponding to the ’stars’ from the left panel. Again this ’eye’ fit gives a reasonable description of the spectra (except of the very low $`m_T`$ part).
Here we have to mention again that the extraction of freeze-out parameters from the experimental data is very sensitive to the details of the thermodynamical model applied as well as to the observables considered. For example, the analysis of SIS data at 1.0 GeV from Ref. gives thermal freeze-out parameters – $`T52`$ MeV and $`\beta 0.4`$. At SPS energies the chemical freeze-out temperature extracted in Ref. from the thermal-analysis of particle ratios is $`T168`$ MeV, whereas the analysis of particle spectra and two-particle correlations (HBT data) provides a much lower thermal freeze-out temperature $`T9095`$ MeV. For a survey different freeze-out parameters the reader is referred to Fig. 4 of Ref. ).
In view of the various uncertainties inherent in the extraction of the thermal freeze-out parameters we conclude that a full, i.e. thermal and chemical, thermodynamical equilibrium at freeze-out cannot be deduced from such an analysis.
## 5 Summary
In this paper we have performed a systematic study of equilibration phenomena and equilibrium properties of ’infinite’ hadronic matter as well as of relativistic nucleus-nucleus collisions using a BUU transport model that contains resonance and string degrees-of-freedom. The ’infinite’ hadron matter is modelled by initializing the system at fixed baryon density, strange density and energy density by confining it in a cubic box with periodic boundary conditions.
We have shown that the equilibration times $`\tau _{eq}`$ for different particles depend on baryon density and energy density. The time $`\tau _{eq}`$ for non-strange particles is much shorter than for particles including strangeness; for kaons and antikaons the equilibration time is found to be larger than $``$ 40 fm/c for all baryon and energy densities considered. The overall abundance of the dominant strange particles (kaons and $`\mathrm{\Lambda }`$’s) being produced and obtained within the BUU cascade model for heavy-ion collisions can therefore not be described by assuming a perfect chemical equilibrium as strangeness is typically still undersaturated to a quite large extent. We mention taht transport model calculations like ours can describe the yield and spectra of the produced nonstrange hadrons as well as $`K^+,K^{},\mathrm{\Lambda }`$ yields quite well at SPS energies . On the other hand, at AGS energies the measured $`K^+/\pi ^+`$ ratio in central Au + Au collisions is underestimated by about 30% . However, we have to point out that the more exotic strange particles (like the measured antihyperon yields of Ref. ) can by far not be explained within such standard hadronic multiple channel reactions. These hadronic data possibly point towards new physics.
We have, furthermore, shown that thermal equilibrium is established quickly, within about 5 fm/c at SIS energies and samewhat larger times at high energies. The inclusion of continuum excitations, i.e. hadron ’strings’, leads to a limiting temperature of $`T_s150`$ MeV in our transport approach which practically does not depend on the baryon density and energy. We have compared our results with the statistical model (SM), which contains the same degrees of freedom and the same spectral functions of particles as our transport model. We found that the limiting temperature behaviour can be reproduced in the statistical model only after including continuum excitations of the Hagedorn type, otherwise the fireball temperature extracted from the particle abundances and spectra is overestimated substantially.
Close to the critical temperature $`T_s`$, the hadronic energy densities can increase to a couple of GeV/fm<sup>3</sup>. From lattice QCD calculations one expects that a phase transition to a potentially deconfined QGP state should occur. Referring to the limiting temperature $`T_s150`$ MeV obtained, a QGP should be revealed and clearly distinguished from a hadronic state of matter if one can unambiguously prove the existence of an equilibrated and thermal phase of strongly interacting matter with temperatures exceeding, e.g., 200 MeV. The best candidates are electromagnetic probes, either direct photons or dileptons. On the other hand these are also ‘contaminated’ by hadronic background and/or preequilibrium physics. So far no thermal electromagnetic source with temperatures larger or equal than 200 MeV has been clearly identified.
We have also studied the expansion of the equilibrated hadronic fireball and found that the slope parameters of the particles after expansion increase with their mass; the pions leave the fireball much faster than nucleons and accelerate heavier hadrons by rescattering (’pion wind’). If the system before expansion is close to the limiting temperature $`T_s`$, the slope parameters for all particles after expansion practically do not depend on energy and baryon density. This is due to the fact that the particle velocity distributions in equilibrium do not change any more for $`TT_s`$. We have fitted the resulting spectra within the blast model of Siemens and Rasmussen. Our analysis shows a strong sensitivity of the $`(\beta ,T)`$ parameters on the spectral shape at low energy (or a low energy cut-off) so that no reliable parameter determination can be reported. However, a global ’eye’ fit with ’average’ $`(\beta ,T)`$ parameters describes the data reasonably well.
Additionally, we have considered the equilibration in realistic nucleus-nucleus collisions of light (C + C) and heavy (Au + Au and Pb + Pb) systems. The $`(\beta ,T)`$ parameters extracted from our calculations for Au + Au at 1 A$``$GeV agree with those extracted from the TAPS collaboration and for Pb + Pb at 160 A$``$GeV with the parameters from Ref. . Here the reaction time is a few 10 fm/c and decreases with the initial energy due to the fast expansion. Since the reaction time is much shorter than the equilibration time for strangeness, a chemical equilibrium of strange particles in heavy-ion collisions is not supported by our transport calculations. Although again simple fits within the blast model provide a decent parametrization of our transport results for the differential particle spectra (Fig. 18) a deduction of global parameters for thermal freeze-out is again found to be rather ambiguous, especially when considering also the lower momentum contributions of the various particle spectra.
## Acknowledgements
The authors are grateful for valuable discussions with V. Metag, H. Oeschler and H. Stöcker. |
warning/0001/astro-ph0001033.html | ar5iv | text | # RADIO EMISSION OF THE GALACTIC X-RAYS BINARIES WITH RELATIVISTIC JETS
## X-rays binaries with jets
The Galactic plane bright X-ray sources are close binaries, in which a neutron star (NS) or black hole (BH) traps matter from a companion, forming a accretion disk. Amongst $`250`$ X-ray binaries only $`30`$ are detected as radio sources (S$`{}_{\nu }{}^{}>1`$ mJy). Many (or all!) such bright sources resolved into the relativistic jets with detectable proper motion of bright blobs. This sample of 10-15 such X-ray binaries with jets (RJXB) is extremely interesting for understanding of the astrophysical relativistic jets phenomenon.
In this paper I discuss mainly spectral and temporal properties of the RJXB radio emission. I would like to address the reader to some remarkable reviews, devoted to properties of RJXB: Fender et al. (1997), Mirabel and Rodriguez (1998), Fender (2000a,b,c), papers from two workshops “Relativistic Jets from Galactic Sources” (1997, Jodrell Bank; 1999, Paris).
Just after the discovery the superluminal sources in 1994, ten Galactic radio-emitting X-ray binaries were unified into a sample of RJXB showing relativistic jets. In such close binaries the accretion disks are created by flowing of material from a normal star, filling its Roche lobe, to a compact one. The rate of accretion of material in such binaries could change sporadically or periodically, in accordance with orbital or precession periods. Then two opposite directed relativistic jets from poles of the accretion disk around a compact object are created because of a high accretion rate. Powerful variable radio emission of RJXB seems to form in relativistic moving matter of jets, and could be a trace of the jets formation in the Galactic sources. The causes and details of relativistic electrons generation are unclear, but the accumulated data confirm the strong shocks formation, spreading within the jets. The circumstellar envelope and strong stellar wind, forming a hot corona must be added in convincing models.
The monitoring is continued by Rossi XTE satellite as All Sky Monitoring (ASM) program during the last 150 weeks. In radio band Ryle’s interferometer at 15 GHz and Green Bank interferometer (GBI) at 2.25 and 8.3 GHz carried out monitoring of RJXB during some years. RATAN-600 could obtain complementary daily radio spectra in a wide frequency range.
## The RATAN-600 radio monitoring of RJXB
In last years we carried out the monitoring program of radio variability of RJXB: GRO J1655$``$40, LS I $`+61^{}303`$, GRS 1915+10, SS433, Cyg X-1, Cyg X-3, and CI Cam. In 1995–1999 long-time observational sets of RJXB the dynamical radio spectra at 1–31 cm wavelength range were detected. We found consistent patterns in spectral variability of RJXB. Some light curves of the RATAN monitoring of the RJXB are in Bursov & Trushkin (1995), Trushkin (1998) and Trushkin & Bursov (1999). The last paper is from the “Odessa” part of GMIC’99.
## GRO J1655-40
The dynamically resolved black hole GRO J1655$``$40 is a second (after GRS 1915+105) superluminal X-ray transient discovered by BATSE in July 1994, with 0.92c jets detected by VLBI and the VLA. This is bright emission lines variable optical object showing an orbital modulation of 2.6 days.
We observed GRO J1655$``$40 only during the decay of most powerful flare in August 1994. Fig.1 shows the light curves at five frequencies and the beginning of the flare from MOST data at 843 MHz. The straight lines are fits by a exponential law of the flare decay. The rates of flux decreasing are similar at different frequencies, but the power law spectra are steeper from 12th to 22th days from the flare beginning as seen in Fig.2.
The RATAN data indicate that GRS J1655$``$40 could be in a weak nonthermal radio shell diameter of $`6^{}`$ or 5.6 $`d/d_{HI}`$ pc (if distance $`d_{HI}=36`$ kpc) with the spectrum $`S_\nu `$\[Jy\]=0.27$`\nu _{GHz}^{0.5}`$. Such a shell does not seem to be a supernova remnant because it is much weaker than a typical Galactic SNR with such a surface brightness and diameter.
## LSI 61 303
The Be star/X-ray source LSI 61 303 was discovered as the variable source GT0236+61 in the patrol Galactic plane survey by Gregory and Taylor (1983). In Fig.3 are shown the light curves of LS I $`61^{}`$ 303 at five frequencies and the GBI data at 2.25 GHz. The x-axis is modified Julian days. We marked the orbital period phases 0.6, during which usually the radio flares (or flux maxima) occur again.
Generally the light curves are well correlation, while the latter maxima of flares come at the lower frequencies. The delay at 2.3 GHz is equal to 1–2 days from the maximum of flux at 11.2 GHz. Daily radio spectra changed quickly day by day and were flat at the beginning of the flare rise, as a consequence of the delay, and then became usual optically-thin synchrotron spectra. Both results could be explained in a model of relativistic jets, moving away from a compact object in the radio-absorbing dense thermal envelope formed by a stellar wind of the binary.
The light curves of LSI $`61^{}`$ 303 in X-ray band 2–12 keV, from RXTE satellite (quick-look data in RXTE ASM program; Levine et al., 1996) show that soft X-ray emission was lower than the detection level during the radio flux maxima, equal to $`5`$ mCrab, and just in period the radio flux minimum (MJD $`51000\pm 3`$) reached a prominent value of 20 mCrab (0.0362 mJy).
During the observational set we detected variable radio emission from the recently discovered X-ray binary CI Cam (XTE J0421+560) at 3.9 GHz. Its nonthermal radio emission was firstly detected at a level of 0.5–1 Jy at the end of March 1998, just during a powerful X-ray flare. After three months the radio flux decreased to $`50`$ mJy and the source showed small daily flux fluctuations to 20–30% around the mean value.
## GRS 1915+105
Mirabel et al. (1994) detected apparent superluminal motions in first Galactic source GRS 1915+105, X-ray and radio transient mapping of which with VLA has revealed to possess relativistic jets with true resolved velocity of 0.92c.
We detected GRS 1915+105 with RATAN-600 observations at 3-4 frequencies simultaneously. Thus we plotted daily spectra, which are often inverse with a positive spectral index of +0.4 – +0.9 during “quiet” periods. In Fig.4 radio light curves of GRS 1915+105 are shown at 2.25 and 8.3 GHz (JD = 2450900.5 is 00UT 28 March 1998). We used a semi-log plot in order to show the laws of flares decay. Unfortunately it is difficult to say which law is present really. The first and second powerful flares are best fitted by power and exponential laws, respectively while a last law is better fit in general. The spectral index changed from $`0.4`$ to $`1.0`$ for the first flare in Fig.4 and, in contrast, for second one from $`0.2`$ to $`+0.5`$.
Possible association GRS 1915+105 with SNR G45.7$``$0.4 is unclear. The X-ray binary is located at the south-west boundary of SNR. The defined size of this SNR ($`22^{}`$) is very indefinite and G45.7$``$0.4 could be wider in X-ray binary direction, as seen from the NVSS map at 1.4 GHz. The distances to SNR and GRS 1915+105 are comparable: d($`\mathrm{\Sigma }D`$)=9-10 kpc, and d(GRS)=10-12.5 kpc.
## Cyg X-3
Since 1972, when Gregory detected a powerful variable radio source, associated with the X-ray binary Cyg X-3, it is regular monitored at radio frequencies. Trushkin (1998) described in details some interesting flaring events during the 80-days monitoring set in summer of 1997. It is coincided with collaborating multi-band monitoring program. Analysis of the light curves in X-ray and radio range (McCollough et al., 1998a,b) confirmed a high correlation of the hard X-ray flux (20–100 keV) and flaring radio flux, and a anticorrelation with soft X-ray emission (2–12 keV) during powerful flaring activity. These light curves are shown by Trushkin & Bursov (1999). For the powerful flare dates of the correlation coefficients are found: $`\rho `$ (RXTE – 11 GHz) = $`0.64\pm 0.04`$, $`\rho `$ (BATSE – 11 GHz) = $`0.54\pm 0.04`$. Then for the post-flare period the picture is changed: $`\rho `$ (RXTE – 3.9 GHz) = $`+0.69\pm 0.01`$ or $`\rho `$ (RXTE – 11 GHz) = $`+0.66\pm 0.02`$. In general the soft and hard X-ray emission show the anticorrelation in total active period: $`\rho `$ (2-10 keV – 20-100 keV) = $`0.64\pm 0.01`$, that also confirms before discovered dependence, that a hardness of the X-ray emission anti-correlates with its brightness in a flare.
Probably a flaring variability is a direct evidence of jets formation and their expanding in thermal shell around Cyg X-3 (Trushkin, 1998). On other hand during the deep minimum of the radio flux (below 10 mJy) the formation of the jets are ceased temporarily, in 1–2 weeks. And we see only weak radio emission with flat spectrum from envelope.
The variability of Cyg X-1, the well-known black hole candidate, was studied at 3.9 GHz during the same set. Non-thermal radio emission weakly fluctuated in a range from 10 to 30 mJy. We could not find any significant radio flux modulation with orbital period 5.6 days, recently detected at high frequencies (Pooley et al., 1999). The orbital modulation of the soft X-ray emission was in detail investigated using two-year monitoring program data with RTXE (Wen et al., 1999).
In the 40-day set of observations in December 1998 – January 1999 flux variability of the Cyg X-3 is monitored during very low soft X-ray flux ($`80`$ mCrab in range 2-12 keV). Then optically thick radio emission has positive spectral index ($`+0.3`$) in the range 2.3–11.2 GHz. A flux changed from 40 to 160 mJy at the frequencies.
## SS433
The X-ray binary and luminous star SS 433 is a persistent bright radio source, detected first in 4C survey. Its active periods are characterized by powerful flares with the flux increasing two-ten times during 1-2 days (e.g. Bursov & Trushkin, 1995). SS433 was resolved into radio blobs on scales from $`0.005^{\prime \prime }`$ to $`3^{\prime \prime }`$, the radio structure location followed a kinematic model, constructed from “moving” emission lines, originated in two opposite directing precessing jets.
Here we can discuss only some characteristics of the powerful flares of SS433 in the last years.
In Fig.5 the radio light curves of the SS433 flares, derived with the continuum radiometric complex of the RATAN-600 radio telescope are given. The usual quiet spectrum, defined in the date range of MJD50210–50215 was subtracted.
In Fig.6 plots of the first flare decays of SS433 are given on different log/semi-log scales with two different fits for comparison. In the figures there are no a distinguished difference of the fits, but the exponential law gives other important and reliable dependence – a power law frequency dependence of a decay rate. In Fig.7 (left) this dependence is shown. The fitting gives $`\tau `$(days)=14.3 $`\nu _{GHz}^{0.4}`$. It allowed me to argue that the flare decay follows the exponential law. In Fig.7 (right) the best power law fitting of the second flare in May 1996 is shown, indicating that these dependences are steeper with increasing frequency, from -0.5 to -1.0 at 0.96 and 11.2 GHz, respectively.
In the 60-day set of observations in April 23 – June 23 1999 the strong flare of SS433 was detected during relative low soft X-ray emission state – 20 mCrab in the range 2–12 keV. The data of GBI monitoring of SS433 show that its X-ray and radio activity increase in last year. That is determined by common increasing quiet level in 1.5 times at all frequencies with constant spectral index ($`0.6`$) and increasing the frequency of powerful X-ray and radio flares. (see Trushkin & Bursov, 1999).
In Fig.5(right) radio light curves of the SS433 flares detected in May 1999. We subtracted the “minimum quiet state” that is determined far long ego and being coincident to fluxes in the beginning of June 1999: $`S_\nu `$\[Jy\] = $`1.15\nu _{\mathrm{G}Hz}^{0.60}`$. Daily radio spectra during the set indicated a negative spectral index. The flare spectrum became steeper during its decay.
A delay of flaring maxima increase with decreasing of the frequency, and is particularly traceable at 0.96 GHz. Decline of the flare (May 10–28) is exactly fitted by exponential law at all frequencies, that is reliable during twenty days at different frequencies. In contrast with the first flare in May 1996 there are no any detectable dependence of rate of decreasing $`\tau `$ ($`S_\nu S_{}exp(t/\tau )`$) upon frequency. This rate of decay seems to be similar $`\tau =6\pm 1`$ days at six frequencies.
If a light curve in the beginning of a flare are characterized by a high absorption in a thermal circumstellar envelope, then a spectrum of $`S_{}`$ does indeed indicate the initially injected distribution of the relativistic electrons. In Fig.8 dependence $`S_{}`$ on frequency is exactly followed to a power law with spectral index equal to $`0.8`$ that is usual for non-thermal cosmic sources, like quasars or radio galaxies, indicating acceleration of relativistic electrons on strong shocks in jets.
In Fig.9 shifts of maximum fluxes dependence, $`\mathrm{\Delta }T`$ on frequency are given. It is exactly followed to a power law with spectral index equal to $`0.64`$ that is usual for non-thermal cosmic sources, like quasars or radio galaxies, indicating acceleration of relativistic electrons on strong shocks in jets. It is commonly for SS433 that $`\mathrm{\Delta }T(\nu )`$ is exactly followed to a power law with spectral index being in range from $``$0.8 to $``$0.4.
The common property of RJXB is anticorrelation a radio flaring radio flux and a level of soft X-ray emission in comparison of these data. Usually X-ray flares coincided with quiet periods in radio light curves and, vice versa, rare and powerful radio flares coincided with periods of low X-ray emission, probably formed in far regions around SS433. The variable hard X-ray emission, originated from inverse Compton scattering of stellar photons off relativistic electrons responsible for the radio flux could strongly increase during flaring events.
## Conclusions
Below we summarize some general radio emission parameters in RJXB:
$``$ Sources Cyg X-3, Cir X-1, CI Cam (XTE 0421+560), 1E1740.7$``$2942, GRS 1758$``$258, GX339$``$4, LSI+$`61^{}303`$ and SS433, are radio jet X-ray binaries. Probably jets play a key role in formation of powerful non-thermal radio emission. The correlation of hard X-ray and radio emission seems to be a common feature of RJXB.
$``$ All RJXB are strongly variable X-ray and IR sources. The Cir X-1 show the neutron star and GRO J1655$``$40 is a dynamically resolved system with a black hole. Other binaries are only probable NS or BH.
$``$ VLBI observations of RJXB show multi-component structure on a scale $`0.0015^{\prime \prime }`$. High velocities, $`0.10.92`$c, are detected from the proper motion of blobs in resolved sources, often showing a apparent superluminal expansion. The relativistic electrons and magnetic fields reserved a large portion of the total power of flare.
$``$ The synchrotron spectra $`S_\nu =S_{}\nu ^\alpha `$ of RJXB are variable with flares up to 1000 times comparing with quiescent state. Also high linear polarization was detected in SS433, Cyg X-3 during the flares. The processes of thermal absorption in a dense envelope are used for explanation of the frequency-dependent delays of flare maxima.
$``$ The basic model of the synchrotron emission evolution is an adiabatic expansion of the blobs moving away from binaries, which contain the relativistic electrons and magnetic fields (Shklovski, 1960, van der Laan, 1966; then Marti et al., 1992). A conical geometry of jets and considerations of the radiative losses, synchrotron radiation and inverse Compton scattering (ICS) are a modification the basic model to satisfy the spectral and temporal dependences during the flare evolution. ICS in the hot corona around a binary could be responsible for correlation of hard X-ray and radio emission.
$``$ Monitoring of radio variability shows that the decay of flaring flux after the maximum follows: a power law $`S_\nu =S_{}t^{2p}`$, as the Shklovski-van der Laan model predicts and a exponential law $`S_\nu =S_{}e^{t/\tau }`$ the cause of which is a geometric structure of jets. Here often (but not always) $`\tau \nu ^\beta `$, where $`\beta `$ ranged from $``$0.8 to $``$0.4, thus a flare decays faster at the higher frequency. Maybe there are two different types of flare in RJXB with or without delays and different laws of decay: “flare of core” and “brightening zone” in SS433.
$``$ Quasi-periodical oscillations are detected at X-ray or radio light curves with characteristic frequencies: 0.01-1000 Hz in the black hole candidate RJXB.
$``$ Non-thermal radio halos are produced by precessing jets. Good examples are SS433 and Cir X-1, probably GRS J1655$``$40 could be in a weak radio shell, and GRS 1915+10 seems to be associated with SNR.
$``$ Similar (magnetic, hydrodynamic, relativistic) models of jet confinement are applied for the Galactic and extragalactic jets despite their different spatial scales.
Acknowledgements. Author is thankful to RFBR for supporting the project of monitoring X-ray binaries, grant N98-02-17577. |
warning/0001/math-ph0001025.html | ar5iv | text | # 1 Existence
## 1 Existence
In the sense of convergence in the space $`\mathrm{S}^{}`$ (distributions),
$$\mathrm{f}_\kappa =\underset{\mathrm{y}0,\mathrm{y}\mathrm{C}^+}{lim}\mathrm{F}_\kappa ^\mathrm{y}(\mathrm{x});\kappa =1,2,\mathrm{},\mathrm{m},$$
with $`\mathrm{F}_\kappa ^\mathrm{y}(\mathrm{x})=\mathrm{f}_\kappa ^+(\mathrm{x}+\mathrm{iy})\mathrm{f}_\kappa ^{}(\mathrm{x}\mathrm{iy})`$, $`(\mathrm{f}_\kappa ^\pm (\mathrm{x})`$ are holomorphic in tabular domains $`\mathrm{T}^{\mathrm{C}_\mathrm{R}^\pm }`$ and satisfy
$$|f(x+iy)|C(R^{},C^{})|y|^\alpha (1+|x|)^\beta $$
(1)
and
$$zR^n+i(C^{}U(0,R^{}))$$
$`\alpha ,\beta 0`$, independent of $`\mathrm{R}^{}`$ and $`\mathrm{C}^{}`$. From this, it follows that there exists in $`\mathrm{S}^{}`$ a unique boundary value
$$\mathrm{f}(\mathrm{x})=\underset{\mathrm{y}0,\mathrm{y}\mathrm{c}}{lim}\mathrm{f}(\mathrm{x}+\mathrm{iy})\mathrm{S}^{(\mathrm{m})};\mathrm{m}=\alpha +\beta +\mathrm{n}+3.)$$
Let us suppose that for arbitrary $`\phi \mathrm{S}`$ there exists a finite limit
$$\underset{y0,yC^+}{lim}F_1^y(x)\mathrm{}F_m^y(x)\phi (x)𝑑x$$
(2)
independent of the sequence $`\mathrm{y}0,\mathrm{y}\mathrm{C}^+`$. Then, since the space $`\mathrm{S}^{}`$ is dense, this limit defines a distribution in $`\mathrm{S}^{}`$ which we call the product $`\mathrm{f}_1.\mathrm{f}_2.\mathrm{}.\mathrm{f}_\mathrm{m}`$ of the distributions $`\mathrm{f}_1,\mathrm{f}_2,\mathrm{},\mathrm{f}_\mathrm{m}`$. Thus
$$f_1.f_2.\mathrm{}.f_m=\underset{y0,yC^+}{lim}F_1^y\mathrm{}F_m^y(\mathrm{in}S^{})$$
(3)
if the limit of the RHS exists and is independent of the sequence $`\mathrm{y}0,\mathrm{y}\mathrm{C}^+`$. This product is obviously commutative and associative. So the set of boundary values that are holomorphic in $`\mathrm{T}^{\mathrm{C}_\mathrm{R}^+}`$ and satisfy (1) constitute a commutative ring with unity, without zero divisors with respect to the multiplication defined above.
We note that the existence of the lim in (2) for $`\phi \mathrm{S}`$ implies the existence of the limit in (3) with respect to the norm of the functional in $`\mathrm{S}^{(\mathrm{N})}`$ for some $`\mathrm{N}`$, which depends on $`\mathrm{f}_1\mathrm{}\mathrm{f}_\mathrm{m}`$ (notice that weak convergence in $`\mathrm{S}^{}`$ implies strong convergence).
## 2 General case
Suppose now that (2) does not exist for all $`\phi \mathrm{S}`$, but that it exists for all $`\phi `$ in a closed subspace $`\mathrm{M}`$ of $`\mathrm{S}^{(\mathrm{N})}`$ for some $`\mathrm{N}`$. (Since $`\mathrm{M}`$ is closed in $`\mathrm{S}^{(\mathrm{N})}`$ it is a Banach space with norm $`_\mathrm{N}`$). From the Banach-Steinhaus theorem, (2) defines a continuous linear functional $`\overline{\mathrm{T}}`$ on $`\mathrm{M}`$. We use now the term ‘product’ $`\mathrm{f}_1.\mathrm{}.\mathrm{f}_\mathrm{m}`$ of the distributions $`\mathrm{f}_1,\mathrm{f}_2,\mathrm{},\mathrm{f}_\mathrm{m}`$ for any continuous linear functional in the space $`\mathrm{S}^{(\mathrm{N})}\mathrm{S}^{}`$ that is a continuation of $`\overline{\mathrm{T}}`$ from $`\mathrm{M}`$ to $`\mathrm{S}^{(\mathrm{N})}`$. According to the Hahn-Banach theorem, such an extension always exists but is not unique in general.
We shall concentrate now on the case of those $`\phi `$ in $`\mathrm{S}^{(\mathrm{N})}`$ that vanish together with all derivatives of order $`\mathrm{p}\mathrm{N}`$ inclusively, at $`\mathrm{x}=0`$. In this case, all continuations $`\mathrm{f}_1.\mathrm{f}_2.\mathrm{}.\mathrm{f}_\mathrm{m}`$ of $`\overline{\mathrm{T}}`$ from $`\mathrm{M}`$ onto $`\mathrm{S}^{(\mathrm{N})}`$ are given by
$$(f_1.f_2.\mathrm{}.f_m,\phi )=(\overline{T},\overline{\phi })+\underset{\kappa p}{}c_\kappa (\delta ^{(\kappa )},\phi )$$
(4)
where
$$\overline{\phi }(\mathrm{x})=\phi (\mathrm{x})\underset{\kappa \mathrm{p}}{}\phi ^{(\kappa )}(\mathrm{o})\omega (\mathrm{x})\frac{\mathrm{x}^\kappa }{\kappa ^!}$$
and $`\omega (\mathrm{x})`$ is an arbitrary function, $`\omega \mathrm{S}`$, identically equal to 1 in a neighbourhood of the point $`\mathrm{x}=0`$; the $`\mathrm{c}_\kappa `$ are arbitrary constants. (Notice that the extension (4) is actually independent of $`\omega (\mathrm{x})`$).
In conclusion, the formula (4) represents the desired result, given at the end of the preamble with $`_{\kappa \mathrm{p}}\mathrm{c}_\kappa \delta ^{(\kappa )}`$ the general solution of $`(\mathrm{f}_1.\mathrm{}.\mathrm{f}_\mathrm{m},\phi )=0`$ and $`(\mathrm{T},\overline{\phi })=(\overline{\mathrm{T}},\mathrm{x}^{\kappa +1}\psi )=(\mathrm{x}^{\kappa +1}\overline{\mathrm{T}},\psi ),\psi \mathrm{S}`$, a particular solution of $`(\mathrm{f}_1.\mathrm{}.\mathrm{f}_\mathrm{m},\phi )`$.
It is therefore shown that the solution (4) is not unique, the $`\mathrm{c}_\kappa `$ being arbitrary constants. |
warning/0001/cond-mat0001302.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In the standard rebound theory of earthquakes, deformation elastic energy is progressively stored in the crust and is suddenly released in an earthquake when a threshold is reached. The Ruina-Dieterich friction laws \[Dieterich, 1972; 1978; Ruina, 1983\] constitute the basic ingredient used to describe the interaction between the two sides of the sliding fault. Friction coefficients based on laboratory experiments \[Byerlee, 1977; Scholz, 1998\] fail to account for modern observations of strain \[Jackson et al., 1997\], stress \[Zoback et al., 1987; Zoback, 1992a; 1992b\] and heat flows \[Henyey and Wasserburg, 1971; Lachenbruch and Sass, 1980; 1988; 1992; Lachenbruch et al., 1995\] (see Sornette for a synthesis).
Resolutions of these paradoxes usually call for additional mechanisms, involving fluids \[National Research Council, 1990\], crack opening modes of slip \[Brune et al., 1993\], dynamical collision effects \[Lomnitz, 1991; Pisarenko and Mora, 1994\], frictional properties of a granular gouge model under large slip \[Scott, 1996\], space filling bearings with compatible kinematic rotations \[Herrmann et al., 1990\], hierarchical scaling \[Schmittbuhl et al., 1996\], etc.
Melosh has recently suggested that the mechanism of “acoustic fluidization” could provide an alternative to theories that invoke pressurized fluids as an explanation for why some faults appear to be weak. Fluidization usually refers to the experimental observation that granular material in the presence of an interstitial fluid can liquidify when shaken sufficiently strongly \[Russo et al., 1995\]. The liquifaction is due to the fact that granular media are first compressive for small deformation leading to an increase of the interstitial fluid pressure. This increase in turn decreases the friction between the grains that can eventually become free to shear. For example, liquifaction of sediments by resonant amplified seismic waves have been proposed to be in part responsible for the damage and collapse of certain buildings during the Michoacan earthquake, 1985 \[Lomnitz, 1987\] and for the damage in the Marina district of San Francisco during the Loma Prieta earthquake \[Bardet, 1990\].
In the acoustic fluidization mechanism \[Melosh, 1979; 1996\], no interstitial fluid is invoked. A fraction $`e`$ of the earthquake energy is released as high-frequency acoustic waves that scatter off and shake the granular gouge leading to the build-up of a local acoustic pressure. According to \[Melosh, 1996\], when this pressure becomes of the order of the overburden lithostatic pressure $`\rho gh`$, the granular gouge becomes essentially free to slip without much residual friction.
The purpose of this note is to show that there is a problem with this mechanism because it predicts a slip velocity during an earthquake more than two orders of magnitudes smaller than the typical meters per second for observed earthquakes, in contradiction to the result of Melosh . The problem stems from a confusion in the definition of dissipation and scattering lengths. We then suggest possible modifications of Melosh’s theory that could resolve this problem and which lead to a richer theory.
## 2 Summary of Melosh’s theory and useful background
### 2.1 Acoustic wave energy transport
The first ingredient is the generation and transport of high-frequency acoustic waves in the core of the fault. Melosh uses the standard diffusion equation (his equation 2) for the elastic transport of acoustic waves \[Ishimaru, 1978; Sornette, 1989a-c\] with a dissipation and a source term :
$$\frac{dE}{dt}=\frac{\xi }{4}^2E\frac{c}{\lambda Q}E+e\frac{\dot{ϵ}\tau }{2},$$
(1)
where $`E`$ is the acoustic wave elastic energy density. The diffusion coefficient $`\xi /4`$, where $`\xi `$ is called the scattering diffusivity by Melosh , can be expressed in terms of the elastic mean free path $`l_e`$ and of the transport acoustic wave velocity $`c`$ at scales below $`l_e`$. The velocity $`c`$ is of the order of the shear wave velocity \[Turner and Weaver, 1996; Van Albada et al., 1991; Van Tiggelen and Lagendijk, 1993\]. We thus get \[Ishimaru, 1978; Sornette, 1989a,c\] :
$$\frac{\xi }{4}\frac{1}{3}cl_e.$$
(2)
We stress that the term “elastic” refers to the fact that $`l_e`$ is the characteristic distance over which an acoustic wave propagates before being scattered in other directions, without any loss of energy.
The l.h.s. and first term of the r.h.s. of (1) give the diffusion equation which describes the transport of an acoustic wave in a multiple-scattering medium. The second term of the r.h.s. of (1) will be shown to describe the presence of a genuine absorption, while the last source term corresponds to the conversion to acoustic waves of a fraction $`e`$ of the mechanical work performed per unit time by the fault motion with strain rate $`\dot{ϵ}`$ and shear stress $`\tau `$.
### 2.2 Diffusive transport
The first two terms of (1)
$$dE/dt=(\xi /4)^2E$$
(3)
gives the standard parabolic diffusion equation which is based on the following processes. Once generated from a source, an acoustic wave propagates roughly ballistically over a typical distance of the order of the elastic (scattering) mean free path $`l_e`$. Over this distance, the equation governing the acoustic wave propagation is the hyperbolic wave equation
$$\frac{^2A}{t^2}=c^2^2A,$$
(4)
where the wave amplitude $`A`$ is related to $`E`$ by $`E=|A|^2`$ and the wave velocity $`c`$ may depend locally on position to reflect the heterogeneity of the medium. Due to this heterogeneity, the wave is scattered off its initial propagation path along the direction $`x`$ and its intensity in this direction $`x`$ decays as $`\mathrm{exp}[x/l_e]`$. This exponential decay of the intensity does not correspond to a genuine absorption but rather reflects the loss of acoustic energy along the direction $`x`$ to all possible scattered waves in all other directions. Mathematically, the exponential decay $`\mathrm{exp}[x/l_e]`$ can be derived from (4) using standard scattering theory \[Ishimaru, 1978\]. The conservation of acoustic energy is ensured by the fact that the sum of wave intensity over all directions of propagation remains constant.
Beyond the distance $`l_e`$, the nature of the transport of the wave intensity crosses over from ballistic (i.e. straight propagation) to diffusive, corresponding to the picture where the acoustic wave can be viewed as a superposition of random walks with typical step length equal to $`l_e`$. This means that the equation for the wave propagation changes from the hyperbolic wave equation for the wave amplitude to the parabolic diffusion equation for the wave intensity given by the first two terms of (1).
One can quantify this by the following example. Consider an acoustic wave of energy $`E_0`$ impinging on a slab of thickness $`L`$ made of heterogeneities that scatter off the acoustic wave, and whose scattering strength is quantified by the elastic mean free path $`l_e`$. Anderson and Sornette \[1989c,d\] have revisited this diffusion equation to get the transmission coefficient in this example, i.e. the acoustic energy which is transmitted to the other size of the slab, as a function of its thickness $`L`$. The result is :
$$E(L)E_0\frac{l_e}{L}.$$
(5)
Note that the decay follows the algebraic $`1/L`$ law rather than an exponential law. Furthermore, the acoustic intensity profile within the slab is linear and not exponentially decreasing :
$$E(z)E_0\frac{L+l_e/3x}{L},\mathrm{for}l_e<x<Ll_e.$$
(6)
These results (5,6) highlight that the diffusive transport of the acoustic energy due to multiple scattering event is very different from the exponential attenuation that a genuine absorption would produce.
### 2.3 Absorption
The third term $`(c/\lambda Q)E`$ of (1) quantifies genuine absorption processes. The parameter $`\lambda `$ is the acoustic wavelength and $`Q`$ is the quality factor. To see that this term reflects absorption, we consider (1) in absence of the spatial derivative $`^2E`$ and of the last source term :
$$\frac{dE}{dt}=\frac{c}{\lambda Q}E.$$
(7)
Its solution is
$$E(t)=E_0\mathrm{exp}\left(\frac{c}{\lambda Q}t\right),$$
(8)
which is very different from the energy decay (5) solely due to diffusion. It is thus clear that the term $`(c/\lambda Q)E`$ is not coming from elastic scattering but solely from genuine absorption, i.e. conversion of acoustic energy into thermal energy.
The usual definition of the quality factor $`Q`$ is \[Knopoff, 1964\]
$$Q2\pi \frac{l_a}{\lambda },$$
(9)
where $`l_a`$ is the absorption length defined by the exponential decay $`\mathrm{exp}\left(x/l_a\right)`$ of a ballistically propagating wave in an absorbing medium.
Melosh introduces a characteristic length $`l_{}`$, which he calls (misleadingly) the “scattering length”, defined by
$$l_{}\sqrt{\frac{\xi \lambda Q}{4c}}.$$
(10)
Using (2) and (9), we get
$$l_{}\sqrt{\frac{2\pi }{3}}\sqrt{l_el_a}.$$
(11)
Calling $`l_{}`$, a “scattering length”, is misleading because $`l_{}`$ is in reality the effective absorption length in the diffusive medium. To see this, we use the standard diffusion relation
$$l_{}^26\frac{\xi }{4}\tau _a$$
(12)
linking the radius of gyration $`l_{}`$ covered by a diffusing process over a time $`\tau _a=l_a/c`$ equal to the time needed for the wave to cover the real distance $`l_a`$, along its convoluted multi-scattered path. The prefactor $`6`$ holds for diffusion in a three dimensional space.
Using (2), we get $`l_{}\sqrt{2l_el_a}`$, which recovers the (11) up to a numerical factor of order unity. The expression (11) can be derived by several other methods \[Sornette, 1989d\]. What is important is that $`l_{}`$ scales as the geometrical mean of $`l_e`$ and $`l_a`$, which comes from the random walk nature of the diffusive process. Physically, in the diffusive regime, the acoustic wave energy is absorbed over the characteristic length $`l_{}`$, which stems from the fact that, to cross the distance $`l_{}`$, the wave follows random walk paths of length $`l_al_{}^2/l_e`$. This reflects that attenuation of a wave in a scattering medium is a function of both absorption of energy and scattering.
### 2.4 Feedback of the acoustic vibrations on the slip rate
The interesting idea of Melosh is that the high-frequency vibrations may shake the fault and unlock it, leading to an easier sliding motion. For this, he proposes the following effective friction equation, relating the strain rate $`\dot{ϵ}`$, the shear stress $`\tau `$ and the normalized acoustic wave energy $`\mathrm{\Psi }=E\rho c^2/(\rho gh)^2`$ :
$$\dot{ϵ}=\frac{\tau }{\rho \lambda c}[\frac{1\mathrm{erf}(\frac{1}{2\sqrt{\mathrm{\Psi }}})}{1+\mathrm{erf}(\frac{1}{2\sqrt{\mathrm{\Psi }}})}].$$
(13)
The main physical phenomenon taken into account in this equation is that, due to the acoustic shaking, the effective viscous friction $`\tau /\dot{ϵ}`$ is a decreasing function of the acoustic wave energy. This mechanism is related to the velocity weakening mechanism induced by collision between asperities that lead to a transfer of momentum from the direction parallel to the fault to the direction transverse to it \[Lomnitz-Adler, 1991; Maveyraud et al., 1998\].
Putting (13) in (1) and looking for stationary modes gives the non-linear ordinary differential equation (21) (for the case $`\eta =1`$, see below), whose analysis leads to the prediction of two rupture modes \[Melosh, 1996\].
## 3 Problem with Melosh’s theory
In order to obtain realistic values, there are some constraints that the model parameters must satisfy. The key parameter is the “regeneration” parameter
$$r\mathrm{\Sigma }^2=\frac{eQ}{2}\left(\frac{\tau }{\rho gh}\right)^2,$$
(14)
where $`\mathrm{\Sigma }=\tau /\rho gh`$ is the normalized shear stress. Melosh finds reasonable solutions only for $`r\mathrm{\Sigma }^2>2.8`$. For a typical fraction $`e0.1`$ of conversion to acoustic waves of the mechanical work performed per unit time by the fault motion and for a ratio $`\tau /\rho gh`$ as low as $`0.1`$ as suggested from observations on the San Andreas fault, this value $`r\mathrm{\Sigma }^2>2.8`$ corresponds to $`Q>5600`$. This estimation may vary by an order of magnitude with the conversion factor and the relative shear stress. However, the message is that the quality factor $`Q`$ measuring the attenuation of the acoustic waves must be high, in the range of $`10^3`$ for the acoustic waves to be self-sustained during the earthquake slip motion. This is the first condition.
On the other hand, Melosh’s theory predict the slip velocity
$$\dot{u}1.4\frac{\tau }{\rho c}\frac{l_{}}{\lambda }$$
(15)
during a typical earthquake. Using a shear stress $`\tau 10MPa`$, a density $`\rho =3000kg/m^3`$ and $`c=4km/s`$ gives $`\dot{u}1.2(l_{}/\lambda )m/s`$. Thus, a realistic slip velocity $`\dot{u}1m/s`$ requires that
$$l_{}\lambda .$$
(16)
Together with (11) and (9), this leads to
$$Q3\frac{\lambda }{l_e}.$$
(17)
This last expression (17) is totally incompatible with the above condition $`Q10^3`$, as this would lead to $`l_e\lambda /100`$ or smaller. This last condition is a physical impossibility : the elastic scattering length is always much larger than or at the extreme limit of the same order as the wavelength. The physical intuition is that a wave is defined over a length scale of the order of the wavelength (otherwise, there are no spatial oscillations) and the scattering process needs at least this scale to operate. The limit $`l_e\lambda `$ is attained only under exceptional circumstances leading to a novel phenomenon, called Anderson localization, in which the acoustic wave do not propagate anymore but oscillate locally. Extraordinary efficient scatterers are needed to reach this regime \[Sornette, 1989c\]. It is thus clear that the condition $`l_e\lambda /100`$ is utterly unphysical.
If in constrast, we put $`Q10^3`$ in (10), we get $`l_{}160\lambda `$, which from (15) leads to a maximum slip velocity $`\dot{u}7.5mm/s`$, using the numerical example of Melosh . This slow sliding velocity is unrealistic for earthquakes.
## 4 Possible remedies
A first remedy is to relax the condition used by Melosh that the acoustic pressure needs to reach the overburden pressure in order to significantly affect the fault friction. We propose that only a small fraction $`\eta `$ of it is enough to liquidify the fault.
Indeed, it is well-established experimentally \[Biarez and Hicher, 1994\] that the elastic modulii of granular media under large cyclic deformations are much lower than their static values. This effect occurs only for sufficiently large amplitudes of the cyclic deformation, typically for strains $`ϵ_a`$ above $`10^4`$. At $`ϵ_a=10^3`$, the elastic modulii are halved and at $`ϵ_a=10^2`$, the elastic modulii are more than five times smaller than their static values. As a consequence, the strength of the granular medium is decreased in proportion. Melosh (private communication) also finds in laboratory experiments that a large decrease in elastic modulus is required to fit the measured flow rate of acoustically fluidized debris. This is consistent with flow in granular material in a completely dilatent state and agrees with measurements reported in the literature for the elastic modulus of highly strained granular materials.
In absence of cohesion forces, the strength of a granular material is solely due to the effect of gravity weight that put grains in contact together and the resistance to shear is governed by Coulomb’s law according to which the shear stress at the threshold for sliding is proportional to the normal stress. The coefficient of proportionality is the friction coefficient. As we have discussed above, in presence of acoustic fluidization, the resistance to shear deformation is decreased as a consequence of the reduction of the effective elastic modulus. Correlatively, the threshold for sliding is also decreased (strength is often proportional to elastic modulus in brittle material). In order to capture these phenomena, we propose, what is maybe the simplest approximation, that the criterion for unlocking the fault is changed from Melosh’s criterion to the condition that the acoustic pressure needs only reach a fraction $`\eta `$ of the overburden pressure. Let us stress that the essence of our argument is that the acoustic energy is fed by the moving fault and thus the acoutic particle velocity adjust to the changing elastic modulus so that the overall acoustic energy is “externally” controlled by the rate of global elastic dissipation. The acoustic fluidization thus controls the sliding threshold rather than solely the acoustic particle velocity.
We need to estimate the strain created by the acoustic field. The acoustic pressure is related to the acoustic particle velocity $`v`$ by
$$p=\rho cv.$$
(18)
Assuming
$$p=\eta \rho gh,$$
(19)
this yields
$$v=\eta gh/c12m/s$$
(20)
for $`p200MPa`$, a density $`\rho =310^3kg/m^3`$, a velocity $`c=4000m/s`$ and $`\eta =0.1`$. This corresponds to an acoustic wave displacement $`u_a=v/2\pi f210^3m`$ at a frequency $`f10^3Hz`$. The corresponding strain $`u_a/w`$ is $`210^3`$ for a gouge width $`w`$ of the order of one meter \[Melosh, 1996\] over which the intense shaking occurs. These estimations suggest that Melosh’s criterion that the acoustic stress fluctuations must approach the overburden stress on the fault for acoustic fluidization to occur is too drastic and smaller shaking can reduce significantly the fault friction.
Persuing this reasoning, we see that the fundamental equation (6) in \[Melosh, 1996\] is changed into
$$\frac{d^2\mathrm{\Psi }}{d\zeta ^2}=\mathrm{\Psi }r\mathrm{\Sigma }^2[\frac{1erf(\frac{\eta }{2\sqrt{\mathrm{\Psi }}})}{1+erf(\frac{\eta }{2\sqrt{\mathrm{\Psi }}})}],$$
(21)
where $`\eta =1`$ recovers the case treated by Melosh. $`\mathrm{\Psi }`$ is the normalized acoustic energy, $`\zeta =z/l_{}`$, $`z`$ is the coordinate perpendicular to the fault, the regeneration parameter is $`r=eQ/2`$ where $`e`$ is the acoustic energy conversion efficiency, $`\mathrm{\Sigma }=\tau /\rho gh`$ and erf$`(x)`$ is the error function. We see that a factor $`\eta <1`$ implies a more effective generation of acoustic waves because the second source term of the r.h.s. is larger. But, since the bracket term saturates to one for large energies and/or small $`\eta `$, this does not lead to a significantly larger slip velocity than found above. This remains a problem of the theory.
This problem might be alleviated by treating $`e`$ self-consistently as a decreasing function of the friction coefficient, and thus of the acoustic energy density. The problem then becomes even more non-linear because it reflects in addition the dependence of the acoustic radiation efficiency of the granular gouge on the amplitude of the acoustic vibrations. In addition, the elastic modulus is also really nonlinear and it is only the tangent modulus that decreases close to the yield in the dilatent region, which suggests that the above linear formalism is not adequate and should be modified.
Further improvement could also take into account that the stochastic acoustic energy may deviate from a Rayleigh distribution \[Ishimaru, 1978; Mirlin et al., 1998\] when the medium is strong heterogeneous. This modifies the functional form of the term in bracket in eq.(21) and thus all numerical estimations.
We ackowledge stimulating discussions and correspondence with H.J. Melosh.
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warning/0001/astro-ph0001262.html | ar5iv | text | # Electron Firehose instability and acceleration of electrons in solar flares
## 1 Introduction
Particle acceleration is a phenomenon occurring at many different sites throughout the universe. An important example of particle acceleration are solar flares, offering a wide range of observations that allow one to probe electron and ion acceleration. It is now widely accepted that the hard X-ray emission observed by various spacecraft reflects the energization of almost all electrons in the flaring plasma to energies up to $`25\mathrm{k}\mathrm{e}\mathrm{V}`$. These observations and the observed magnetic fields encompassing the solar flare suggest that most of the dissipated energy is released by restructuring the magnetic field, e.g. magnetic reconnection events.
During the impulsive phase of the flare, when the most powerful energization takes place, electrons must be accelerated to mean energies of about $`25\mathrm{k}\mathrm{e}\mathrm{V}`$ at a rate of about $`10^{36}`$ electrons per second in order to sustain the observed intensity of the hard X-ray bursts. Taking the impulsive phase of a flare to last about $`10\mathrm{s}`$ and assuming an electron density of about $`10^{10}\mathrm{cm}^3`$ (Moore &Fung moorefung1972 (1972); Vaiana & Rosner vaianarosner1978 (1978)), the bulk energization must process a coronal volume of at least $`10^{27}\mathrm{cm}^3`$. Thus energization must affect a large fraction of the electron population in the flaring region.
In view of this background we want to briefly describe the processes which may be responsible for particle acceleration in solar flares. For a detailed review of possible acceleration processes in impulsive solar flares see e.g. Miller et al. (milleretal1997 (1997)) and Cargill (cargill1999 (1999)).
1) Shock Acceleration: There are two types of shock acceleration. The one referred to as ’shock drift acceleration’ involves the shock electric field that reflects and accelerates the particles moving along the shock surface. Since this mechanism is only effective when the shock normal approaches an angle of $`90^{}`$ to the background magnetic field, either the gained energy or the particle flux is very limited. It seems to be unlikely that this mechanism is responsible for the large number of accelerated particles in solar flares. The second kind of shock acceleration is called ’diffusive shock acceleration’. In this process the particles cross the shock-front several times, interacting with scattering centers on both sides of the shock. In the rest frame of the shock these centers approach each other and the particles systematically gain energy. This kind of acceleration process requires a certain initial velocity in order to become effective. The ion velocity has to exceed the Alfvén speed $`V_\mathrm{A}=\mathrm{B}_0/\sqrt{4\pi \rho }`$ while the electrons must have velocities at least above $`\sqrt{m_\mathrm{i}/m_\mathrm{e}}V_\mathrm{A}`$. This has been called the ’injection problem’.
2) Acceleration by parallel electric fields: Direct acceleration by electric fields depend on its strength compared to the Dreicer field $`E_D=(2\pi e^3n_\mathrm{e}\mathrm{ln}\mathrm{\Lambda })/(k_\mathrm{B}T)`$. If $`E>E_\mathrm{D}`$ most electrons and ions gain energy. If $`E<E_\mathrm{D}`$ only electrons in the high energy tail of the velocity distribution function will be accelerated. The limitation in both cases is the maintenance of overall neutrality of charge and pre-existing current in the acceleration region.
3) MHD turbulence: This acceleration mechanism occurs when particles interact many times with randomly moving MHD waves. Due to a slight overplus in head-on collisions the interaction results in an energy gain for the particle. As in the shock acceleration model, the acceleration by MHD turbulence suffers from an ’injection problem’: Thermal ions and electrons cannot resonate with MHD waves for typical solar pre-flare conditions.
A solution for this problem is the assumption of MHD turbulent cascades that channel the energy residing in the MHD turbulence to smaller scales and into the region where interaction with thermal particles is possible. A realization of this scenario is proposed in Miller (miller1991 (1991)), Miller & Roberts (millerroberts1995 (1995)) and Miller (miller1997 (1997)). An MHD turbulent cascade transfers the energy from large scale MHD waves to smaller scales where the energy may be absorbed by the particles. The mechanism that dissipates the wave energy into the particles is transit-time damping (Fisk fisk1976 (1976); Stix stix1992 (1992)). It is is basically a resonant Fermi acceleration of second order. Only particles in resonance with a low-amplitude MHD wave are affected. The resonance condition is the usual $`l=0`$ (or Landau) resonance given by $`\omega k_{}v_{}0`$. As the particles can only gain energy in the direction parallel to the background magnetic field, the temperature in parallel direction increases.
A preference for acceleration along the background magnetic field is a common feature of the acceleration models mentioned above. The velocity distribution thus becomes more and more anisotropic during acceleration. If energization in parallel direction is from a thermal level of some 0.1 keV to 20 keV or more but the perpendicular temperature remains constant, the anisotropy is substantial. The free energy residing in parallel direction may give rise to growth of plasma waves.
For $`T_{}^\mathrm{e}>T_{}^\mathrm{e}`$ and high beta plasmas, Hollweg & Völk (hollwegvolk1970 (1970)) and Pilipp & Völk (pilippvolk1971 (1971)) have proposed the Electron Firehose instability. This instability is an extension to higher frequencies of the (MHD) Firehose instability, originally mentioned by Parker (parker1958 (1958)). While the Firehose instability is of a completely non-resonant nature, the Electron Firehose instability involves non-resonant electrons but resonant protons.
For large anisotropy of the electron distribution, the electrons become also resonant. This instability is described in Pilipp & Benz (pilippbenz1977 (1977)) and is called the Resonant Electron Firehose instability.
Having been applied to a variety of problems, the Electron Firehose instability has not been considered to occur during electron acceleration in solar flares.
Here we investigate the threshold for growth of plasma modes resulting from acceleration and infer a prediction for the evolution of the distribution function in velocity space with respect to conditions expected to occur in solar flares.
We assume that no significant instability of Langmuir waves occurs. This is suggested by the following argument: If a large fraction of the available energy normally did go into Langmuir waves, we would expect always a radio signature orders of magnitude higher than ever observed during the impulsive phase (Benz & Smith benzsmith1987 (1987)).
Section 2 describes the techniques used to solve the dispersion equations. In the following section the results obtained are shown and the thresholds of the instability are presented. In section 4 we discuss the effect on the acceleration of electrons and conclude this work.
## 2 Method
Consider electromagnetic transverse waves of the form $`\mathrm{exp}(ikxi\omega t)`$ propagating in the direction of the background magnetic field in an electron-proton plasma. The plasma dispersion equation can then be written as
$`det\left((\text{1}k^2kk){\displaystyle \frac{c^2}{\omega ^2}}ϵ(\omega ,k)\right)=0,`$ (1)
where, according to linearized kinetic theory, the dielectric tensor $`ϵ(\omega ,k)`$ is given by
$`ϵ(\omega ,k)=`$
$`\text{1}{\displaystyle \frac{\omega _\mathrm{p}^2}{\omega ^2}}\{\text{1}{\displaystyle \underset{j}{}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle }v{\displaystyle }\times `$ (2)
$`{\displaystyle \frac{\frac{n\mathrm{\Omega }_j}{v_{}}\frac{}{v_{}}+k_{}\frac{}{v_{}}}{\omega k_{}v_{}n\mathrm{\Omega }_j}}f_j^0\}.`$
The gyrofrequency of the $`j`$th species is given by $`\mathrm{\Omega }_j=q_jB/(cm_j)`$ and $`\omega _\mathrm{p}`$ denotes the plasma frequency defined as $`\omega _\mathrm{p}=\sqrt{_j\omega _{\mathrm{p}_j}^2}=\sqrt{_j4\pi n_jq_j^2/m_j}`$. The tensor $``$ is given by the matrix
$`{\displaystyle }=`$ (3)
$`\left(\begin{array}{ccc}\left(\frac{n\mathrm{\Omega }_j}{k_{}}J_n\right)^2& i\frac{n\mathrm{\Omega }_j}{k_{}}v_{}J_nJ_n^{}& \frac{n\mathrm{\Omega }_j}{k_{}}v_{}J_n^2\\ i\frac{n\mathrm{\Omega }_j}{k_{}}v_{}J_nJ_n^{}& \left(v_{}J_n^{}\right)^2& iv_{}v_{}J_nJ_n^{}\\ \frac{n\mathrm{\Omega }_j}{k_{}}v_{}J_n^2& iv_{}v_{}J_nJ_n^{}& \left(v_{}J_n\right)^2\end{array}\right)`$ (7)
where the argument of the Bessel function $`J_n`$ is $`k_{}v_{}/\mathrm{\Omega }_j`$. $`f_j^0(v_{},v_{})`$ in equation (2) denotes the zero order distribution function in velocity space of the particle species $`j`$. In order to obtain full solutions of this equation, the computer code WHAMP (Rönnmark roennmark1982 (1982)) has been applied to the problem. The usage of this code has been facilitated by programming an interface for the programming language IDL. It is called IDLWhamp and provides the user with a comfortable tool to input parameters to the code and for management and visualization of the results.
According to the capability of the WHAMP code, the most general form of the particle distribution function $`f_j^0(v_{},v_{})`$ for each species $`j`$ is given by
$`f_j^0(v_{},v_{})=`$
$`{\displaystyle \frac{1}{\sqrt{2\pi }v_{j_{\mathrm{th}}}}}\mathrm{exp}(({\displaystyle \frac{v_{}}{2v_{j_{\mathrm{th}}}}}v_{dj})^2)\times `$ (8)
$`[{\displaystyle \frac{\mathrm{\Delta }_j}{\alpha _{1j}}}\mathrm{exp}({\displaystyle \frac{v_{}^2}{2\alpha _{1j}v_{j_{\mathrm{th}}}^2}})+{\displaystyle \frac{1\mathrm{\Delta }_j}{\alpha _{1j}\alpha _{2j}}}\times `$
$`\{\mathrm{exp}({\displaystyle \frac{v_{}^2}{2\alpha _{1j}v_{j_{\mathrm{th}}}^2}})\mathrm{exp}({\displaystyle \frac{v_{}^2}{2\alpha _{2j}v_{j_{\mathrm{th}}}^2}})\}].`$
This is the original notation used in Rönnmark (roennmark1982 (1982)) beside the choice of the thermal speed to be $`v_{j_{\mathrm{th}}}=\sqrt{(k_\mathrm{B}T_j/m_j)}`$. The $`\alpha _{1j}`$ parameter is the anisotropy $`\alpha _{1j}=T_{}^j/T_{}^j`$ of the $`j`$-th distribution function. $`\mathrm{\Delta }_j`$ and $`\alpha _{2j}`$ define the depth and size of a possible loss-cone.
We assume that the electron velocity distribution function can be described by a bi-maxwellian with different temperatures in parallel and perpendicular direction with respect to the background magnetic field. Hence for our problem the parameters $`\mathrm{\Delta }_j`$ and $`\alpha _{2j}`$ were set to unity in equation (2).
Taking into account the uncertainties in the acceleration region including a possible pre-heating mechanism altering the pre-flare plasma conditions, we do not want to restrict our work to the parameters of a particular scenario. According to Pallavicini et al. (pallaviciniseriovaiana1977 (1977)) reasonable ranges in the acceleration region of an impulsive solar flare would be $`100500\mathrm{G}`$ for the background magnetic field, $`10^910^{11}\mathrm{cm}^3`$ for the number density and $`10^610^7\mathrm{K}`$ for the temperature of electrons and protons.
## 3 Results
### 3.1 Electron Firehose Instability
The only mode exhibiting significant growth rates in our calculations is a lefthand circularly polarized wave which was identified to be the Electron Firehose instability. This mode evolves out of a stable righthand polarized whistler wave at small anisotropy. With increasing anisotropy, the frequency $`\omega _\mathrm{r}`$ is shifted so that, with the convention $`\omega _\mathrm{r}>0`$, the mode becomes lefthand circularly polarized at $`k||\mathrm{B}_0`$ in the unstable regime, cf. section 7 in Gary (gary1993 (1993)). A typical dispersion relation is plotted in Fig. 1.
By introducing the resonance factor
$`\zeta _j^\pm {\displaystyle \frac{\omega _\mathrm{r}\pm \mathrm{\Omega }_j}{\sqrt{2}|k_{}|v_{j_{\mathrm{th}}}}},`$ (9)
the values for the protons and electrons are found to be $`|\zeta _\mathrm{p}^{}|1`$ and $`|\zeta _\mathrm{e}^{}|1`$, demonstrating resonance for the protons and non-resonance for the electrons.
According to Hollweg & Völk (hollwegvolk1970 (1970)) there are also right hand circularly polarized modes, which can become unstable for this extension of the Firehose instability. These modes have been found, but the growth rates are smaller than the growth rates of the lefthand polarized modes described above.
As the instability first appears, the phase velocities of the resonant waves are near the peak of the proton distribution. Fig. 2 displays a plot of the proton resonance factor (9) for the fastest growing modes versus electron anisotropy.
With increasing anisotropy less protons are resonant and the resonance factor increases. The change in the fraction of resonant protons is mirrored in the excited frequency range. As depicted in Fig. 3 the unstable frequency range grows to a maximum value at an anisotropy of about $`\frac{T_{}}{T_{}}12`$, coinciding with the maximum value of the resonance factor at $`|\zeta _\mathrm{p}^{}|0.57`$ of the protons (cf. Fig. 2). As the resonance factor decreases again, the excited frequency range becomes narrow around $`\omega _r/|\mathrm{\Omega }_\mathrm{p}|1`$. This narrow range is in itself evidence for the resonant character of the instability.
### 3.2 Instability Threshold
In this section we present the calculated threshold for linear growth of L-mode waves excited by the Electron Firehose instability.
The initial plasma is assumed to be maxwellian with temperatures $`T_0=T_0=T_0`$ perpendicular and parallel to the background magnetic field for both plasma species, the electrons and the protons. Taking into account an acceleration mechanism for the electrons acting only in parallel direction, the perpendicular temperature remains constant throughout the whole acceleration process, i.e. $`T_{}^\mathrm{e}=T_0^\mathrm{e}`$. In order to investigate the condition of the pre-flare plasma for the Electron Firehose instability to occur during the acceleration process, the initial plasma parameters have to be connected to the actual plasma parameters during the acceleration. With the assumptions above, this can be done by defining an initial parallel plasma beta, $`\beta _0^\mathrm{e}`$, via the perpendicular plasma beta
$`\beta _0^\mathrm{e}\beta _{}^\mathrm{e}={\displaystyle \frac{8\pi n_\mathrm{e}k_\mathrm{B}T_{}^\mathrm{e}}{\mathrm{B}_{0}^{}{}_{}{}^{2}}},`$ (10)
and the usual parallel plasma beta by
$`\beta _{}^\mathrm{e}{\displaystyle \frac{8\pi n_\mathrm{e}k_\mathrm{B}T_{}^\mathrm{e}}{\mathrm{B}_{0}^{}{}_{}{}^{2}}}=\beta _0^\mathrm{e}{\displaystyle \frac{T_{}^\mathrm{e}}{T_{}^\mathrm{e}}},`$ (11)
where the connection between these two quantities is given by the temperature anisotropy $`T_{}^\mathrm{e}/T_{}^\mathrm{e}`$.
According to Hollweg & Völk (hollwegvolk1970 (1970)) the instability criterion for the Electron Firehose instability may be approximated by
$`1\beta _{}^\mathrm{e}A_\mathrm{e}<0,`$ (12)
where the anisotropy factor is defined by $`A_\mathrm{e}=1T_{}^\mathrm{e}/T_{}^\mathrm{e}`$.
As one can see from inequality (12), the instability threshold does not depend directly on the parameters $`n_\mathrm{e},T_{}^\mathrm{e},\mathrm{B}_0`$, but only on the resulting $`\beta _{}^\mathrm{e}`$. This independence is also reproduced with the numerically obtained data. For our purpose, the plasma is therefore fully described by the plasma beta.
In Fig. 4 the function $`\gamma _{\mathrm{max}}(T_{}^\mathrm{e}/T_{}^\mathrm{e})`$ is plotted for five different values of $`\beta _0^\mathrm{e}`$. The maximum growth rate of the instability steeply raises at the threshold of the instability and flattens for larger anisotropies, where $`\gamma _{\mathrm{max}}/|\mathrm{\Omega }_\mathrm{p}|`$ approaches unity.
From these results, the contour of zero growth rate in the $`A_\mathrm{e}\beta _{}^\mathrm{e}`$ plane has been derived (cf. Fig. 5). The discrepancy between the analytically derived relation (12) and the numerically obtained values is due to the approximation used in Hollweg & Völk (hollwegvolk1970 (1970)).
According to inequality (12), instability cannot occur for a parallel beta smaller than unity. Due to the deviations of the approximation mentioned above, this limit is shifted to a value of $`1.6`$ (cf. Fig. 5).
In order to investigate the necessary properties of the pre-flare plasma for the Electron Firehose instability to occur, it is the initial plasma beta that is of interest. Fig. 6 depicts the same plot as Fig. 5 but this time the anisotropy factor $`A_\mathrm{e}`$ has been plotted versus the initial plasma beta $`\beta _0^\mathrm{e}`$. The dotted line in both figures represents a fit to the numerically obtained values and is an extrapolation to a broader range of beta values. The negative $`\beta _0^\mathrm{e}`$ at the $`A_\mathrm{e}1`$ limit is an artifact of this extrapolation.
The values of the initial plasma beta for the Electron Firehose instability to occur at considerable values of $`T_{}^\mathrm{e}/T_{}^\mathrm{e}`$ are well within the range of usually assumed pre-flare plasma parameters. For example, an initial plasma beta of $`\beta _0^\mathrm{e}0.05`$ can be realized by assuming pre-flare plasma parameters of $`n_\mathrm{e}=510^{10}\mathrm{cm}^3`$, $`T_0^\mathrm{e}=310^6K`$ and $`\mathrm{B}_0=100\mathrm{G}`$. This plasma becomes unstable at an anisotropy of $`T_{}/T_{}32`$.
According to the acceleration model via transit-time damping, this is a reasonable value for the anisotropy to occur during the acceleration process (Lenters & Miller lentersmiller1998 (1998)).
### 3.3 Influence of Anisotropic Protons
If we assume the protons to be heated by the same or a similar mechanism, it is to be expected that they will grow anisotropic in the same way the electrons do. Hence, we also have investigated the influence of anisotropic protons and briefly discuss the effect of an additional proton anisotropy on the instability.
Consider a plasma with anisotropic electrons and isotropic protons that is already unstable to the Electron Firehose instability. When the protons are anisotropized by increasing the temperature in parallel direction, more and more become resonant to the L-waves, non-resonantly excited by the electrons. As shown by Hollweg & Völk (hollwegvolk1970 (1970)), the protons are damping these waves. Hence, it is to be expected that the resulting growth rate of the L-waves decreases as the proton anisotropy is increased. This expectation has been verified by numerical calculation.
Moreover, the protons are heated by absorption of the excited waves at the expense of the electrons (Pilipp & Völk pilippvolk1971 (1971)). According to Kennel & Petschek (kennelpetschek1966 (1966)) this scatters the protons to higher perpendicular velocities and hence, destroys or inverts the parallel proton anisotropy. It inhibits the growth of the electron anisotropy and may complicate the acceleration to higher energies, but increases the bulk energy of the protons. The energy transfer from the electrons to the protons via the Electron Firehose instability could be responsible for the proton energization, which is a problem in the transit-time damping scenario (Miller miller1998 (1998)).
If the protons are anisotropic, there is an additional righthand polarized wave mode. This mode is the extension of the lefthand Electron Firehose mode to negative frequencies. According to Hollweg & Völk (hollwegvolk1970 (1970)) this mode has real frequencies mainly below the proton gyrofrequency and is in resonance with the protons. As the anisotropy of the electrons increases, this righthand polarized mode becomes less and less significant.
The proton cooling through the instability of the righthand mode competes with the heating by the lefthand mode and it is not yet clear what the net energization of the protons will be.
### 3.4 Oblique Propagation: Preliminary Results
Sample calculations in oblique directions indicate an unstable branch of solutions that grows faster than the parallel Electron Firehose mode. This mode is stable at parallel propagation and is also lefthand circularly polarized. Fig. 7 depicts frequency and maximum growth rate versus the angle $`\mathrm{\Theta }`$ of both modes. The dashed lines represent the branch, that is excited by the Electron Firehose instability at parallel propagation. At a propagation angle of about $`10^{}`$ with respect to the background magnetic field, the growth rate of the oblique mode becomes larger than the growth rate of the parallel mode. Hence, the determining mode for instability thresholds is the oblique mode rather than the parallel Electron Firehose mode.
As calculations have shown, not only the growth rate of the oblique mode is larger than in the parallel case, but also the instability threshold may be lower with respect to anisotropy. Plasmas being stable with respect to the parallel mode exhibited instability to the oblique mode. Therefore, the thresholds for instability derived in the section above can be considered as upper boundaries.
Lefthand circularly polarized oblique modes seem to have never been considered in connection with the Electron Firehose instability. The further investigation of these modes is the subject of ongoing work.
## 4 Discussion and Conclusion
Numerical solutions of the dispersion equation for lefthand circularly polarized electromagnetic waves, propagating parallel to the background magnetic field, have shown that the Electron Firehose instability, usually considered as a ’high-beta plasma’ instability, must be expected in coronal plasmas which are processed by an acceleration mechanism with a preference in parallel direction. The distribution function of the electrons in velocity space has been represented by a bi-maxwellian with temperatures $`T_{}^\mathrm{e}`$ and $`T_{}^\mathrm{e}`$, perpendicular and parallel with respect to the background magnetic field. The protons have been assumed to be isotropic and in thermal equilibrium.
Considering the uncertainty in the pre-flare conditions, we have investigated instabilities in a broad range of plasma temperature $`T`$, density $`n`$ and background magnetic field $`\mathrm{B}_0`$. For these plasmas, it was found that the Electron Firehose instability occurs at anisotropies that must be expected for an acceleration mechanism acting predominantly in parallel direction and being capable of producing the observed electron bulk energization.
The unstable parallel modes that have been found are lefthand circularly polarized and non-resonantly excited by the electrons, but partially cyclotron resonant with the protons, which absorb the wave energy. Hence, energy is transferred from the electrons to the protons. Assuming the density and the magnetic field to be constant during the acceleration process, there is a limiting electron temperature in parallel direction that cannot be exceeded without loosing energy to the protons via the Electron Firehose instability. The Electron Firehose instability may thus inhibit the electron acceleration process and limit the reachable energies.
At angles $`\mathrm{\Theta }0`$ with respect to the background magnetic field, an oblique mode has been found, that exhibits even larger growth rate than the parallel Electron Firehose instability. This mode is also lefthand circularly polarized. The properties of the oblique mode open new aspects on the Electron Firehose instability and its thresholds.
Taking anisotropic protons into account, the lefthand mode may extend to negative frequencies. They correspond to the righthand circularly polarized mode resonantly excited by the protons. Due to the cooling effect of this instability it is not yet clear to what the net energization of the electrons and the protons will amount. This topic will be the subject of future particle simulations.
###### Acknowledgements.
The authors thank S. Peter Gary and James A. Miller for their helpful advice and Kjell Rönnmark for providing them with a copy of the KGI report describing the original WHAMP code. The authors also want to acknowledge Gérard Belmont and Laurence Rezeau who gave them free access to their improved version of the WHAMP code, which has become the mathematical core of IDLWhamp. This work was financially supported by the Swiss National Science Foundation (grant No. 20-53664.98). |
warning/0001/cond-mat0001080.html | ar5iv | text | # THEORY OF THE SPIN BATH
## I Introduction
Many problems in quantum physics can be discussed using a model in which one or more mesoscopic or even macroscopic coordinates $`Q,Q^{}`$, etc., interact with a background environment (one coordinate might also represent an experimental probe, or even an observer). In such models (which have a long history ) all variables, including the environmental ones, are treated quantum-mechanically. The aim is to find the behaviour of $`Q,Q^{}`$, etc., after averaging over the environmental variables in some way.
It is certainly not obvious that one can discuss the real world in this way, given the complexity of $`N`$-body systems. However we now know that many (but not all) mesoscopic or macroscopic systems can be described at low energies by a few ”canonical models”, where a simple ”central system” (eg., a 2-level system, or an oscillator) couples to an environment. Remarkably, there appear to be only two types of environment, describable as baths of either oscillators or spins. One way of trying to justify such models is the ”renormalisation group” viewpoint , which maintains that most physical systems fall into a few ”universality classes”, each scaling to its own ”fixed point” in the space of possible Hamiltonians. All systems in a given university class use the same canonical Hamiltonian- the differences between different systems lying in the different values of the relevant couplings in this Hamiltonian.
Although this ”hard RG” philosophy clearly fails in some cases, it is a useful starting point for the present article, in which the quantum environment is modelled by a ”spin bath” (usually of 2-level systems, or ”spin-$`1/2`$” systems). The finite Hilbert space of each bath spin makes the spin bath appropriate for describing the low energy dynamics of a set of localised environmental modes. We concentrate on one particular ”central spin” model in which the central system itself reduces to a 2-level system; but we also discuss cases where the central system is a one-dimensional coordinate (a ”particle”) moving through a field of spins.
Another well-known set of canonical models describes the environment as a set of uncoupled oscillators- these include the ”spin- boson” model and the ”Caldeira-Leggett” model . The spin-boson model couples a central 2-level system to the oscillators, and is thus the analogue of the central spin model; and the Caldeira-Leggett model couples a tunneling particle to the oscillators. These oscillator models all derive from a scheme proposed by Feynman and Vernon , to describe a central system coupled weakly to $`N`$ environmental modes; as they showed, the mapping to an oscillator bath can only be made rigourously if the coupling is weak. Oscillator models are thus best adapted to $`N`$ delocalised environmental modes (where the coupling is automatically $`1/N^{1/2}`$, and thus small for large $`N`$).
However, readers familiar with low-temperature physics will know that at low energies, the entropy and heat capacity of almost all real physical systems are dominated instead by local modes such as defects, impurity spins, and nuclear spins . Typically these relax very slowly at low $`T`$ because little phase space is available in their coupling to any delocalised modes (or to each other). However they often couple strongly to any mesoscopic or macroscopic collective coordinate, which then easily perturbs them. This coupling is of course independent of $`N`$.
Unfortunately, even though spin bath models have been studied sporadically for many years , the results have often been misleading, either because they treat some weak-coupling limit (sometimes made by arbitrarily multiplying the coupling to each of the $`N`$ bath spins by a factor $`1/N`$, for no good physical reason), or because they drop some of the important couplings to the bath spins, in order to solve the model. In the weak-coupling limit, spin bath models can be mapped to oscillator baths (in accordance with the original remarks of Feynman and Vernon ). However one is often nowhere near the weak-coupling limit, and the mapping to the oscillator bath then fails in general . This demands a new approach, which is the subject of this review.
It may be useful to mention why many physicists are interested in models of this kind. Here are some of the reasons:
(i) Very rapid progress in work on intrinsically quantum processes (interference, tunneling, etc.) occurring at the nanoscopic and mesoscopic scales , plus speculations about the coming ”nanotechnological revolution”. Perhaps the most exciting idea in this area is that of making ”quantum computers” using nanoscopic superconductors , semiconductors , or nanomagnets. Needless to say, the technological repercussions of this work will be enormous, provided the crucial problem of decoherence can be overcome.
(ii) Physicists need to understand the mechanisms of decoherence and quantum dissipation in nature, and the crossover to (or ”emergence” of) classical behaviour from quantum physics as either size, temperature $`T`$, external fields, or couplings to the environment are increased. These issues are not only relevant to quantum device design, but also to problems in quantum gravity, and to the infamous ”quantum measurement” problem . The existence of low-$`T`$ canonical models, going beyond the phenomenology of stochastic or master equations to work with closed Hamiltonians, is invaluable here. Recent examples include the analysis of quantum spin glass relaxation , quantum relaxation in nanomagnets and tunneling superconductors ), and the study of quantum chaotic systems, stochastic resonance and dissipative tunneling in AC fields . Earlier such models have been used for decades to discuss relaxation in fields like quantum chemistry or nuclear physics.
(iii) Both oscillator and spin bath models map to many important models in quantum field theory. Thus the ”spin-boson” model mentioned above maps, for specific parameter values, to the Kondo model, the Thirring and Sine-Gordon models, and various other 2-dimensional field theories. Although similar mappings have yet to be exploited in great detail for the spin bath, they will obviously be very useful for, eg., lattice spin models.
Most work in these areas has used oscillator bath representations of the environment, with the tacit assumption that delocalised environmental modes dominate the physics. However, experiments on quantum nanomagnetic systems , on glasses , and on mesoscopic conductors , as well as theoretical debate about the mechanisms of decoherence in nature, clearly require a more general point of view. In fact we shall see that spin baths behave very differently from oscillator baths. For example, the oft-discussed connection between dissipation and decoherence which exists for oscillator bath environments is absent here- one can even have decoherence with no dissipation at all , because of the quantum phase associated with the spin bath dynamics.
Although the existence of a quantum phase associated with spin is obvious (it is a quantum variable), it was not until Haldane and Berry discussed its topological properties that physicists realised its practical importance, in terms of the path traced out on the spin sphere. Just as in the usual Aharonov-Bohm effect, the ”flux” enclosed by a path (given here by $`\omega S`$, where $`S`$ is the spin and $`\omega `$ is the enclosed solid angle on the unit spin sphere) is equal to a dynamical phase- but now these are both in spin space, not in real space. These ideas (and related experiments) were the centre of enormous interest in the mid-late 1980’s, in almost all fields of physics (and were extensively reviewed then ).
In this article we will be interested in the spin phase of the environment. We stress here that the environmental ”spin” variables may not necessarily refer to real spins (they can describe defects, or other such ”2 level systems”), but they will still have an associated dynamic topological phase, which can be described by an effective spin bath variable. The environmental spin phase interacts with the phase of the central system, causing phase decoherence in its dynamics . From the point of view of measurement theory, this environmental phase decoherence comes from a ”phase measurement” made by the spin environment , in a kind of ”inverse Stern-Gerlach” setup (where the spins, instead of being measured, are themselves doing the measuring!). Such phase decoherence also arises from oscillator baths, in a rather different (and much less effective) way . In fact, the relevant phases involve both an adiabatic ”Berry” term and a second term coming from transitions between different bath states (section 3.A). There are also other decoherence mechanisms associated with the bath spins, coming from both the temporal fluctuations of the bias on the central system caused by the spin bath (section 3.D), and from the precession of the spins in the spin bath (with their associated phase) in between transitions of the central system (section 3.C). Thus the question of how the bath spin dynamics influences the central system is not simply a question of looking at Berry phases.
Practical application of the theory to, eg., SQUIDs, or nanomagnets, or ”qubits” (section 5), must include all mechanisms properly (section 3.E). The tactic adopted in this article is to focus on a ”Central Spin” model (sections 3 and 4), in which the role of each term is exposed rather clearly; after this one sees how things generalise to other models. This model is directly relevant to qubits, and to the observation of mesoscopic or ”macroscopic” quantum coherence- indeed we maintain that any practical design of such devices must involve the elimination, by one means or another, of all decoherence mechanisms from the relevant spin bath .
We begin the article (section 2) by showing how both oscillator and spin bath models arise as the low-energy truncated versions of higher energy Hamiltonians. We give several examples, both magnetic and superconducting, to illustrate this. Then, in section 3 we explain how, mathematically, one averages over the spin bath variables to find the behaviour of the central system. This is done pedagogically- we use the example of the Central Spin model (and compare it with the spin-boson model). Various simple limits are introduced, and solved, before the general technique is given at the end of section 3. In Section 4 we give some results for the dynamics of the Central Spin model, in various regimes, to show how the physics is influenced by the spin bath; and we also show how the system reduces to the spin-boson system in the weak-coupling limit. For those readers interested in the mathematical details, these are sketched in 2 Appendices.
Finally, in section 5 we return to physical applications, particularly to quantum magnetic systems and superconductors. We then discuss decoherence, and show how this should persist even in the $`T0`$ limit. We finish by discussing the important application of the Central Spin model to decoherence in qubits and in quantum computation.
## II The Low-Energy Effective Hamiltonian
In studying the low-energy dynamics of a central quantum system coupled to an environment, we begin by “truncating out” the unwanted high-energy physics, to produce a low-energy effective Hamiltonian. This is of course a quite general technique in physics, and one way to approach it is illustrated in Fig.1. Typically one has a reasonably accurately known “high-energy” or “bare” Hamiltonian (or Lagrangian) for a quantum system, valid below some “ultraviolet” upper energy cut-off energy $`E_c`$, and having the form
$$\stackrel{~}{H}_{\text{Bare}}(E_c)=\stackrel{~}{H}_o(\stackrel{~}{P},\stackrel{~}{Q})+\stackrel{~}{H}_{\text{int}}(\stackrel{~}{P},\stackrel{~}{Q};\stackrel{~}{p},\stackrel{~}{q})+\stackrel{~}{H}_{\text{env}}(\stackrel{~}{p},\stackrel{~}{q})(E<E_c),$$
(1)
where $`\stackrel{~}{Q}`$ is an $`\stackrel{~}{M}`$-dimensional coordinate describing that part of the system we are interested in (with $`\stackrel{~}{P}`$ the corresponding conjugate momentum), and $`(\stackrel{~}{p},\stackrel{~}{q})`$ are $`\stackrel{~}{N}`$-dimensional coordinates describing all other degrees of freedom which may couple to $`(\stackrel{~}{P},\stackrel{~}{Q})`$. Conventionally one refers to $`(\stackrel{~}{p},\stackrel{~}{q})`$ as environmental coordinates. $`\stackrel{~}{H}_{\text{Bare}}`$ is of course a low-energy form of some other even higher-energy Hamiltonian, in a chain extending ultimately back to quarks, leptons, and perhaps strings.
If, however, one is only interested in physics below a much lower energy scale $`\mathrm{\Omega }_o`$, then the question is - can we find a new effective Hamiltonian, of form
$$H_{\text{eff}}(\mathrm{\Omega }_o)=H_o(P,Q)+H_{\text{int}}(P,Q;p,q)+H_{\text{env}}(p,q)(E<\mathrm{\Omega }_o),$$
(2)
in the truncated Hilbert space of energies below $`\mathrm{\Omega }_o`$? In this $`H_{\text{eff}}`$, $`P`$ and $`Q`$ are generalised $`m`$-dimensional coordinates of interest, and $`p,q`$ are $`N`$-dimensional environmental coordinates coupled to them. Since we have truncated the total Hilbert space, we have in general that $`M<\stackrel{~}{M}`$ and $`N<\stackrel{~}{N}`$.
Why do we make this truncation (after all, its inevitable effect will be to generate new couplings between the low-energy modes)? Essentially because in many cases the truncation pushes the new $`H_{\text{eff}}`$ towards some low-energy “fixed point” Hamiltonian; and many different physical systems may flow to the same fixed point. This allows us to speak of “universality classes” of quantum environment, and of a small number of ”canonical” effective Hamiltonians. All physical systems in the same universality class will be described by the same form for $`H_{\text{eff}}`$, albeit with different values for the couplings. As one varies the UV cut-off $`\mathrm{\Omega }_o`$, the couplings change and any given system moves in the ”coupling” or ”effective Hamiltonian” space; but they all move towards the same fixed point (or fixed line, in some cases) as $`\mathrm{\Omega }_0`$ is reduced. The various coupling terms in $`H_{eff}(\mathrm{\Omega }_o)`$, simply parametrise the path it takes as $`\mathrm{\Omega }_o0`$ (Fig.1).
Explicit derivations of $`H_{\text{eff}}`$ for particular systems are lengthy; see eg. for general discussions, ref. for comparision of spin and oscillator bath systems, and for specific examples. In this article we will go directly to the canonical models, giving some examples of each so that readers can see how the high-energy Hamiltonians are related to the models for some real systems. To warm up we recall the basic structure of the oscillator bath effective Hamiltonians, and then move on immediately to discuss various canonical models involving spin baths.
### A Oscillator Bath models
For models in the general ”universality class” of oscillator bath environments, $`H_{\text{eff}}`$ takes the form :
$`H_{\text{eff}}(\mathrm{\Omega }_o)=H_o(P,Q)+{\displaystyle \underset{q=1}{\overset{N}{}}}\left[F_q(P,Q)x_q+G_q(P,Q)p_q\right]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{q=1}{\overset{N}{}}}\left({\displaystyle \frac{p_q^2}{m_q}}+m_q\omega _q^2x_q^2\right),`$ (3)
where the generalised bath coordinates $`(q_k,p_k)`$ are now oscillator displacements $`x_q`$ and momenta $`p_q`$; these describe delocalised modes. The couplings $`F_q(P,Q)`$ and $`G_q(P,Q)`$ are $`O(N^{1/2})`$, so that in the ”thermodynamic limit” $`N1`$, appropriate to a macroscopic environment of delocalised oscillators, these couplings are small (a number of studies have also shown how higher-order couplings can be absorbed into linear but $`T`$-dependent couplings ). A special case of (3) is the Feynman-Vernon bilinear coupling form :
$`H_{\text{eff}}(\mathrm{\Omega }_o)=H_o(P,Q)+{\displaystyle \underset{q=1}{\overset{N}{}}}c_qx_qQ+{\displaystyle \frac{1}{2}}{\displaystyle \underset{q=1}{\overset{N}{}}}\left({\displaystyle \frac{p_q^2}{m_q}}+m_q\omega _q^2x_q^2\right),`$ (4)
where the couplings $`c_qO(N^{1/2})`$ as well. In recent years great attention has been given to problems where $`H_o(P,Q)`$ describes a tunneling system (the ”Caldeira-Leggett” model ); there have also been extensive studies of an oscillator coupled to oscillators , of free particles coupled to oscillators and of band particles coupled to oscillators .
Suppose now the potential $`V(Q)`$ has a 2-well form, with small oscillation frequencies $`\mathrm{\Omega }_o`$, and a ”bias” energy difference between the two minima $`<\mathrm{\Omega }_o`$. Then for energies $`<\mathrm{\Omega }_o`$, one further truncates to the celebrated “spin-boson” model :
$`H_{\text{SB}}(\mathrm{\Omega }_o)=\mathrm{\Delta }(\mathrm{\Phi }_o)\widehat{\tau }_x+\xi _H\widehat{\tau }_z+{\displaystyle \underset{q=1}{\overset{N}{}}}[c_q^{}\widehat{\tau }_z+(c_q^{}\widehat{\tau }_{}+h.c.)]x_q+{\displaystyle \frac{1}{2}}{\displaystyle \underset{q=1}{\overset{N}{}}}({\displaystyle \frac{p_q^2}{m_q}}+m_q\omega _q^2x_q^2),`$ (5)
where the two-level central system (with tunneling amplitude $`\mathrm{\Delta }(\mathrm{\Phi }_o)`$ and longitudinal bias $`\xi _H`$) is described by the Pauli matrix vector $`\widehat{\stackrel{}{\tau }}`$, coupled to background oscillators having energies $`\omega _q<\mathrm{\Omega }_o`$. We have introduced a topological phase $`\mathrm{\Phi }_o`$ for the central system, which depends in general on an external field; the simplest and best-known example is a form $`\mathrm{\Delta }(\mathrm{\Phi }_o)=2\mathrm{\Delta }_o\mathrm{cos}\mathrm{\Phi }_o`$, arising from the interference between 2 paths of amplitude $`\mathrm{\Delta }_oe^{\pm i\mathrm{\Phi }_o}`$ in the motion of the central system. This kind of ”Aharonov-Bohm” interference is well-known in superconductors (where the phase is just the superconducting order parameter phase), and in magnets (where it is the topological spin phase ). One can have a more complicated dependence of $`\mathrm{\Delta }(\mathrm{\Phi }_o)`$ on $`\mathrm{\Phi }_o`$ (eg., using multiple SQUIDs, or a nanomagnet with $`m`$-fold rotation symmetry), but we will stick with the simple $`2\mathrm{cos}\mathrm{\Phi }_o`$ dependence in this article.
For consistency we must assume $`\xi _H<\mathrm{\Omega }_o`$, otherwise higher levels will be involved. Typically $`c_q^{}`$ is dropped, because its effects are down on those of $`c_q^{}`$ by a factor $`(\mathrm{\Delta }_o/\mathrm{\Omega }_o)^2`$ in tunneling rates; but sometimes $`c_q^{}=0`$ (for reasons of symmetry), and then $`c_q^{}`$ must be retained. The tunneling amplitude $`\mathrm{\Delta }_o\mathrm{\Omega }_oe^{A_o}`$, where $`A_o`$ is the tunneling action.
The fame of the spin-boson model partly arises because many well-known problems in condensed matter physics can be mapped to it- this is a good example of the “universality” mentioned above. Because the effect of each oscillator on the central system (and vice-versa) is very small, it may be entirely incorporated in second-order perturbation theory (ie., to order $`(F_q^2/\omega _q),(G_q^2/\omega _q)`$) for the general form (3), or to order $`(c_q^2/\omega _q)`$ for the bilinear forms (4),(5). In the latter case this immediately encapsulates all environmental effects in the spectral function
$$J_\alpha (\omega )=\frac{\pi }{2}\underset{q=1}{\overset{N}{}}\frac{|c_q^\alpha |^2}{m_q\omega _q}\delta (\omega \omega _q),$$
(6)
where $`\alpha =,`$. In general $`J_\alpha (\omega )`$ also depends on $`T`$, even in the low-$`\omega `$ limit . The case where $`J_\alpha (\omega )\omega `$ is referred to as ”Ohmic” . Because the $`c_q^\alpha N^{1/2}`$, the $`J_\alpha (\omega )`$ are independent of $`N`$ and have the usual ”response function” form.
### B Two examples of Spin-Boson systems
We give just 2 examples here of how a spin-boson model can arise, in the description of mesoscopic systems at low energies. Both truncations ignore the presence of spin bath modes (for which see sections 2.E and 2.F, where we return to these 2 examples).
(i) Nanomagnet coupled to phonons or electrons: The electronic spin dynamics of nanomagnets are often described by a ”giant spin” Hamiltonian $`H_o(\stackrel{}{S})`$, describing a quantum rotator with spin quantum number $`S=|\stackrel{}{S}|1`$. This model assumes the individual electronic moments are locked together by strong exchange interactions $`J_{ij}`$ into a monodomain giant spin, with $`\stackrel{}{S}=_j\stackrel{}{s}_j`$ (summed over local moment sites). This only only works below a UV cut-off energy $`E_c`$ considerably less than $`J_{ij}`$. However we are interested in the quantum dynamics for energies $`<\mathrm{\Omega }_o`$, where $`\mathrm{\Omega }_o`$ is controlled mostly by the single-ion magnetic anisotropy; in real nanomagnets $`\mathrm{\Omega }_o0.110K`$.
Any real nanomagnet has couplngs to a spin bath of nuclear and paramagnetic spins, and to oscillator baths of phonons and electrons, which we now describe.
The ”high-energy”coupling between phonons and $`\stackrel{}{S}`$ is described by terms like :
$$_2^\varphi \mathrm{\Omega }_oU(\widehat{\stackrel{}{S}})\left(\frac{m_e}{M_a}\right)^{1/4}\underset{\stackrel{}{q}}{}\left(\frac{\omega _q}{\mathrm{\Theta }_D}\right)^{1/2}[b_\stackrel{}{q}+b_\stackrel{}{q}^{}],$$
(7)
where m<sub>e</sub> is the electron mass, $`M_a`$ the mass of the molecule, and $`\mathrm{\Theta }_Dc_sa^1`$ is the Debye temperature (with $`a`$ the relevant lattice spacing, and $`c_s`$ the sound velocity). The interaction $`U(\stackrel{}{S})S`$ and dimensionless; a typical example is the non-diagonal term $`(\widehat{S}_x\widehat{S}_z/S)`$, which causes phonon emission or absorption with a change $`\pm 1`$ in $`S_z`$ (since $`S_x=\frac{1}{2}(S_++S_{}))`$. One also has diagonal terms in which $`S_zS_x`$ is replaced by, eg. $`S_z^2`$; and there are also higher couplings to, eg, pairs of phonons.
Truncation to the ”quantum regime” then gives the spin-boson model (5), with a dominant non-diagonal coupling $`c_q^{}S\mathrm{\Omega }_o|\stackrel{}{q}|^{1/2}`$, coming from terms like (7). In the absence of external fields in Hamiltonian $`H_o^{eff}(\stackrel{}{S})`$, the diagonal coupling $`c_q^{}`$ is actually zero (because of time-reversal symmetry). The Caldeira-Leggett spectral function for the system has the form $`J_{}(\omega )S^2(\mathrm{\Omega }_o^2/\rho c_s^5)\omega ^3`$ where $`B_{}(S^2\mathrm{\Omega }_o^2/\mathrm{\Theta }_D)`$; here $`\rho `$ is the density of the medium supporting Debye phonons, and $`\mathrm{\Theta }_D^4\rho c_s^5`$.
If now we couple electrons to the giant spin, it is the diagonal coupling which dominates . The electronic coupling to $`\stackrel{}{S}`$ depends on the type of magnetism. Some details have been worked out for Kondo interactions with conduction electrons- the coupling to $`\stackrel{}{S}`$ is
$$H_{int}^{GK}=\frac{1}{2}\overline{J}\widehat{\stackrel{}{S}}.\widehat{\stackrel{}{\sigma }}^{\alpha \beta }\underset{\stackrel{}{k}\stackrel{}{q}}{}F_qc_{\stackrel{}{k}+\stackrel{}{q}_i\alpha }^+c_{\stackrel{}{k}_i\beta }$$
(8)
where $`\overline{J}`$ is the mean value of the Kondo couplings to each individual electronic spin in the nanomagnet, and $`\stackrel{}{S}F_q=(d^3r/V_o)\stackrel{}{s}(\stackrel{}{r})e^{i\stackrel{}{q}.\stackrel{}{r}}`$ is a ”form factor” integrating the localised electron spin density $`\stackrel{}{s}(\stackrel{}{r})`$ over the nanomagnetic volume $`V_o`$. At low energies the corresponding spin-boson model has an ”Ohmic” diagonal spectral function $`J_{}(\omega )=\pi \alpha _\kappa \omega `$. The size of $`\alpha _\kappa `$ depends on how the electrons permeate the nanomagnet ; if they permeate freely, $`\alpha _\kappa g^2S^{4/3}`$, where $`g=\overline{J}N(0)`$ is the mean dimensionless Kondo coupling, and $`N(0)`$ the Fermi surface density of states. Typically $`g0.1`$, so $`\alpha _\kappa `$ can be big.
(ii) RF SQUID (Flux coupled to electrons): We briefly recall one well-known application of the spin-boson model, to an RF SQUID coupled to both normal electrons (in shunts, etc.), and Bogoliubov quasiparticles . The flux $`\varphi `$ passing through a superconducting ring with a weak link moves in a multiwell potential, which can be adjusted so that the lowest 2 wells (each with small oscillation or ”Josephson plasma” frequency $`\mathrm{\Omega }_o2\pi [E_J/\pi C]^{1/2}/\varphi _o`$, where $`E_J`$ is the Josephson weak link coupling energy) are almost degenerate, and dominate the low-energy properties. The high-energy coupling between the flux and the electronic quasiparticles has the form
$$H_{int}=\{\mathrm{cos}(\pi \varphi /\varphi _o)\underset{q}{}t_qU_q^S(a_q+a_q^{})+i\mathrm{sin}(\pi \varphi /\varphi _o)\underset{q}{}t_qU_q^A(a_qa_q^{})\}$$
(9)
where $`q(\stackrel{}{k},\stackrel{}{k}^{})`$ labels oscillator states describing a quasiparticle pair $`|\stackrel{}{k}\stackrel{}{k}^{}`$ with energy $`\omega _q=E_k+E_k^{}`$, $`t_q`$ is the relevant junction tunneling matrix element, $`U_q^{S/A}`$ the symmetric/antisymmetric BCS coherence factor, and $`\varphi _o`$ is the flux quantum. Thus we have a coupling to both the momenta and coordinates of the oscillators, which can also be written as a coupling to 2 independent oscillator baths . The $`T`$-dependence of the coherence factors (coming from the BCS gap dependence) as well as the gap structure in their energy dependence, gives a complex structure in $`J(\omega ,T)`$. The reduction to the spin-boson model is now trivial , the minima in $`\varphi `$-space of the effective potential corresponding to the 2 eigenstates of $`\widehat{\tau }_z`$.
### C Spin Bath Environments
Now suppose we have a high-energy Hamiltonian of form (1), but where the environmental coordinates $`(\stackrel{~}{p},\stackrel{~}{q})`$ are a set of $`N`$ spin-$`1/2`$ variables $`\{\widehat{\stackrel{}{\sigma }}_k\}`$, (i.e., two-level systems); and we assume the interspin couplings to be weak. Then, instead of (3), we have
$$H=H_o(P,Q)+H_{\text{int}}(P,Q;\{\widehat{\stackrel{}{\sigma }}\})+H_{\text{env}}(\{\widehat{\stackrel{}{\sigma }}\});$$
(10)
$$H_{\text{int}}(P,Q;\{\widehat{\stackrel{}{\sigma }}\})=\underset{k=1}{\overset{N}{}}[F_k^{}(P,Q)\widehat{\sigma }_k^z+[F_k^{}(P,Q)\widehat{\sigma }_k^{}+h.c.]];$$
(11)
$$H_{\text{env}}(\{\widehat{\stackrel{}{\sigma }}\})=\underset{k=1}{\overset{N}{}}\stackrel{}{h}_k\widehat{\stackrel{}{\sigma }}_k+\underset{k=1}{\overset{N}{}}\underset{k^{}=1}{\overset{N}{}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta ,$$
(12)
for energy scales $`E<E_c`$. Thus we now have a central system coupled to a ”spin bath”, described by $`H_{\text{env}}(\{\widehat{\stackrel{}{\sigma }}\})`$ in (12). The couplings $`F_k^{}(P,Q)`$ and $`F_k^{}(P,Q)`$, between the central system and the bath spins, are usually much greater than the interspin couplings $`V_{kk^{}}^{\alpha \beta }`$; this means that the dynamics of each spin is largely ”slaved” to that of the central system.
Unlike oscillator baths (whose modes typically represent delocalised environmental degrees of freedom), the $`\{\widehat{\stackrel{}{\sigma }}_k\}`$ represent localised modes (whose weak spatial overlap explains why the $`V_{kk^{}}^{\alpha \beta }`$ are small). This fact underlies a crucial difference between oscillator and spin bath environments- the couplings $`F_k^{}(P,Q)`$ and $`F_k^{}(P,Q)`$ are independent of the number $`N`$ of bath spins. Thus the larger is $`N`$, the larger is the total effect of the spin bath on the central system- there is no strict thermodynamic limit in the system, and it is not meaningful to let $`N\mathrm{}`$. We emphasize also that we see no justification in general for spin bath models in which $`F_k^{},F_k^{}O(N^{1/2})`$, or even $`O(1/N)`$ (although one can certainly invent artificial models of this kind). Thus, if we add more localised environmental modes to our environment, it is clear that the different modes are approximately independent (as they will be if quasi-localised), so that their individual couplings to the central system will be hardly affected, ie., will depend only weakly on $`N`$.
There is nothing to stop generalisation of this model to include bath spins $`\{\stackrel{}{I}_k\}`$, with $`I_k=|\stackrel{}{I}_k|>1/2`$; the $`(2I_k+1)`$ states then represent the degrees of freedom of, eg., a defect, or a spin (again, localised). This introduces tensor (eg., quadrupolar) couplings to the bath spins , and thereby complicates the mathematics- but does not alter the basic physics. We will not discuss this here (for the relevant formalism, and its application to the $`Fe`$-8 molecular nanomagnet, see refs. ).
### D Particle moving through a spin bath
A particle moving through a spin bath is described by (10), in which $`P`$ and $`Q`$ describe the momentum and position of the moving particle. The diagonal term $`F_k^{}(P,Q)`$ is analogous to the ”position” oscillator coupling $`F_q(P,Q)`$ in (3), and likewise $`F_k^{}(P,Q)`$ to corresponds to $`G_q(P,Q)`$. However both forms can be altered by canonical transformation, corresponding to a rotation between the different coordinates. The most common problems involve a diagonal coordinate coupling $`F_k^{}(Q)`$ and a non-diagonal momentum coupling $`F_k^{}(P)`$. Then bath transitions (spin flips) are induced by the motion of the particle, whereas a stationary particle sees a ”potential” $`U(Q,\{\sigma _k^z(t)\})=_kF_k^{}(Q)\sigma _k^z(t)`$, in general time-dependent.
A nice mesoscopic example of this is a large magnetic soliton coupled to background spins . In many realistic cases the most important such coupling will be to paramagnetic impurities, but here we consider the simpler case of a hyperfine coupling to a set $`\{\widehat{\stackrel{}{\sigma }}\}`$ of $`N`$ spin-$`1/2`$ nuclear spins. In this case $`V_{kk^{}}^{\alpha \beta }`$ describes the extremely weak internuclear dipolar coupling; typically $`|V_{kk^{}}^{\alpha \beta }|10^7K`$; and $`\stackrel{}{h}_k`$ is any external field that might unfluence these nuclei.
The ”high-energy” Hamiltonian for such a wall is usually determined as an integral over the magnetisation density $`𝐌(𝐫)`$ and its spacetime gradients . ¿From this one eliminates the details of the wall profile altogether, to produce a “bare” Hamiltonian (ie., neglecting the environment) for the wall coordinate; in simple cases where the wall demagnetisation field keeps the wall almost flat, this gives
$$H_w=\frac{1}{2}M_w\dot{Q}^2V(Q)2S_w\mu _BM_0H_eQ$$
(13)
for a wall with centre of mass coordinate $`Q`$ and surface area $`S_w`$. The ”pinning potential” $`V(Q)=V_0\text{sech}^2(Q/\lambda _w)`$, provided the pinning centre is much smaller than the wall width $`\lambda _w`$. The term linear in $`Q`$ comes from an external magnetic field $`𝐇_e`$.
What now of the environment? In the literature there is extensive discussion of the effects of magnons (ie., spin waves), electrons , and phonons on the wall dynamics- these are all oscillator baths. However at low $`T`$ spin bath effects, coming from nuclear and paramagnetic spins, will completely dominate. Even in Ni (where only $`1\%`$ of the nuclei have spins, with a tiny hyperfine coupling $`\omega _0=28.35\text{MHz}1.4\text{mK}`$), all real samples have an important concentration of paramagnetic spins (caused by Oxygen in the sample) as well as many defects. In rare earths, the hyperfine coupling $`\omega _k110\text{GHz}`$ ($`0.050.5\text{K}`$), and hyperfine effects alone are quite massive. Thus we must modify $`H_w`$ above to
$$H=H_w+\underset{k=1}{\overset{N}{}}\omega _k^{\alpha \beta }s_k^\alpha I_k^\beta +\underset{k}{}\underset{k^{}}{}V_{kk^{}}^{\alpha \beta }I_k^\alpha I_k^{}^\beta $$
(14)
in which the electronic spins $`𝐬_k`$ couple locally to $`N`$ nuclear spins $`𝐈_k`$ at positions $`𝐫_k`$ ($`k=1,2,3,\mathrm{}N`$), via a hyperfine coupling $`\omega _k^{\alpha \beta }`$ (and also in general to paramagnetic spins). The internuclear coupling $`|V_{kk^{}}^{\alpha \beta }|1100\text{kHz}`$ ($`0.055\mu \text{K}`$), ie., $`\omega _k`$, but it gives the spin bath its own dynamics.
To write the Hamiltonian in the form (11), we write the continuum magnetisation $`𝐌(𝐫)=𝐌_o(𝐫)+𝐦(𝐫)`$, where $`𝐌_o(𝐫)`$ is the slowly-varying part descibing the wall profile and $`𝐦(𝐫)`$ describes fluctuations around this. Then we rotate the spin quantisation axis to be locally parallel to $`𝐌_o(𝐫)`$, and get
$$H=H_w+\underset{k=1}{\overset{N}{}}\frac{d^3r}{\gamma _g}\delta (𝐫𝐫_k)\left[\omega _k^{}M_z(𝐫)I_k^z+\omega _k^{}[m_x(𝐫)I_k^x+m_y(𝐫)I_k^y]\right]+\frac{1}{2}\underset{k}{}\underset{k^{}}{}V_{kk^{}}^{\alpha \beta }I_k^\alpha I_k^{}^\beta $$
(15)
displaying explicitly the longitudinal and transverse couplings. The ”particle” moves through a slowly fluctuating ”random walk” potential field $`U_{}(Q)`$ (coming from the sum over couplings $`\omega _k^{}`$ to randomly oriented spins). The transverse coupling (independent of $`Q`$ but not of $`P`$) causes ”spin flip” transitions in the spin bath when the wall moves, even if the bath is at $`T=0`$.
One may also discuss problems in superconductors and normal metals involving nuclear and paramagnetic spins, and other ”defects”, which can also be mapped to the same model of a particle moving through a spin bath (sections 2.F, 5.B).
### E The Central Spin Model
Is there a low-energy effective Hamiltonian, analogous to the spin-boson model, in which a ”central” 2-level system couples instead to a spin bath? The answer is yes, but the effective Hamiltonian does not look quite so simple as the spin-boson one. In the absence of any external field, the analogue of the spin-boson form in (5) for a spin bath is actually
$`H_{\text{CS}}(\mathrm{\Omega }_o)`$ $`=`$ $`\{2\stackrel{~}{\mathrm{\Delta }}\widehat{\tau }_{}\mathrm{cos}[\mathrm{\Phi }{\displaystyle \underset{k}{}}\stackrel{}{V}_k\widehat{\stackrel{}{\sigma }}_k]+H.c.\}`$ (16)
$`+`$ $`\widehat{\tau }_z{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _k^{}\stackrel{}{l}_k\widehat{\stackrel{}{\sigma }}_k+{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _k^{}\stackrel{}{m}_k\widehat{\stackrel{}{\sigma }}_k+{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \underset{k^{}=1}{\overset{N}{}}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta .`$ (17)
where $`\stackrel{}{\tau }`$ describes the central spin, and the $`\sigma _k`$ the spin bath degrees of freedom. This form is not the most general one- apart from dropping external field effects (for which see below) we have also restricted the central spin phase in a simple $`\mathrm{cos}\mathrm{\Phi }`$ form (cf. Introduction). As discussed below, both $`\stackrel{~}{\mathrm{\Delta }}`$ and $`\mathrm{\Phi }`$ incorporate spin bath renormalisation effects. The factor of 2 in front of $`\stackrel{~}{\mathrm{\Delta }}`$ is somewhat arbitrary (if the cosine is one, then the actual ”tunnel splitting” coming from (17), in the absence of spin bath effects, will be $`4\stackrel{~}{\mathrm{\Delta }}`$).
The basic form of (17) is actually fairly easy to understand. The extra phase in the first ”non-diagonal” term (adding to $`\mathrm{\Phi }`$) comes from the topological phase of the bath spins as they make transitions . There are also diagonal terms, and a weak interaction $`V_{kk^{}}`$ between the bath spins. This form assumes the diagonal couplings $`\omega _k^{},\omega _k^{}`$ are $``$ the UV cutoff $`\mathrm{\Omega }_o`$, and that $`V_{kk^{}}\omega _k^{},\omega _k^{}`$. The ratio $`V_{kk^{}}/\stackrel{~}{\mathrm{\Delta }}`$ is arbitrary. In sections 4 and 5 we shall see there is a weak coupling limit to this model, in which it reduces to a spin-boson system.
Before this model was derived in this general form and then solved, a number of special cases had already been looked at . In particular, Shimshoni and Gefen included only the diagonal terms $`\omega _k^{}`$ and $`\omega _k^{}`$, and examined the results in weak coupling in the presence of an AC field (see also ); they clearly recognised that the problem was different from an oscillator bath one.
Let us now discuss the different terms in (17), in the order they appear (cf. Fig. 2).
(i) Non-Diagonal terms: That an extra phase term should exist, coming from the spin bath, is obvious on general grounds (cf. Introduction). One can understand its algebraic form in the following way. Notice that $`H_{eff}`$ in (17) operates on both the central spin and the spin bath; and the effect of a single giant spin transition on $`\widehat{\stackrel{}{\sigma }}_k`$ can always be written as a transition between initial and final nuclear states, in the form $`\chi _k^{fin}=\widehat{T}_k\chi _k^{in}`$, where
$$\widehat{T}_k=e^{i{\scriptscriptstyle 𝑑\tau H_{int}(\tau )}}=e^{[\delta _k+i(\stackrel{}{V}_k\stackrel{}{\sigma }_k+\varphi _k)]}$$
(18)
The integral over $`H_{int}`$ is only defined once we know the trajectory of the central system during the tunneling event. It is in this sense that we say that the instanton has become an ”operator” in the space of the spin bath modes. Notice that in general the central spin phase $`\mathrm{\Phi }`$ and splitting $`\mathrm{\Delta }`$ are renormalised by the bath couplings :
$$\mathrm{\Phi }=\mathrm{\Phi }_o+\underset{k}{}\varphi _k,\stackrel{~}{\mathrm{\Delta }}=\mathrm{\Delta }_o\mathrm{exp}\{\underset{k}{}\delta _k\}.$$
(19)
The $`\{\varphi _k\}`$ are ”Berry phase” terms coming from the bath spin dynamics during a central spin flip, and the $`\{\delta _k\}`$ come from high-frequency modifications of the original high-energy potential (”barrier fluctuations”). One can show that both $`\varphi _k`$ and $`\delta _k`$ are $`O(\omega _k^2/\mathrm{\Omega }_o^2)`$, where $`\omega _k`$ is the larger of $`\omega _k^{},\omega _k^{}`$, and we will ignore these terms from now on.
Expanding out the cosine in (17) gives a series of terms like $`\widehat{\tau }_\pm \mathrm{\Gamma }_{\alpha \beta \gamma \mathrm{}}\widehat{\sigma }_{k_1}^\alpha \widehat{\sigma }_{k_2}^\beta \widehat{\sigma }_{k_3}^\gamma \mathrm{}`$ in which the instanton flip of the central spin couples simultaneously to many different bath spins- a single central spin transition can cause multiple transitions in the bath (Fig. 2). Later we will introduce a parameter $`\lambda `$ which measures the average number of bath spins flipping during each instanton.
(ii) Diagonal terms: These act between transitions of the central spin (Fig. 2), and are also easy to understand. Formally, one starts by considering the ”initial” and ”final” fields (ie., before and after a transition of $`\stackrel{}{\tau }`$) acting on $`\stackrel{}{\sigma }_k`$. Calling these fields $`\stackrel{}{\gamma }_k^{(1)}`$ and $`\stackrel{}{\gamma }_k^{(2)}`$ respectively (Fig. 3), we define the sum and the difference terms as
$`\omega _k^{}\stackrel{}{l}_k`$ $`=`$ $`(\stackrel{}{\gamma }_k^{(1)}\stackrel{}{\gamma }_k^{(2)})/2`$ (20)
$`\omega _k^{}\stackrel{}{m}_k`$ $`=`$ $`(\stackrel{}{\gamma }_k^{(1)}+\stackrel{}{\gamma }_k^{(2)})/2.`$ (21)
where $`\stackrel{}{l}_k`$ and $`\stackrel{}{m}_k`$ are unit vectors. Then the truncated diagonal interaction takes the form
$$H_{eff}^D=\underset{k=1}{\overset{N}{}}\left\{\stackrel{}{\gamma }_k^{(1)}\frac{1+\widehat{\tau }_z}{2}+\stackrel{}{\gamma }_k^{(2)}\frac{1\widehat{\tau }_z}{2}\right\}\widehat{\stackrel{}{\sigma }}_k\widehat{\tau }_z\underset{k=1}{\overset{N}{}}\omega _k^{}\stackrel{}{l}_k\widehat{\stackrel{}{\sigma }}_k+\underset{k=1}{\overset{N}{}}\omega _k^{}\stackrel{}{m}_k\widehat{\stackrel{}{\sigma }}_k,$$
(22)
i.e., one term which changes during a transition of the central system, and one which does not.
The longitudinal coupling $`\widehat{\tau }_z\omega _k^{}\sigma _k^z`$ determines the gross structure of the bath states in energy space: it also determines an ”internal bias field” $`ϵ(\{\sigma _k^z\})=\omega _k^{}\sigma _k^z`$ acting on $`\tau _z`$. The 2 levels of bath spin $`\stackrel{}{\sigma _k}`$ are split by energy $`\omega _k^{}`$, depending on whether $`\sigma _k^z`$ is parallel or antiparallel to $`\tau _z`$. The effect of this on the $`2^N`$ fold multiplet of bath states surrounding each central spin state is shown in Fig. 4. Suppose we classify these states by their ”polarisation group”; all bath states whose total longitudinal polarisation $`\sigma _k^z=M`$ are in polarisation group $`M`$. Since the $`\omega _k^{}`$ vary from one bath spin to another, states in polarisation group $`M`$ are spread over an energy range $`\stackrel{~}{\mathrm{\Gamma }}_M`$ ; and the entire manifold of states, comprising all polarisation groups, is spread over a larger energy range $`E_o`$. Let us define normalised densities of states $`G_M(ϵ)`$ and $`W(ϵ)`$ for these 2 distributions, so that
$$W(ϵ)=(1/2^N)\underset{M}{}C_N^{(N+M)/2}G_M(ϵ)$$
(23)
where $`C_n^m=n!/m!(nm)!`$. In almost any physical case one will have strongly overlapping polarisation groups, so that for all but very small values of $`N`$, or except in the extreme wings of the distributions, one has
$$G_M(ϵ)(2/\pi \stackrel{~}{\mathrm{\Gamma }}_M^2)^{1/2}e^{2ϵ^2/\stackrel{~}{\mathrm{\Gamma }}_M^2}.$$
(24)
$$W(ϵ)(2/\pi E_o^2)^{1/2}e^{2ϵ^2/E_o^2}.$$
(25)
The simplest case is where the $`\omega _k^{}`$ cluster around a single central value $`\omega _o`$ with variance
$$\delta \omega _o=\sqrt{\frac{1}{N}\underset{k}{}(\omega _k^{}\omega _o)^2}$$
(26)
For this case we define a parameter $`\mu =N^{1/2}\delta \omega _o/\omega _o`$, characterising the degree of polarisation group overlap; overlap is complete when $`\mu >1`$. Then $`\stackrel{~}{\mathrm{\Gamma }}_M2N^{1/2}\delta \omega _o`$ and $`E_o=2N^{1/2}\omega _o`$ (so that $`\stackrel{~}{\mathrm{\Gamma }}_M/E_o\delta \omega _o/\omega _o=N^{1/2}\mu `$). In the extremely unlikely case where $`\mu 1`$, $`W(ϵ)`$ can no longer be treated as Gaussian- however there is an intrinsic lower limit to the linewidth of each polarisation group set by the interspin interaction $`V_{kk^{}}`$. This ”intrinsic linewidth” $`\mathrm{\Gamma }_oN^{1/2}V_o`$, where $`V_o`$ is a typical value of $`V_{kk^{}}`$; in any physically realistic case this is usually enough by itself to cause complete overlap of all groups (the essentially non-interacting case where $`\stackrel{~}{\mathrm{\Gamma }}=0`$, ie., $`\omega _k=\omega _o`$ for all bath spins, and $`V_{kk^{}}=0,\stackrel{}{V}_k=0,\omega _k^{}=0`$, was actually studied by Garg ).
The ”transverse” couplings $`\omega _k^{}`$ arise when the fields before/ after a transition, acting on the $`\{\stackrel{}{\sigma }_k\}`$, are not exactly parallel or antiparallel. This can happen in many ways, either because of external fields which couple to the bath spins, or because of a lack of symmetry in the underlying dynamics of the central system, or in its coupling to the bath modes. Thus they are non-zero in any realistic situation.
(iii) Internal Spin Bath dynamics: Finally, the interaction $`V_{kk^{}}`$ is usually so weak that it does not change under truncation. If this term is absent, the spin bath will have no ”intrinsic” dynamics, and remains inert between transitions of $`\stackrel{}{\tau }`$. Thus even if small, $`V_{kk^{}}`$ is important, since it allows the bath state to evolve during these intervals. The most important effect of $`V_{kk^{}}`$ is that it allows the longitudinal bias field $`ϵ(\{\sigma _k^z\})`$ to fluctuate in time, between and during transitions of $`\stackrel{}{\tau }`$. Notice, however, that with the Hamiltonian in the form (17), only fluctuations within the same polarisation group are allowed. In NMR language, only $`T_2`$ processes occur in the intrinsic dynamics of the spin bath- changes in $`M`$ can only occur via the interaction with the central spin. This will only be true at low $`T`$\- at higher $`T`$ longitudinal relaxation (ie., ”$`T_1`$ processes”, in NMR language) between different polarisation groups should be included in a realistic model. Such processes (which arise from the interaction of the spin bath modes with other environmental modes, or with thermally excited higher modes of the central system, above the UV cutoff energy $`\mathrm{\Omega }_o`$) are almost always very slow when $`kT\mathrm{\Omega }_o`$. On the other hand $`T_2`$ fluctuations will persist until $`kT\mathrm{\Gamma }_o`$; in the physical examples studied so far this means they persist down to $`\mu K`$ temperatures or below. We will return briefly to this very low $`T`$ regime at the end of the article (section 5.C).
External Field Effects: In general all of the parameters in (17) will depend on any external field $`𝐇_o`$, because it changes the high-energy dynamics of both central system and bath- however we can make low-field expansions and separate out the most important terms. Defining the ”Zeeman” coupling energies $`\mathrm{\Omega }_{H_o}`$ and $`\omega _k^{H_o}`$ of central and bath spins to this field, it is easy to see that under the conditions $`\mathrm{\Omega }_{H_o}/\mathrm{\Omega }_o<1`$ and $`\omega _k^{H_o}/\omega _k<1`$, the principal changes to (17) will be (i) the addition of an obvious longitudinal coupling $`\xi _{H_o}\widehat{\tau }_z`$ to the central spin, and (ii) Qthe changes $`\mathrm{\Phi }_o\mathrm{\Phi }_{H_o}`$, $`\stackrel{}{V}_k\stackrel{}{V}_k^{H_o}`$, and $`\stackrel{}{\omega }_k^{}\omega _k^{}\widehat{\stackrel{}{m}}_k\stackrel{}{\omega }_k^{}(𝐇_o)`$, where up to linear order in $`H_o`$ one has
$`\mathrm{\Phi }_{H_o}`$ $`=`$ $`\mathrm{\Phi }_o+\psi (H_o);\psi 2\pi \mathrm{\Omega }_{H_o}/\mathrm{\Omega }_o`$ (27)
$`\stackrel{}{V}_k^{H_o}`$ $`=`$ $`\stackrel{}{V}_k+\stackrel{}{v}_k(H_o);|\stackrel{}{v}_k(H_o)|/|\stackrel{}{V}_k|\omega _k^{H_o}/\omega _k`$ (28)
$`\stackrel{}{\omega }_k^{}(𝐇_o)`$ $`=`$ $`\stackrel{}{\omega }_k^{}+𝐝_k(𝐇_o);|𝐝_k(𝐇_o)|/\omega _k^{}\omega _k^{H_o}/\omega _k.`$ (29)
Thus even at low fields one has an important change in all the topological phases in the problem, and also to the transverse diagonal coupling (which itself arises from internal fields). In general $`𝐇_0`$ will also change the interspin bath couplings $`V_{kk^{}}`$; in a way which depends on the specific details of the problem. In the next 2 sections we see how this works for both magnetic and superconducting systems.
### F Nanomagnet coupled to nuclear and paramagnetic spins
If we start from the ”giant spin” model introduced above for a nanomagnet, then a simple isotropic contact hyperfine coupling to nuclear spins will lead to a Hamiltonian (for $`E<E_c`$) like:
$$H(\stackrel{}{S};\{\widehat{\stackrel{}{\sigma }}\})=H_o(\stackrel{}{S})+\frac{1}{S}\underset{k=1}{\overset{N}{}}\omega _k\stackrel{}{S}\widehat{\stackrel{}{\sigma }}_k+H_{\text{env}}(\{\widehat{\stackrel{}{\sigma }}\});$$
(30)
where $`H_o(\stackrel{}{S})`$ is the “giant spin” Hamiltonian, and $`H_{\text{env}}(\{\widehat{\stackrel{}{\sigma }}\})`$ is the same as in (12). The generalisation of this simple Hamiltonian to include dipolar hyperfine interactions, as well as to higher spin nuclei and to paramagnetic spins (with tensor and quadrupolar couplings) can be used if necessary . Here we will assume for simplicity that $`\stackrel{}{I}_k=I=\frac{1}{2}`$, and write $`\stackrel{}{I}_k\sigma _k`$, i.e. the nuclear spins will be described by spin-$`\frac{1}{2}`$ Pauli matrices. In fact in many cases even if $`I\frac{1}{2}`$, the low-energy nuclear spin dynamics is well described by a 2-level system.
The truncation of $`H(\stackrel{}{S};\{\widehat{\stackrel{}{\sigma }}\})`$ to a central spin Hamiltonian $`H_{\text{eff}}(\stackrel{}{\tau };\{\widehat{\stackrel{}{\sigma }}\})`$ has been discussed in several papers . As an example we quote the result for a simple easy axis-easy plane nanomagnet (for which $`H_o(\stackrel{}{S})=(1/S)[K_2^{}\widehat{S}_z^2+K_2^{}\widehat{S}_y^2]`$, and give it a physical interpretation. In this case, assuming $`\omega _k\mathrm{\Omega }_o`$ and also a weak external field $`𝐇_o`$, the effective Hamiltonian is :
$`H_{eff}(\mathrm{\Omega }_o)`$ $`=`$ $`\{2\mathrm{\Delta }_o\widehat{\tau }_{}\mathrm{cos}[\pi Si{\displaystyle \underset{k}{}}\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k\beta _o𝐧_o.𝐇_o]+H.c.\}`$ (31)
$`+`$ $`\widehat{\tau }_z\left[\xi _H+{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _k^{}\widehat{\sigma }_k^z\right]+{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \underset{k^{}=1}{\overset{N}{}}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta .`$ (32)
ie., a special case of the general form (17), with the parameters $`\xi _H=g\mu _BS_zH_o^z`$, $`\stackrel{}{l}_k=\widehat{\stackrel{}{z}}`$, $`\omega _k^{}=\omega _k`$, and $`\omega _k^{}=0`$. The vectors $`\alpha _k\stackrel{}{n}_k`$ and $`\beta _o𝐧_o`$ for this easy axis-easy plane case turn out to be (again assuming $`\omega _k\mathrm{\Omega }_o`$, and small $`𝐇_o`$):
$$\alpha _k\stackrel{}{n}_k=\frac{\pi \omega _k}{\mathrm{\Omega }_o}(\widehat{\stackrel{}{x}},i\sqrt{K_{}/K_{}}\widehat{\stackrel{}{y}});\beta _o𝐧_o=\frac{\pi g\mu _BS}{\mathrm{\Omega }_o}(\widehat{\stackrel{}{x}},i\sqrt{K_{}/K_{}}\widehat{\stackrel{}{y}})$$
(33)
In this example there are 2 tunneling trajectories (clockwise and counterclockwise in the easy plane), giving a result $`e^{\pm i\alpha _k\stackrel{}{n}_k\stackrel{}{\sigma }_k}`$ for the ”transfer matrix” $`\widehat{T}_k`$ (in zero applied field). The resulting vector $`\alpha _k\stackrel{}{n}_k`$ is the ”average hyperfine field” acting on $`\stackrel{}{\sigma }_k`$ during the tunneling event. To understand its orientation (and why it is complex) we note that the nuclear spin itself exercises a torque on $`\stackrel{}{S}`$ while it is tunneling, and this pushes $`\stackrel{}{S}`$ away from the easy plane. Consequently (a) the average field acting on $`\stackrel{}{\sigma }_k`$ has a component out of the easy plane, in the $`y`$-direction, and (b) $`\stackrel{}{S}`$ no longer moves exactly along the easy-plane path, while tunneling, that it would in the absence of $`\stackrel{}{\sigma }_k`$ (and so its action increases, via the imaginary part of $`\alpha _k\stackrel{}{n}_k`$).
The contribution $`\beta _o𝐧_o.𝐇_o`$ to the topological phase comes from the area swept out by the giant spin on the spin sphere (cf. Introduction), which changes as the field $`𝐇_o`$ changes; it is essentially an ”Aharonov-Bohm” contribution to this phase from the external field , which leads to spectacular oscillations (Fig. 5) in the effective tunneling amplitude $`\stackrel{~}{\mathrm{\Delta }}2\mathrm{\Delta }_o\mathrm{cos}[\pi S+i\beta _o𝐧_o.𝐇_o]`$ for a field perpendicular to $`\stackrel{}{z}`$ . Very recently oscillations in the tunneling amplitude of $`Fe`$-8 magnetic molecular crystals were seen which are related to this , although the presence of both nuclear spins and dipolar fields seriously complicates their interpretation (see section 5.A).
We recall from section II.B that the giant spin model truncates for nanomagnets to a 2-level system for energies $`\mathrm{\Omega }_o110K`$. Contact hyperfine couplings are in the range $`1.330mK`$ (transition metals) or $`40500mK`$ (rare earths); on the other hand the internuclear couplings $`V_{kk^{}}10^810^5K`$. In the $`Fe`$-8 system just mentioned the hyperfine interactions are actually dominated by dipolar couplings between the 8 $`Fe^{+3}`$ ions and the 120 protons in the molecule; these couplings are in the range $`1100MHz`$ ($`0.055mK`$). When the hyperfine couplings are this weak we must also take into account the effect of external fields on the nuclear dynamics (section 5.A). The values of $`\mathrm{\Delta }`$ vary over a huge range, but typically $`\mathrm{\Delta }\omega _k`$ (in $`Fe`$-8, $`\mathrm{\Delta }10^7K`$).
### G SQUID coupled to nuclear and paramagnetic spins
We consider again the RF SQUID, but now concentrate on the coupling of the flux $`\mathrm{\Phi }`$ to the spin bath of nuclear and paramagnetic spins which are within a penetration depth of the surface of the superconductor. This example is very instructive in understanding the weak-coupling limit of the central spin model (the following discussion is based on refs. ).
Suppose to start with we consider a ”cubic geometry” , in which a cube of side $`L=1cm`$ has a hole of radius $`R=0.2cm`$ through it, with in addition a slit connecting the hole to the exterior, spanned by a cylindrical junction of length $`l=10^4cm`$ and diameter $`d=2\times 10^5cm`$. The magnetic field inside the hole corresponding to a half-flux quantum is $`B_o=(\pi \mathrm{}c/e\pi R^2)=2\times 10^6G`$, whereas the magnetic field in the junction is as high as $`B_jB_oL/d=10^1G`$.
There are both nuclear spins and paramagnetic impurities in the spin bath. Consider first the nuclear spins; assuming all nuclei have spins, we find that in the bulk of the ring, within a penetration depth of the surface, there are $`N_r2\pi R\lambda _LL\times 10^{23}5\times 10^{17}`$ nuclear spins coupling to the ring current; and in the junction itself, a number more like $`N_j=(\pi dl\lambda _L)\times 10^{23}3\times 10^9`$. Thus each ring nuclear spin couples to the SQUID with a diagonal coupling $`\omega _r^{}\mu _nB_o2\times 10^{13}K`$; on the other hand for junction spins this coupling is $`\omega _j^{}=\mu _nB_j10^8K`$. Notice we have ignored any coupling to substrate spins (assume, eg., the ring is in superfluid He-4!), which might have a much larger coupling to the current.
At any temperature such that $`kT\omega _r^{},\omega _j^{}`$, the typical polarisation of these spin baths will be $`\sqrt{N_r},\sqrt{N_j}`$ respectively, giving a distribution of longitudinal bias energies with typical values $`E_o^j\omega _j^{}\sqrt{N_j}5\times 10^4K`$, and $`E_o^r\omega _r^{}\sqrt{N_r}10^4K`$ acting on the tunneling flux coordinate $`\mathrm{\Phi }`$. If we now add paramagnetic impurities to the ring, with concentration $`n_{pm}`$, and coupling $`\omega _{pm}^{}2\times 10^3\omega _r4\times 10^{10}K`$ to the current, this gives a typical longitudinal bias energy $`E_o^{pm}n_{pm}^{1/2}\times 0.2K`$. This longitudinal term is obviously bigger than the nuclear contribution, unless the superconductor is very pure indeed!
However, there is another much stronger transverse term, because each spin feels the dipolar fields from the other spins. This field is $`1G`$ (much higher near to the paramagnetic spins), and for the nuclear spins has an associated energy $`\omega _k^{}10^7K`$, which is $`\omega _r^{},\omega _j^{}`$. Physically, when the SQUID flips, the field on each nuclear spin hardly changes its direction, being dominated by the more slowly varying (but much stronger) nuclear dipolar field. For the paramagnetic spins the analogous coupling $`\omega _{pm}^{}`$ is $`>10^3`$ times larger, which in the absence of nuclear fields would give an inter-paramagnetic ”flip-flop” rate $`V_{kk^{}}^{pm}10^9n_{pm}`$ Hz, except that in pure samples these processes will themselves be blocked by the local dipolar coupling between the impurity and nearby nuclear spins (of strength $`\omega _{pm}^{}10^4K`$); this will happen once $`n_{pm}10^3`$.
We can write down an effective Hamiltonian for this system, valid over timescales considerably greater than $`\mathrm{\Delta }^1`$; we will use this later to analyse the effect of the spins on the SQUID dynamics (section 5.B). We will assume that $`\mathrm{\Delta }V_{kk^{}}`$ (the only case of experimental interest); then we can treat the internuclear dipolar fields as slowly-varying in time. The effective Hamiltonian can then be derived, to give :
$$H\text{eff}(\mathrm{\Omega }_o)=\{\mathrm{\Delta }_o(\mathrm{\Phi }_o)\widehat{\tau }_+e^{i_k\stackrel{}{\alpha }_k.\widehat{\stackrel{}{\sigma }}_k}+H.c.\}+\xi _H\widehat{\tau }_z+\underset{k=1}{\overset{N}{}}[\widehat{\tau }_z\omega _k^{}\widehat{\sigma }_k^z+\omega _k^{}\widehat{\sigma }_k^x]$$
(34)
where $`\omega _k^{}=\omega _r`$, $`\omega _j`$ or $`\omega _{pm}^{}`$, depending on the spin, and $`\omega _k^{}`$ has just been discussed; and microscopic analysis shows that $`|\stackrel{}{\alpha }_k|\omega _k^{}/\mathrm{\Omega }_o`$, where $`\mathrm{\Omega }_o`$ is the Josephson plasma frequency (section 2.B). Notice that $`\omega _k^{}/\omega _k^{}1`$, which is the opposite limit considered to that for the giant spin! Notice further that these couplings are far less than $`\mathrm{\Delta }_o`$ ($`E_o^r,E_o^j,E_o^{pm}\mathrm{\Delta }_o`$ only because there are so many spins involved). Thus the spin bath is no longer ”slaved” to the central system. In section 4 we see how this allows a mapping to an oscillator bath, coupled to $`\stackrel{}{\tau }`$ (ie., a spin-boson model).
Finally we note that the external field also acts on the nuclear and paramagnetic spin dynamics, via the Zeeman coupling (which simply adds to $`\omega _k^{}`$). This can help to suppress decoherence effects, by freezing out the spin bath dynamics- for more details see ref..
### H General Canonical Models
We now recall our assertion that almost any mesoscopic ”central” system, coupled to its environment, may be described at low energies by (2), with the environment being written as a sum of an oscillator bath term (3) and a spin bath term ((12), or a higher spin generalisation). The simplest example is of a single central spin coupled to both oscillator and spin baths. Such a model seems forbidding but in fact a fairly complete analysis has been given of its dynamics \- we recall some of the results at the end of section 4. One can also consider a much more complicated model in which a macroscopic array of central spins $`\{\stackrel{}{\tau }_j\}`$, at positions $`\{𝐫_j\}`$, couples to both oscillator and spin baths. The effective Hamiltonian is then an obvious generalisation of what has gone before:
$`H_{\text{CS}}(\mathrm{\Omega }_o)`$ $`=`$ $`{\displaystyle \underset{j}{}}\{\mathrm{\Delta }_j\widehat{\tau }_j^{}\mathrm{cos}[\mathrm{\Phi }_ji{\displaystyle \underset{k}{}}\alpha _{jk}\stackrel{}{n}_{jk}\widehat{\stackrel{}{\sigma }}_k]+H.c.\}+{\displaystyle \underset{i<j}{}}V(𝐫_i𝐫_j)\widehat{\tau }_i^z\widehat{\tau }_j^z`$ (35)
$`+`$ $`{\displaystyle \underset{j}{}}\left\{\widehat{\tau }_j^z{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _{jk}^{}\stackrel{}{l}_{jk}\widehat{\stackrel{}{\sigma }}_k+{\displaystyle \underset{k=1}{\overset{N}{}}}\omega _{jk}^{}\stackrel{}{m}_{jk}\widehat{\stackrel{}{\sigma }}_k\right\}+{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \underset{k^{}=1}{\overset{N}{}}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta `$ (36)
$`+`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{q}{}}[c_{jq}^{}\widehat{\tau }_j^z+(c_{jq}^{}\widehat{\tau }_j^{}+h.c.)]x_q+{\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}({\displaystyle \frac{p_q^2}{m_q}}+m_q\omega _q^2x_q^2),`$ (37)
where $`V(𝐫_i𝐫_j)\widehat{\tau }_i^z\widehat{\tau }_j^z`$ is a ”high-energy” diagonal coupling between the various ”central spin” systems. If we throw away the spin bath we get a set of 2-level systems coupling to an oscillator bath, of which the simplest example is the ”PISCES” model (in which there are only two 2-level systems ).
Such models seem impossibly complicated, but actually one can solve for their dynamics in many important regimes! What is crucial is the separation of the 2 baths. Often (as with nuclear spins) their mutual interaction is very weak (and can be parametrised by a time $`T_1(T)`$ which may be very long at low $`T`$); in this case this separation is a good one. If there are certain spin bath modes that interact strongly with the oscillators, then typically we can simply absorb these modes into an ”augmented” oscillator bath by a canonical transformation. An obvious example arises with electronic spins in a metallic host (the Kondo or Kondo lattice problems); one rewrites the bath to include the ”Kondo resonance” in the oscillator bath spectrum.
A proof that one may do this in all cases seems rather difficult- in any case the usefulness of these models tends to be established by their application. Models like (37) describe mesoscopic systems like coupled SQUIDs or coupled nanomagnets , Quantum Spin Glasses and low-$`T`$ dipolar glasses , as well as coupled anisotropic coupled Kondo spins and Kondo lattices, coupled nuclear spin systems , superconducting arrays, or coupled defects in solids. They are also useful for analysing purely theoretical questions about relaxation, dissipation, decoherence and quantum measurements in quantum systems- many questions remain unanswered, having only been studied thoroughly in restricted models such as the spin-boson model or the PISCES model . We return to experimental and theoretical applications in section 5.
## III Averaging over the Spin Bath
To extract useful information from the low-energy canonical models, we must calculate their dynamical properties. Since we are typically not interested in the environment (one usually has little control over it), one performs a statistical average over the environment. This procedure is fraught with danger, because of ”memory” effects in the environment, and because assumptions such as ”self-averaging” in the environmental correlation functions may not strictly be valid.
In this section we show how the spin bath may be ”integrated out” by means of 4 different statistical averages, each involving an integration over a particular variable. The end result is a description of the time evolution of the ”reduced” density matrix for the central system- provided we can ignore memory effects in the environment. The starting point is no different from that involved in functional averaging over oscillators ; both begin with a path integral form for the propagator of the reduced central system density matrix, written as
$$K(1,2)=_{Q_1}^{Q_2}𝑑Q_{Q_1^{}}^{Q_2^{}}𝑑Q^{}e^{i/\mathrm{}(S_o[Q]S_o[Q^{}])}[Q,Q^{}],$$
(38)
where $`S_o[Q]`$ is the free central system action, and $`[Q,Q^{}]`$ is the famous “influence functional” , defined in general by
$$[Q,Q^{}]=\underset{k}{}\widehat{U}_k(Q,t)\widehat{U}_k^{}(Q^{},t),$$
(39)
Here the unitary operator $`\widehat{U}_k(Q,t)`$ describes the evolution of the $`k`$-th environmental mode, given that the central system follows the path $`Q(t)`$ on its ”outward” voyage, and $`Q^{}(t)`$ on its ”return” voyage; and $`[Q,Q^{}]`$ acts as a weighting function, over different possible paths $`(Q(t),Q^{}(t^{}))`$.
For a central 2-level system, the paths $`Q(t),Q^{}(t)`$ are simple (recall Fig. 2):
$$Q_{(n)}(s)=1\underset{i=1}{\overset{2n}{}}\left[sgn(st_{2i1})+sgn(t_{2i}s)\right],$$
(40)
where $`sgn(x)`$ is the sign-function, and $`n`$ is the number of transitions of the central system, occuring at times $`t_1,t_2,\mathrm{},t_{2n}`$ (for definiteness we assume trajectories starting and ending in the same state, and use the convention that $`Q=\pm 1`$ corresponds to $`\tau _z=\pm 1`$). The goal is to find the central spin density matrix; in this article we give results for the ”return probability” $`P_{11}(t)`$ for the system to be in the same state $`|`$ at time $`t`$ as it was at $`t=0`$. Using (40) this can be written as an ”instanton expansion” over flips of the central spin (Appendix A):
$$P_{11}(t)=\underset{nm}{\overset{\mathrm{}}{}}(i\mathrm{\Delta }_o)^{2(n+m)}_0^t𝑑t_1\mathrm{}_{t_{2n1}}^t𝑑t_{2n}_0^t𝑑t_1^{}\mathrm{}_{t_{2m1}^{}}^t𝑑t_{2m}^{}[Q_{(n)},Q_{(m)^{}}]$$
(41)
Further simplification arises if the environmental modes are uncoupled- then $`[Q,Q^{}]`$ factorises, and we can write $`[Q,Q^{}]=\mathrm{exp}(i\mathrm{\Phi }[Q,Q^{}])=\mathrm{exp}(i_{k=1}^N\varphi _k[Q,Q^{}])`$, where the complex phase $`\varphi _k[Q,Q^{}]`$ contains both real (reactive), and imaginary (damping) contributions.
Now for an oscillator bath one simplifying feature is crucial, viz., the very weak coupling to each oscillator. This allows one to evaluate each $`\varphi _k[Q,Q^{}]`$ up to 2nd order only in these couplings, in terms of a spectral function for the unperturbed oscillator dynamics (compare $`J(\omega ,T)`$ in (6)). Even though the paths $`Q(t)`$ and $`Q^{}(t)`$ may be complicated, the calculation of $`[Q,Q^{}]`$ is often tractable .
However because the coupling to each spin bath mode is not necessarily weak, it will in general strongly alter their dynamics, often slaving them to the motion of the central system. Thus we cannot start from the unperturbed spin bath dynamics\- the problem is fundamentally non-perturbative in the $`\{\omega _k\}`$. However it is not intractable, because one can rather easily deal with the dominant longitudinal terms $`\{\omega _k^{}\}`$. The other terms can then be dealt with perturbatively (and sometimes even non-perturbatively). It is the separation of the effects of the various terms in the Hamiltonian which leads to not one, but 4 different averages. What is quite remarkable is that these averages can be evaluated analytically in most cases (section 4).
We begin by explaining the 4 different averaging integrals required for a general spin bath. This is done pedagogically, by solving for 4 different limiting cases of the central spin Hamiltonian (sections III.A-III.D), each of which requires only one of the 4 averaging integrals. Then the general procedure (combining the 4 averages) is given in III.E. One reason for going through these averages one by one is that each corresponds to a different physical mechanism of decoherence- we return to this in section 5 (note that more detailed results for the 3 limiting cases discussed in sections III.B-III.D are given in refs. ).
### A Phase averaging: Topological decoherence
Formally the case of pure topological decoherence applies to the following special case of $`H_{eff}`$, in which only non-diagonal terms are included:
$$H_{\text{eff}}^{top}=2\mathrm{\Delta }_o\{\widehat{\tau }_{}\mathrm{cos}[\mathrm{\Phi }_oi\underset{k=1}{\overset{N}{}}\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k]+H.c.\},$$
(42)
Since $`\omega _k^{}=\omega _k^{}=0`$ are zero, all the $`2^N`$ environmental states are degenerate, and there is no exchange of any energy between $`\stackrel{}{\tau }`$ and the $`\{\widehat{\sigma }_k\}`$. The only thing that is exchanged is phase; the phase $`\mathrm{\Phi }_o`$ of $`\stackrel{}{\tau }`$ becomes entangled with that of the $`\{\widehat{\sigma }_k\}`$, during the transitions of $`\stackrel{}{\tau }`$ between $``$ and $``$, so that the initial and final states of the spin environment are different. Physically (42) would arise if the original high-energy Hamiltonian contained only ”transverse” couplings (to $`\{\sigma _k^\pm \}`$), which only act while the central system is making transitions (in the case of a moving particle coupled to a spin bath, they would be ”velocity couplings”, only acting when the particle is moving).
We wish to determine $`P_{11}(t)`$ for this case. For pedagogical purposes let us begin by assuming that $`i\alpha _k`$ is real, ie., $`\alpha _k`$ is pure imaginary; then we have added a pure environmental phase term to the free 2-level Hamiltonian. Then, writing $`i\alpha _k\stackrel{~}{\alpha }_k`$, the formal solution to this problem can be written immediately as
$$P_{11}(t)=\frac{1}{2}\left\{1+\mathrm{cos}\left[4\mathrm{\Delta }_ot\mathrm{cos}\left(\mathrm{\Phi }_o+\underset{k=1}{\overset{N}{}}\stackrel{~}{\alpha }_k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k\right)\right]\right\},$$
(43)
where the brackets $``$ trace over the spin bath. By writing this as an instanton expansion over central spin flips (see App. A and refs. ), we transform it to a weighted integration over topological phase:
$`P_{11}(t)`$ $`=`$ $`{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}F(m){\displaystyle _0^{2\pi }}{\displaystyle \frac{d\phi }{2\pi }}e^{i2m(\mathrm{\Phi }\phi )}\left\{{\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2}}\mathrm{cos}(2\mathrm{\Delta }_o(\phi )t)\right\}`$ (44)
$`=`$ $`{\displaystyle \frac{1}{2}}\left\{1+{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}(1)^mF(m)e^{i2m\mathrm{\Phi }}J_{2m}(4\mathrm{\Delta }_ot)\right\},`$ (45)
where $`\mathrm{\Delta }_o(\phi )=2\stackrel{~}{\mathrm{\Delta }}_o\mathrm{cos}\phi `$ is a phase-dependent tunneling amplitude, $`J_{2m}(z)`$ is a Bessel function, and
$`F(m)={\displaystyle \underset{k=1}{\overset{N}{}}}e^{2im\stackrel{~}{\alpha }_k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k}={\displaystyle \underset{k=1}{\overset{N}{}}}\mathrm{cos}(2m\stackrel{~}{\alpha }_k).`$ (46)
For small $`\alpha _k`$ we may approximate the product in (46) as
$$F_\lambda (\nu )=e^{4\lambda \nu ^2};\lambda =\underset{k=1}{\overset{N}{}}\stackrel{~}{\alpha }_k^2/2.$$
(47)
Notice that $`\lambda `$ is just the mean number of environmental spins that are flipped, each time $`\tau `$ flips.
Clearly phase decoherence is important if $`\lambda >1`$, in which case $`F_\lambda (\nu )=\delta _{\nu ,0}+smallcorrections`$. Then, rather surprisingly, we get a universal form (shown in Fig. 6), in the intermediate coupling limit, for $`P_{11}(t)`$:
$$P_{11}(t)\frac{1}{2}\left[1+J_0(4\mathrm{\Delta }_ot)\right]\frac{d\phi }{2\pi }P_{11}^{(0)}(t,\mathrm{\Phi }=\phi )(intermediatecoupling)$$
(48)
with a phase integration over the free central spin propagator $`P_{11}^{(0)}(t,\mathrm{\Phi })`$ \[cf. eqtn. (A9)\]. Thus random phases arise because successive flips of $`\stackrel{}{\tau }`$ cause, in general, a different topological phase to be accumulated by the spin environment. In fact, the universal behaviour comes from complete phase phase randomisation , so that all possible phases contribute equally to the answer! The final form shows decaying oscillations , with an envelope $`t^{1/2}`$ at long times, which can also be understood by noting that the ”zero phase” trajectories contributing to $`P_{11}`$ constitute a fraction $`(2s)!/(2^ss!)^2s^{1/2}`$ of all possible trajectories, where $`s\mathrm{\Delta }_ot`$.
In the strong coupling limit of this model, where the bath spins rotate adiabatically with the the central spin, one has $`\alpha _k\pi /2`$, so that $`F(m)=(1)^m`$ and $`P_{11}(t)=\frac{1}{2}[1+\mathrm{cos}(4\mathrm{\Delta }_ot\mathrm{cos}\stackrel{~}{\mathrm{\Phi }})]`$, where $`\stackrel{~}{\mathrm{\Phi }}=\mathrm{\Phi }_o+N\pi /2`$, i.e., the Haldane/Kramers phase is now $`\stackrel{~}{\mathrm{\Phi }}`$, since the $`N`$ bath spins are forced to rotate with $`\stackrel{}{\tau }`$.
The results for complex $`\alpha _k`$ are given in Appendix A; the basic ideas behind them (and the techniques for their calculation) are simple elaborations of the above.
### B Average over longitudinal fields: Degeneracy blocking
We now consider the effective Hamiltonian
$$H_{\text{eff}}=2\mathrm{\Delta }_o\tau _x+\tau _z\{\xi _H+\underset{k=1}{\overset{N}{}}\omega _k^{}\widehat{\sigma }_k^z\};$$
(49)
To solve this we will assume the model discussed previously, in which all $`\{\omega _k^{}\}`$ cluster around a single central value $`\omega _o`$ (cf. eqtns (26),(25), and Fig. 4). Since the bath now just acts as an extra longitudinal field, we are dealing with the trivial case of a biased two-level system, with bias energy $`(\xi _H+ϵ)`$, where $`ϵ=_{k=1}^N\omega _k^{}\sigma _k^z`$. The only question is how to average over the internal bias- this depends on whether we deal with a single central system, or a statistical ensemble of them (corresponding to either an average over many measurements on a single system, or a single measurement on a large number of non-interacting systems- interactions are discussed in section 5.A).
For a single central spin, the dynamics of $`\stackrel{}{\tau }`$ in this model are completely trivial- one has $`P_{11}(t)=[1(\mathrm{\Delta }_o^2/E^2)\mathrm{sin}^2Et]`$, where $`E=\xi _H+ϵ`$, ie., resonant tunneling of an isolated spin in a longitudinal bias feld (compare (A5)).
For an ensemble of central spins, we must average over the whole bias range. In what follows let us assume for definiteness a spin bath at some equilibrium temperature $`T=1/\beta `$; then the ensemble average is just a weighted average over bias:
$$𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}$$
(50)
where $`Z(\beta )`$ is the appropriate partition function.
$$P_{11}(t)=𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}\left[1\frac{2\mathrm{\Delta }_o^2}{(ϵ+\xi _H)^2+4\mathrm{\Delta }_o^2}\left(1\mathrm{cos}(2t\sqrt{(ϵ+\xi _H)^2+4\mathrm{\Delta }_o^2})\right)\right]$$
(51)
The physical interpretation is obvious ; only a very small fraction $`A(\xi _H)\mathrm{\Delta }_o/E_o`$ of central spins in the ensemble in the ”resonance window”, ie.,, having total bias $`|(\xi _H+ϵ)|\mathrm{\Delta }_o`$, can make transitions- this selects states with internal bias around $`ϵ\xi _H`$. All other states lack the near-degeneracy between initial and final energies required for resonant tunneling- they are ”degeneracy blocked” . The resulting correlation function is then
$$P_{11}(t)=12A(\xi _H)\underset{k=0}{\overset{\mathrm{}}{}}J_{2k+1}(4\mathrm{\Delta }_ot).$$
(52)
where $`A(\xi )=\mathrm{\Delta }_oW(\xi )`$ (except in the unphysical case where the polarisation group linewidths $`\stackrel{~}{\mathrm{\Gamma }}_M<\mathrm{\Delta }_o`$. For the usual case where all polarisation groups overlap, and $`W(ϵ)`$ has the Gaussian form (25), one has $`A(\xi )/(2\pi )^{1/2}=(\mathrm{\Delta }_o/E_o)\mathrm{exp}(2\xi ^2/E_o^2)`$. It is not surprising to find that the spectral absorption function $`\chi ^{\prime \prime }(\omega )=Im𝑑tP_{11}(t)`$, corresponding to (52), has the ”BCS” form
$$\chi ^{\prime \prime }(\omega )=A(\xi _H)\frac{8\mathrm{\Delta }_o}{\omega \sqrt{\omega ^216\mathrm{\Delta }_o^2}}\eta (\omega 4\mathrm{\Delta }_o),$$
(53)
Finally, let us note that one may imagine a case where one has an ensemble of systems in which, although the bath state is not fixed, the polarisation group is known to be equal to $`M`$. In this case we must replace (50) by
$$𝑑ϵG_M(ϵ)\frac{e^{\beta ϵ}}{Z_M(\beta )}$$
(54)
in (51) for $`P_{11}(t)`$, where $`Z_M(\beta )`$ is the partition function for the $`M`$-th polarisation group. If we then recalculate $`\chi ^{\prime \prime }(\omega )`$ in the same way we find almost zero absorption unless $`|M\omega _o+\xi _H|<Max(\mathrm{\Delta }_o,\stackrel{~}{\mathrm{\Gamma }}_M)`$, where $`\stackrel{~}{\mathrm{\Gamma }}_M`$ is again the linewidth of the $`M`$-th polarisation group (cf. (24).
### C Average over transverse fields: Orthogonality blocking
Until now we have ignored the ”transverse field” part $`_{k=1}^N\omega _k^{}\stackrel{}{m}_k\widehat{\stackrel{}{\sigma }}_k`$ of the diagonal term in the effective Hamiltonian (17). To study this let us consider again an effective Hamiltonian which has no non-diagonal terms apart from the ”bare” tunneling, but having all diagonal terms:
$$H_{\text{eff}}=2\mathrm{\Delta }_o\tau _x+\widehat{\tau }_z\omega _o^{}\underset{k=1}{\overset{N}{}}\stackrel{}{l}_k\widehat{\stackrel{}{\sigma }}_k+\underset{k=1}{\overset{N}{}}\omega _k^{}\stackrel{}{m}_k\widehat{\stackrel{}{\sigma }}_k,$$
(55)
We will assume $`\omega _o^{}\omega _o^{}`$, ie., that the transverse ”orthogonality blocking” part of the diagonal interaction is much smaller than the longitudinal part. To make things as simple as possible we drop all degeneracy blocking effects, ie., we assume all $`\omega _k^{}`$ are equal (ie., $`\mu =0`$). It then follows that the spin bath spectrum is split by $`\omega _o^{}`$ into $`2N`$ ”polarization groups” of degenerate lines, with $`C_N^{(N+M)/2}`$ degenerate states in polarisation group $`M`$ (cf. eqtn. (23).
At first glance it seems that the $`\{\stackrel{}{\sigma }_k\}`$ in (55) simply act on the central spin $`\stackrel{}{\tau }`$ as an external field. However this is wrong; it ignores their role as dynamic quantum variables. The dynamics come because the $`\{\stackrel{}{\sigma }_k\}`$ can precess in the fields $`\{\stackrel{}{\gamma }_k\}`$ acting on them, and these change each time $`\stackrel{}{\tau }`$ makes a transition (cf. eqtn. (21)). Quantum mechanically, the precession caused by $`\omega _k^{}`$ is equivalent to saying that some bath spins are flipped when the central spin $`\stackrel{}{\tau }`$ flips (in general a different number of them during each transition of $`\stackrel{}{\tau }`$). To see this formally, recall that $`\omega _k^{}`$ exists when the initial and final fields $`\stackrel{}{\gamma }_k^{(1)}`$ and $`\stackrel{}{\gamma }_k^{(2)}`$ acting on $`\stackrel{}{\sigma }_k`$ are not exactly equal and opposite (cf. eqtn (21)). Defining the small angle $`\beta _k`$ by $`\mathrm{cos}2\beta _k=\stackrel{}{\gamma }_k^{(1)}\stackrel{}{\gamma }_k^{(2)}/|\stackrel{}{\gamma }_k^(1)||\stackrel{}{\gamma }_k^(2)|`$ (recall Fig. 3), we see that the initial and final states of $`\stackrel{}{\sigma }_k`$ are related by
$$\stackrel{}{\sigma }_k^f=\widehat{U}_k\stackrel{}{\sigma }_k^{in}=e^{i\beta _k\widehat{\sigma }_k^x}\stackrel{}{\sigma }_k^{in}.$$
(56)
Suppose now the initial spin bath state belongs to polarisation group $`M`$. If, when $`\stackrel{}{\tau }`$ flips, bath spins also flip so that $`MM`$, then since $`E_{}(M)=E_{}(M)`$, resonance is still preserved, a transition is possible- indeed $`\stackrel{}{\tau }`$ cannot flip at all unless there is a change in polarisation state of magnitude $`2M`$. For this change in polarisation of $`2M`$, at least $`M`$ spins must flip; moreover, for resonant transitions to continue (incoherently), the bath polarisation state must change by $`\pm 2M`$ each time $`\stackrel{}{\tau }`$ flips.
Let us therefore define $`P_M(t)`$ as the correlation function $`P_{11}(t)`$ restricted to systems for which the bath polarisation is $`M`$. For a thermal ensemble,
$$P_{11}(t;T)=\underset{M=N}{\overset{N}{}}w(T,M)P_M(t),$$
(57)
with a weighting $`w(T,M)=Z^1C_N^{(N+M)/2}e^{M\omega _o^{}/k_BT}`$, where $`Z`$ is the partition function.
In Appendix A.2 we calculate $`P_M(t)`$ (see Eq. (A40) as a weighted average over an orthogonality variable $`x`$:
$`P_M(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑xxe^{x^2}\left(1+\mathrm{cos}[4\mathrm{\Delta }_\mathrm{\Phi }J_M(2x\sqrt{\kappa })t]\right)`$ (58)
$`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑xxe^{x^2}P_{11}^{(0)}(t,\mathrm{\Delta }_M(x)),`$ (59)
where $`P_{11}^{(0)}(t,\mathrm{\Delta }_M)`$ is just the usual free 2-level correlator (Eq.(A1)), but now with an an $`x`$-dependent tunneling amplitude $`\mathrm{\Delta }_M(x)=2\mathrm{\Delta }_\mathrm{\Phi }J_M(2x\sqrt{\kappa })`$; and $`\kappa `$ is the ”orthogonality exponent”, defined by
$$e^\kappa =\underset{k=1}{\overset{N}{}}\mathrm{cos}\beta _ke^{(1/2)_k\beta _k^2},$$
(60)
The orthogonality blocking term $`\beta _k`$ is analogous to the topological decoherence term $`\alpha _k`$. It is important to understand why we must introduce the average over $`x`$. Mathematically, it comes from the restriction to a single polarisation group (see Appendix A.2). Physically, it corresponds to a phase average just like that in (45) (compare the Bessel functions), but now the phase is that accumulated between transitions of the central spin (rather than during these transitions, as in topological decoherence). This phase accumulates if $`\omega _k^{}0`$, because then the field $`\stackrel{}{\gamma }_k`$ on the $`k`$-th bath spin does not exactly reverse when the central spin flips, and so this spin must start precessing in the new field. It is random simply because the waiting time between flips is a random variable, in the path integral.
We shall not give full details for the dynamics of this limiting case (for which see ref ), but just enough to understand the physics. First, note that the terms $`P_M(t)`$ are easily verified to be incoherent, and so $`P_{11}(t)fP_{M=0}(t)+incoherent`$, where $`f=\sqrt{2/\pi N}`$. Even the small fraction $`f`$ of systems in an ensemble having $`M=0`$ will only have coherent dynamics if $`\kappa 1`$. The easiest way to see this is to again calculate the spectral absorbtion function from (58), to get :
$$\chi _{M=0}^{\prime \prime }(\omega )=\frac{\pi f}{2\mathrm{\Delta }_o\kappa ^{1/2}}\underset{j}{}\frac{x_je^{x_j^2}}{|J_1(2\kappa ^{1/2}x_j)|}|_{J_0(2\kappa ^{1/2}x_j)=\pm (\omega /2\mathrm{\Delta }_o)}$$
(61)
which leads, as $`\kappa `$ increases, to an ever-increasing number of square-root singularities in $`\chi _{M=0}^{\prime \prime }(\omega )`$. For $`\kappa <O(1)`$ only a single root $`x_1[(1\omega /2\mathrm{\Delta }_o)/\kappa ]^{1/2}`$ enters, and $`\chi _{M=0}^{\prime \prime }(\omega )(\pi f/2\mathrm{\Delta }_o\kappa )e^{[(1\omega /2\mathrm{\Delta }_o)/\kappa ]}`$ for $`\omega <2\mathrm{\Delta }_o`$, ie., a fairly sharp asymmetric peak at the resonant frequency of the free spin. For larger $`\kappa `$ the multiple peaks and tails mix, and $`\chi _{M=0}^{\prime \prime }(\omega )`$ shows no obvious peak- moreover, its spectral weight is shifted to ever-lower frequencies:
$$\begin{array}{c}P_{M=0}(t)=1/2[1+\mathrm{cos}(4\mathrm{\Delta }_{eff}t)]\hfill \\ \mathrm{\Delta }_{eff}=2\mathrm{\Delta }_\mathrm{\Phi }e^\kappa \hfill \end{array}\}\kappa 1$$
(62)
$$\begin{array}{c}P_0(t)=14\mathrm{\Delta }_{eff}^2t^2+O(\mathrm{\Delta }_{eff}^4t^4)\hfill \\ \mathrm{\Delta }_{eff}=\mathrm{\Delta }_\mathrm{\Phi }/(\pi \kappa )^{1/4}\hfill \end{array}\}\kappa 1$$
(63)
Notice that this reduction of the transition rate is much slower than the usual polaronic or Anderson/Hopfield orthogonality catastrophe, relevant to oscillator baths, which gives exponential suppression of $`\mathrm{\Delta }_{eff}`$ for strong coupling. This is understood as follows. In oscillator bath models, band narrowing comes essentially without any bath transitions (most of the polaron ”cloud” is in virtual high-frequency modes) - it is adiabatic. Here, however, roughly $`\kappa `$ spins flip each time $`\stackrel{}{S}`$ flips (the probability of $`r`$ flips is $`\kappa ^re^\kappa /r!`$, which peaks at $`r\kappa `$), even though we only consider $`P_{M=0}(t)`$, i.e., even though $`\mathrm{\Delta }M=0`$ (just as many bath spins flip one way as the other).
A further examination of the correlation function shows that the structure of $`P_{M=0}(t)`$ is exceedingly bizarre- it was described in detail in ref. . In section 4.A and Fig. 8 we will return to the physics of orthogonality blocking, but including other mechanisms as well (see also Fig. 7 below).
### D Averaging over spin bath fluctuations
The previous 3 averages assume no intrinsic spin bath dynamics- the bath acquired its dynamics from the central spin. Consider now a Hamiltonian
$$H_{\text{eff}}=\mathrm{\Delta }\tau _x+\xi \tau _z+\tau _z\underset{k=1}{\overset{N}{}}\omega _k^{}\widehat{\sigma }_k^z+\underset{k=1}{\overset{N}{}}\underset{k^{}=1}{\overset{N}{}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta ;$$
(64)
in which we assume $`|V_{kk^{}}|\omega _k^{}`$, but arbitrary $`V_{kk^{}}/\mathrm{\Delta }`$. The addition of $`V_{kk^{}}`$, to what would have been a simple degeneracy blocking Hamiltonian, gives the spin bath its dynamics, and causes 2 changes (as noted previously in our introductory presentation of the central spin model). First, a polarisation group $`M`$ acquires an ”intrinsic linewidth” $`\mathrm{\Gamma }_oV_oN^{1/2}`$, where $`V_o`$ is a typical value of $`|V_{kk^{}}|`$ for the $`N`$ bath spins (of course normally $`\mathrm{\Gamma }_o\stackrel{~}{\mathrm{\Gamma }}_M`$, unless the $`\{\omega _k^{}\}`$ happen to be extremely tightly bunched together). Second, the transverse part of $`V_{kk^{}}^{\alpha \beta }`$ causes pairwise flipping amongst bath spins (eg., transitions $`|_k_k^{}|_k_k^{}`$), at a rate $`NT_2^1`$, where $`T_2^1V_o`$. This ”spin diffusion” in the bath causes the internal bias $`ϵ`$ to fluctuate in time, inside the energy range of the polarisation group $`M`$, with a random walk correlation $`[ϵ(t)ϵ(t^{})]^2=\mathrm{\Lambda }^3|tt^{}|`$, where $`\mathrm{\Lambda }^3=\stackrel{~}{\mathrm{\Gamma }}^2T_2^1`$, for timescales $`\mathrm{\Lambda }|tt^{}|1`$.
It thus follows that for a given single central spin, with its surrounding spin bath in polarisation group $`M`$, the problem reduces to calculating the dynamics of a 2-level system in a longitudinal bias field $`(\xi +ϵ(t))`$, where $`\xi `$ is the applied bias, and the internal field is $`ϵ(t)=M\omega _o+\delta ϵ(t)`$. The correlation properties of $`\delta ϵ(t)`$ are those just described- our task is simply to functionally average over these fluctuations in the calculation of $`P_{11}(t)`$. In doing this we will make the physically sensible assumption of ”fast diffusion” of $`\delta ϵ(t)`$, such that the time $`\mathrm{\Delta }t`$ it takes for the bias to diffuse across the ”resonance window”, of energy width $`\mathrm{\Delta }`$, satisfies $`\mathrm{\Delta }t1/\mathrm{\Delta }`$. Then the system has no time to tunnel coherently, but can only make an incoherent ”Landau-Zener” transition. Since the bias changes by $`\delta ϵ\delta \omega _o(N/(T_2\mathrm{\Delta }))^{1/2}`$ in a time $`1/\mathrm{\Delta }`$, this formally requires that
$$\mathrm{\Delta }^3\mathrm{\Lambda }^3\stackrel{~}{\mathrm{\Gamma }}^2T_2^1(fastdiffusion).$$
(65)
This problem is solved in Appendix A, by performing a weighted average over dynamic bias fluctuations, with the restriction that these only occur inside polarisation group $`M`$; the relevant average is
$$𝒟ϵ(t)e^{\frac{1}{2}{\scriptscriptstyle 𝑑t_1𝑑t_2K(t_1t_2)ϵ(t_1)ϵ(t_2)}};$$
(66)
where $`2K^1(t_1t_2)=\mathrm{\Lambda }^3(|t_1|+|t_2||t_1t_2|)`$ is the correlator of the dynamic spin bath fluctuations (see Appendix A). One finds that $`P_{11}(t)`$ decays as a simple exponential $`P_{11}(t)=e^{t/\tau _M(\xi )}`$, where
$$\tau _M^1=2\pi ^{1/2}\frac{\mathrm{\Delta }^2}{\stackrel{~}{\mathrm{\Gamma }}}e^{(\xi +M\omega _o)^2/\stackrel{~}{\mathrm{\Gamma }}^2}$$
(67)
This result is easily understood- the bias fluctuations can cause the system to pass briefly through resonance (allowing a transition of the central spin), but only if the net static bias $`\xi +M\omega _o`$ is not greater than the range $`\stackrel{~}{\mathrm{\Gamma }}`$ of the fluctuations. By comparing with the case of pure degeneracy blocking we see that the important role of the bath dynamics is (i) to unblock the central spin dynamics, by helping it to find resonance, now over an energy window of width $`\stackrel{~}{\mathrm{\Gamma }}`$ (instead of $`\mathrm{\Delta }_o`$) around zero bias (recall Fig. 4), and (ii) to change the central spin dynamics from coherent to incoherent tunneling.
Note that in a model like (64) we have eliminated bath fluctuations between different polarisation groups. The basic assumption is that any ”$`T_1`$ processes” in the intrinsic bath dynamics, which could change $`M`$ in the absence of the central spin, are very slow (At low $`T`$, $`T_1`$ for nuclear or paramagnetic spins does become extremely long). However this is not always realistic- we return to this point in section 5.A. If $`T_1`$ is short, one must make a dynamical average over 2 kinds of fluctuation, usually with quite different time correlation, viz., the intra-polarisation group fluctuations, described by (66) and with correlation time $`NT_2^1`$, and the inter-polarisation group fluctuations, occuring on a timescale $`T_1`$ (cf. ).
### E Averaging over the Spin Bath: General Results
We now turn to the problem of averaging over the spin bath for the general form of the Central Spin Hamiltonian given in eqtn. (17). This can be given in the form of a marvellously simple prescription- one simply applies the 4 averages we have just seen, to the problem of a simple biased 2-level system! We begin by giving the explicit prescription (whose proof is given in Appendix B), and make a few comments on it.
The prescription begins with the following 4 averages (all of which we have seen in the preceding 4 sub-sections):
$$\text{ (a) A ”topological phase average” }\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(\nu )\frac{d\phi }{2\pi }e^{i2\nu (\mathrm{\Phi }\phi )};$$
(68)
$$\text{ (b) An ”orthogonality average” }\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}_0^{\mathrm{}}𝑑xxe^{x^2};$$
(69)
$$\text{ (c) A ”bias average” }𝑑ϵG_M(ϵ)\frac{e^{\beta ϵ}}{Z_M(\beta )}OR𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}\underset{M}{};$$
(70)
$$\text{ (d) A ”bath fluctuation average” }𝒟ϵ(t)e^{\frac{1}{2}{\scriptscriptstyle 𝑑t_1𝑑t_2K(t_1t_2)ϵ(t_1)ϵ(t_2)}};$$
(71)
As before, we assume a thermal distribution over spin bath biases, with a corresponding partition function and $`Z(\beta )`$. All averages are normalized to unity. The weighting function $`F_\lambda ^{}(\nu )=e^{4\lambda ^{}\nu ^2}`$ in (68) is a generalisation of that in (47), to allow for arbitrary directions of the unit vector $`\widehat{\stackrel{}{n}}_k`$; we now define
$$\lambda =\frac{1}{2}\underset{k}{}|\alpha _k|^2(1(n_k^z)^2),\lambda ^{}=\frac{1}{2}\underset{k}{}\alpha _k^2(n_k^z)^2,$$
(72)
Now, suppose we want to calculate $`P_{11}(t)`$. The prescription is fairly obvious in the light of the results given above for the 4 limiting cases. One follows the following steps:
(i) Begin with the quantity
$$P_{11}^{(0)}(t;\mathrm{\Delta }_M(\phi ,x);ϵ)=1\frac{\mathrm{\Delta }_M^2(\phi ,x)}{E_M^2(\phi ,x)}\mathrm{sin}^2(E_M(\phi ,x)t),$$
(73)
which is just the free central spin correlator (cf. (A5)) having tunneling matrix element $`\mathrm{\Delta }_M(\phi ,x)`$, and in an ”internal field” bias $`ϵ`$. The energy splitting $`E_M`$ is given by $`E_M^2(\phi ,x)=\mathrm{\Delta }_M^2(\phi ,x)+ϵ^2`$, and the matrix element $`\mathrm{\Delta }_M`$ is
$$\mathrm{\Delta }_M(\phi ,x)=2\stackrel{~}{\mathrm{\Delta }}_o|\mathrm{cos}(\phi )J_M(2x\sqrt{\gamma })|,$$
(74)
$$\gamma =\{\begin{array}{cc}\lambda \hfill & \text{if }\lambda \kappa \text{ (topological decoherence regime)}\hfill \\ \kappa \hfill & \text{if }\kappa \lambda \text{ (orthogonality blocking regime)}\hfill \end{array}.$$
(75)
We defined $`\kappa `$ previously (eqtn. (60). We will not give results for the case $`\kappa \lambda `$; they are extremely complex, do not appear to add new physics, and seem unlikely to be realised in practice.
(ii) Now carry out the averages over topological phase \[Eq. (68)\] and orthogonality \[Eq. (69)\], to give an expression $`P_M(t,ϵ)`$ describing the central spin dynamics in a bias $`ϵ`$, coming from a bath in polarisation state $`M`$:
$$P_M(t;ϵ)=2_0^{\mathrm{}}𝑑xxe^{x^2}\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(\nu )\frac{d\phi }{2\pi }e^{i2\nu (\mathrm{\Phi }\phi )}\left[1\frac{\mathrm{\Delta }_M^2(\phi ,x)}{E_M^2(\phi ,x)}\mathrm{sin}^2(E_M(\phi ,x)t)\right],$$
(76)
where the weighting function is $`F_\lambda ^{}(\nu )=e^{4\lambda ^{}\nu ^2}`$ over winding number $`\nu `$ (recall eqtn. (47));
(iii) Then, carry out the bias average \[Eq. (70)\]. We will assume in the following for definiteness an ensemble average over all polarisation groups, thereby ensuring a summation over $`M`$, to give
$$P_{11}(t;T)=1𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}\underset{M=N}{\overset{N}{}}\left(1P_M(t,ϵM\omega _o)\right);$$
(77)
This result summarizes the central spin dynamics in the case where the spin bath has no dynamics of its own, and only acquires dynamics through its interaction with the central system. In some cases there will be no intrinsic bath dynamics, and this will be the final answer. If we wish to apply the theory to a single central system, or for some reason we can fix the polarisation group to be a definite value $`M`$, then we drop the summation over $`M`$ in (77), and replace $`W(ϵ)`$ by $`G_M(ϵ)`$.
(iv) When the interaction term $`V_{kk^{}}`$ plays a role, we apply the 4th average (71) to (77), as described in Appendix A (cf. also the discussion in ref. ). This gives the completely incoherent form
$$P_{11}(t)=\underset{M}{}w(T,M)x𝑑xe^{x^2}\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}\frac{\phi }{2\pi }F_\lambda ^{}(\nu )e^{i2\nu (\mathrm{\Phi }\phi )}\left[1+e^{t/\tau _M(x,\phi )}\right],$$
(78)
where the relaxation rate $`\tau _M^1(x,\phi )`$ is given by
$$\tau _M^1(x,\phi )=2\mathrm{\Delta }_M^2(x,\phi )𝑑ϵG_\mu (ϵ)_0^{\mathrm{}}𝑑se^{iϵs}e^{\mathrm{\Lambda }^3s^3/6}=2\mathrm{\Delta }_M^2(x,\phi )_0^{\mathrm{}}𝑑se^{(\mu \omega _o)^2s^2/4}e^{\mathrm{\Lambda }^3s^3/6};$$
(79)
with $`\mathrm{\Delta }_M(x,\phi )`$ given by (74), and where $`G_\mu (ϵ)`$ is a Gaussian of width $`\stackrel{~}{\mathrm{\Gamma }}=\mu \omega _o`$. Since $`\stackrel{~}{\mathrm{\Gamma }}T_2^1`$ we have also $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Lambda }`$, and so we get
$$\tau _M^1(x,\phi )=\frac{2\mathrm{\Delta }_M^2(x,\phi )}{\pi ^{1/2}\stackrel{~}{\mathrm{\Gamma }}},$$
(80)
This result is the most general one for the dynamics of the central spin, if all 4 bath averages are included- it is generally valid, with only the single restriction that the diffusion of the fluctuating bath bias in energy space be fast (cf. eqtn. (65). In the absence of such fluctuations we go back to (77).
In essentially all physically realistic cases the different polarisation groups strongly overlap. It is then simpler to transform the sum over $`M`$ in (77) or (80) into an integral over energy bias $`\xi `$ , using the change of variables $`_M𝑑\xi /2\omega _o`$, and then integrate over $`\xi `$. One way to do this (using steepest descents) was detailed in ref. (compare eqtns (4.41)-(4.47) in that paper). For exact answers one can use the identity
$$_0^{\mathrm{}}x𝑑xe^{x^2}J_M^2(2x\sqrt{\gamma })\mathrm{cos}^2\varphi =I_M(2\gamma )e^{2\gamma }\mathrm{cos}^2\varphi $$
(81)
to evaluate either (77) or (78).
In the next section we will evaluate $`P_{11}(t)`$ and its Fourier transform for a number of different parameter regimes. But even before doing the integrals, the qualitative behaviour is obvious. Relaxation is only occurring for central spins which happen to be within a bias $`\xi _o`$ of exact resonance. The width $`\xi _o`$ of this ”resonance window” is coming from the energy which the bath spins can provide to the central spin, by flipping up to $`\gamma `$ bath spins; hence $`\xi _o\gamma \omega _o`$ (formally this is obvious from the properties of the Bessel functions in (74) and (81), which fall off very fast once $`M>\gamma `$). We show graphically the relaxation of different groups in Fig. 7(a); again one sees how only groups with $`M\gamma `$ relax quickly. The $`T_2`$ bath fluctuations help this process by bringing the a central spin in polarisation group $`M`$ to its resonance window (of width $`\mathrm{\Delta }_M`$). Only transitions of systems having $`M=0`$ can show (partial) coherence; all transitions with $`M0`$ are essentially incoherent. Note that the resonance window will not be visible in a resonant absorption experiment (Fig. 7(b)); higher $`M`$ groups contribute only a very low frequency contribution to this. This nicely demontrates that one is very far from any linear-response regime in the present system (so that, e.g., the fluctuation-dissipation theorem is somewhat irrelevant here).
Without the bath, transitions of the central spin would be coherent, but over the far smaller resonance window of width $`\mathrm{\Delta }_o`$. If we only had $`T_2`$ bath fluctuations, but $`\gamma 1`$ (ie., no bath spins flipped during the transitions of the central spin), then we would again get incoherent relaxation, but this time with $`\xi _o\stackrel{~}{\mathrm{\Gamma }}_{M=0}`$ (the width of the $`M=0`$ polarisation group). One can also imagine a situation in which $`T_1`$ is very short, so that all polarisation groups are involved in the relaxation, and the resonance window is just the distribution $`W(\xi )`$, with $`\xi _o=E_o=N^{1/2}\omega _o`$.
These results thus tell us that in the presence of a spin bath, any ensemble of central spins, initially spread over a range of biases, will start relaxing by digging a ”hole” of width $`\xi _o`$ around zero bias . This hole reflects the intrinsic central spin dynamics (ie., it is not being produced by interaction with some external resonant signal- it should not be confused in any way with the ”spectral hole-burning” done by experimenters working on glasses or in optics, using an external source). As discussed in a number of papers , evaluation of (78), using either steepest descents or other means, shows that the system under most conditions relaxes incoherently with a relaxation rate (after summing over all polarisation groups, and doing the orthogonality and phase integrals) given approximately as a function of bias $`\xi `$ by
$$\tau ^1(\xi )=\tau _o^1e^{|\xi |/\xi _o}\frac{2\mathrm{\Delta }^2}{\pi ^{1/2}\stackrel{~}{\mathrm{\Gamma }}}e^{|\xi |/\xi _o}$$
(82)
All of this is in complete contrast to how inelastic tunneling works in the presence of an oscillator bath ; there the relaxation rate typically increases as one moves away from resonance, usually as a power in bias ($`\tau ^1(\xi )\xi `$ for diagonal coupling to phonons, $`\xi ^3`$ for non-diagonal coupling to phonons, and $`\xi ^{2\alpha 1}`$ for diagonal coupling to Ohmic baths like electrons via a dimensionless coupling $`\alpha `$). Thus one does not expect hole-digging for oscillator bath-mediated quantum relaxation, except over a very narrow region of width $`\mathrm{\Delta }_0`$.
Finally, we note that one may also give a formal prescription for the case where some central spin couples simultaneously to an oscillator bath and a spin bath. We do not give the details here- they are discussed fairly exhaustively (along with the results for the central spin dynamics) in ref. (see also end of section 4).
## IV Dynamics of the Central Spin
Given the large number of parameters entering into the 4 averages just described, we see little point in an exhaustive description of $`P_{11}(t)`$ over the whole parameter domain (for more extensive results see refs. ). Instead we concentrate on 3 points. First, we show how in the strong coupling regime, coherence is destroyed, leaving incoherent quantum relaxation; this regime applies to almost all mesoscopic or nanoscopic magnetic systems, because of their coupling to nuclear spins and to paramagnetic impurities. Second, we discuss the physics of the weak coupling regime (applicable to, eg., SQUIDs), and how in one limit of this regime one may formally map the spin bath to an oscillator bath. Finally, and very briefly, we comment on the results obtained when one couples simultaneously to a spin bath and an oscillator bath.
### A Strong coupling regime
As already explained, the strong-coupling regime is defined by the condition $`\omega _k^{}`$ and/or $`\omega _k^{}\mathrm{\Delta }_o`$. This condition applies to virtually all situations in which the couplings are hyperfine ones to nuclear spins, or exchange couplings to paramagnetic spins; and also when one has dipolar couplings to paramagnetic impurities or defects.
Almost all interesting physical examples in this regime fall either into the catagory of “strong orthogonality blocking” (when $`\kappa \lambda ^{}`$) or strong “phase decoherence” (when $`\lambda \kappa `$). In both cases the central system makes transitions accompanied by flips in the bath spins- so that even if the isolated central system is not in resonance, it can ”find resonance” by using the flipped bath spins to make up the energy difference. If the couplings are such that roughly $`\kappa `$ bath spins flip, the range of energy bias over which transtions can occur is extended to roughly $`2\kappa \omega _o`$. The central system is helped in this task by the fluctuations in bath bias caused by the interspin interactions $`V_{kk^{}}`$.
In what follows we concentrate on the physics of decoherence in this regime, with an eye to the physics of ”qubits” and of ”macroscopic quantum coherence”. We also look at the form of the relaxation. We begin by explaining the results without the bath fluctuations, and then show what happens on adding these.
(i) Results without bath fluctuations: In this case we must evaluate (77) and (76), suppressing either the orthogonality average or the topological average. In what follows we look at each case in turn, focussing particularly on the $`M=0`$ polarisation group contribution.
(a) Orthogonality Blocked regime ($`\kappa \lambda `$). In this regime only the orthogonality average $`2x𝑑xe^{x^2}`$, and the average over bias $`ϵ`$, are relevant- the phase average is approximated by a delta-function. The presence of the $`x`$-dependent transition matrix element $`\mathrm{\Delta }_M(x)=2\mathrm{\Delta }_oJ_M(2x\sqrt{\kappa })`$ means that polarisation groups with $`M\kappa `$ play a dominant role. Suppose however that we are interested in any coherent dynamics of the central spin- what will be found? It is obvious that transitions with finite $`M`$ will be essentially incoherent, so we concentrate on central spins for which $`M=0`$. Thus we simply integrate $`P_M(t,ϵ)`$ over $`ϵ`$, in the weighted bias average, to get
$`P_{11}^{M=0}(t)`$ $`=`$ $`14A{\displaystyle 𝑑xxe^{x^2}J_0(2x\sqrt{\kappa })\underset{k=0}{\overset{\mathrm{}}{}}J_{2k+1}\left[4\mathrm{\Delta }_oJ_0(2x\sqrt{\kappa })t\right]}`$ (83)
$`=`$ $`12{\displaystyle 𝑑xxe^{x^2}2A(x)\underset{k=0}{\overset{\mathrm{}}{}}J_{2k+1}\left[2\mathrm{\Delta }_0(x)t\right]},`$ (84)
where the $`x`$-dependent spectral weight is $`A(x)=A|J_0(2x\sqrt{\kappa })|`$. Notice we have just done an ”orthogonality average” over a ”biased averaged” expression for the free system with $`x`$-dependent tunneling frequency $`\mathrm{\Delta }_0(x)`$. A Fourier transform to frequency space (which is essentially a picture of the relaxation rate, as a function of energy bias $`\xi `$ for this system) gives the absorption spectrum
$$\chi _{M=0}^{\prime \prime }(\omega )=\frac{1}{\omega }𝑑xxe^{x^2}4A(x)\frac{\mathrm{\Delta }_o(x)}{[\omega ^24\mathrm{\Delta }_0^2(x)]^{1/2}}\eta (\omega 2\mathrm{\Delta }_o(x)),$$
(85)
Fig. 8 shows some representative plots for this ”coherent” part of $`\chi ^{\prime \prime }(\omega )`$; it is in fact almost completely incoherent, with total spectral weight
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}(d\omega /2\pi )\chi ^{\prime \prime }(\omega )`$ $`=`$ $`2A{\displaystyle 𝑑xxe^{x^2}J_0(2\sqrt{\kappa }x)}`$ (86)
$`=`$ $`{\displaystyle \frac{2\mathrm{\Gamma }(3/4)}{\pi ^{3/2}}}{\displaystyle \frac{A}{\kappa ^{1/4}}};`$ (87)
a result which is very accurate even for $`\kappa 0.02`$. Note that the shape of $`\chi ^{\prime \prime }(\omega )`$ will change once we include all other (incoherently relaxing) polarization sectors $`M0`$, and its total weight increases - in fact the total weight is $`A\kappa ^{1/4}`$ for large $`\kappa `$, since $`\kappa ^{1/2}`$ different polarization sectors contribute. The absorption in $`\chi ^{\prime \prime }(\omega )`$ from these higher $`M`$ groups will be concentrated at frequencies $`\mathrm{\Delta }_o`$ (compare also Fig 7(b)). Note however that relaxation itself (not described by the linear response function $`\chi ^{\prime \prime }(\omega )`$) will be spread incoherently over a frequency range $`\xi _o2\kappa \omega _o`$, ie., the ”hole-digging” in the relaxation occurs over a window of width $`\xi _o2\kappa \omega _o`$.
(b) Phase decoherence regime ($`\lambda \kappa `$). Let us now suppose the transverse field terms $`\{\omega _k^{}\}`$ are negligible compared to the $`\{\stackrel{}{\alpha }_k\}`$. We then deal with an effective Hamiltonian
$$H_{\text{eff}}=2\mathrm{\Delta }_o\{\widehat{\tau }_{}\mathrm{cos}[\mathrm{\Phi }i\underset{k=1}{\overset{N}{}}\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k]+H.c.\}+\widehat{\tau }_z\underset{k=1}{\overset{N}{}}\omega _k^{}\widehat{\sigma }_k^z.$$
(88)
Since $`\omega _k^{}\mathrm{\Delta }_o`$ by assumption, energy conservation requires that environmental spins flip with the central spin, just as in our discussion of pure orthogonality blocking. Thus in this case we must also keep the orthogonality average, to enforce this constraint, ie., we must perform the full average embodied in eqtns. (76),(77). The full answer, including both the real and imaginary parts of $`\alpha _k`$, is rather complicated, and is presented in Appendix B. Here we will consider the more transparent answer one gets when $`\alpha _k`$ is purely imaginary and adds to directly as a random variable to the central spin phase.
Let us again start with only the $`M=0`$ contribution to $`P_{11}(t)`$. Then we have
$$P_{11}^{M=0}(t)=𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}P_0(t,ϵ);$$
(89)
with $`P_0(t,ϵ)`$ given by (77) with $`M=0`$. We may now carry out the integration in (89), assuming that $`W(ϵ)`$ is given by the usual Gaussian form (25), to get
$$P_{11}^{M=0}(t)=12_0^{\mathrm{}}𝑑xxe^{x^2}\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(\nu )\frac{d\phi }{2\pi }e^{i2\nu (\mathrm{\Phi }\phi )}2A(\phi ,x)\underset{k=0}{\overset{\mathrm{}}{}}J_{2k+1}\left[2\mathrm{\Delta }_0(\phi ,x)t\right];$$
(90)
with $`A(\phi ,x)=A\mathrm{cos}\phi J_0(2x\sqrt{\lambda })`$. The corresponding absorption $`\chi ^{\prime \prime }(\omega )`$ is
$$\chi _{M=0}^{\prime \prime }(\omega )=\frac{2}{\omega }𝑑xxe^{x^2}\underset{\nu =\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(\nu )\frac{d\phi }{2\pi }e^{i2\nu (\mathrm{\Phi }\phi )}\frac{A(\phi ,x)\mathrm{\Delta }_0(\phi ,x)}{[\omega ^24\mathrm{\Delta }_0^2(\phi ,x)]^{1/2}}.$$
(91)
It is possible to write analytic expressions starting from (91), but in this somewhat pedagogical presentation we simply discuss the case when $`\mu =0`$, i.e., zero degeneracy blocking, when $`\omega _k=\omega _o`$ for all nuclei. The integration over bias is then absent (since $`W(ϵ)`$ is now just a set of $`\delta `$-function peaks, ie., $`W(ϵ)_{M=N}^NC_N^{(N+M)/2}\delta (ϵM\omega _o)`$), and we get
$$P_{11}(t)=\underset{M}{}w(T,M)P_M(t);w(T,M)=C_N^{(N+M)/2}e^{M\omega _o/T}/Z(\beta ),$$
(92)
where $`P_M(t)`$ now describes the dynamics in zero bias; it is given by exactly the same weighted average over phase as in (45):
$`P_M(t)`$ $`=`$ $`{\displaystyle 𝑑xxe^{x^2}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(m)\frac{d\phi }{2\pi }e^{i2m(\mathrm{\Phi }\phi )}\left\{1+\mathrm{cos}[4\mathrm{\Delta }_otJ_M(2x\sqrt{\lambda })\mathrm{cos}\phi ]\right\}}`$ (93)
$`=`$ $`{\displaystyle 𝑑xxe^{x^2}\left\{1+\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}(1)^mF_\lambda ^{}(m)e^{i2m\mathrm{\Phi }}J_{2m}[4\mathrm{\Delta }_otJ_M(2x\sqrt{\lambda })]\right\}},`$ (94)
This can be interpreted either as an orthogonality-blocked expression, with frequency scale $`\mathrm{\Delta }_M(\phi ,x)=2\mathrm{\Delta }_o\mathrm{cos}(\phi )J_M(2x\sqrt{\lambda })`$ which is then averaged over $`\phi `$, to give phase randomisation; or as an integration $`𝑑x`$ over an already topologically decohered function having frequency scale $`\mathrm{\Delta }_M(x)=2\mathrm{\Delta }_oJ_M(2x\sqrt{\lambda })`$. It is intuitively obvious (and easily demonstrated) that only $`P_0(t)`$ may behave coherently, with a fractional weight $`\sqrt{2/\pi N}`$ in an ensemble.
There are various interesting cases of (94) for $`M=0`$. If $`\lambda =0`$ (i.e., $`\stackrel{}{n}_k`$ is parallel to $`\widehat{\stackrel{}{z}}`$), then we go back to pure topological decoherence - the projection operator then commutes with the cosine operator. On the other hand if $`\lambda ^{}=0`$, we have pure orthogonality blocking as stated earlier, and in fact when $`\lambda ^{}=0`$, the parameter $`\lambda `$ plays the role of $`\kappa `$ in (58). Notice that whereas the case $`\lambda =0`$ can only occur accidentally, $`\lambda ^{}=0`$ is quite common - indeed it pertains to the model in Eqs. (30).
We really begin to see the analogy between orthogonality blocking and topological decoherence when $`\lambda ,\lambda ^{}1`$; just as with pure topological decoherence, $`F_\lambda ^{}(m)`$ collapses to a Kronecker delta, and we get the universal projected topological decoherence form:
$$P_{11}^{M=0}(t)𝑑xxe^{x^2}\left[1+J_0[4\mathrm{\Delta }_otJ_0(2x\sqrt{\lambda })]\right](\mu =0);$$
(95)
$$\chi _{M=0}^{\prime \prime }(\omega )\sqrt{\frac{2}{\pi N}}𝑑xxe^{x^2}\frac{4}{[16\mathrm{\Delta }_o^2J_0^2(2x\sqrt{\lambda })\omega ^2]^{1/2}}\eta (4\mathrm{\Delta }_oJ_0(2x\sqrt{\lambda })\omega ),$$
(96)
which generalizes the result of (48) for pure topological decoherence. Eq.(95) should be compared to (58).
We show in Fig.9 some results for $`\chi ^{\prime \prime }(\omega )`$ for selected values of $`\lambda `$. The results are startling; even a very small value of $`\lambda `$ significantly washes out pure topological decoherence; but for any large value of $`\lambda ^{}`$, we never get back the pure orthogonality blocking spectrum.
The results in the case where $`\alpha _k`$ is real, and $`\widehat{\stackrel{}{n}}_k`$ is along the $`\widehat{\stackrel{}{x}}`$-direction (see Eq. (33)) are obtained by simply converting the Bessel functions $`J_m`$ to Bessel functions of imaginary argument $`I_m`$.
Finally, note again that the above discussion of the $`M=0`$ polarisation group is irrelevant to the real experimental lineshape- an evaluation of the full expression (77), summing over all $`M`$, for the $`\lambda ,\lambda ^{}1`$ regime, simply gives incoherent relaxation , spread over a frequency range $`\xi _o\lambda \omega _o`$.
(ii) Including Bath Fluctuations: The modification of the above results, occasioned by the intrinsic spin bath fluctuations, was given in detail in Prokof’ev and Stamp . The fluctuations in bias allow the central system to cycle rapidly through the whole range of biases within a given polarisation group (transitions between different polarisation groups can occur through $`T_1`$ processes- usually much slower). Here we simply recall the main result, which is obtained by summing the relaxation forms from each polarisation group in an ensemble (cf. eqtns. (77) and (78)), and assuming that the spin bath $`T_1`$ is longer than all experimental times scales.
For a single central system, coupled to a spin bath in equilibrium at temperature $`kT\omega _o`$, one finds that after an initial short-time transient, the relaxation is roughly logarithmic over a very long period; in fact one finds for the strong coupling regime that
$$1P(t,\xi _H=0)\sqrt{\frac{1}{2\pi N}}\frac{\mathrm{ln}(t/\tau _o)}{\mathrm{ln}\left[\frac{1}{e\sqrt{\gamma }}\mathrm{ln}(t/\tau _o)\right]};(tt_c),$$
(97)
for times $`tt_c`$, where $`t_c\tau _o(2N/e^2\gamma )^{\sqrt{2N}}`$, and $`\tau _o=2\mathrm{\Delta }^2/\pi ^{1/2}\stackrel{~}{\mathrm{\Gamma }}`$ is the relaxation time of the $`M=0`$ polarisation group (compare eqtn. (80)); thus $`t_c`$ is extremely long! For $`tt_c`$ the system settles down to a rather different behaviour . This logarithmic behaviour can be roughly understood as coming from a distribution of barrier heights, for the different polarisation groups, which are then summed over- as discussed in ref. , section 4.3(b), the final result looks basically the same as that shown in Fig. 7(a). The fastest relaxation comes from the those polarisation groups in the ”resonance window” (recall the discussion at the end of section 3).
Two cautionary notes are in order here. First, (97) applies to a single relaxing system- but in the case of nanomagnetic systems, all experiments until now have been done on large numbers of nanomagnets, coupled together via long-range dipolar forces, which drastically changes the relaxation (see section 5.A below). Second, (97) should not be applied uncritically to experiments, even on single quantum systems. This is because in a real experiment there will also be (i) couplings to oscillator baths, and (ii) the relaxation will change once $`tT_1`$. In superconductors or metals, electronic oscillator baths will often dominate the relaxation , even at short times (for their effect on coherence, see section 5.C below). Even in insulating systems, one eventually expects the coupling to phonons to take over at long times , since this causes exponential decay- even if very slow, this will eventually become faster than the spin bath-mediated logarithmic decay in (97).
### B The weak coupling regime; relation to the oscillator bath
A question of considerable theoretical (and practical) interest is the transition to the weak coupling regime, where the perturbation on the central system dynamics by a single bath spin is small (even though the net effect of all bath spins, measured by parameters like $`\lambda `$ or $`\kappa `$, may still be large if $`N`$ is very large). The weak-coupling regime is thus defined by the condition $`\omega _k\mathrm{\Delta }_o`$, ie., both $`\omega _k^{}`$ and $`\omega _k^{}`$ are $`\mathrm{\Delta }_o`$. Again, a variety of cases is possible depending on how large are ratios like $`V_{kk^{}}/\mathrm{\Delta }_o`$ and $`\omega _k^{}/\omega _k^{}`$, or parameters like $`\lambda `$ and $`\kappa `$. In the following we will not be exhaustive, but simply consider two theoretically interesting cases, in which $`V_{kk^{}}`$ is assumed negligible, and we look at the limiting behaviour arising when either $`\omega _k^{}/\omega _k^{}1`$ or $`\omega _k^{}/\omega _k^{}1`$. We will also assume $`N1`$, otherwise the problem is trivial (the bath has little effect at all).
There are 2 ways to solve for the dynamics in this regime. One is to use the averaging already developed above- this simplifies considerably in the weak- coupling regime. The other is to map the problem onto an oscillator bath one, and then use standard techniques to solve this. We demonstrate the 2 methods by solving one problem with each.
(i) Longitudinally dominated case ($`\omega _k^{}/\omega _k^{}1`$): We assume the same Hamiltonian as in the discussion of the phase decoherence regime (eqtn. (88)), but now we can drop the orthogonality average- the projection to a single polarisation group is not required since $`\omega _k^{}\mathrm{\Delta }_o`$. Again, for simplicity we consider the case when $`\alpha _k`$ is imaginary.
We shall solve this using the techniques previously developed; we let $`x=0`$ in (74), and hence use a matrix element $`\mathrm{\Delta }(\phi )=2\mathrm{\Delta }_o\mathrm{cos}\phi `$ (which is independent of $`M`$). All polarisation groups overlap, and so we simply average over bias and topological phase:
$`P_{11}(t)`$ $`=`$ $`1{\displaystyle 𝑑ϵW(ϵ)\frac{e^{\beta ϵ}}{Z(\beta )}\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda (m)\frac{d\phi }{2\pi }e^{i2m(\mathrm{\Phi }\phi )}\left\{1\frac{\mathrm{\Delta }_0^2(\phi )}{ϵ^2+\mathrm{\Delta }_0(\phi )}\left(1\mathrm{cos}\left[2t\sqrt{ϵ^2+\mathrm{\Delta }_0^2(\phi )}\right]\right)\right\}}`$ (98)
$`=`$ $`1{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}F_\lambda (m){\displaystyle \frac{d\phi }{2\pi }e^{i2m(\mathrm{\Phi }\phi )}2A(\phi )\underset{k=0}{\overset{\mathrm{}}{}}J_{2k+1}[2\mathrm{\Delta }_0(\phi )t]},`$ (99)
with $`\mathrm{\Delta }_0(\phi )=2\mathrm{\Delta }_o\mathrm{cos}\phi `$ as before, and $`A(\phi )=A\mathrm{cos}\phi `$. This gives an absorption form
$$\chi ^{\prime \prime }(\omega )=\frac{2A}{\omega }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda (m)\frac{d\phi }{2\pi }e^{i2m(\mathrm{\Phi }\phi )}\frac{\mathrm{cos}^2\phi }{[(\omega /4\mathrm{\Delta }_o)^2\mathrm{cos}^2\phi ]^{1/2}}\eta (\omega /4\mathrm{\Delta }_o\mathrm{cos}\phi ),$$
(100)
which for large $`\lambda `$ simplifies to
$$\chi ^{\prime \prime }(\omega )=\frac{2A}{\omega }\frac{d\phi }{2\pi }\frac{\mathrm{cos}^2\phi }{[(\omega /4\mathrm{\Delta }_o)^2\mathrm{cos}^2\phi ]^{1/2}}\eta (\omega /4\mathrm{\Delta }_o\mathrm{cos}\phi ),$$
(101)
and can be expressed in terms of Elliptic functions.
Notice that since $`\omega _k\mathrm{\Delta }_o`$, the number $`N`$ of environmental spins must be very large to have a noticeable effect, i.e., for $`\lambda `$ to be appreciable. Thus if $`\lambda N\alpha _k^21`$, since $`\alpha _k\omega _k/\mathrm{\Omega }_o`$, we have $`N(\mathrm{\Omega }_o/\omega _k)^2(\mathrm{\Omega }_o/\mathrm{\Delta }_o)^2`$. This not only implies that $`\mu =N^{1/2}\delta \omega _k/\omega _o1`$ (ie., very strongly overlapping polarisation groups) but also a Gaussian half-width $`N^{1/2}\omega _o\mathrm{\Delta }_o`$, so that the internal bias $`ϵ\mathrm{\Delta }_o`$. The reason why the two mechanisms (topological decoherence and degeneracy blocking) are so easily combined is just because this bias is produced by all of the environmental spins, whereas only a few of them are actually flipped and their contribution to the bias field is small.
(ii) Transverse dominated case ($`\omega _k^{}/\omega _k^{}1`$): Consider now the case described by
$$H\text{eff}=\mathrm{\Delta }_o\widehat{\tau }_x+\underset{k=1}{\overset{N}{}}(\omega _k^{}\widehat{\sigma }_k^z+\widehat{\tau }_z\omega _k^{}\widehat{\sigma }_k^x),$$
(102)
with $`\omega _k^{}/\omega _k^{}1`$. One can map this problem to the spin-boson problem, in a way similar to that given in the paper by Caldeira et al. . We begin by noting that the time evolution operator for the $`k`$-th bath spin, under the influence of a central spin moving along a path $`Q_{(n)}(t)`$, is just
$$\widehat{U}_k(Q_{(n)},t)=T_\tau \mathrm{exp}\left\{i_0^t𝑑s\left[\omega _k^{}\widehat{\sigma }_k^z+Q_{(n)}(s)\omega _k^{}\widehat{\sigma }_k^x\right]\right\},$$
(103)
Weak coupling means we can expand the time-ordered exponent to second order in $`\omega _k^{}`$; completing the average in Eq. (41), and exponentiating the answer, one derives the influence functional for this problem . The ”high temperature” result $`kT\omega _k^{}`$ is readily found:
$$[Q,Q^{}]=\mathrm{exp}\left\{\underset{k}{}\frac{(\omega _k^{})^2}{2}\left|_0^t𝑑se^{2i\omega _k^{}s}[Q_{(n)}(s)Q_{(m)}^{}(s)]\right|^2\right\},$$
(104)
The evaluation of this depends on the distribution of couplings. In the ”strong decoherence” case where the bath states are spread over an energy range $`E_o`$ (defined as usual by $`E_o^2=2_k(\omega _k^{})^2`$) which is large (ie., $`E_o\omega _k^{},\mathrm{\Delta }_o`$), then $`[Q,Q^{}]`$ is simply approximated by its form for a single kink - anti-kink trajectory. With kink at $`t=t_1`$ and anti-kink at $`t_2`$, one only has a contribution if $`\omega _k^{}(t_2t_1)1`$, which gives $`[Q,Q^{}]=e^{\frac{1}{4}E_o^2(t_2t_1)^2}`$, ie., a decay on a time scale $`E_o^1`$. It then follows that kink/ anti-kink transitions are bound in closed pairs, with phase correlations decaying over a time $`\mathrm{\Delta }_o^1`$, and the leading terms sum to give exponential relaxation, ie., $`P_{11}(t)=1/2(1e^{t/\tau _R})`$, with a relaxation rate
$$\tau _R^1=2\mathrm{\Delta }_o_0^{\mathrm{}}𝑑te^{\frac{1}{4}E_o^2t^2}2\frac{\sqrt{\pi }\mathrm{\Delta }_o^2}{E_o}.$$
(105)
On the other hand if $`E_o\mathrm{\Delta }_o`$, we will get coherent oscillations of the system over a time scale $`\tau _\varphi \mathrm{\Delta }_o/E_o^2`$, ie., over roughly $`\mathrm{\Delta }_o^2/E_o^2`$ oscillation periods. We will use this result below, in discussing the possibility of seeing true mesoscopic quantum coherence effects in SQUIDs and nanomagnets (section 5.C).
It is of course obvious from the original discussion of Feynman and Vernon that in the weak-coupling regime, where an expansion to 2nd order in all bath couplings is sufficient, a mapping to an oscillator bath model must be possible . The method given above is of course not the only one- others are explained in refs. .
Central Spin coupled to oscillator and spin baths: For reference we also briefly note results for the case where a central spin is coupled simultaneously to a spin bath and an oscillator bath; the Hamiltonian is that in eqtn. (37), but without the sum over $`j`$, (and of course dropping the coupling term $`V(𝐫_i𝐫_j)`$ between different central spins). Some of the dynamic properties of this model were studied in ref. \- a complete study is still lacking. Quite generally one expects that at very short times the dynamics will be controlled by the spin bath, but at longer times incoherent oscillator bath- mediated transitions will take over. However these transitions still contain spin bath effects, via an integration over spin bath bias distribution, to give a result
$$P_{11}(t,\xi )=AZ^1(\beta )𝑑ϵW(ϵ\xi )e^{\beta (ϵ\xi )}\{f(T,ϵ)+[1f(T,ϵ)]e^{t/\tau (ϵ,T)}\}$$
(106)
where we assume that the ensemble of central spins is in equilibrium with a bath of oscillators at temperature $`T=1/\beta `$, in external bias $`\xi `$ (for more general cases see ref. ). In this case $`f(ϵ)=e^{\beta ϵ}/2\mathrm{cosh}(\beta ϵ)`$, and $`\tau (ϵ;T)`$ is the oscillator bath- mediated relaxation time; $`A`$ is a ”renormalisation” constant $`O(1)`$ coming from the spin bath dynamics. Most features of the results are obvious from (106). Thus, the oscillator bath unblocks transitions for a central spin way off resonance (ie., having $`\xi E_o`$); these spins relax much as they would without the spin bath. When $`\xi /leqE_o`$ there is a wide range of relaxation times, and the physical relaxation depends in a complex way on the oscillator bath spectrum as well as the bias. In any experiment one would see a crossover between purely spin-bath mediated relaxation at short times and this more complex behaviour- this crossover has not yet been studied in any detail.
## V Physical application of spin bath models
In this section we briefly review some recent work, mainly experimental, in which interactions with a spin bath play a role. We begin with magnetic systems where such interactions are strong, and where clear evidence exists for their influence on tunneling phenomena. We then discuss superconducting and normal systems. Finally, we discuss the difficult and controversial problem of decoherence and the mechanisms which govern it in nature. In this discussion the results of the Central Spin model are crucial- they describe the effect of a spin bath on a ”qubit”, or on a SQUID, or a magnetic macromolecule, which is trying to show coherent oscillations.
### A Magnetic Systems
(i) Magnetic Solitons: Magnetically ordered systems support a wide variety of soliton excitations, depending on the symmetry of the order parameter. These couple to various environmental excitations, which strongly affect their dynamics. These include both linear and non-linear couplings to “quasiparticle” excitations such as magnons , phonons , electrons , and photons , all oscillator baths. More serious effects come from any localized modes coupling to the soliton, most notably paramagnetic impurities and nuclear spins. Both of these can be understood theoretically using a model in which a moving particle couples to a spin bath , as in section 2.C.
Most experimental work in this area has looked at domain wall tunneling in ferromagnets. Early experiments looked at the dynamics of multi-wall systems . More recently attention has focussed on the tunneling of single domain walls, whose position can be monitored in various ways . These walls propagate along magnetic nanowires (of diameter 300-600 Å)and are thus mesoscopic objects (the material of choice is usually Ni, for which the domain wall thickness $`\lambda _w700\AA `$). Although some predictions of the tunneling theory have been verified (eg., the square root dependence of the tunneling exponent on the pinning field $`H_c`$; cf. ), the crossover temperatures to tunneling are usually much higher than calculated values . A clue to the reason for this may be found in the microwave resonance experiments of Hong & Giordano , which show extremely broad resonances, even down to $`1.5K`$. The coupling to phonons or electrons is far too small to give this; however, Oxygen impurities, which will act as paramagnetic impurities, would act as a strong time-varying potential on any wall, and broaden the line . It is clear that much more experimental work will be necessary to really understand the quantum dynamics of magnetic domain walls.
(ii) Magnetic Macromolecules: On the other hand, a number of spectacular quantum effects have unquestionably been seen in crystalline arrays of magnetic macromolecules, in one of the most important developments in magnetism in recent years. Experiments on tunneling phenomena have had striking success in the ”Mn-12” and ”Fe-8” molecules (each of which behaves at low $`T`$ as a ”giant spin” of spin 10). There have also been highly-publicised experiments on the large ”ferritin” molecule, which claim the observation of ”macroscopic quantum coherence”; these are discussed in section 5.C below.
The early molecular work demonstrated resonant tunneling (coming when spin levels of each molecule are brought into resonance), at relatively high temperatures . More recent experiments have gone into the quantum regime, in which only the 2 lowest levels of each giant spin are occupied. The classical-quantum crossover is clearest in the case of the Fe-8 system (which is very conveniently described by the easy axis/easy plane model discussed in section 2.F); below a temperature $`T_c0.4K`$, the dynamics is completely independent of $`T`$ (the energy gap to the next spin level in this system is $`5K`$). In Mn-12 there are effects arising from ”rogue” molecules which make things more complicated . At such low temperatures phonons are utterly irrelevant (experiments have now been pursued to $`T30mK`$), and the only dynamic environment left is the nuclear spin bath. The system thus seems at first to be an ideal realisation of the central spin model discussed in sections 3 and 4. The main complication is that experiments are done on a crystalline array of molecules, which interact via strong magnetic dipolar interactions. However, once one knows that the dynamics of a single central spin is going to be incoherent relaxation (cf. sections 3.E, 4.A), these interactions are rather easy to deal with . The system is then described by an obvious generalisation of the central spin Hamiltonian, viz.
$$H=\underset{j}{}H_j^{CS}+\underset{ij}{}V(𝐫_i𝐫_j)\widehat{\tau }_i^z\widehat{\tau }_j^z,$$
(107)
in zero applied field (compare eqtn. (37), without the oscillators). Here $`H_j^{CS}`$ is the central spin Hamiltonian (17) for the molecule at lattice position $`𝐫_j`$, and $`V(𝐫_i𝐫_j)`$ is the magnetic dipolar coupling between molecules $`i`$ and $`j`$.
To solve for the dynamics of such an interacting array of systems, one begins by defining a distribution function $`P_\alpha (\xi ,𝐫,t)`$ for a molecule at position $`𝐫`$ to be in a longitudinal bias $`\xi `$, with polarisation $`\tau _z=\alpha =\pm 1`$. It is then trivial to write down a BBGKY-like hierarchy of kinetic equations for $`P_\alpha (\xi ,𝐫,t)`$ and its multimolecular generalisations $`P^{(2)}(1,2)P_{\alpha _1,\alpha _2}^{(2)}(\xi _1,\xi _2;\stackrel{}{r}_1,\stackrel{}{r}_2;t)`$, and $`P^{(3)}(1,2,3)`$, etc., of which the first member is
$`\dot{P}_\alpha (\xi ,\stackrel{}{r})=`$ $`\tau _N^1(\xi )[P_\alpha (\xi ,\stackrel{}{r})P_\alpha (\xi ,\stackrel{}{r})]`$ (108)
$``$ $`{\displaystyle \underset{\alpha ^{}}{}}{\displaystyle \frac{d\stackrel{}{r}^{}}{\mathrm{\Omega }_0}\frac{d\xi ^{}}{\tau _N(\xi ^{})}\left[P_{\alpha \alpha ^{}}^{(2)}(\xi ,\xi ^{};\stackrel{}{r},\stackrel{}{r}^{})P_{\alpha \alpha ^{}}^{(2)}(\xi \alpha \alpha ^{}V(\stackrel{}{r}\stackrel{}{r}^{}),\xi ^{};\stackrel{}{r},\stackrel{}{r}^{})\right]},`$ (109)
in which relaxation is driven by the nuclear spin-mediated relaxation rate $`\tau _N^1(\xi )(\mathrm{\Delta }^2/\stackrel{~}{\mathrm{\Gamma }})e^{|\xi |/\xi _o}`$ (cf eqtn (82), in conjunction with the dipolar interactions. Under general conditions we must also solve for $`P^{(2)}`$ in terms of $`P^{(3)}`$, etc.; but if the initial experimental state is either polarised or annealed then $`P^{(2)}`$ factorises at $`t=0`$ (ie., $`P^{(2)}(1,2)=P(1)P(2)`$), and (109) can be solved. This led to the following predictions :
i) Relaxation should only occur for molecules having $`|\xi |\xi _0`$; consequently a “hole” rapidly appears in $`M(\xi ,t)=P_+(\xi ,t)P_{}(\xi ,t)`$ with time, with initial width $`\xi _o`$ determined entirely by the nuclear dynamics (and of course typically $`\xi _o\mathrm{\Delta }`$). In fact microscopic calculations of $`\xi _o`$ for the $`Fe`$-8 molecule can be done (recall Fig. 5(b)), since the hyperfine couplings are essentially dipolar (the relevant nuclei include 120 protons, 8 $`Br`$ nuclei, and 18 $`N`$ nuclei). One finds that even the weak molecular dipolar fields strongly distort the hyperfine fields, mixing up the different polarisation groups, and also giving a large value of $`\kappa `$. The final value $`\xi _o`$ will obviously depend sensitively on any nuclear isotopic substitution . The subsequent evolution of the hole depends on sample shape; this has been studied theoretically using Monte Carlo simulations .
(ii) The short-time relaxation for the total magnetisation $`M(t)=𝑑\xi M(\xi ,t)`$ should have a ”square-root” time form
$$M(H_0,t)=M_0[1(t/\tau _Q(H_0))^{1/2}],$$
(110)
where $`M_0`$ is the initial magnetization (or with appropriate modifications for other protocols, such as a zero-field cooling followed by relaxation in a field ), and $`\tau _Q^1(H_0)=c(\xi _o/E_D)\mathrm{\Delta }^2M(\xi =g\mu _BSH_0,t=0)`$, where $`W_D^1`$ is the ”density of states” of the distribution. The constant $`c`$ is dimensionless, depending on sample shape, and can be evaluated analytically or numerically. Notice that the existence of the square root does not depend on sample shape, and indeed its persistence over fractional relaxations of $`0.1`$ is clearly demonstrated in Monte Carlo simulations for different shapes .
This result implies that by varying $`H_0`$, one can measure $`M(\xi )`$, by extracting $`\tau _Q^1(H_0)`$ at successive values of $`H_0`$. If one knows $`M(\xi )`$ (as one does in an annealed sample - it should be Gaussian) then one may then extract $`\mathrm{\Delta }`$ from measurements of $`\tau _Q^1`$. The $`\sqrt{t}`$ dependence of $`M(t)`$ has since been reported in quite a few experiments, on both Fe(8) crystals and Mn(12) crystals (although the situation in Mn-12 is seriously complicated by ”impurities” ). The Fe-8 experiments have produced remarkable ”maps” of $`M(\xi )`$, and its time variation . The ”hole-digging” has been found in both Fe-8 and Mn-12 , with an ”intrinsic” short-time intrinsic linewidth which is roughly that expected from the hyperfine interactions, provided one takes account of the effect of internal fields , which make $`T_1`$ very short (so that $`\kappa `$ is effectively $`\sqrt{N}`$).
Wernsdorfer et al also used this technique to extract the value of $`\mathrm{\Delta }`$ for Fe-8. In a remarkable experiment, Wernsdorfer and Sessoli extended these measurements of $`\mathrm{\Delta }`$ to include a transverse applied field $`H_{}`$; as noted in section 2.F, the topological giant spin phase $`\mathrm{\Phi }`$ should vary with field $`H_{}`$ , producing Aharonov-Bohm oscillations in $`\mathrm{\Delta }(H_{})`$. These oscillations were found, both as oscillations in the relaxation rate $`\tau _Q^1(H_{})`$, and using a quite different AC absorption (”Landau-Zener”) technique; these independent techniques agreed rather well in the measurement of $`\mathrm{\Delta }`$. Notice, incidentally, that neither method can properly measure $`\mathrm{\Delta }`$ near its nodes (ie., where $`\mathrm{\Delta }0`$); this is because of the distribution of internal transverse fields. Incidentally, we should strongly emphasize that these experiments (even the Aharonov-Bohm ones) do not demonstrate coherent tunneling- indeed they show exactly the opposite! This is because the experiments are inherently relaxational (this is why all rates are $`|\mathrm{\Delta }|^2`$, and not $`\mathrm{\Delta }`$). Readers puzzled about how an Aharonov-Bohm effect can occur in a relaxation rate, are encouraged to think about other examples in physics where phase interference shows up in irreversible quantum phenomena. In the present case we may loosely define a ”decoherence time” $`\tau _\varphi `$ for the molecular spins (one should not push this too far, given the non-exponential nature of the central spin relaxation!), and one finds that $`\tau _\varphi \mathrm{\Delta }_o1`$ (no coherent oscillations possible) but $`\tau _\varphi \mathrm{\Omega }_o1`$ (ie., coherence is maintained during a single very rapid tunneling transition).
In a very recent experiment , deliberate modification of the nuclear isotopes in an $`Fe`$-8 crystal has shown a dramatic modification of the bulk magnetic relaxation- this is the most direct evidence so far for the role of the nuclear spins in mediating the quantum relaxation. It will be interesting to see a quantitative comparison with the calculated results for different isotopes .
(iii) Quantum Spin Glasses: The work on nanomagnets probes our understanding of collective phenomena in many magnetic systems, ranging from spin chains to quantum spin glasses. This latter example is very closely related to the molecular nanomagnets just described, since the model Hamiltonian is just (107), with the dipolar interaction replaced by a set $`\{J_{ij}\}`$ of ”frustrating” interactions, usually long-ranged. A typical experimental example is provided by the disordered dipolar spin system LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> (where the Ho moments interact via dipolar interactions , and the $`\mathrm{\Delta }`$ are induced by a transverse field). Until now most theory has ignored the environment in this problem, but from our discussion above, this is clearly a mistake if one wishes to discuss dynamics. This is also the view expressed in recent papers of Cugliandolo et al. , who have included coupling to an oscillator bath environment. This is presumably correct in metallic glasses, where the coupling to electrons dominates (compare the discussion following (8)), but in insulators the coupling to nuclear spins dominates (in Li<sub>1-x</sub>Ho<sub>x</sub>F<sub>4</sub> this is clear- the hyperfine coupling even modifies the phase diagram!). An observation of strong hole-digging in $`M(\xi ,t)`$ in a spin glass would be consistent with this, since oscillator baths typically give most rapid relaxation away from resonance ($`cf.`$ discussion after eqtn. (82). Note also that one of the main theoretical questions concerning quantum spin glasses, viz., what happens when $`\mathrm{\Delta }J_{ij}`$, could be examined experimentally using molecular magnets like Fe-8, or the LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> system, in a strong transverse field, where one might reasonably expect to see coherent propagation of tunneling events from one site to another if $`\mathrm{\Delta }`$ is sufficiently large.
### B Conductors and Superconductors
In mesoscopic conductors the standard weak localisation theory evaluates a ”decoherence time” $`\tau _\varphi (T)`$ in terms of the electron-electron scattering rate, which is itself strongly influenced at low $`T`$ by elastic impurity scattering. This is essentially an oscillator bath problem- the relevant oscillators being diffusons, Cooperons, and phonons. Curiously, given the known importance of scattering off dynamic ”2-level” fluctuators in these systems , there has been little theory on the effect of these on $`\tau _\varphi `$, apart from the pioneering work of Altshuler and Spivak, and Feng et al. . This is of course a spin bath problem, with the spins representing defects, paramagnetic impurities, etc., in the environment. It will be interesting to see if the ”saturation” in $`\tau _\varphi (T)`$ reported at low $`T`$ may at least be partially explained by such scattering. The physics of this saturation depends on the more fundamental question of how decoherence behaves in the low-$`T`$ limit in conductors, and has caused considerable debate in the recent theoretical literature .
In the case of superconductors the situation is similar, in that almost all theoretical work on dissipation and other environmental effects has looked at the effects of electronic quasiparticle modes or photons, ie., delocalised modes which can be mapped directly to oscillator modes. The pioneering papers led to a massive subsequent literature, both experimental and theoretical . In most work on tunneling there is no question that theory and experiment correspond very well . However the situation is more delicate for coherence, discussed below. In spite of the extensive theory of spin impurity effects in superconductors we are aware of only 2 theoretical papers (and no experiments) examining their effect on the flux dynamics (in particular tunneling) of SQUIDs.
### C Coherence and decoherence; and ”qubits”
An understanding of decoherence mechanisms is central to the exploitation of mesoscopic systems in quantum devices, as well as to general questions about how quantum mechanics applies on the large scale, , and the quantum measurement problem . It has become particularly important now that efforts are being made to construct ”qubit” devices, with a view to making quantum computers.
The last 15 years have seen a total transformation in how such questions are discussed- instead of vague analyses in terms of ”measurements” by the environment, we now have precise and generally applicable models, which can be tested in many experiments. The basic issues are (i) whether phase coherence can be preserved in the reduced density matrix of the system of interest and (ii) what are the decoherence mechanisms destroying it. Here we briefly review studies of superconducting and magnetic systems, and then examine things from a more general theoretical standpoint.
(i) Decoherence in superconductors: This has been discussed intensively ever since the theoretical predictions of Leggett et al. concerning ”macroscopic quantum coherence” in SQUIDs, and subsequent proposals for experimental searches . To date no experimental success has been reported (although there is good evidence for resonant ”one passage” tunneling transitions between near degenerate levels in 2 wells ). Almost all microscopic analyses of this problem have assumed environments of electronic excitations which can be mapped onto oscillator baths (see, eg., ). In our opinion, as discussed in section 4.B, the basic problem is simply that the main source of decoherence in most systems (including SQUIDs) at low $`T`$ will not be any oscillator bath, but the spin bath of paramagnetic and nuclear spins. As discussed in section 2.G and 4.B, the low-energy scale of this spin bath means it will not usually have a big effect on SQUID tunneling, but its effect on macroscopic coherence or on superconducting qubits will be rather large.
Although a superconducting qubit has not yet been built, experiments may be getting rather close . To see how big spin bath effects on coherence might be, let us recall that the effect of paramagnetic impurities is to create a Gaussian multiplet of spin bath states of width $`E_o\mu _BB_o\sqrt{N_{pm}}`$ in energy for a SQUID containing a total of $`N_{pm}`$ paramagnetic impurities interacting with the supercurrent, where $`B_o`$ is the change in field on each paramagnetic impurity caused by the change in flux state of the SQUID. To see coherence it is necessary, from the discussion of section 4.B, that $`\mathrm{\Delta }_oE_o`$, because the decoherence time coming from the spin bath is $`\tau _\varphi \mathrm{\Delta }_o/E_o^2`$; this essentially sets a lower bound for $`\mathrm{\Delta }_o`$. As discussed in some detail in a recent paper , this turns out to be a rather stringent requirement on real SQUIDs; in fact in the experiments of the Lukens group one infers a value $`E_o0.4K`$ from their resonant linewidths. Obviously this value could be reduced a great deal by careful attention to the nuclear and paramgnetic spin impurity composition in the system (as well as to the sample geometry).
(ii) Decoherence in Magnets: The most dramatic claims for the observation of macroscopic coherence have been made by Awschalom et al. , working on randomly oriented dilute ensembles of ferritin macromolecules (which order antiferromagnetically, but carry an excess moment of somewhat random size; the antiferromagnetic ”Néel” moment is $`23,000\mu _B`$). In an effort to make the molecule size as uniform as possible, these authors filtered them magnetically. They also artificially engineered molecules of smaller size. The essential result was the observation of an absorption peak at $`MHz`$ frequencies, whose frequency varies approximately exponentially with the size of the molecules. This was interpreted as a signature of coherent tunneling between ”up” and ”down” states of the Neel vector. There have been widespread objections to this interpretation, both on theoretical and experimental grounds , and so far no other group has succeeded in confirming the experiments. Note that the Awschalom group saw similar resonances (also with an exponential dependence of resonant frequency on size) in large FeCo<sub>5</sub> particles, but did not attribute this to tunneling .
As we discussed in the previous sub-section, there is now very extensive evidence that nanomagnetic molecules in macroscopically ordered crystals tunnel incoherently in the low-$`T`$ quantum regime . There is thus an apparent contradiction between the ferritin work and that done in the Mn-12 and Fe-8 systems (particularly since the ferritin molecules are much larger and certainly contain a lot of spin disorder). Thus in the very well-characterised $`Fe`$-8 molecular crystals used by the Florence and Grenoble groups, the parameter $`\kappa `$ characterising decoherence from the orthogonality blocking mechanism varies , even in an ideal sample, between $`\kappa 615`$ in zero applied field (depending on the annealing-dependent spread $`W_D`$ in the intermolecular dipolar fields), to $`\kappa 80`$ when $`H_x0.2T`$ (where the first zero in $`\mathrm{\Delta }`$ is supposed to occur ). Recalling from sections 3.C and 4.A that coherence is practically eliminated unless $`\kappa 1`$, it is hard to see how experiments on these particular molecules in low fields will stand much chance of seeing it.
On the other hand it is clear that future experiments on single nanomagnets in the quantum regime might have a chance of seeing coherence iff one could raise $`\mathrm{\Delta }_o`$ to values $``$ the hyperfine couplings $`\omega _k`$ (presumably using a large external transverse field), thereby making $`\kappa 1`$ and so removing decoherence (and also reducing the problem to a straightforward spin-boson model- see section 4.B). Another possibility, which could be realised in, eg., the LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> system at high transverse fields, would be to see coherent propagation of spin flips (ie., spin waves) in a lattice of spins, by making $`\mathrm{\Delta }>W_D`$ (here one could also make $`\omega _kkT`$, thereby freezing the nuclear dynamics!). Conceivably the same could be done in an $`Fe`$-8 crystal (now with the inequality $`\mathrm{\Delta }>W_D\omega _k`$ operating). We see no reason why such experiments could not be done in the next few years.
(iii) Decoherence as $`T0`$: Let us now consider the general question of decoherence effects at low $`T`$. Decoherence is often (particularly in conductors) characterised by a ”decoherence time” $`\tau _\varphi `$, for the phase dynamics of the degree of freedom of interest. If and when $`\tau _\varphi `$ is meaningful, it may be much shorter than the energy relaxation time $`\tau _E`$ (cf. the example of a single oscillator coupled to an oscillator bath , or the examples of topological decoherence given for the spin bath in sections 3.A and 4.A above)). Coherence exists if $`\tau _\varphi \mathrm{\Delta }_o1`$, where $`\mathrm{\Delta }_o`$ is the characteristic frequency of the system’s phase dynamics.
Notice that what allows us to discuss this problem with any generality at all is the assumption, discussed in sections 1 and 2, that a few canonical models describe the low-$`T`$ behaviour of most physical systems. Extensive study of the relevant canonical oscillator bath models (in particular, the spin-boson model and the ”oscillator on oscillators” model ) show that with a power-law form $`J(\omega )\omega ^n`$, decoherence disappears as $`T=0`$ for $`n>1`$; for the Ohmic form $`J(\omega )=\pi \alpha \omega `$ decoherence is finite at $`T=0`$, but can be made small if $`\alpha 1`$. If the electronic spectrum is gapped the Ohmic dissipation falls off exponentially in the low $`T`$ limit (thus for superconductors one has $`J(\omega ,T)\omega e^{\mathrm{\Delta }_{BCS}/kT}`$ for $`\omega <2\mathrm{\Delta }_{BCS}`$, and for magnetic solitons one has $`J(\omega ,T)\omega (kT/\mathrm{\Delta }_m)e^{\mathrm{\Delta }_m/kT}`$ for $`\omega <3\mathrm{\Delta }_m`$, where $`\mathrm{\Delta }_{BCS}`$ and $`\mathrm{\Delta }_m`$ are the BCS and magnon gaps respectively). Thus, if one believes the oscillator bath models, coherence ought to be easily observable at low $`T`$; the condition $`\tau _\varphi (T)\mathrm{\Delta }_o1`$ is clearly satisfied for temperatures well below the gap energy.
If we examine the canonical spin bath models we find a very different story . Consider first the central spin model as $`T0`$; we will go to such a low temperature that all of the spins in the spin bath order in the field of the central spin (ie., $`T1\mu K`$ in some cases), ie., all intrinsic fluctuational dynamics of the bath is frozen out. Does decoherence disappear? No, because the mechanisms of topological decoherence (induced bath spin flip) and orthogonality blocking (precession of the bath spins in between central spin flips) still exist- the bath can still acquire dynamics from the central spin. We emphasize here that this physics cannot be described by an oscillator bath model. From the results in section 4 we see there is a residual constant decoherence as $`T0`$, coming from the spin bath, which in an experiment would be signalled by a saturation of $`\tau _\varphi `$ once oscillator bath effects had disappeared. The extent of this decoherence is characterised by the behaviour of $`P_{M=0}(t)`$, (cf. section 4.A), and we saw that unless both $`\kappa `$ and $`\lambda `$ were $`1`$, decoherence was strong. In the case where all bath spins are polarised by the central spin, transitions are blocked anyway- there are no spins in the $`M=0`$ polarisation group!
We conclude that for any system described by the central spin model (ie., where the central system reduces to a 2-level system at low energies), a general consequence of the coupling to a spin bath will be a loss of coherence, via either the topological decoherence or orthogonality blocking mechanisms, even in the $`T0`$ limit. The residual coherence (if any) will depend on the strength of the couplings to the spin bath, in a way discussed quantitatively in sections 4 and 5.
We may generalise these considerations to models in which a ”particle” moves through a spin bath (the same model also describes a network of spins, or of mesoscopic superconductors, etc., coupled to a spin bath), and get the same result. Consider, eg., a particle hopping from site to site on a $`D`$-dimensional hypercubic lattice, whilst coupled to a spin bath , and described by a Hamiltonian
$`H^{latt}(\mathrm{\Omega }_o)`$ $`=`$ $`\mathrm{\Delta }_o\{{\displaystyle \underset{<ij>}{}}[c_i^{}c_j\mathrm{cos}[\mathrm{\Phi }_{ij}+{\displaystyle \underset{k}{}}\stackrel{}{V}_k^{ij}.\stackrel{}{\widehat{\sigma }}_k]+H.c.]`$ (111)
$`+`$ $`{\displaystyle \underset{k}{}}[^{}\omega _k^{ij}(c_i^{}c_ic_j^{}c_j)\stackrel{}{\widehat{\sigma }}_k^z+\omega _k^{}\widehat{\sigma }_k^x]\}+{\displaystyle }_{k,k^{}}V_{kk^{}}^{\alpha \beta }\widehat{\sigma }_k^\alpha \widehat{\sigma }_k^{}^\beta ;`$ (112)
where $`<ij>`$ sums over nearest neighbour sites. The couplings $`{}_{}{}^{}\omega _{k}^{ij}`$ and $`\stackrel{}{V}_k^{ij}`$ to the bath spin $`\stackrel{}{\sigma }_k`$ usually depend upon which ”lattice rung” $`(i,j)`$ the particle happens to be, because the coupling normally has a finite range and the spin $`\stackrel{}{\sigma }_k`$ has some position with respect to the lattice (one can also add a longitudinal coupling to the $`\{\widehat{\sigma }_k^z\}`$ on each site, as in eqtn. (37)). In the continuum limit (112) reduces to (11) and (12), after dropping the external field $`\stackrel{}{h}_k`$ and the dependence of $`F_k^{}(P,Q)`$ on $`Q`$. To isolate the decoherence effects let us assume $`V_{kk^{}}=0`$, ie., we now study the analogue of the strong-coupling regime in section 4.A. Without loss of generality we may then concentrate on the ”phase decoherence” regime, ie., on an effective Hamiltonian $`H_{eff}=_Mw_MH_M^{eff}`$, where
$$H_M^{eff}=\mathrm{\Delta }_o\underset{<ij>}{}[c_i^{}c_j\widehat{𝒫}_Me^{i_k\alpha _k\widehat{\sigma }_k}\widehat{𝒫}_M+H.c.]$$
(113)
This is just a generalisation of (88) to the lattice; the dependence of the coupling on the lattice position is dropped because it is inessential to what follows.
Coherence, if it exists, will appear in the function $`P_{n0}(t)`$ (the probability to start at site $`0`$, and be at site $`n`$ a time $`t`$ later). Perfect coherence (ie., no bath) yields
$$P_{\stackrel{}{n}0}^{(0)}(t)=\underset{\stackrel{}{p}\stackrel{}{p}^{}}{}e^{i[(\stackrel{}{p}\stackrel{}{p}^{}).\stackrel{}{n}(E_\stackrel{}{p}E_\stackrel{}{p}^{})t]}=\underset{\mu =1}{\overset{D}{}}J_{n_\mu }^2(2\mathrm{\Delta }_ot)$$
(114)
for which $`P_{00}^{(0)}(t)1/(\mathrm{\Delta }_ot)^D`$ at long times; moreover, the 2nd moment $`(n(t)^2=_\stackrel{}{n}n^2P_{\stackrel{}{n}0}^{(0)}(t)D(\mathrm{\Delta }_ot)^2`$ at long times. Both are characteristic of coherent ”ballistic” band motion.
If we now go to the interacting case one finds, by a straightforward generalisation of the calculations in sections 3 and 4, some rather interesting results . Consider first a bath with all spin states equally populated. Then at long times $`(n(t)^2D(\mathrm{\Delta }_ot)^2`$ but $`P_{00}^{(0)}(t)(1/\mathrm{\Delta }_ot)`$, independent of $`D`$! Moreover, if we start with an initial Gaussian wave-packet of width $`\mathrm{\Delta }n(t=0)=R_o`$, one finds $`P_{00}^{(0)}(t)(1/R_o\mathrm{\Delta }_ot)`$ as $`t\mathrm{}`$. These results show that the naive inference of ballistic propagation from the second moment result is wrong- in reality one has strongly anomalous diffusion (with an energy-dependent diffusion coefficient, demonstrated by the $`R_o`$-dependence of the results). In fact the probability function $`P_{\stackrel{}{n}0}^{R_o}(t)`$ decays like $`R_o/(\mathrm{\Delta }_ot)|\stackrel{}{n}|^{(D1)}`$ at long times (ie., $`\mathrm{\Delta }_otR_o`$) and distances $`|\stackrel{}{n}|R_o`$, out to a ”ballistic” distance $`l(t)\mathrm{\Delta }_ot`$, for dimension $`D>1`$. Thus there is an advance ”ballistic front” which decays in amplitude with increasing distance/time from the origin, but only as a power law; and it is followed by a much larger anomalously diffusive (and of course incoherent) contribution. For sufficiently long times (in fact for $`\lambda \mathrm{\Delta }_ot1`$) this result is independent of $`\lambda `$. Similar results apply for the orthogonality- blocked case as a function of $`\kappa `$.
Now let’s take the limit $`T0`$, meaning that we allow the bath spins to order in the field of the particle. It is clear that if $`\omega _k^{}\mathrm{\Delta }_o`$ for all spins, and provided there are no non-diagonal momentum couplings to the bath, we cannot get decoherence by the same mechanism as above, since there is only one state in the relevant polarisation group- the particle will then move freely without disturbing the spins in any way. However in any other case phase will be exchanged with the bath in the same way as above, with or without dissipation (which will certainly arise in the weak-coupling limit $`\omega _k^{}<\mathrm{\Delta }_o`$).
We therefore conclude that a finite decoherence in the $`T0`$ limit is a generic consequence of the existence of spin bath environments.
(iv) Qubits and Quantum Computation: What is the impact of these results on hopes for quantum computation? A Quantum computer is an information processing device which can be imagined as an assembly of 2-state qubits, these being none other than the spin-boson and Central spin models discussed in this article . Such a computer has yet to be built, and papers and books on this topic tend to divide into 2 classes. The first simply ignores decoherence (apart from occasionally referring to it as the main stumbling block preventing the construction of a quantum computer!), whereas the second regards decoherence as the crucial problem, and either tries to treat it theoretically (eg., ) or maintains that a quantum computer will never be built because of it (eg., ).
The basic problem here is the lack of serious theory on the effects of decoherence, starting from realistic models which can be tested quantitatively by experiment. The analysis of the Central Spin model in sections 3 and 4 can be used for a single qubit- it will be clear that the main task is to reduce diagonal couplings to the spin bath (and also any oscillator bath) as far as possible. However this is only the beginning of the problem- the operation of a quantum computer involves multi-qubit wave-function entanglement, and thus one wants to understand the behaviour of a decoherence time $`\tau _\varphi ^{(M)}(\xi _1,\xi _2,..\xi _M)`$, governing loss of $`M`$-spin phase correlations, in the presence of coupling to spin and oscillator baths. We are aware of no studies of this problem (or even recognition that it is a problem) in the literature. Studies of mutual coherence and decoherence in the problem of 2 spins coupled to an oscillator bath give some feeling for what might happen, as do the studies of lattice systems. However we still have no answer to, for example, the question of how $`\tau _\varphi ^{(M)}`$ behaves for large $`M`$. If it ends up having a generic decrease $`e^{aM}`$ then the whole quantum computing enterprise will be in very serious trouble!
It is clear that in the near future research on quantum computers will have to proceed on both practical designs (with serious attention given to the decoherence characteristics of the relevant materials), and also on general studies of decoherence for systems of coupled qubits, themselves coupled to spin and oscillator baths. This promises to be one of the great challenges in condensed matter physics during the next decade.
## VI Acknowledgements
We would like to thank NSERC and the CIAR in Canada, and the Russian Foundation for Basic Research (grant 97-02-16548) and European Community (grant INTAS-2124), for financial support of our research. We also thank I. Affleck, Y. Aharonov, B. Barbara, J. Clarke, M. Dubé, W. Hardy, A.J.Leggett, T. Ohm, C. Paulsen, R. Penrose, G. Rose, I.Tupitsyn, W.G. Unruh, I.D. Vagner, and W. Wernsdorfer for helpful discussions about the spin bath, going back to 1987.
## VII Note Added in Proof (Feb 22, 2000)
A number of papers touching upon the present subject (particularly on decoherence) have appeared since this article was written.
On the highly controversial question of zero-temperature decoherence in mesoscopic conductors (which according to some authors throws the entire conventional theory of metals into doubt), a large number of suggestions have appeared. A list of them appears in a short experimental review by Mohanty , which concludes that only the original suggestion, of electronic coupling to zero-point fluctuations of the EM field , can explain the experimental saturation of $`\tau _\varphi `$. This conclusion is hotly disputed by other authors, both on experimental and theoretical grounds . The idea that the decoherence might be coming from ”two-level systems” in the sample (ie., from a spin bath) has been explored by Imry et al. and Zawadowski et al. , but rejected by Mohanty, mainly because it would imply very large low-frequency noise levels in the sample. In our opinion this requires further work- some of the decoherence mechanisms discussed in the present article do not appear in these papers, and would not necessarily show up in the low-frequency noise. In any case this controversy shows clearly how the theory of transport, including weak localisation theory, depends crucially on a correct understanding of decoherence mechanisms.
The most interesting recent experimental progress concerning the mechanisms of decoherence may be work in magnetic systems. The most recent work by Wernsdorfer et al. , which looked at the dependence of tunneling rates and ”hole widths” in the resonant quantum relaxation of crystals of $`Fe`$-8 molecules, has been supplemented by further work on the same system . Taken together these experiments give rather strong evidence that the tunneling is mediated and controlled by the nuclear spins in the system. No direct test has yet appeared of the theory of the decoherence coming from these same nuclear spins, however (which depends on calculations of the parameter $`\kappa `$; cf. section 2.6, including Fig. 5, and sections 3.3, 4.1, 5.1, and 5.3, including Fig. 8, and refs. ). Experiments looking for coherence in $`Fe`$-8 have also been done at higher fields ; however we are unable to see why the reported results give evidence either for or against coherence, since they simply show very broad and rather weak peaks in the ESR spectrum as a function of field. We emphasize a point here which has often been made, viz., that any demonstration of coherence requires direct observation of combinations of multi-time correlation functions- even a very sharp peak in, say, the AC absorption is not enough to demonstrate coherence. In this connection the reader is referred to Figs. 7-9 in the present paper, which show peaks in $`\chi ^{\prime \prime }(\omega )`$ even when there is no coherence at all. Several reviews of the many different experiments in tunneling nanomagnets have also recently appeared . One topic not covered in these is the LiHo<sub>x</sub>Y<sub>1-x</sub>F<sub>4</sub> system (section 5.1.3); some interesting new results on this compare the thermal and quantum annealing, showing the efficiency of the latter in quantum optimisation.
More general discussions of coherence and decoherence appearing recently include several theoretical reviews (although these do not really discuss spin bath environments). There has also been interesting experimental work on systems other than superconductors and magnets; see, eg., Wiseman et al. for coupled quantum dots, and Myatt et al. for decoherence in quantum optical systems.
## A Derivations for limiting cases
In this appendix we give the derivations of some key formulae in sections 3 and 4. Instanton methods are used to handle the spin environment. In Appendix A.1 we give details of the fairly trivial calculations required to deal with averaging over bias, topological phase, and fluctuations in bias. Then in Appendix A.2 we discuss the more lengthy derivation of the expression (58) involved in orthogonality blocking.
### 1 Topological phase, bias, and bias fluctuation effects
We wish to evaluate $`P_{11}(t)`$, the probability for the central spin to return to an initial state $`|`$ after time $`t`$, in the presence of a static bias $`\xi `$ and a noisy bias $`ϵ(t)`$. We use ssh physics.ubc.ca standard instanton techniques , because they easily generalise to include the spin bath.
We begin by ignoring the topological phase of the central spin since its effects are trivial to add. Then the amplitude for a 2-level system to flip in a time $`dt`$ is $`i\mathrm{\Delta }_odt`$, and the return probability is given by summing over even numbers of flips:
$$P_{11}^o(t)=\frac{1}{2}\left\{1+\underset{s=0}{\overset{\mathrm{}}{}}\frac{(2i\mathrm{\Delta }_ot)^{2s}}{(2s)!}\right\}=\frac{1}{2}[1+\mathrm{cos}(2\mathrm{\Delta }_ot)],$$
(A1)
Now consider the modification introduced by a longitudinal static bias $`\xi `$. In instanton language we begin with the return amplitude $`A_{11}(t,\xi )`$, and Laplace transform it; using the action $`e^{\pm i\xi dt}`$ in bias $`\xi `$, over time $`dt`$, we get
$$A_{11}(t,\xi )=\underset{n=0}{\overset{\mathrm{}}{}}(i\mathrm{\Delta }_o)^{2n}_0^t𝑑t_{2n}\mathrm{}_0^{t_2}𝑑t_1e^{i(\xi (tt_{2n})\xi (t_{2n}t_{2n1})+\mathrm{}\xi t_1)}=_i\mathrm{}^i\mathrm{}e^{pt}A_{11}^{^{TLS}}(p,\xi ),$$
(A2)
$$A_{11}^{^{TLS}}(p,\xi )=\frac{1}{pi\xi }\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{(i\mathrm{\Delta }_o)^2}{p^2+\xi ^2}\right)^n=\frac{1}{pi\xi }\frac{p^2+\xi ^2}{p^2+E_o^2}.$$
(A3)
where $`E_o=[\xi ^2+\mathrm{\Delta }_o^2]^{1/2}`$ ; then the standard answer for $`P_{11}(t,\xi )`$ is just:
$`P_{11}(t,\xi )`$ $`=`$ $`{\displaystyle _i\mathrm{}^i\mathrm{}}𝑑p_1𝑑p_2{\displaystyle \frac{e^{(p_1+p_2)t}}{(p_1i\xi )(p_2i\xi )}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{(i\mathrm{\Delta }_o)^2}{p_1^2+\xi ^2}}\right)^n\left({\displaystyle \frac{(i\mathrm{\Delta }_o)^2}{p_2^2+\xi ^2}}\right)^m`$ (A4)
$`=`$ $`[1{\displaystyle \frac{\mathrm{\Delta }_o^2}{2E^2}}(1\mathrm{cos}2Et)]=[1{\displaystyle \frac{\mathrm{\Delta }_o^2}{E^2}}\mathrm{sin}^2Et]`$ (A5)
This representation avoids the difficulty in the usual representation coming from the square root $`E=[\xi ^2+\mathrm{\Delta }_o^2]^{1/2}`$ in the cosine.
We now add a fluctuating bias $`ϵ(t)`$ to $`\xi `$, with the correlation $`[ϵ(t)ϵ(t^{})]^2=\mathrm{\Lambda }^3|tt^{}|`$ for short times. Noise averages are then given by the Gaussian average
$$F[ϵ(t)]=𝒟ϵ(t)F[ϵ]e^{\frac{1}{2}{\scriptscriptstyle 𝑑t_1𝑑t_2ϵ(t_1)K(t_1t_2)ϵ(t_2)}}$$
(A6)
with a noise correlator $`2K^1(t_1t_2)=\mathrm{\Lambda }^3(|t_1|+|t_2||t_1t_2|)`$. Averaging over a phase function $`F(t_1,t_2)=e^{i_{t_1}^{t_2}𝑑sϵ(s)}`$ then gives the standard result
$$F(t_1,t_2)=e^{iϵ(t_1)(t_2t_1)}e^{\mathrm{\Lambda }^3(t_2t_1)^3/6},$$
(A7)
In applying this to (A2) we assume fast diffusion (see text), and thus only expand to $`O(\mathrm{\Delta }^2)`$. This gives $`P_{11}(t)=e^{t/\tau (\xi )}`$, with
$$\tau ^1(\xi )=\mathrm{\hspace{0.33em}2}\mathrm{\Delta }^2𝑑ϵG_\mu (ϵ)_0^{\mathrm{}}𝑑se^{i(ϵ+\xi )s}e^{\mathrm{\Lambda }^3s^3/6}=\mathrm{\hspace{0.33em}\hspace{0.33em}2}\pi ^{1/2}\frac{\mathrm{\Delta }^2}{\mathrm{\Gamma }}e^{(\xi /\mathrm{\Gamma })^2}$$
(A8)
where $`G(ϵ)=(2/\pi \mathrm{\Gamma }^2)^{1/2}e^{2(ϵ/\mathrm{\Gamma })^2}`$ is the probability $`ϵ(t)`$ takes value $`ϵ`$ (ie., it is the lineshape of the polarisation group).
Adding the the topological phase $`\mathrm{\Phi }_o=\pi S`$ to these calculations changes the flip amplitude to $`i\mathrm{\Delta }_o\mathrm{exp}\{\pm i\mathrm{\Phi }_o\}dt`$; one then sums over all paths with an even number of flips and over all combinations of $`\pm `$ (clockwise and counterclockwise) flips; thus (A1) becomes
$$P_{11}^{(0)}(t)=\frac{1}{2}\left\{1+\underset{s=0}{\overset{\mathrm{}}{}}\frac{(2i\mathrm{\Delta }_ot)^{2s}}{(2s)!}\underset{n=0}{\overset{2s}{}}\frac{(2s)!}{(2sn)!n!}e^{i\mathrm{\Phi }_o(2s2n)}\right\}=\frac{1}{2}[1+\mathrm{cos}(4\mathrm{\Delta }_o\mathrm{cos}\mathrm{\Phi }_o)t],$$
(A9)
with similar obvious modifications to include bias. The generalization to include the phase from the bath spins (topological decoherence) is now obvious for both imaginary and full complex $`\alpha _k`$ (see text, section 3.A).
### 2 Orthogonality blocking effects
We consider the situation described in section IV.B., where the ”initial” and ”final” fields on the bath spin $`\stackrel{}{\sigma }_k`$ are $`\stackrel{}{\gamma }_k^{(1)}`$ and $`\stackrel{}{\gamma }_k^{(2)}`$, related by an angle $`\beta _k`$, which is assumed small, and is defined by $`\mathrm{cos}2\beta _k=\stackrel{}{\gamma }_k^{(1)}\stackrel{}{\gamma }_k^{(2)}/|\stackrel{}{\gamma }_k^{(1)}||\stackrel{}{\gamma }_k^{(2)}|`$. We choose axes in spin space such that the initial and final spin bath wave-functions are related by
$$\{\stackrel{}{\sigma }_k^f\}=\underset{k=1}{\overset{N}{}}\widehat{U}_k\{\stackrel{}{\sigma }_k^{in}\}=\widehat{U}\{\stackrel{}{\sigma }_k^{in}\}.$$
(A10)
where $`\widehat{U}_k=e^{i\beta _k\widehat{\sigma }_k^x}`$ (compare eqtn. (56).
In general the initial spin bath state will belong to some polarisation group $`M_o`$ (not necessarily $`M_o=0`$), where the polarisation is defined along some direction defined by the central system (for example, in a nanomagnetic problem, one could define it as the direction of initial orientation of the nanomagnetic spin). As explained in the text, during a central system transition energy conservation requires the polarisation to change from $`M_o`$ to $`M_o2M`$ (and back, for further transitions); for ”pure” orthogonality blocking (ie., when no other terms are involved in the Hamiltonian), $`M_o=M`$. In what follows we calculate the correlation function $`P_{M_o,M}(t)`$, the central spin correlator defined under the restriction that the spin bath transitions are between subspaces defined by $`\widehat{𝒫}_{k=1}^N\widehat{\sigma }_k^z=M_o`$ and $`\widehat{𝒫}=M_o2M`$ subspaces, which are supposed to be in resonance. The statistical weight of states with $`M_o>N^{1/2}`$ is negligible, so we will assume that $`M_o,M<N`$.
We enforce the restriction to a polarisation group $`M`$ using the projection operator
$$\widehat{\mathrm{\Pi }}_M=\delta (\underset{k=1}{\overset{N}{}}\widehat{\sigma }_k^zM)=_0^{2\pi }\frac{d\xi }{2\pi }e^{i\xi (_{k=1}^N\widehat{\sigma }_k^zM)}.$$
(A11)
We can now write down an expression for the amplitude (not the probability!) $`A_{M_o,M}^{11}(t)`$ for the central spin $`\stackrel{}{\tau }`$ to stay in state $`|`$ during a time $`t`$:
$$A_{M_o,M}^{11}(t)=\left\{\underset{n=0}{\overset{\mathrm{}}{}}\frac{(i\mathrm{\Delta }_o(\mathrm{\Phi })t)^{2n}}{(2n)!}\underset{i=1}{\overset{2n}{}}\frac{d\xi _i}{2\pi }e^{iM_o(\xi _{2n}+\xi _{2n1}+\mathrm{}+\xi _1)}e^{2iM(\xi _{2n1}+\xi _{2n3}+\mathrm{}+\xi _1)}\widehat{T}_{2n}\right\}\{\stackrel{}{\sigma }_k^{in}\},$$
(A12)
where $`\widehat{T}_{2n}`$ is
$$\widehat{T}_{2n}=\left[e^{i\xi _{2n}_{k=1}^N\widehat{\sigma }_k^z}\widehat{U}^{}e^{i\xi _{2n1}_{k=1}^N\widehat{\sigma }_k^z}\widehat{U}\mathrm{}\widehat{U}^{}e^{i\xi _1_{k=1}^N\widehat{\sigma }_k^z}\widehat{U}\right].$$
(A13)
¿From (A12) we can now write the full correlation function $`P_{M_o,M}^{11}(t)`$ as
$`P_{M_o,M}(t)R_{M_o,M}^{}(t)R_{M_o,M}(t)`$ (A14)
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\mathrm{\Delta }_o(\mathrm{\Phi })t)^{2(n+m)}}{(2n)!(2m)!}}{\displaystyle \underset{i=1}{\overset{2n}{}}}{\displaystyle \underset{j=1}{\overset{2m}{}}}{\displaystyle \frac{d\xi _i}{2\pi }\frac{d\xi _j^{}}{2\pi }e^{iM_o(_i^{2n}\xi _i_j^{2m}\xi _j^{})}e^{2iM(_{i=odd}^{2n1}\xi _i_{j=odd}^{2m1}\xi _j^{})}\widehat{T}_{2m}^{}\widehat{T}_{2n}}.`$ (A15)
We now use the assumption that the $`\beta _k`$ are small; more precisely we assume that the orthogonality exponent $`\kappa `$, defined previously by $`e^\kappa =\mathrm{cos}\beta _k`$ (cf. eqtn (60)), can be approximated by the perturbative expansion $`\kappa \frac{1}{2}\beta _k^2`$. This assumption makes it much easier to calculate the average in (A15). We consider first the problem with only one environmental spin $`\stackrel{}{\sigma }_k`$, and calculate the average $`\widehat{T}_{2m}^{}\widehat{T}_{2n}_k`$ in this case; since $`\widehat{T}_{2m}^{}\widehat{T}_{2n}`$ is a product of operators acting separately on each $`\stackrel{}{\sigma }_k`$, the average over all spins is also the product of single spin results.
We only need consider processes with $`0,\mathrm{\hspace{0.33em}1}`$, or $`2`$ flips of the environmental spin, i.e., we expand in powers of $`\beta _k`$, and stop at $`\beta _k^2`$. Then it is clear that, if the initial state of $`\stackrel{}{\sigma }_k`$ is $`_k`$
$`\widehat{T}_{2n}^{(k)}_k`$ $`=`$ $`e^{i\xi _{2n}\widehat{\sigma }_k^z}e^{i\beta _k\widehat{\sigma }_k^x}\mathrm{}e^{i\beta _k\widehat{\sigma }_k^x}e^{i\xi _1\widehat{\sigma }_k^z}e^{i\beta _k\widehat{\sigma }_k^x}_k`$ (A16)
$`=`$ $`e^{i_{i=1}^{2n}\xi _i}[(1n\beta _k^2)_k+i\beta _k_k{\displaystyle }_{l=1}^{2n}(1)^{l+1}e^{2i_{i=l}^{2n}\xi _i}`$ (A17)
$`\beta _k^2_k{\displaystyle \underset{l^{}=l+1}{\overset{2n}{}}}{\displaystyle \underset{l=1}{\overset{2n1}{}}}(1)^{l^{}l}e^{2i_{i=l}^{l^{}1}\xi _i}+O(\beta _k^3)],`$ (A18)
where the first term arises from the sequence $`[11\mathrm{}11]`$, the second from the sequence $`[11\mathrm{}\mathrm{}]`$, with a flip when $`j=l`$; and so on. In the same way we find
$`_k(\widehat{T}_{2m}^{(k)})^{}\widehat{T}_{2n}^{(k)}_k`$ $`=`$ $`e^{i(_{i=1}^{2n}\xi _i_{j=1}^{2m}\xi _j^{})}[1\beta _k^2[(n+m)+{\displaystyle \underset{l^{}=l+1}{\overset{2n}{}}}{\displaystyle \underset{l=1}{\overset{2n1}{}}}(1)^{l^{}l}e^{2i_{i=l}^{l^{}1}\xi _i}`$ (A19)
$`+`$ $`{\displaystyle \underset{p^{}=p+1}{\overset{2m}{}}}{\displaystyle \underset{p=1}{\overset{2m1}{}}}(1)^{p^{}p}e^{2i_{j=p}^{p^{}1}\xi _j^{}}`$ (A20)
$``$ $`{\displaystyle \underset{p=1}{\overset{2m}{}}}{\displaystyle \underset{l=1}{\overset{2n}{}}}(1)^{l+p}e^{2i(_{i=l}^{2n}\xi _i_{j=1}^{p1}\xi _j^{})}]],`$ (A21)
to order $`\beta _k^2`$. The sequence $`_k(\widehat{T}_{2m}^{(k)})^{}\widehat{T}_{2n}^{(k)}_k`$ will have a similar expression, but with reversed signs coming from the $`e^{i\xi _j\widehat{\sigma }_k^z}`$ factors.
We now observe that the state with polarisation $`M_o`$ consists of $`(N+M_o)/2`$ spins up and $`(NM_o)/2`$ spins down. Consequently, for each $`\stackrel{}{\sigma }_k`$, we add $``$ or $``$ averages like (A21), and then take the product
$$\widehat{T}_{2m}^{}\widehat{T}_{2n}=\underset{k=1}{\overset{N_{}}{}}(\widehat{T}_{2m}^{(k)})^{}\widehat{T}_{2n}^{(k)}\underset{k^{}=1}{\overset{N_{}}{}}(\widehat{T}_{2m}^{(k)})^{}\widehat{T}_{2n}^{(k)}$$
(A22)
Substituting (A21) into this expression we get
$$\widehat{T}_{2m}^{}\widehat{T}_{2n}=e^{iM_o(_i^{2n}\xi _i_j^{2m}\xi _j^{})}\mathrm{exp}\left\{K_{nm}^{eff}(\xi _i,\xi _j^{},M_o)\right\},$$
(A23)
where the ”effective action” $`K_{nm}^{eff}(\xi _i,\xi _j,M_o)`$ has two contributions $`K^{eff}=K_1+K_2`$:
$`K_1=2\kappa (1{\displaystyle \frac{M_o}{N}})\{(n+m)`$ $`+`$ $`{\displaystyle \underset{l^{}>l}{}}(1)^{l^{}l}\mathrm{cos}[2{\displaystyle \underset{i=l}{\overset{l^{}1}{}}}\xi _i]+{\displaystyle \underset{p^{}>p}{}}(1)^{p^{}p}\mathrm{cos}[2{\displaystyle \underset{j=p}{\overset{p^{}1}{}}}\xi _j^{}]`$ (A24)
$``$ $`{\displaystyle \underset{p=1}{\overset{2m}{}}}{\displaystyle \underset{l=1}{\overset{2n}{}}}(1)^{l+p}\mathrm{cos}[2{\displaystyle \underset{i=l}{\overset{2n}{}}}\xi _i{\displaystyle \underset{j=1}{\overset{p1}{}}}\xi _j^{}]\},`$ (A25)
$`K_2=2\kappa {\displaystyle \frac{M_o}{N}}\{`$ $`{\displaystyle \underset{l^{}>l}{}}(1)^{l^{}l}\mathrm{exp}[2i{\displaystyle \underset{i=l}{\overset{l^{}1}{}}}\xi _i]+{\displaystyle \underset{p^{}>p}{}}(1)^{p^{}p}\mathrm{exp}[2i{\displaystyle \underset{j=p}{\overset{p^{}1}{}}}\xi _j^{}]`$ (A26)
$``$ $`{\displaystyle \underset{p=1}{\overset{2m}{}}}{\displaystyle \underset{l=1}{\overset{2n}{}}}(1)^{l+p}\mathrm{exp}[2i{\displaystyle \underset{j=1}{\overset{p1}{}}}\xi _j^{}2i{\displaystyle \underset{i=l}{\overset{2n}{}}}\xi _i]\},`$ (A27)
We recall now that $`M_oN^{1/2}N`$, which allows us to neglect the contribution due to $`K_2`$ and drop the correction $`M_o/N`$ to the coefficient $`\kappa `$ in $`K_1`$. Notice also that the phase factor in front of $`\mathrm{exp}\{K^{eff}\}`$ in (A22) cancels exactly the phase proportional to $`M_o`$ in the formula (A15) for $`P_{M_o,M}(t)`$. Thus, quite surprisingly, we find the correlation function to be independent of $`M_o`$ in this limit:
$$P_M(t)=\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{\mathrm{}}{}}\frac{(i\mathrm{\Delta }_o(\mathrm{\Phi })t)^{2(n+m)}}{(2n)!(2m)!}\underset{i=1}{\overset{2n}{}}\underset{j=1}{\overset{2m}{}}\frac{d\xi _i}{2\pi }\frac{d\xi _j^{}}{2\pi }\mathrm{exp}\left\{2iM(\xi _{2n1}+\xi _{2n3}+\mathrm{}+\xi _1)K_{nm}^{eff}(\xi _i,\xi _j^{})\right\},$$
(A28)
We can render this expression more useful by changing variables; first we consider the whole sequence $`\xi _\alpha =(\xi _1,\mathrm{},\xi _{2n},\xi _1^{},\mathrm{},\xi _{2m}^{})`$ together, and then define new angular variables
$$\chi _\alpha =\underset{\alpha ^{}=\alpha }{\overset{2(n+m)}{}}2\xi _\alpha ^{}+\pi \alpha ,$$
(A29)
so that now
$`P_M(t)=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(i\mathrm{\Delta }_o(\mathrm{\Phi })t)^{2(n+m)}}{(2n)!(2m)!}}\left({\displaystyle \underset{\alpha =1}{\overset{2(n+m)}{}}}{\displaystyle \frac{d\chi _\alpha }{2\pi }}\right)`$ (A30)
$`\times `$ $`\mathrm{exp}\left\{iM{\displaystyle \underset{\alpha }{}}(1)^{\alpha +1}\chi _\alpha 2\kappa \left[(n+m)+{\displaystyle \underset{\alpha ^{}>\alpha }{}}\mathrm{cos}(\chi _\alpha \chi _\alpha ^{})\right]\right\}.`$ (A31)
Thus we have mapped our problem onto the partition function of a rather peculiar system of spins, interacting via infinite range forces, with interaction strength $`2\kappa `$.
To deal with this partition function , we define ”pseudo-spins” $`\stackrel{}{s}_\alpha =(\mathrm{cos}\chi _\alpha ,\mathrm{sin}\chi _\alpha )`$ and $`\stackrel{}{𝒮}`$, such that
$$\stackrel{}{𝒮}=\underset{\alpha =1}{\overset{2(n+m)}{}}\stackrel{}{s}_\alpha ,\underset{\alpha ^{},\alpha }{}\mathrm{cos}(\chi _\alpha \chi _\alpha ^{})=\stackrel{}{𝒮}^2,$$
(A32)
We can think of $`\stackrel{}{s}_\alpha `$ as rotating in our fictitious angular space defined by the projection operator (A11). Now consider the term $`G(\stackrel{}{𝒮})`$ in (A31) defined by
$`G(\stackrel{}{𝒮})`$ $`=`$ $`\left({\displaystyle \underset{\alpha =1}{\overset{2(n+m)}{}}}{\displaystyle \frac{d\chi _\alpha }{2\pi }e^{iM(1)^{\alpha +1}\chi _\alpha }}\right)\mathrm{exp}\left\{\kappa {\displaystyle \underset{\alpha ^{},\alpha }{}}\mathrm{cos}(\chi _\alpha \chi _\alpha ^{})\right\}`$ (A33)
$`=`$ $`\left({\displaystyle \underset{\alpha =1}{\overset{2(n+m)}{}}}{\displaystyle \frac{d\chi _\alpha }{2\pi }e^{iM(1)^{\alpha +1}\chi _\alpha }}\right)e^{\kappa \stackrel{}{𝒮}^2}.`$ (A34)
This is easily calculated, viz.,
$`G(\stackrel{}{𝒮})`$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{𝒮}e^{\kappa \stackrel{}{𝒮}^2}\underset{\alpha =1}{\overset{2(n+m)}{}}\frac{d\chi _\alpha }{2\pi }e^{iM(1)^{\alpha +1}\chi _\alpha }\delta (\stackrel{}{𝒮}\underset{\alpha }{}\stackrel{}{s}_\alpha )}`$ (A35)
$`=`$ $`{\displaystyle \frac{d\stackrel{}{z}}{2\pi }𝑑\stackrel{}{𝒮}e^{\kappa \stackrel{}{𝒮}^2+i\stackrel{}{z}\stackrel{}{𝒮}}\left(_0^{2\pi }\frac{d\chi _\alpha }{2\pi }e^{i\stackrel{}{z}\stackrel{}{s}_\alpha +iM\chi _\alpha }\right)^{2(n+m)}}`$ (A36)
$`=`$ $`{\displaystyle \frac{1}{2\kappa }}{\displaystyle 𝑑zze^{z^2/4\kappa }J_M^{2(n+m)}(z)},`$ (A37)
where $`J_M(\lambda )`$ is the $`M`$th-order Bessel function. Using
$$\underset{l=0}{\overset{\mathrm{}}{}}\frac{\delta _{2(n+m),2s}}{(2m)!(2n)!}=\frac{\delta _{s,0}+2^{2s}}{2(2s)!},$$
(A38)
to reorganize the sum over $`n`$ and $`m`$ in (A31) and changing the integration variable $`z2x\sqrt{\kappa }`$, we then find
$`P_M(t)`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑xxe^{x^2}{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{[2it\mathrm{\Delta }_o(\mathrm{\Phi })J_M(2x\sqrt{\kappa })]^{2s}}{(2s)!}}\right)`$ (A39)
$`=`$ $`P_M(t)={\displaystyle _0^{\mathrm{}}}𝑑xxe^{x^2}\left(1+\mathrm{cos}[2\mathrm{\Delta }_o(\mathrm{\Phi })J_M(2x\sqrt{\kappa })t]\right)2{\displaystyle _0^{\mathrm{}}}𝑑xxe^{x^2}P_{11}^{(0)}(t,\mathrm{\Delta }_M(x)).`$ (A40)
$$\mathrm{\Delta }_M(x)=\mathrm{\Delta }_o(\mathrm{\Phi })J_M(2x\sqrt{\kappa }).$$
(A41)
Here we come to the crucial point in our derivation. Eq.(A40) gives the final answer as a superposition of non-interacting correlation functions for effective tunneling rates $`\mathrm{\Delta }_M(x)`$ with the proper weighting. For $`M=0`$ this is the form quoted in Eq.(58) of the text.
It is worth noting that non-zero $`M`$ enters this calculation as the overall phase factor which we can follow from (A15) up to (A37), where we finally integrate over $`\{\chi _\alpha \}`$ to produce the Bessel function of order $`M`$. This observation allows one to generalise any calculation done for $`M=0`$ to finite $`M`$ by simply replacing $`J_0J_M`$ in the final answer - the prescription which we use in other Appendices.
## B Derivations for the generic case
We outline here the derivations for section 4, in which topological decoherence, degeneracy blocking, and orthogonality blocking are all simultaneously incorporated (the average over bias fluctuations being essentially trivial- see Appendix A.1). We have demonstrated in section 3 and Appendix A how each different term in the effective Hamiltonian (17) influence the central spin dynamics. From these limiting cases we learned that static (or diagonal) terms in the Hamiltonian can be partly absorbed into a redefinition of the transition amplitude between states with equal initial and final energies. If we now deal with the full central spin Hamiltonian, we can still write the instanton expansion in central spin transitions in the form (see (A5)):
$$P_M(t)=_i\mathrm{}^i\mathrm{}𝑑p_1𝑑p_2\frac{e^{(p_1+p_2)t}}{(p_1iϵ)(p_2iϵ)}\underset{n=0}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{\mathrm{}}{}}\left(\frac{(i\mathrm{\Delta }_o)^2}{p_1^2+ϵ^2}\right)^n\left(\frac{(i\mathrm{\Delta }_o)^2}{p_2^2+ϵ^2}\right)^mB_{nm}(M),$$
(B1)
$$B_{nm}(M)=\underset{\{g_l=\pm \}}{}e^{i\mathrm{\Phi }_{l=1}^{2(n+m)}g_l}\widehat{T}_{2m}^{}(M,g_l)\widehat{T}_{2n}(M,g_l),$$
(B2)
where the sum over $`\{g_l=\pm \}`$ with $`1l2(n+m)`$ describes all possible clockwise and anticlockwise transitions, and the operator product is defined as
$$\widehat{T}_{2n}(M,g_l)=\widehat{U}_M(g_1)\widehat{U}_M^{}(g_2)\widehat{U}_M(g_3)\mathrm{}\widehat{U}_M^{}(g_{2n}),$$
(B3)
$$\widehat{U}_M(g)=\widehat{\mathrm{\Pi }}_Me^{ig_k\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k}e^{i\beta _k\widehat{\sigma }_k^x}\widehat{\mathrm{\Pi }}_M.$$
(B4)
Here $`\widehat{\mathrm{\Pi }}_M`$ as before projects on the polarisation state $`\mathrm{\Delta }N=M`$, and $`\beta _k`$ describes the mismatch between the initial and final nuclear spin states. If all the couplings were equal ($`\omega _k=\omega _o\mathrm{\Delta }_o`$) then the above set of equations would be the complete solution of the $`M`$ polarisation group dynamics. One may then further average over different grains in the ensemble by summing over different polarisation groups with proper weigthing. If there is a small spread in the nuclear hyperfine couplings, it will produce an internal bias field acting on the grain, as described in the text (section IV.C. The final answer for the ensemble of grains is obtained then by averaging (B1) over the bias field. This bias field is due to all environemntal spins interacting with $`\stackrel{}{S}`$, and hardly changes when a few nuclei (of order $`\lambda `$) flip with $`\stackrel{}{S}`$. For this reason there is no back influence of the induced nuclear spin flips on the bias field, at least during the time scale set by the damping of coherent oscillations (when many environemntal spins are flipped the coherence is obviously already lost).
The crucial observation is that if the sum over the clockwise and anticlockwise trajectories and the average of the operator product in (B2) can be presented as some weighted average and/or sum, of form
$$B_{nm}(M)=𝑑x_1𝑑x_2\mathrm{}𝑑x_a\underset{k_1k_2\mathrm{}k_b}{}Z_M(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b)R_M^{2(n+m)}(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b),$$
(B5)
with fixed integer values $`a`$ and $`b`$, then the problem may be considered as solved because the instanton summation then reduces to that of a coherent (non-interacting) dynamics with the renormalized tunneling amplitude
$$\mathrm{\Delta }_o\mathrm{\Delta }_M(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b)=\mathrm{\Delta }_oR_M(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b).$$
(B6)
and the final answer acquires a form
$$P_M(t,ϵ)=dx_1dx_2\mathrm{}dx_a\underset{k_1k_2\mathrm{}k_b}{}Z_M(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b)P_{11}^{TLS}(t,ϵ,\mathrm{\Delta }_M(x_1,\mathrm{},x_a;k_1,\mathrm{},k_b),$$
(B7)
where $`P_{11}^{TLS}(t,ϵ,\mathrm{\Delta }_M)`$ is described by Eq. (A5).
We have already seen that Eq. (B5) is indeed valid for the cases of pure topological decoherence and pure orthogonality blocking- we now prove that it also holds when we combine the effects of topological decoherence with the projection on a given polarisation state. Here we evaluate the $`M=0`$ contribution; the result for $`P_M(t)`$ then follows from the generalisation explained at the end of the previous Appendix.
Introducing as before the spectral representation for the projection operator \[see Eq. (A11)\] we write
$$B_{nm}(M)=\underset{\{g_l=\pm \}}{}e^{i\mathrm{\Phi }_{l=1}^{2(n+m)}g_l}\underset{\rho =1}{\overset{2(n+m)}{}}\frac{d\xi _\rho }{2\pi }\mathrm{exp}\left\{K_{nm}^{eff}(\{g_l\},\{\xi _\rho \})\right\}.$$
(B8)
$$\mathrm{exp}\left\{K_{nm}^{eff}(\{g_l\},\{\xi _\rho \})\right\}=e^{i\xi _1\widehat{𝒫}}e^{ig_1_k\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k}e^{i\xi _2\widehat{𝒫}}e^{ig_2_k\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k}\mathrm{}e^{i\xi _{2(n+m)}\widehat{𝒫}}e^{ig_{2(n+m)}_k\alpha _k\stackrel{}{n}_k\widehat{\stackrel{}{\sigma }}_k}.$$
(B9)
With the usual assumption that the individual $`\alpha _k`$ are small (but not necessarily $`\lambda `$), the ”effective action” $`K_{nm}^{eff}`$ can be written as (compare Eq.(A31))
$$K_{nm}^{eff}(\{\xi _\rho \})=\lambda ^{}\underset{\rho ^{},\rho }{\overset{2(n+m)}{}}g_\rho g_\rho ^{}+\lambda \underset{\rho ^{},\rho }{\overset{2(n+m)}{}}\mathrm{cos}(\chi _\rho \chi _\rho ^{})g_\rho g_\rho ^{},$$
(B10)
which generalizes from orthogonality blocking; the $`\chi _\rho `$ are defined as in (A29), and
$$\lambda =\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k^2(1(n_k^z)^2);\lambda ^{}=\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k^2(n_k^z)^2.$$
(B11)
as before. We use the same trick of introducing “pseudo-vectors” $`\stackrel{}{s}_\rho =(\mathrm{cos}\chi _\rho ,\mathrm{sin}\chi _\rho )`$, and $`\stackrel{}{𝒮}^2=\left(_{\rho =1}^{2(n+m)}g_\rho \stackrel{}{s}_\rho \right)^2=_{\rho ^{},\rho }g_\rho g_\rho ^{}\mathrm{cos}(\chi _\rho \chi _\rho ^{})`$, to decouple integrals over the new variables $`\chi _\rho =\chi _\rho +\pi g_\rho /2`$. After some lengthy, but straightforward algebra we get $`P_{11}(t)`$ in the form
$$P_0(t)=2𝑑xxe^{x^2}\frac{d\phi }{2\pi }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}F_\lambda ^{}(m)e^{i2m(\mathrm{\Phi }\phi )}P_{11}^{(0)}(t,\mathrm{\Delta }_o(\phi ,x))$$
(B12)
where $`\mathrm{\Delta }_o(\phi ,x)=2\mathrm{\Delta }_o\mathrm{cos}\phi J_0(2x\sqrt{\lambda })`$ as before.
The case of complex $`\alpha _k`$ (even assuming $`\omega _k^{}=0`$ in the effective Hamiltonian) is more subtle technically, but goes through in exactly the same way. Here we just outline the key steps; a more detailed derivation may be found in . The effective action now has the form
$`K_{nm}^{eff}`$ $`=`$ $`{\displaystyle \underset{\rho ,\rho ^{}=1}{\overset{2(n+m)}{}}}g_\rho g_\rho ^{}\{[\lambda ^{}\eta ^{}(1)^{\rho +\rho ^{}}i\gamma ^{}((1)^\rho +(1)^\rho ^{})]`$ (B13)
$`+`$ $`\mathrm{cos}(\chi _\rho \chi _\rho ^{})[(\lambda \lambda ^{})(1)^{\rho +\rho ^{}}(\eta \eta ^{})i(\gamma \gamma ^{})((1)^\rho +(1)^\rho ^{})]\},`$ (B14)
where the constants are defined by:
$$\lambda =\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k^2;\lambda ^{}=\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k^2(n_k^z)^2;$$
(B15)
$$\eta =\frac{1}{2}\underset{k=1}{\overset{N}{}}\xi _k^2;\eta ^{}=\frac{1}{2}\underset{k=1}{\overset{N}{}}\xi _k^2(v_k^z)^2;$$
(B16)
$$\gamma =\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k\xi _k\stackrel{}{n}_k\stackrel{}{v}_k;\gamma ^{}=\frac{1}{2}\underset{k=1}{\overset{N}{}}\alpha _k\xi _kn_k^zv_k^z;$$
(B17)
As before we change variables according to $`\chi _\rho =\chi _\rho +\pi `$ when $`g_\rho =1`$, to introduce odd and even spin fields
$$\stackrel{}{𝒮}_o=\underset{\rho =odd}{\overset{2(n+m)1}{}}\stackrel{}{s}(\chi _\rho );\stackrel{}{𝒮}_e=\underset{\rho =even}{\overset{2(n+m)}{}}\stackrel{}{s}(\chi _\rho ).$$
(B18)
which are used to decouple integrations over $`\chi _\rho `$ with the final goal to get the answer in the form of Eq. (B5). This indeed can be done, and the final answer reads
$`P_{11}(t)={\displaystyle \frac{d\phi _1}{2\pi }\frac{d\phi _2}{2\pi }}`$ $`{\displaystyle \underset{m_1=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m_2=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle 𝑑x_1𝑑x_2}`$ (B19)
$`\times `$ $`𝒵(\phi _1,\phi _2,x_1,x_2,m_1,m_2)P_{11}^{(0)}[t,\stackrel{~}{\mathrm{\Delta }}_o(x_1,x_2,\phi _1,\phi _2)],`$ (B20)
and (B20) has an obvious generalisation to include the bias integration $`𝑑ϵ`$. The weight is given by
$`𝒵`$ $`=`$ $`e^{2i[m_1(\mathrm{\Phi }\phi _1)m_2\phi _2+4m_1m_2\gamma ^{}]}e^{4(\eta ^{}m_2^2\lambda ^{}m_1^2)}`$ (B21)
$`\times `$ $`{\displaystyle \frac{x_1x_2}{8(abc^2)}}I_0\left({\displaystyle \frac{(a+b)x_1x_2}{8(abc^2)}}\right)\mathrm{exp}\left\{{\displaystyle \frac{(ab+2ic)x_1^2+(ab2ic)x_2^2}{16(abc^2)}}\right\},`$ (B22)
and the renormalized tunneling splitting equals
$$\stackrel{~}{\mathrm{\Delta }}_o^2(x_1,x_2,\phi _1,\phi _2)=4\stackrel{~}{\mathrm{\Delta }}_o^2\mathrm{cos}(\phi _1+\phi _2)\mathrm{cos}(\phi _1\phi _2)J_0(x_1)J_0(x_2),$$
(B23)
where
$$a=\lambda \lambda ^{};b=\eta \eta ^{};c=\gamma \gamma ^{}.$$
(B24)
FIGURE CAPTIONS
Figure 1 The flow of a class of effective Hamiltonians describing a central system coupled to a background environment, in coupling constant space, as the UV cutoff in the joint Hilbert space is reduced from $`E_c`$ to $`\mathrm{\Omega }_o`$. Here we show flow to a fixed point FP, in a simplified 2-dimensional space of couplings $`\alpha _1,\alpha _2`$, but one may also have fixed lines or more complex topologies.
Figure 2 A typical path for a 2-level central system (solid line) coupled to environmental modes (wavy lines) as a function of time, showing the couplings which exist in both the spin-boson and central spin models. We show both diagonal couplings D to $`\tau _z`$ and non-diagonal couplings ND to $`\tau _\pm `$ (in the central spin model these are strong enough to lead to multiple excitation of environmental modes).
Figure 3 Definition of the longitudinal and transverse parts of the diagonal coupling to a bath spin in the Central Spin Hamiltonian, in terms of the initial and final fields $`\stackrel{}{\gamma }^{(1)}`$ and $`\stackrel{}{\gamma }^{(2)}`$ acting on this spin- this also defines the angle $`\beta `$, and the mutually perpendicular unit vectors $`\widehat{l}`$ and $`\widehat{m}`$ (see text).
Figure 4 Classifying the states of the Central spin Hamiltonian. Each level of $`\stackrel{}{\tau }`$ is associated with a $`2^N`$-fold multiplet of bath states (Fig 6(a)). These are classified into polarisation groups $`\{M\}`$ (where $`M`$ is the total polarisation along $`\widehat{z}`$), separated by energy $`\omega _o`$ and with width $`\stackrel{~}{\mathrm{\Gamma }}_M`$; Fig 6(b) shows the density of states $`G_M(\xi )`$ of the separate groups, and Fig. 6(c) their sum $`W(\xi )`$. We show $`W(\xi )`$ for 2 different values of the parameter $`\mu =\stackrel{~}{\mathrm{\Gamma }}_M/\omega _o`$; in realistic cases $`\mu 1`$ (ie., the polarisation groups strongly overlap), and $`W(\xi )`$ is Gaussian. Longitudinal transitions between 2 different polarisation groups $`M_1`$ and $`M_2`$ go at a rate $`T_1^1`$; transitions within a polarisation group at a rate $`T_2^1`$.
Figure 5 Example of the application of the Central Spin model to a magnetic macromolecule (the $`Fe`$-8 molecule, further descibed in section 5). In (a) we show the effective tunneling matrix element $`|\stackrel{~}{\mathrm{\Delta }}_{eff}|=|\mathrm{\Delta }\mathrm{cos}(\pi S+i\beta _o𝐧_o.𝐇_o|)`$, for this easy axis/easy plane nanomagnet in the presence of a field $`𝐇_o=\widehat{x}H_x`$ in the $`x`$-direction (transverse to the easy axis), assuming an angle $`\phi `$ between $`\widehat{x}`$ and the magnetic ”hard axis” (perpendicular to the easy plane). Aharonov-Bohm oscillations appear when $`\phi `$ is small, so that the action of the 2 relevant paths on the spin sphere have similar magnitudes, but almost opposite phase. For larger $`\phi `$, one path dominates over the other and oscillations are suppressed. In (b) we show a histogram of the $`\omega _k^{}`$ for this system- the main figure shows the protons and the lower inset the $`N`$ and $`O`$ contributions. The upper inset in (b) shows the variation of $`E_o`$ and $`\xi _o`$ with $`H_x`$ (the parameter $`\xi _o`$ is discussed in sections 4 and 5). These figures are adapted from Ref. .
Figure 6 Behaviour of $`P_{11}(t)`$ in the case of pure topological decoherence. We show $`P_{11}(t)1/2`$ for intermediate coupling, for which $`P_{11}(t)`$ takes the ”universal form” discussed in the text.
Figure 7 The effect of relaxation on a statistical ensemble of central spins, each interacting with a spin bath. In (a) we assume that $`\lambda =0,\kappa =5`$ and $`N=1000`$, and show the normalised time dependence of 3 different contributions $`P_M(t)`$ to the total relaxation function $`P_{11}(t)`$; they sum to give a roughly logarithmic time dependence for the total function $`P_{11}(t)`$. The small $`M`$ contributions relax quickly (up to $`M\kappa `$), so the effect on an initial ensemble distributed over bias $`\xi `$, at short times, is to dig a hole around zero bias, of width $`\kappa \omega _o`$. In (b) we show the spectral absorption function $`\chi ^{\prime \prime }(\omega )`$ for $`\kappa =2`$, dividing this into the $`M=0`$ contribution and the contributions from higher $`M`$ groups (which relax more slowly and thus peak at lower $`\omega `$).
Figure 8 The spectral absorption function $`\chi _{M=0}^{\prime \prime }(\omega )`$ for several values of $`\kappa `$, for an ensemble of central spins in the $`M=0`$ polarisation group, in the case where orthogonality blocking dominates, and degeneracy blocking effects (ie., a bias average) are also incorporated. Contributions from higher polarisation groups $`M0`$ are not shown; they are spread over a range $`M\omega _o`$, up to $`\kappa \omega _o`$. Contributions from groups with $`M>\kappa `$ are negligible.
Figure 9 Graphs of $`\chi _{M=0}^{\prime \prime }(\omega )`$ for ”projected topological decoherence” (ie., including a bias average over an ensemble in which topological decoherence dominates), for several different values of the parameter $`(\lambda \lambda ^{})`$. Contributions from higher polarisation groups, which are spread over an energy range $`\lambda \omega _o`$, are not shown. |
warning/0001/gr-qc0001024.html | ar5iv | text | # The inverse scattering method for cylindrical gravitational waves
## 1 Introduction
In this paper we study the cylindrical gravitational waves model with two polarization modes. The space-time of the model possesses two commuting one-parameter isometry groups of which one is isomorphic to $`R`$ and the other to $`SO(2)`$, and is non-stationary solution of the vacuum Einstein equations. The generators of these groups are called Killing vector fields. In the case of one polarization (both Killing fields are hypersurface orthogonal) this model is also known as Einstein-Rosen waves. The theory of Einstein-Rosen waves is linear, however, still physically interesting and even in recent time provides us with new important results both in classical and quantum scope Ref. where one finds additional discussion and reference list. The polarized model is non-linear and hence far more difficult to study.
In two Killing reduction of vacuum or electro-vacuum general relativity Einstein equations belong to the class of integrable equations. By the other words they can be written as a compatibility condition of an auxiliary linear system. Although in this paper we use the linear system proposed in Ref. it is worth mentioning that there exits essentially different approach Ref..
The best tool to study integrable equations is the Inverse Scattering Method (ISM). There are many works where the ISM is applied for constructing exact solutions and there are a few papers where the physically relevant models are considered. The author knows only Ref. where the model of colliding gravitational plane waves was investigated and Rev. where the solutions with disconnected horizon were analyzed.
The main goal of the present paper is to develop a framework to study the cylindrical gravitational waves model.
## 2 Field equations and boundery conditions
The metric of the space-time with two commuting space-like Killing fields can always be chosen locally in the following form
$$ds^2=ds_{II}^2+Xd\varphi ^2+2Wd\varphi dz+Vdz^2,$$
(1)
where
$$ds_{II}^2=\gamma _{ab}dx_adx_b$$
is the metric of an orthogonal to both Killing vectors two-surface. The variables $`X,W,V`$ depend only on the coordinates $`x_a`$ adapted to this surface.
Let
$$g=\left(\begin{array}{cc}V& W\\ W& X\end{array}\right),\rho ^2=detg.$$
Then, the first group of Einstein equations can be written as
$$d\rho dgg^1=0.$$
(2)
Here $``$ is the Hodge operator with respect to the metric $`ds_{II}^2`$. The function $`\rho `$ has a space-like gradient hence it can be used as a radial coordinate. Let the time coordinate be dual to space one, $`dt=d\rho `$. Coordinate chart constructed is often called the Weyl canonical coordinates. We employ it throughout the paper. In these coordinates the 2-metric $`ds_{II}^2`$ is given by
$$ds_{II}^2=f(\rho ,t)\left(dt^2+d\rho ^2\right),f(\rho ,t)=\frac{X}{\rho ^2}e^{2\gamma }$$
(3)
The last formula is a definition of the function $`\gamma `$. Then the system (2) takes the form
$$(\rho g_{,t}g^1)_{,t}+(\rho g_{,\rho }g^1)_{,\rho }=0.$$
(4)
We impose the following boundary conditions on $`g`$
$$g=\left(\begin{array}{cc}V& \rho ^2\widehat{W}\\ \rho ^2\widehat{W}& \rho ^2\widehat{X}\end{array}\right)$$
(5)
where $`V,\widehat{W}`$ and $`\widehat{X}`$ are smooth functions not equal to zero and
$$V=1+J/\rho +O(1/\rho ^2),\widehat{X}=1J/\rho +O(1/\rho ^2),\widehat{W}=O(1/\rho ^3)$$
at $`\rho =\mathrm{}`$. Here $`J`$ is a constant independent on $`t`$.
An alternative formulation of Einstein equations (2) exists. For this, let us introduce matrix potential
$$dH=\rho dgg^1,Y=H_{12}.$$
(6)
Here $`Y`$ is the Ernst potential that is determined by the transition Killing vector field. The functions $`Y`$ and $`V`$ are independent dynamical variables which set a solution of (2) completely. Potentials $`Y`$ and $`V`$ satisfy the Ernst equation which we will not use in this paper.
Using (5) one can easily prove that
$$\rho g_{,\rho }g^1=\left(\begin{array}{cc}0& _tY\\ 0& 2\end{array}\right),\rho g_{,t}g^1=0$$
(7)
at $`\rho =0`$ and
$$\rho g_{,\rho }g^1\left(\begin{array}{cc}0& 0\\ 0& 2\end{array}\right)=\left(\begin{array}{cc}O(1/\rho )& O(1/\rho ^3)\\ O(1/\rho )& O(1/\rho )\end{array}\right),$$
(8)
$$\rho g_{,t}g^1=\left(\begin{array}{cc}O(1/\rho )& O(1/\rho ^3)\\ O(1/\rho )& O(1/\rho )\end{array}\right),$$
(9)
at $`\rho =\mathrm{}`$.
The second group of Einstein equations allows one to determine the coefficient $`\gamma `$ from the matrix $`g`$. From (5) it is possible to show that $`_t\gamma =0`$ as $`\rho =0`$. The solution is regular at the symmetry axis iff $`\gamma =0`$ for $`\rho =0`$. However, in this case $`\gamma \gamma _00`$ as $`\rho \mathrm{}`$ and the conical singularity appears at the space infinity. Its angle of deficit can be treated as the energy of the system . It is interesting to note that the deficit angle of the conical singularity at the symmetry axis of the stationary solutions with disconnected horizon has the sense of the force between the black holes (see Ref. and references therein). In the present paper, we restrict ourselves to the study of the system (4).
## 3 Auxiliary linear problem
System of equations (4) is the compatibility condition for the following pair of matrix linear differential equations:
$$D_1\mathrm{\Psi }=U(\omega ,\rho ,t)\mathrm{\Psi },U(\omega ,\rho ,t)=\frac{\rho ^2g_{,\rho }g^1+\omega \rho g_{,t}g^1}{\rho ^2\omega ^2},$$
(10)
$$D_2\mathrm{\Psi }=V(\omega ,\rho ,t)\mathrm{\Psi },V(\omega ,\rho ,t)=\frac{\rho ^2g_{,t}g^1+\omega \rho g_{,\rho }g^1}{\rho ^2\omega ^2},$$
(11)
Here $`D_1`$ and $`D_2`$ are the commuting differential operators
$$D_1=_\rho \frac{2\omega \rho }{\omega ^2\rho ^2}_\omega ,D_2=_t\frac{2\omega ^2}{\omega ^2\rho ^2}_\omega ,$$
(12)
and $`\omega `$ is a complex parameter that does not depend on $`\rho ,t`$. We also use the $`UV`$-pair representation in which $`\omega `$ is a dependent parameter. To be more precise, let $`\omega `$ be a root of the equation
$$\omega ^2+2(tk)\omega +\rho ^2=0$$
(13)
where $`k`$ is an independent parameter. Determined through (13), $`\omega `$ satisfies
$$_\rho \omega =\frac{2\omega \rho }{\rho ^2\omega ^2},_t\omega =\frac{2\omega ^2}{\rho ^2\omega ^2}.$$
(14)
Passing from $`\mathrm{\Psi }(\omega )`$ to $`\mathrm{\Psi }^{}(k)=\mathrm{\Psi }(\omega (k))`$, one obtains from (10, 11) that
$$_\rho \mathrm{\Psi }^{}=U(\omega (k),\rho ,t)\mathrm{\Psi }^{},_t\mathrm{\Psi }^{}=V(\omega (k),\rho ,t)\mathrm{\Psi }^{}.$$
(15)
We will denote the solution of (15) by $`\mathrm{\Psi }(k)`$ missing the prime for brevity. It is worth mentioning that solving (15) with the fixed branch of the root one finds the solution of (10, 11) only in the analyticity domain of $`\omega (k)`$.
We will follow the general scheme for investigating integrable equations Ref. omitting many technical details.
Eq. (13) is invariant w. r. t. the transformation $`\omega \rho ^2/\omega `$. Now suppose $`\omega _\pm (k)`$ are root’s branches such that $`\mathrm{Im}\omega _+>0`$ and $`\mathrm{Im}\omega _{}<0`$. Then the cuts of $`\omega _\pm `$ are half-lines $`(\mathrm{},t\rho ]`$ and $`[t+\rho ,\mathrm{})`$. Before proceed we mention some useful properties of $`\omega _\pm `$. Note that
$$\omega _+(k)=\frac{\rho ^2}{\omega _{}(k)},\overline{\omega }_+(\overline{k})=\omega _{}(k).$$
(16)
As $`\rho \mathrm{}`$ the functions $`\omega _\pm (k,\rho ,t)`$ satisfy uniformly in $`k,t`$ the following estimate
$$\omega _\pm (k,\rho ,t)=\pm i\rho +(kt)+O(1/\rho ).$$
(17)
Furthermore, introduce the monodromy matrixes $`T_\pm (k,\rho ,\rho ^{})`$. By definition, they are solutions to
$$_\rho T_\pm (\rho ,\rho ^{})=U(\omega _\pm (k),\rho ,t)T_\pm (\rho ,\rho ^{}),T_\pm (\rho ,\rho )=I.$$
(18)
Besides, zero caveture condition,
$$U_{,t}V_{,\rho }+[U,V]=0,$$
reveals that
$$_tT_\pm (\rho ,\rho ^{})=V(\omega _\pm (k),\rho )T_\pm (\rho ,\rho ^{})T_\pm (\rho ,\rho ^{})V(\omega _\pm (k),\rho ^{}).$$
(19)
For symmetric real matrix $`g`$ the system (10, 11) remains invariant under the transformations
$$\mathrm{\Psi }(\omega ,\rho ,t)g\stackrel{~}{\mathrm{\Psi }}^1(\frac{\rho ^2}{\omega },\rho ,t),\mathrm{\Psi }(\omega ,\rho ,t)\overline{\mathrm{\Psi }}(\overline{\omega },\rho ,t)$$
(20)
where the tilde denotes transposition. Taking into account (16) one gets
$$T_+(k,\rho ,\rho ^{})=g(\rho )\stackrel{~}{T}_{}^1(k,\rho ,\rho ^{})g^1(\rho ^{}),T_+(k,\rho ,\rho ^{})=\overline{T}_{}(\overline{k},\rho ,\rho ^{})$$
(21)
Note that we regard the parameter $`k`$ as a parameter on the complex plane, not on the Riemann surface of the root, and understand equations (18) as two different ones. Let us remark also that $`T_+(\omega ,\rho ,\rho ^{})`$ and $`T_{}(\omega ,\rho ,\rho ^{})`$ are solutions of (10) as $`\mathrm{Im}\omega 0`$ and $`\mathrm{Im}\omega 0`$ respectively.
Let us determine the Jost functions,
$$\mathrm{\Psi }_\pm (k,\rho )=\underset{\rho ^{}\mathrm{}}{lim}T_\pm (k,\rho ,\rho ^{})e_\pm (k,\rho ^{}),e_\pm (k,\rho ^{})=\left(\begin{array}{cc}1& 0\\ 0& \omega _\pm (k,\rho ^{})\end{array}\right).$$
(22)
They are analytic in the $`k`$-plane with the cuts $`(\mathrm{},t\rho ]`$ and $`[t+\rho ,\mathrm{})`$. Moreover, from equation (19) we derive that they are solutions of compatible system of equations (15). Since $`lim_\rho \mathrm{}e_{}(\rho )g^1(\rho )e_+(\rho )=I`$ we see also that
$$\mathrm{\Psi }_+(k)=g\stackrel{~}{\mathrm{\Psi }}_{}^1(k),\mathrm{\Psi }_+(k)=\overline{\mathrm{\Psi }}_{}(\overline{k}).$$
(23)
It is worth mentioning that $`det\mathrm{\Psi }_\pm (k)=\omega _\pm (k)`$.
Recall that functions $`\mathrm{\Psi }_+(\omega )`$ and $`\mathrm{\Psi }_{}(\omega )`$ are analytic solutions of (10), (11) as $`\text{Im}\omega >0`$ and as $`\text{Im}\omega <0`$ respectively. Hence for real $`\omega `$ there is a matrix $`T(k)`$ such that
$$\mathrm{\Psi }_+(\omega )=\mathrm{\Psi }_{}(\omega )T(k),k=t+\frac{\omega ^2+\rho ^2}{2\omega }.$$
(24)
Since $`\mathrm{\Psi }_\pm (\omega )`$ are coupled by transformation (20) and $`det\mathrm{\Psi }_\pm (\omega )=\omega `$, $`T(k)`$ satisfies
$$T(k)=\stackrel{~}{T}(k),\overline{T}(\overline{k})=T^1(k),detT(k)=1,$$
(25)
by the other words it is a symmetric unitary matrix. It is interesting to note that if initial data $`(V1,\widehat{X}1,V_{,t},\widehat{X}_{,t})`$ are of compact support then $`TI`$ is also of compact support. In general case we have $`T(k)=I+O(1/|k|)`$ as $`|k|\mathrm{}`$.
Let
$$\chi _\pm (\omega )=\mathrm{\Psi }_\pm (\omega )e^1(\omega ),e(\omega )=\left(\begin{array}{cc}1& 0\\ 0& \omega \end{array}\right).$$
(26)
It is possible to prove that $`\chi _\pm `$ are regular at $`\omega =0`$ and $`\chi _\pm (\mathrm{})=I`$. The last condition and symmetry (20) reveal that
$$g=\chi _\pm (0)\left(\begin{array}{cc}1& 0\\ 0& \rho ^2\end{array}\right)$$
(27)
Since $`\chi _{}(0)=\chi _+(0)`$ from (24) we derive that $`T_{12}(k)=T_{21}(k)=O(1/|k|^2)`$ at $`k=\mathrm{}`$.
Boundary condition (7) yields
$$V(\omega )|_{\rho =0}=\frac{1}{\omega }\left(\begin{array}{cc}0& _tY\\ 0& 2\end{array}\right)$$
Solving (11) one obtains
$$\mathrm{\Psi }_\pm (\omega )|_{\rho =0}=\left(\begin{array}{cc}1& Y(t)\\ 0& \omega \end{array}\right)T_\pm (k),T_{}(k)=\overline{T}_+(\overline{k})$$
(28)
Here $`T_+`$ is analytic in $`k`$ as $`\text{Im}k>0`$. Inserting (28) in (24) we come to identities:
$$T_{}^1(k)T_+(k)=T(k)=\stackrel{~}{T}_+(k)\stackrel{~}{T}_{}^1(k)$$
(29)
It follows from normalization of $`\chi _\pm `$ at $`\omega =\mathrm{}`$ that
$$\underset{\omega \mathrm{}}{lim}e(\omega )T_+(k)e^1(\omega )=I$$
(30)
As a function of $`\omega `$, $`T_+(k(\omega ))`$ is analytic in domains: $`\text{Im}\omega >0,|\omega |>\rho `$ and $`\text{Im}\omega <0,|\omega |<\rho `$. Then using (24) and (29) one can easily check that
$$\mathrm{\Phi }_{}(\omega )=\{\begin{array}{cc}\mathrm{\Psi }_+(\omega )T_+^1(k)\text{Im}\omega >0,|\omega |>\rho & \\ \mathrm{\Psi }_{}(\omega )T_{}^1(k)\text{Im}\omega <0,|\omega |>\rho & \end{array}$$
(31)
is analytic as $`|\omega |>\rho `$ while
$$\mathrm{\Phi }_+(\omega )=\{\begin{array}{cc}\mathrm{\Psi }_+(\omega )\stackrel{~}{T}_{}(k)\text{Im}\omega >0,|\omega |<\rho & \\ \mathrm{\Psi }_{}(\omega )\stackrel{~}{T}_+(k)\text{Im}\omega <0,|\omega |<\rho & \end{array}$$
(32)
is analytic as $`|\omega |<\rho `$. On the circle $`|\omega |=\rho `$, $`\mathrm{\Phi }_\pm `$ satisfy the conjugation condition,
$$\mathrm{\Phi }_+(\omega )=\mathrm{\Phi }_{}(\omega )G(k),G(k)=T_+(k)\stackrel{~}{T}_{}(k)=T_{}(k)\stackrel{~}{T}_+(k).$$
(33)
Set $`\widehat{\chi }_\pm =\mathrm{\Phi }_\pm e^1`$. The matrices $`\mathrm{\Phi }_+`$ and $`\mathrm{\Phi }_{}`$ are still coupled by transformation (20). Using it and the fact that $`\widehat{\chi }_{}(\mathrm{})=I`$ which follows from (30), we come to
$$g=\widehat{\chi }_+(0)\left(\begin{array}{cc}1& 0\\ 0& \rho ^2\end{array}\right)$$
(34)
Due to the boundary condition (5) the limit, $`lim_{\rho 0}\widehat{\chi }_+(0)`$, exists. On the other hand we have from (33) that
$$\underset{\rho 0}{lim}\widehat{\chi }_+(0)=\underset{\omega 0}{lim}\left(\begin{array}{cc}1& Y(t)\\ 0& \omega \end{array}\right)G(k)e^1(\omega ).$$
The above limit exits if and only if $`Y(t)=G_{12}(t)/G_{22}(t)`$. Eventually,
$$\left(\begin{array}{cc}V(t)& \widehat{W}(t)\\ 0& \widehat{X}(t)\end{array}\right)=\left(\begin{array}{cc}\frac{1}{G_{22}(t)}& \frac{2_tG_{12}(t)G_{22}(t)2G_{12}(t)_tG_{22}(t)}{G_{22}(t)}\\ 0& G_{22}(t)\end{array}\right)$$
(35)
where unimodularity of $`G(k)`$ was used as well.
We summarize the results of this section as follow. For any solution of the initial-value problem posed in Section 2 there exits analytic in upper $`k`$-plane matrix, $`T_+(k)`$, that satisfies the conditions,
$$T_+T_+^{}=T_+^{}T_+,\text{Im}T_+T_+^{}=0$$
for real $`k`$. Here $``$ is the Hermitian conjugate. It uniquely determines symmetric unitary matrix, $`T(k)`$, and symmetric real matrix, $`G(k)`$. Then the inverse problem is reduced to the Riemann-Hilbert problem (RHP) (24) or (33). In addition, there is one to one correspondence between $`G(k)`$ and the data on the symmetry axis, for example $`Y(t)`$ and $`V(t)`$ (35).
## 4 Einstein-Rosen waves
In this section we apply the results of the previous section to a space-time with diagonal metric ($`W=0`$). Set $`X=\rho ^2e^{2\psi }`$. Then the system (4) reduces to one linear equation, viz
$$\frac{^2\psi }{t^2}+\frac{^2\psi }{\rho ^2}+\frac{1}{\rho }\frac{\psi }{\rho }=0$$
(36)
The above equation is the compatibility condition of the following pair of scaler linear equations
$$_\rho F=\frac{2\rho ^2_\rho \psi +2\rho \omega _\pm _t\psi }{\rho ^2\omega _\pm ^2}F,_tF=\frac{2\rho ^2_t\psi +2\rho \omega _\pm _\rho \psi }{\rho ^2\omega _\pm ^2}F.$$
(37)
The solution of the system (15) is given by
$$\mathrm{\Psi }_\pm (k,\rho ,t)=\left(\begin{array}{cc}e^{\theta _\pm (k,\rho ,t)}& 0\\ 0& \omega _\pm e^{\theta _\pm (k,\rho ,t)}\end{array}\right)$$
with
$$\theta _\pm (k,\rho ,t)=\underset{\rho }{\overset{\mathrm{}}{}}𝑑\rho ^{}\left(\frac{2\rho ^2_\rho \psi }{\rho ^2\omega _\pm ^2}+\frac{2\rho ^{}\omega _\pm _t\psi }{\rho ^2\omega _\pm ^2}\right).$$
The function $`k(\omega )=t+\frac{\omega ^2+\rho ^2}{2\omega }`$ maps the upper half-plane(lower half-plane) into $`k`$-plane with the cuts $`(\mathrm{},t\rho ]`$ and $`[t+\rho ,\mathrm{})`$. The function $`\omega _+(k,\rho ^{})(\omega _{}(k,\rho ^{}))`$ maps the $`k`$-plane with the cuts $`(\mathrm{},t\rho ^{}]`$ and $`[t+\rho ^{},\mathrm{})`$ into the upper half-plane(lower half-plane). Analyzing these mappings in the case when $`\rho \rho ^{}`$ we conclude that
$$\omega _+(k(\omega +i0),\rho ^{})\omega _{}(k(\omega i0),\rho ^{})=2i\sqrt{\rho ^2(kt)^2}$$
as $`|kt|\rho ^{}`$ and zero otherwise. Note that $`|kt|\rho `$ for any real $`\omega `$ and $`\omega _\pm (k,\rho ^{})=kt\pm i\sqrt{\rho ^2(kt)^2}`$ as $`\rho |kt|\rho ^{}`$ and real $`k`$. Therefore
$$\theta _+(k(\omega +i0),\rho ,t)\theta _{}(k(\omega i0),\rho ,t)=2i\mathrm{\Delta }(k)$$
(38)
where
$$\mathrm{\Delta }(k)=\underset{|kt|}{\overset{\mathrm{}}{}}𝑑\rho ^{}\frac{(kt)_\rho \psi +\rho ^{}_t\psi }{\sqrt{\rho ^2(kt)^2}}.$$
The piece-wise analytic function $`\theta (\omega )`$, $`\theta (\omega )=\theta _{}(\omega )`$ as $`\text{Im}\omega <0`$ and $`\theta (\omega )=\theta _+(\omega )`$ as $`\text{Im}\omega >0`$, is a unique solution of the RHP (38) with $`\theta (\mathrm{})=0`$. This solution can be written as
$$\theta (\omega )=\frac{1}{\pi }\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\lambda \frac{\mathrm{\Delta }(k(\lambda ))}{\lambda \omega }.$$
(39)
Then we can represent the solution of (36) as
$$\psi (t,\rho )=\frac{1}{2}\theta (0)=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\lambda \frac{\mathrm{\Delta }(k(\lambda ))}{\lambda }.$$
(40)
Changing the variable in the above integral one has
$$\psi (t,\rho )=\frac{1}{\pi }\underset{\mathrm{}}{\overset{t\rho }{}}𝑑k\frac{\mathrm{\Delta }(k)}{\sqrt{(kt)^2\rho ^2}}+\frac{1}{\pi }\underset{t+\rho }{\overset{\mathrm{}}{}}𝑑k\frac{\mathrm{\Delta }(k)}{\sqrt{(kt)^2\rho ^2}}$$
(41)
## 5 Conclusions
In this paper we outline the ISM scheme for the cylindrical symmetric waves model. However, to complete this scheme one needs to analyze the smoothness property of the matrices $`T(k)`$ and $`G(k)`$. We assume that for smooth initial data they are also smooth. The relation of the matrix $`G`$ with the data on the symmetry axis allows to determine the behavior of the solution at the time infinities through the study of the asymptotic properties of $`G`$. The matrix $`T`$ defines the behavior of the solution at the null infinities and we assume that it is uniquely restored from the data on the null infinities. An important open question is the derivation of the formula for the angle of deficit at the space infinity in terms of the matrix $`T`$ or $`G`$. We plan to consider the problems mentioned in separate papers without promising to do it soon.
## Acknowledgements
This work was partly supported by RFBR grant No 98-01-01063. |
warning/0001/quant-ph0001026.html | ar5iv | text | # Generalized Affine Coherent States: A Natural Framework for Quantization of Metric-like Variables
## 1 Introduction
It has been suggested that the affine group may have a significant role to play in a quantum theory of gravity. The (one-dimensional) affine Lie algebra is in some sense the simplest possible non-abelian Lie algebra; it may be considered to be generated by two self adjoint operators $`\kappa `$ and $`\sigma `$ whose commutation relation reads
$`[\sigma ,\kappa ]=i\sigma .`$ (1)
There exists a representation of (1) in which the spectrum of $`\sigma `$ is strictly positive. It has been noted by Klauder that a multidimensional generalization of the algebra (1) exists in which $`\sigma `$ is replaced by a matrix operator corresponding to a symmetric, positive definite matrix degree of freedom. Such an object is clearly well suited to the description of the spacial part of a metric tensor, and attempts have been made to construct an affine field theory in the context of quantum gravity. It has been argued by Isham and Kakas that an affine algebra arises naturally in an attempt to quantize a nonlinear phase space such as that which exists in general relativity. The strong-coupling limit of such a theory of gravity has previously been discussed by Pilati .
The coherent state representation associated with the one-dimensional affine algebra (1) is fairly well known . Our present aim is to generalize this analysis to a multi-dimensional affine algebra. Our discussion is limited largely to kinematics; furthermore, we make no attempt to build a field theory. These questions are considered by Klauder .
## 2 The One-Dimensional Affine Algebra and its Associated Coherent States
The algebra in (1) generates a 2-parameter Lie group known as the affine group (also known as the $`ax+b`$” group , with $`a>0`$) - the group of translations and dilations of the real line. The operators $`\kappa `$ and $`\sigma `$ induce the following transformations:
$`e^{iB\kappa }\sigma e^{iB\kappa }`$ $`=`$ $`e^B\sigma ,`$ (2)
$`e^{iF\sigma }\kappa e^{iF\sigma }`$ $`=`$ $`\kappa +F\sigma .`$ (3)
Here $`B`$ and $`F`$ are real parameters.
It is well known that there exist three, faithful, inequivalent, irreducible unitary representations of the affine algebra, characterized respectively by the operator $`\sigma `$ posessing positive, negative, and null spectra . The positive representation is of particular interest in this article since the positivity of $`\sigma `$ will be generalized to a positive definite matrix degree of freedom (see section 5). In this representation, the operators $`\sigma `$ and $`\kappa `$ are represented by the following operators:
$`\sigma `$ $`=`$ $`k,`$ (4)
$`\kappa `$ $`=`$ $`\frac{1}{2}[\theta k+k\theta ],`$ (5)
with
$`\theta =i{\displaystyle \frac{}{k}},`$ (6)
where the representation space is the space of $`L^2`$ functions on the open interval $`(0,\mathrm{})`$, equipped with the inner product
$`\varphi |\psi `$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\varphi (k)^{}\psi (k)𝑑k,`$ (7)
for any two functions $`\varphi `$ and $`\psi `$ in the space. The operator $`\kappa `$ acts to generate unitary dilations in the representation space:
$`e^{iB\kappa }\psi (k)`$ $`=`$ $`e^{B/2}\psi (e^Bk),`$ (8)
$`e^{iB\kappa }|\psi `$ $`=`$ $`|\psi .`$ (9)
It is important to note that $`\theta `$, unlike its canonical conjugate $`\sigma `$, is not self adjoint in this representation - nor does it possess self-adjoint extensions. Its interpretation as an observable is therefore not possible.
We now define a family of unitary operators $`U(F,B)`$ via
$`U(F,B)=e^{iF\sigma }e^{iB\kappa }.`$ (10)
The composition rule for the operators in (10) is
$`U(F^{},B^{})U(F,B)=U(F^{}+e^B^{}F,B^{}+B).`$ (11)
This family of unitary operators may be used to constuct a set of coherent states,
$`|F,B=U(F,B)|\eta .`$ (12)
Here $`|\eta `$ is an as yet unspecified normalized fiducial vector in the representation space. The coherent states in (12) admit a resolution of unity in the form
$`N^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑F{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑Be^B|F,BF,B|`$ $`=`$ $`11,`$ (13)
the validity of which hinges on the fiducial vector admissibility criterion
$`N`$ $`=`$ $`2\pi {\displaystyle _0^{\mathrm{}}}𝑑kk^1|\eta (k)|^2<\mathrm{}.`$ (14)
For a more general discussion regarding coherent states, see .
An example of a normalized admissible fiducial vector is provided by
$`\eta _1(k)=C_1(\alpha ,\beta )k^\alpha e^{\beta k},`$ (15)
where $`\alpha `$ and $`\beta `$ are positive real coefficients, and the normalization constant $`C_1(\alpha ,\beta )`$ is given by
$`C_1(\alpha ,\beta ){\displaystyle \frac{(2\beta )^{\alpha +\mathrm{\hspace{0.33em}1}/2}}{\sqrt{\mathrm{\Gamma }(2\alpha +1)}}}.`$ (16)
It may be verified that this fiducial vector satisfies (14), with $`N=2\pi \beta /\alpha `$. The associated coherent state representative of a general vector $`|\psi `$ in the Hilbert space, $`F,B|\psi `$, exhibits the property that
$`\left[ie^B{\displaystyle \frac{}{F}}\beta ^1{\displaystyle \frac{}{B}}\beta ^1(\alpha +\frac{1}{2})\right]F,B|\psi =0.`$ (17)
The relation (17) defines a complex polarization of the relevant function space, and its form is inherited directly from the property of the fiducial vector (15) described by
$`\left[\sigma i\beta ^1\kappa \beta ^1(\alpha +\frac{1}{2})\right]|\eta =0.`$ (18)
This property, when re-written in the form
$`(\sigma \sigma )|\eta `$ $`=`$ $`i\beta ^1(\kappa \kappa )|\eta ,`$ (19)
where
$`\sigma `$ $``$ $`\eta |\sigma |\eta =\beta ^1(\alpha +\frac{1}{2}),`$ (20)
$`\kappa `$ $``$ $`\eta |\kappa |\eta =\mathrm{\hspace{0.33em}\hspace{0.33em}0},`$ (21)
reveals that the fiducial vectors in (15) form a continuous family (over the parameters $`\alpha `$ and $`\beta `$) of minimum uncertainty states, satisfying equality in the general affine uncertainty relation
$`\mathrm{\Delta }\sigma ^2\mathrm{\Delta }\kappa ^2`$ $``$ $`\sigma ^2/4.`$ (22)
It is straightforward to show that the associated coherent state overlap function is given by
$`F^{},B^{}|F,B=\left[{\displaystyle \frac{e^{(B^{}+B)/2}}{(e^B^{}+e^B)/2+i(F^{}F)/2\beta }}\right]^{2\alpha +1}.`$ (23)
For future comparison, it is helpful to restate our results in terms of the natural variable $`Ge^B`$. The overlap function then takes the form
$`F^{},G^{}|F,G=\left[{\displaystyle \frac{G_{}^{}{}_{}{}^{1/2}G^{1/2}}{(G_{}^{}{}_{}{}^{1}+G^1)/2+i(F^{}F)/2\beta }}\right]^{2\alpha +1},`$ (24)
while the resolution of unity may be written as an integral over a suitable flat measure:
$`N^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑F{\displaystyle _0^{\mathrm{}}}𝑑G|F,GF,G|=\mathrm{\hspace{0.33em}11}.`$ (25)
The polarization property now appears as
$`\left[iG^1{\displaystyle \frac{}{F}}\beta ^1G{\displaystyle \frac{}{G}}\gamma _1\right]F,G|\psi =\mathrm{\hspace{0.33em}0},`$ (26)
where we have introduced
$`\gamma _1\beta ^1(\alpha +\frac{1}{2}).`$ (27)
## 3 Path Integral for 1-D Coherent States
Before building a path integral for the coherent state overlap, we introduce two relevant geometrical objects - the symplectic potential and the ray metric. Suppose that a set of normalized vectors $`\{|l\}`$ defines a continuous curve in a Hilbert space, the label $`l`$ being a continuous real parameter along the curve. Then the overlap of two nearby vectors $`|ldl/2`$ and $`|l+dl/2`$ may be expanded as
$`l+dl/2|ldl/2=\mathrm{exp}(id\theta d\mathrm{\Sigma }^2/2+\mathrm{})`$ (28)
where $`d\theta `$ and $`d\mathrm{\Sigma }^2`$ are called, respectively, the symplectic potential and the ray metric.
The expression (28) may be used to expand the overlap of two of the coherent states from (12),
$`F+dF/2,B+dB/2|FdF/2,BdB/2`$
$`=\mathrm{exp}\left\{i(\alpha +\frac{1}{2})\beta ^1e^BdF(\alpha +\frac{1}{2})[\beta ^2e^{2B}(dF)^2+(dB)^2]\right\}.`$ (29)
The symplectic potential and the ray metric may now be read off,
$`d\theta `$ $`=`$ $`\gamma _1e^BdF`$ (30)
$`=`$ $`\gamma _1GdF`$
$`d\mathrm{\Sigma }^2`$ $`=`$ $`\gamma _1[\beta ^1e^{2B}(dF)^2+\beta (dB)^2]`$ (31)
$`=`$ $`\gamma _1[\beta ^1G^2(dF)^2+\beta G^2(dG)^2].`$
We now introduce a general time-independent Hamiltonian operator $``$ with “upper symbol”, $`H(F,G)`$, defined by
$`H(F,G)=F,G||F,G,`$ (32)
and “lower symbol”, $`h(F,G)`$, defined implicitly by
$`=N_1^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑F{\displaystyle _0^{\mathrm{}}}𝑑Gh(F,G)|F,GF,G|,`$ (33)
and outline the standard construction of the coherent state path integral for the propagator
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})F^{\prime \prime },G^{\prime \prime }|e^{iT}|F^{},G^{}.`$ (34)
The procedure starts with the insertion of $`M`$ resolutions of unity,
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})`$
$`=N^M{\displaystyle \underset{j=1}{\overset{M}{}}dF_jdG_j\underset{k=0}{\overset{M}{}}F_{k+1},G_{k+1}|e^{i\frac{t}{M+1}}|F_k,G_k},`$ (35)
where
$`|F_0,G_0`$ $``$ $`|F^{},G^{},`$ (36)
$`|F_{M+1},G_{M+1}`$ $``$ $`|F^{\prime \prime },G^{\prime \prime }.`$ (37)
The $`M\mathrm{}`$ limit is then taken. It is customary (though not rigorously justified) at this point to interchange the order of the limit and the integrations, and to write the integrand in the form it would take for continuous and differentiable paths. One is then led, with the aid of the symplectic potential (30), to the strictly formal expression
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})={\displaystyle \mathrm{exp}\left\{i_0^T𝑑t[\gamma _1G\dot{F}H(F,G)]\right\}𝒟F𝒟G},`$ (38)
where $``$ represents a suitable normalization.
An alternative path integral representation may be given using the technique of continuous time regularization . This procedure involves the insertion of an appropriate Wiener measure into the integral, and leads to
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})=\underset{\nu \mathrm{}}{lim}_\nu {\displaystyle \mathrm{exp}\left\{i_0^T𝑑t[\gamma _1G\dot{F}h(F,G)]\right\}}`$
$`\times \mathrm{exp}\left\{(\gamma _1/2\nu ){\displaystyle _0^T}𝑑t[\beta ^1G^2\dot{F}^2+\beta G^2\dot{G}^2]\right\}𝒟F𝒟G.`$ (39)
Note that the lower symbol for the Hamiltonian is involved in this formulation.
## 4 A Generalized Affine Algebra
We now construct an $`n`$-dimensional generalization of (1) via the introduction of the set of $`n^2`$, $`GL(n,R)`$ generators $`\kappa _a^b`$ along with their $`\frac{1}{2}n(n+1)`$ symmetric affine conjugates $`\sigma _{jk}`$ ($`=\sigma _{kj}`$), satisfying the commutation relations
$`[\kappa _a^b,\kappa _j^k]`$ $`=`$ $`i\frac{1}{2}\left(\delta _a^k\kappa _j^b\delta _j^b\kappa _a^k\right),`$ (40)
$`[\sigma _{jk},\kappa _a^b]`$ $`=`$ $`i\frac{1}{2}\left(\delta _j^b\sigma _{ak}+\delta _k^b\sigma _{aj}\right),`$ (41)
$`[\sigma _{ab},\sigma _{jk}]`$ $`=`$ $`0,`$ (42)
where all the indices take on values from $`1`$ to $`n`$. The operators $`\kappa _a^b`$ and $`\sigma _{ab}`$ may be contracted with sets of constant coefficients $`B_b^a`$ and $`F^{ab}`$ ($`=F^{ba})`$, respectively, and exponentiated to generate the following transformations:
$`e^{iB_b^a\kappa _a^b}\kappa _j^ke^{iB_b^a\kappa _a^b}`$ $`=`$ $`\left(S^1\right)_j^p\kappa _p^qS_q^k,`$ (43)
$`e^{iB_b^a\kappa _a^b}\sigma _{jk}e^{iB_b^a\kappa _a^b}`$ $`=`$ $`\left(S^1\right)_j^p\sigma _{pq}\left(S^1\right)_k^q,`$ (44)
$`e^{iF^{ab}\sigma _{ab}}\kappa _j^ke^{iF^{ab}\sigma _{ab}}`$ $`=`$ $`\kappa _j^k+F^{kp}\sigma _{jp},`$ (45)
$`e^{iF^{ab}\sigma _{ab}}\sigma _{jk}e^{iF^{ab}\sigma _{ab}}`$ $`=`$ $`\sigma _{jk}.`$ (46)
Here the matrix $`S`$ is defined by
$`S=e^{\frac{1}{2}B},`$ (47)
and clearly we have
$`detS>0`$ (48)
for all values of $`B`$. Observe that (43) and (44) have the flavor of coordinate transformations of tensors of the appropriate valences.
## 5 Unitary Representation for the Generalized Affine Algebra
We now construct a unitary representation for the group generated by the generalized affine algebra described in section 4. As a representation space we choose the space of square integrable functions of a symmetric, positive-definite matrix variable $`\underset{¯}{k}\{k_{ab}\}`$, endowed with the inner product definition
$`\varphi |\psi ={\displaystyle _+}{\displaystyle \underset{ab}{}}dk_{ab}\varphi (\underset{¯}{k})^{}\psi (\underset{¯}{k}),`$ (49)
where the “+” as a limit to the integral indicates that the domain of integration extends only over those regions in which $`\{k_{ab}\}`$ is positive definite. The algebra may be represented as follows:
$`\sigma _{ab}`$ $`=`$ $`k_{(ab)},`$ (50)
$`\kappa _a^b`$ $`=`$ $`i\frac{1}{2}[^{(bp)}k_{(pa)}+k_{(pa)}^{(bp)}]`$ (51)
$`=`$ $`i[k_{(ap)}^{(bp)}+\frac{1}{4}(n+1)\delta _a^b],`$
where
$`k_{(ab)}`$ $``$ $`\frac{1}{2}(k_{ab}+k_{ba}),`$ (52)
$`^{(ab)}`$ $``$ $`\frac{1}{2}\left(/k_{ab}+/k_{ba}\right).`$ (53)
It follows from (51) that
$`e^{iB_q^p\kappa _p^q}\psi (\underset{¯}{k})`$ $`=`$ $`(detS)^{(n+1)/2}e^{B_q^pk_{pr}^{qr}}\psi (\underset{¯}{k})`$ (54)
$`=`$ $`(detS)^{(n+1)/2}\psi (S^T\underset{¯}{k}S),`$
where
$`(S^T\underset{¯}{k}S)_{ab}S_a^pk_{pq}S_b^q.`$ (55)
It may be verified that with the choice of measure in (49), the operators representing $`\sigma _{jk}`$ and $`\kappa _a^b`$ are self adjoint and that the representation of the relevant group is thus rendered unitary:
$`e^{iB_q^p\kappa _p^q}|\psi ^2`$ $`=`$ $`|\psi ^2;`$ (56)
$`e^{iF^{jk}\sigma _{jk}}|\psi ^2`$ $`=`$ $`|\psi ^2.`$ (57)
## 6 Generalized Affine Coherent States
Following the procedure in section (2), we now define a family of unitary operators $`U(F,S)`$ via
$`U(F,S)=e^{iF^{jk}\sigma _{jk}}e^{iB_b^a\kappa _a^b}.`$ (58)
The composition rule for the operators in (58) is
$`U(F^{},S^{})U(F,S)=U(F^{}+S^TFS^{},S^{}S).`$ (59)
where
$`(S^TFS^{})^{jk}`$ $``$ $`S_{}^{}{}_{p}{}^{j}F^{pq}S_q^k,`$ (60)
$`(S^{}S)_a^b`$ $``$ $`S_{}^{}{}_{a}{}^{p}S_p^b.`$ (61)
This family of unitary operators may be used to construct a set of coherent states:
$`|F,SU(F,S)|\eta .`$ (62)
Here $`|\eta `$ is an as yet unspecified normalized fiducial vector in the $`n`$-dimensional representation space.
A resolution of unity in terms of the coherent states in (62) may be established in the usual way, namely, by integrating coherent state projection operators weighted with the appropriate group invariant measure:
$`N^1{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{jk}{}}dF^{jk}{\displaystyle _{detS>0}}{\displaystyle \underset{a,b}{}}dS_a^b(detS)^{(2n+1)}|F,SF,S|=11,`$ (63)
the validity of which hinges on the fiducial vector admissibility criterion
$`N`$ $``$ $`(2\pi )^{n(n+1)/2}{\displaystyle _{detS>0}}{\displaystyle \underset{a,b}{}}dS_a^b(detS)^n|\eta (SS^T)|^2`$ (64)
$`=`$ $`2^1(2\pi )^{n(n+1)/2}{\displaystyle \underset{a,b}{}dS_a^b|detS|^n|\eta (SS^T)|^2}<\mathrm{},`$
where
$`(SS^T)_{ab}{\displaystyle \underset{p}{}}S_a^pS_b^p.`$ (65)
We now choose as a fiducial vector a natural generalization of the one-dimensional vector in (15), that is,
$`\eta (\underset{¯}{k})`$ $`=`$ $`C_n(\alpha ,\beta )det(\underset{¯}{k}^\alpha e^{\beta \underset{¯}{k}})`$ (66)
$`=`$ $`C_n(\alpha ,\beta )(det\underset{¯}{k})^\alpha e^{\beta \mathrm{tr}\underset{¯}{k}},`$
where $`\alpha `$ and $`\beta `$ are positive real coefficients, and the constant $`C_n(\alpha ,\beta )`$ is chosen to be
$`C_n(\alpha ,\beta )={\displaystyle \frac{(2\beta )^{\alpha n+n(n+1)/4}}{\sqrt{K_n(2\alpha )}}}.`$ (67)
Here $`K_n(2\alpha )`$ is defined by
$`K_n(2\alpha ){\displaystyle _+}{\displaystyle \underset{ab}{}}dk_{ab}(det\underset{¯}{k})^{2\alpha }e^{\mathrm{tr}\underset{¯}{k}}.`$ (68)
The integral in (68) may be reduced to a Gaussian integral via a change of variables involving the replacement
$`k_{ab}={\displaystyle \underset{p}{}}Q_a^pQ_b^p.`$ (69)
The result is:
$`K_n(2\alpha )=2^{n1}\mathrm{\Omega }_n^1{\displaystyle \underset{a,b}{}dQ_a^b|detQ|^{4\alpha +1}\mathrm{exp}[\underset{a,b}{}(Q_a^b)^2]}.`$ (70)
Here $`\mathrm{\Omega }_n`$ is the group volume of $`SO(n)`$, which can be expressed as a product of the surface volumes of $`j`$-spheres :
$`\mathrm{\Omega }_n={\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{2\pi ^{j/2}}{\mathrm{\Gamma }(j/2)}}.`$ (71)
The change of variables necessary to obtain (70) is used repeatedly throughout this paper - we refer the reader to the appendix for the details. Clearly the existence (convergence) of the integral expression for $`K_n`$ is manifest in the form (70). The choice (67) ensures that all the coherent states are normalized,
$`F,S|F,S=\eta |\eta =1.`$ (72)
The admissibility of the fiducial vector (66) may be verified by demonstrating the existence of the integral in (64). Again we refer the reader to the appendix for the change of variables necessary to perform this type of integral; the result is
$`N`$ $`=`$ $`2^n(4\pi \beta )^{n(n+1)/2}\mathrm{\Omega }_n{\displaystyle \frac{K_n\left(2\alpha (n+1)/2\right)}{K_n(2\alpha )}}<\mathrm{}.`$ (73)
The overlap of two coherent states based on the fiducial vector (66) may be expressed as
$`F^{},S^{}|F,S`$
$`=[C_n(\alpha ,\beta )]^2[det(S^{}S)]^{(n+1)/2+\mathrm{\hspace{0.33em}2}\alpha }{\displaystyle _+}{\displaystyle \underset{ab}{}}dk_{ab}(det\underset{¯}{k})^{2\alpha }e^{\mathrm{tr}(X\underset{¯}{k})},`$ (74)
where the complex symmetric matrix $`X`$ is defined by
$`X\beta (S^TS^{}+S^TS)+i(F^{}F),`$ (75)
with
$`(S^TS)^{ab}`$ $`=`$ $`{\displaystyle \underset{p}{}}S_p^aS_p^b,`$ (76)
$`(S^TS^{})^{ab}`$ $`=`$ $`{\displaystyle \underset{p}{}}S_{}^{}{}_{p}{}^{a}S_{}^{}{}_{p}{}^{b}.`$ (77)
The $`X`$-dependence may be extracted from the integral in (74) to leave
$`F^{},S^{}|F,S=[C_n(\alpha ,\beta )]^2\left[{\displaystyle \frac{det(S^{^{}}S)}{detX}}\right]^{(n+1)/2+\mathrm{\hspace{0.33em}2}\alpha }K_n(\alpha )`$
$`=\left\{{\displaystyle \frac{det(S^{}S)}{det[(S^TS^{}+S^TS)/2+i(F^{}F)/2\beta ]}}\right\}^{2\alpha +(n+1)/2}.`$ (78)
We appeal to analytic continuation to give the final result in (78) a well-defined meaning.
It will be noticed that the overlap function (78) only depends on $`S`$ through the symmetric combination $`S^TS`$. It is therefore invariant under a transformation
$`SMS,`$ (79)
where $`M`$ is any $`SO(n)`$ matrix. It is appropriate, then, to view the $`SO(n)`$ degrees of freedom as “gauge” degrees of freedom and factor them out of the representation. To this end, we define a new symmetric matrix variable $`G`$ via the relations
$`G^{ab}`$ $``$ $`(S^TS)^{ab}{\displaystyle \underset{p}{}}S_p^aS_p^b,`$ (80)
$`G_{ap}G^{pb}`$ $``$ $`\delta _a^b,`$ (81)
$`G`$ $``$ $`\{G_{ab}\},`$ (82)
and label our coherent states with the parameters $`F`$ and $`G`$. The overlap function (78) then reads
$`F^{},G^{}|F,G`$ $`=`$ $`\left\{{\displaystyle \frac{(detG^{})^{1/2}(detG)^{1/2}}{det\left[(G_{}^{}{}_{}{}^{1}+G^1)/2+i(F^{}F)/2\beta \right]}}\right\}^{2\alpha +(n+1)/2}.`$ (83)
The $`SO(n)`$ variables may be integrated out of the resolution of unity (again, see the appendix for a discussion of the required change of variables), the result being
$`2^nN^1\mathrm{\Omega }_n{\displaystyle \underset{jk}{}dF^{jk}_+\underset{ab}{}dG_{ab}|F,GF,G|}`$ $`=`$ $`11.`$ (84)
We note that the removal of the $`SO(n)`$ degrees of freedom from the representation is only appropriate if the dynamics in question is governed by a classical Hamiltonian which is a function of $`F`$ and $`G`$; this is the point of view we shall take in the remainder of this article. It is, however, easy in principle to envisage Hamiltonians where spinor-like variables couple directly to $`S`$, in which case it would of course be necessary to retain the label $`S`$.
The polarization property analogous to (26) for the one-dimensional case may be written as
$`\left\{iG^{ap}{\displaystyle \frac{}{F^{pb}}}+\beta ^1G^{ap}{\displaystyle \frac{}{G^{pb}}}\gamma \delta _b^a\right\}F,G|\psi =\mathrm{\hspace{0.33em}0},`$ (85)
where we have written
$`\gamma \beta ^1[(n+1)/4+\alpha ].`$ (86)
## 7 Path Integral for the Propagator
The procedure in section 3 may be followed to build a formal path integral expression for the $`n`$-dimensional propagator associated with a general time-independent Hamiltonian $``$ with upper symbol
$`H(F,G)=F,G||F,G.`$ (87)
We first construct the relevant symplectic potential and ray metric. The identity
$`det(1+dA)`$ $`=`$ $`e^{\mathrm{tr}\mathrm{ln}(1+dA)}`$ (88)
$`=`$ $`e^{\mathrm{tr}(dAdA^2/2+\mathrm{})}`$
where $`dA`$ is any infinitesimal matrix, may be used to expand the overlap of two neighboring coherent states as follows:
$`F+dF/2,S+dS/2|FdF/2,SdS/2=e^{id\theta d\mathrm{\Sigma }^2/2},`$ (89)
where the 1-form $`d\theta `$ is given by
$`d\theta =\gamma \mathrm{tr}(GdF),`$ (90)
and the ray metric $`d\mathrm{\Sigma }^2`$ by
$`d\mathrm{\Sigma }^2=\gamma \{\beta \mathrm{tr}[(G^1dG)^2]+\beta ^1\mathrm{tr}[(GdF)^2]\}.`$ (91)
The propagator may then be written as
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})F^{\prime \prime },G^{\prime \prime }|e^{iT}|F^{},G^{}`$
$`={\displaystyle \underset{jk}{}𝒟F^{jk}\underset{ab}{}𝒟G_{ab}\mathrm{exp}\left\{i_0^T[\gamma \mathrm{tr}(G\dot{F})H(F,G)]𝑑t\right\}},`$ (92)
a strictly formal expression to which the remarks immediately preceding (38) again apply.
An alternative representation for the propagator which uses a continuous-time regularization and the lower symbol is given by
$`J_T(F^{\prime \prime },G^{\prime \prime };F^{},G^{})`$
$`=\underset{\nu \mathrm{}}{lim}_\nu {\displaystyle \underset{jk}{}𝒟F^{jk}\underset{ab}{}𝒟G_{ab}\mathrm{exp}\left\{i_0^T[\gamma \mathrm{tr}(G\dot{F})h(F,G)]𝑑t\right\}}`$
$`\times \mathrm{exp}\left\{(\gamma /2\nu ){\displaystyle _0^T}𝑑t\{\beta \mathrm{tr}[(G^1\dot{G})^2]+\beta ^1\mathrm{tr}[(G\dot{F})^2]\}\right\}.`$ (93)
## 8 Conclusion
In this article we have constructed a framework for the quantization of a positive-definite matrix degree of freedom $`\{\sigma _{ab}\}`$. Specifically, we have demonstrated that complementing the operators $`\sigma _{ab}`$ with affine conjugates (40-42) leads to a representation in which the spectrum of the matrix operator $`\{\sigma _{ab}\}`$ is strictly positive definite. Such an approach appears far more satisfactory than the standard use of canonical commutation relations, where the positivity of $`\{\sigma _{ab}\}`$ can only be insured by rather artificial means, if at all. As demonstrated, the generalized affine algebra leads to a set of group-defined coherent states that have been used to construct two versions of path integral expressions for the propagator. Finally, we suggest that the affine construction we have outlined is well suited to the quantization of the spacial part of the metric tensor of general relativity, a program already initiated in .
## Appendix - <br>Jacobian associated with change of variables from $`S`$ to $`G`$
We have repeatedly found it useful to factor out the $`SO(n)`$ degrees of freedom from a real $`n`$-dimensional non-degenerate matrix with positive determinant $`S`$ and display the remaining degrees of freedom as elements of the matrix $`GS^TS`$. We now derive the form of the Jacobian associated with this change of variables.
We first calculate the Jacobian associated with the polar decomposition of $`S`$. Such a decomposition may be defined by
$`S=MT,`$ (94)
where $`MSO(n)`$ and $`T`$ is a positive $`n`$-dimensional upper triangular matrix. We first consider the case $`n=2`$, writing (94) explicitly as
$`\left(\begin{array}{cc}S_1^1& S_1^2\\ S_2^1& S_2^2\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{cc}T_1^1& T_1^2\\ 0& T_2^2\end{array}\right).`$ (95)
It is straightforward to show that in this case,
$`dS_1^1dS_1^2dS_2^1dS_2^2=T_1^1d\theta dT_1^1dT_1^2dT_2^2.`$ (96)
We now generalize (96) for $`n>2`$. It is expeditious to express the $`SO(n)`$ matrix $`M`$ as a product of $`\frac{1}{2}n(n1)`$, $`SO(n)`$ matrices $`R_{ij}=R_{ij}(\theta _{ij})`$, $`ni>j`$, each of which represents a rotation about the $`ij`$ axis through an angle $`\theta _{ij}`$:
$`M=(R_{21}R_{31}\mathrm{}R_{n1})(R_{32}R_{42}\mathrm{}R_{n2})`$
$`\times \mathrm{}\times (R_{(n1)(n2)}R_{n(n2)})(R_{n(n1)}).`$ (97)
It will be noticed that $`R_{n(n1)}`$, whose explicit form is
$`R_{n(n1)}=\left(\begin{array}{cccccc}1& 0& 0& \mathrm{}& \mathrm{}& 0\\ 0& 1& 0& \mathrm{}& \mathrm{}& 0\\ 0& 0& 1& \mathrm{}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& & \\ 0& 0& 0& & \mathrm{cos}\theta _{n(n1)}& \mathrm{sin}\theta _{n(n1)}\\ 0& 0& 0& & \mathrm{sin}\theta _{n(n1)}& \mathrm{cos}\theta _{n(n1)}\end{array}\right),`$ (104)
only affects the bottom right $`2\times 2`$ block of $`T`$. It is therefore responsible for the introduction of a factor $`T_{(n1)}^{(n1)}`$ into the Jacobian. Similarly, each $`R_{ij}`$ only acts on the bottom $`j\times j`$ block of $`T`$, introducing a factor $`T_j^j`$. Building up the entire Jacobian in this way, we find that
$`dS_1^1\mathrm{}dS_n^n=(T_1^1)^{n1}(T_2^2)^{n2}\mathrm{}(T_{n1}^{n1})d\mathrm{\Omega }dT_1^1\mathrm{}dT_n^n`$ (105)
where $`d\mathrm{\Omega }`$ is the invariant measure on $`SO(n)`$.
Having separated the matrix $`S`$ into its “radial” part $`T`$ and “angular” part $`M`$, we are now in a position to make a change of variable from $`T`$ to $`G`$ (recall that both of these matrices possess $`\frac{1}{2}n(n+1)`$ degrees of freedom):
$`G=S^TS=T^TT.`$ (106)
Inspection of the elements of $`G`$ quickly reveals the form of the Jacobian associated with this change of variable:
$`dG^{11}\mathrm{}dG^{nn}`$
$`=\mathrm{\hspace{0.33em}2}^n(T_1^1)^n(T_2^2)^{n1}\mathrm{}(T_{n1}^{n1})^{\mathrm{\hspace{0.33em}2}}(T_n^n)dT_1^1\mathrm{}dT_n^n.`$ (107)
Combining (99) and (107), and noting that $`detS=detT`$, we obtain our central result:
$`d\mathrm{\Omega }dG^{11}\mathrm{}dG^{nn}`$ $`=`$ $`2^n(detS)dS_1^1\mathrm{}dS_n^n`$ (108)
We now integrate a general function $`f(G)`$ against the measure in (108):
$`{\displaystyle _+}{\displaystyle \underset{ab}{}}dG^{ab}f(G)=\mathrm{\hspace{0.33em}2}^n\mathrm{\Omega }_n^1{\displaystyle _{detS>0}}{\displaystyle \underset{a,b}{}}dS_a^b(detS)f(S^TS),`$ (109)
where $`\mathrm{\Omega }_n`$ is the group volume of $`SO(n)`$ given in (71). |
warning/0001/astro-ph0001475.html | ar5iv | text | # Investigations of the Local Supercluster Velocity Field
## 1 Introduction
Study of the local extragalactic velocity field has a considerable history. Rubin (Rubin88 (1988)) pinpoints the beginning of the studies concerning deviations from the Hubble law to a paper of Gamow (Gamow46 (1946)) where Gamow asked if galaxies partake of a large-scale systematic rotation in addition to the Hubble expansion. The pioneer works by Rubin (Rubin51 (1951)) and Ogorodnikov (Ogorodnikov52 (1952)) gave evidence that the local extragalactic velocity field is neither linear nor isotropic. De Vaucouleurs (deVaucouleurs53 (1953)) then interpreted the distribution of bright galaxies and proposed rotation in terms of a flattened local supergalaxy. This short but remarkable paper did not yet refer to differential expansion, introduced by de Vaucouleurs (deVaucouleurs58 (1958)) as an explanation of the “north-south anisotropy” which he stated was first pointed out by Sandage (Humason et al. Humason56 (1956)). Differential expansion was a milder form of Hubble’s “the law of redshifts does not operate within the Local Group” and de Vaucouleurs pondered that “in condensed regions of space, such as groups or clusters, the expansion rate is greatly reduced…”. Though there was a period of debate on the importance of the kinematic effects claimed by de Vaucouleurs and even on the reality of the local supergalaxy (presently termed as the Local Supercluster, LSC), already for two decades the reality of the differential peculiar velocity field around the Virgo cluster has been generally accepted. However, its amplitude and such details as the deviation from spherical symmetry and possible rotational component, are still under discussion.
A theoretical line of research related to de Vaucouleurs’ differential expansion, has been motivated by the work on density perturbations in Friedmann cosmological models, resulting in infall models of matter (Silk Silk74 (1974)) which predict a connection between the infall peculiar velocity at the position of the Local Group towards the Virgo cluster and the density parameter of the Friedmann universe. Later on, Olson & Silk (Olson79 (1979)) further developed the formalism in a way which was found useful in Teerikorpi et al. (Teerikorpi92 (1992); hereafter Paper I). The linearized approximation of Peebles (Peebles76 (1976)) has been often used for describing the velocity field and for making routine corrections for systemic velocities.
Using Tolman-Bondi model (Tolman Tolman34 (1934), Bondi Bondi47 (1947)) Hoffman et al. (Hoffman80 (1980)) calculated the expected velocity dispersions along line-of-sight as a function of angular distance from a supercluster and applied the results to Virgo. They derived a gravitating mass of about $`4\times 10^{14}M_{\mathrm{}}\times 100/h_0`$ inside the cone of $`6\mathrm{°}`$. The Tolman-Bondi (TB) model is the simplest inhomogeneous solution to the Einstein’s field equations. It describes the time evolution of a spherically symmetric pressure-free dust universe in terms of comoving coordinates. For details of the TB-model cf. Ekholm et al. (1999a ; hereafter Paper II).
Then, following the course of Hoffman et al. (Hoffman80 (1980)), Tully & Shaya (Tully84 (1984)) calculated the expected run of radial velocity vs. distance at different angular distances from Virgo and for different (point) mass-age models. Comparison of such envelope curves with available galaxy data agreed with the point mass having roughly the value of Virgo’s virial mass ($`7.5\times 10^{14}M_{\mathrm{}}\times 75/h_0`$) for reasonable Friedmann universe ages.
The Hubble diagram of Tully & Shaya contained a small number of galaxies and did not very well show the expected behaviour. With a larger sample of Tully-Fisher measured galaxies and attempting to take into account the Malmquist bias, Teerikorpi et al. (Teerikorpi92 (1992)) were able to put in evidence the expected features: an initial steeply rising tight velocity-distance relation, the local maximum in front of Virgo and the final ascending part of the relation, expected to approach asymptotically the undisturbed Hubble law. Looking from the Virgo centre the zero-velocity surface was clearly seen around $`r/R_{\mathrm{Virgo}}0.5`$. Using either a continuous mass model or a two-component model, the conclusions of Tully & Shaya (Tully84 (1984)) were generally confirmed and it was stated that “Various density distributions, constrained by the mass inside the Local Group distance (required to produce $`V_{\mathrm{Virgo}}`$), agree with the observations, but only if the mass within the Virgo $`6\mathrm{°}`$ region is close to or larger than the standard Virgo virial mass values. This is so independently of the value of $`q_0`$, of the slope of the density distribution outside of Virgo, and of the values adopted for Virgo distance and velocity”.
It is the aim of the present paper to use the available sample of galaxies with more accurate distances from Cepheids and Tully-Fisher relation to study the virgocentric velocity field. In Paper II galaxies with Cepheid-distances were used to map the velocity field in front of Virgo, here we add galaxies with good Tully-Fisher distances in order to see both the frontside and backside behaviour and investigate how conclusions of Paper I should be modified in the light of new data. It should be emphasized that also our Tully-Fisher distances are now better, after a programme to study the slope and the Hubble type dependence of the zero-point (see Theureau et al. Theureau97 (1997)).
This paper is structured as follows. In Sect. 2 we shortly review the basics of the use of the direct Tully-Fisher relation, give the relation to be used and describe our sample and the restrictions put upon it. In Sect. 3 we examine our sample in terms of systemic velocity vs. distance diagrams and see which distance to Virgo will bring about best agreement between the TB-predictions and the observations. In Sect. 4 we try to answer the question whether we have actually found the Virgo cluster at the centre of the TB-metric. In Sect. 5 we re-examine our sample from a virgocentric viewpoint and compare our results from the TF-distances with the sample of galaxies with distances from the extragalactic Cepheid $`PL`$-relation. In Sect. 6 we shortly discuss the mass estimate and our density profile and, finally, in Sect. 7 we summarize our results with some conclusive remarks.
## 2 The sample based on direct B-band Tully-Fisher relation
The absolute magnitude $`M`$ and the logarithm of the maximum rotational velocity $`\mathrm{log}V_{\mathrm{max}}`$ of a galaxy (for which also a shorthand $`p`$ is used) are related as:
$$M=a\mathrm{log}V_{\mathrm{max}}+b.$$
(1)
The use of this kind of relation as a distance indicator was suggested by Gouguenheim (Gouguenheim69 (1969)). Eq. 1 is known as Tully-Fisher (TF) relation after Tully & Fisher (Tully77 (1977)).
It is nowadays widely acknowledged that the distance moduli inferred using Eq. 1 are underestimated because of selection effects in the sampling. We can see how this Malmquist bias affects the distance determination by considering the observed average absolute magnitude $`M_p`$ at each $`p`$ as a function of the true distance $`r`$. The limit in apparent magnitude, $`m_{\mathrm{lim}}`$, cuts off progressively more and more of the distribution function of $`M`$ for a constant $`p`$. This means that the observed mean absolute magnitude $`M_p`$ is overestimated by the expectation value $`E(M|p)=ap+b`$:
$$M_pE(M|p),$$
(2)
This inequality gives a practical measure of the Malmquist bias depending primarily on $`p`$, $`r`$, $`\sigma _M`$ and $`m_{\mathrm{lim}}`$. The equality holds only when the magnitude limit cuts the luminosity function $`\mathrm{\Phi }(M)`$ insignificantly. For our present purposes it is also important to note that for luminous galaxies, which are also fast rotators (large $`p`$) the effect of the magnitude limit is felt at much larger distances than for intrinsically faint galaxies which rotate slowly. Hence by limiting $`p`$ to large values one expects to add to the sample galaxies which suffer very little from the Malmquist bias within a restricted distance range. For this kind of bias the review by Teerikorpi (Teerikorpi97 (1997)) suggested the name Malmquist bias of the $`2^{\mathrm{nd}}`$ kind, in order to make a difference from the classical Malmquist bias (of the $`1^{\mathrm{st}}`$ kind).
Following Paper I we selected galaxies towards Virgo by requiring $`\mathrm{log}V_{\mathrm{max}}`$ to be larger than $`2.1`$. At the time Paper I was written this value was expected to bring about nearly unbiased TF distance moduli up to twice the Virgo distance. With the present, much deeper sample the limit chosen is much safer. Also, we allow an error in B-magnitude to be at maximum $`0.2^{\mathrm{mag}}`$. We also require the axis ratio to be $`\mathrm{log}R_{25}>0.07`$. Because the maximum amplitude of systemic velocities near Virgo can be quite large, we first restricted the velocities by $`V_{\mathrm{obs}}<3V_{\mathrm{Virgo}}^{\mathrm{cosm}}\mathrm{cos}\mathrm{\Theta }`$, where $`\mathrm{\Theta }`$ is the angular distance from the adopted centre ($`l=284\mathrm{°}`$, $`b=74.5\mathrm{°}`$) and the cosmological velocity of the centre is following Paper II $`V_{\mathrm{Virgo}}^{\mathrm{cosm}}=1200\mathrm{km}\mathrm{s}^1`$. After this the derived TF-distances were restricted by $`R_{\mathrm{TF}}<60\mathrm{Mpc}`$.
With these criteria we found 96 galaxies within $`\mathrm{\Theta }<30\mathrm{°}`$ tabulated in Table 1, where in columns (1) and (2) we give the PGC number and name (the superscript after some galaxies will be explained in Sect. 4). In columns (3) and (4) the galactic coordinates $`l,b`$ in degrees are given. In column (5) we give the morphological type code $`T`$ and in column (6) we give the logarithm of the axis ratio at $`25^{\mathrm{mag}}/\mathit{}\mathrm{}`$, $`\mathrm{log}R_{25}`$. The total B-magnitude corrected according to RC3 (de Vaucouleurs et al. deVaucouleurs91 (1991))<sup>1</sup><sup>1</sup>1Except for galactic extinction which is adopted from RC2 (de Vaucouleurs et al. deVaucouleurs76 (1976)) and the corresponding weighted mean error are given in columns (7) and (8). In columns (9) and (10) we give the logarithm of the maximum rotational velocity $`\mathrm{log}V_{\mathrm{max}}`$ with the weighted mean error. In column (11) we give the observed velocity $`V_{\mathrm{obs}}`$ by which – as in Paper II – we mean the mean observed heliocentric velocity corrected to the centroid of the Local Group according to Yahil et al. (Yahil77 (1977)). Finally, in columns (12) and (13) we have the angular distance $`\mathrm{\Theta }_{\mathrm{gal}}`$ in degrees between a galaxy and the centre and the distance $`R_{\mathrm{gal}}`$ in Mpc from us calculated using the direct TF-relation given below. The data in columns (1) – (11) were extracted from the Lyon-Meudon extragalactic database LEDA.
Our direct TF-parameters for the B-band magnitudes were taken from Theureau et al. (Theureau97 (1997)). The slope for the relation is $`a=5.823`$ and the zero-points corrected for the type-effect are given in Table 2.
The calibration of the zero-points was based on a sample of galaxies with Cepheid distances given in Table 1 in Theureau et al. (Theureau97 (1997)). This calibration corresponded to a Hubble constant $`H_055\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$.
Finally we comment on our notation on velocities. We use systemic velocity in the same sense as in Paper II, i.e. the systemic velocity is a combination of the cosmological velocity and the velocity induced by Virgo with the assumption that the virgocentric motions dominate. When we refer to observed systemic velocity we call it $`V_{\mathrm{obs}}`$ and when to model prediction, $`V_{\mathrm{pred}}`$. If we make no distinction, we use $`V_{\mathrm{sys}}`$.
## 3 The $`V_{\mathrm{sys}}`$ vs. $`R_{\mathrm{gal}}`$ diagram for the TF-sample
In Paper II we found a TB-solution using a simple density law $`\rho (R)=\rho _{\mathrm{bg}}(1+kR^\alpha )`$, which fitted data quite well. Here $`R`$ is the distance from the origin of the TB-metric, $`\alpha `$ is the density gradient and $`k`$ the density contrast. Because an Einstein-deSitter universe was assumed, the background density $`\rho _{\mathrm{bg}}`$ equals the critical cosmological density, $`\rho _\mathrm{c}`$. The relevant quantity, the mass within a radius $`d`$, the radius $`R`$ in units of Virgo distance, was expressed as $`M(d)=M(d)_{\mathrm{EdS}}\times (1+k^{}d^\alpha )`$ (cf. Eq. 9 in Paper II). Here $`k^{}`$ is the mass excess within a sphere having a radius of one Virgo distance.
Unfortunately, the sample of galaxies with distances from the extragalactic Cepheid $`PL`$-relation did not reach well enough behind the LSC. Our present sample is clearly deep enough to reveal the backside infall signal. In Paper I it was well seen how in the front the differences between different TB-models were not large, in contrast to the background, where the model predictions progressively deviate from each other.
In the formalism developed by Ekholm (Ekholm96 (1996)) and adopted in Paper II, the quantity given by Eq. 8 in Paper II, $`A(d,q_0)`$, which is needed for solving the development angle, is no longer an explicit function of $`H_0`$. There are – however – still rather many free parameters, which we shortly discuss below:
1. The deceleration parameter $`q_0`$. In Paper II we considered $`q_0`$ given, restricting our analysis to the Einstein-deSitter universe ($`q_0=0.5`$). In Paper I it was concluded that $`q_0`$ has a minor influence on the $`V_{\mathrm{pred}}`$ vs. $`R_{\mathrm{gal}}`$ curves and on the Virgo mass (though it has a large effect on total mass inside the LG sphere).
2. The density gradient $`\alpha `$ and the relative mass excess at $`d=1`$, $`k^{}`$. We remind that $`k^{}`$ in our formalism does not depend on $`\alpha `$ but only on the amount by which the LG’s expansion velocity with respect to centre of LSC has slowed down. In our two-component model (Sect. 5) $`k^{}`$ will depend also on $`\alpha `$.
3. The velocities $`V_{\mathrm{LG}}^{\mathrm{in}}`$, $`V_{\mathrm{Virgo}}^{\mathrm{obs}}`$ and $`V_{\mathrm{Virgo}}^{\mathrm{cosm}}`$. As in Papers I and II, we presume Virgo to be at rest with cosmological background: $`V_{\mathrm{Virgo}}^{\mathrm{cosm}}=V_{\mathrm{LG}}^{\mathrm{in}}+V_{\mathrm{Virgo}}^{\mathrm{obs}}`$. We feel that our choices for the infall velocity of the Local Group $`V_{\mathrm{LG}}^{\mathrm{in}}=220\mathrm{km}\mathrm{s}^1`$ and for the observed velocity of Virgo $`V_{\mathrm{Virgo}}^{\mathrm{obs}}=980\mathrm{km}\mathrm{s}^1`$ are relatively safe.
We would also like to remind that our solutions in Paper II had an implicit dependence on the Hubble constant $`H_0`$, because we fixed our distance to Virgo kinematically from $`R_{\mathrm{Virgo}}=V_{\mathrm{Virgo}}^{\mathrm{cosm}}/H_0`$ by adopting $`H_0=57\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. This global value was based on SNe Ia (Lanoix Lanoix99 (1999)) and agrees also with the more local results of the KLUN (Kinematics of the Local Universe) project (Theureau et al. Theureau97 (1997); Ekholm et al. 1999b ) and with the findings of Federspiel et al. (Federspiel98 (1998)). Here we allow the distance of Virgo, or equivalently the Hubble constant $`H_0`$, vary keeping the cosmological velocity of Virgo fixed. This choice is justified because even though the estimates for $`H_0`$ have converged to $`60\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ the reported $`1\sigma `$ errors are not small and the different values are still scattered ($`50`$-$`70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$).
In this section we examine how well the present TB-sample agrees with the Model 1 of Paper II, which constitutes of a density excess embedded in a FRW universe with $`q_0=0.5`$ and $`H_0=57\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. The model parameters are $`k^{}=0.606`$ and $`\alpha =2.85`$, which predict for the Virgo cluster ($`\mathrm{\Theta }<6\mathrm{°}`$) a mass $`1.62\times M_{\mathrm{virial}}`$, where $`M_{\mathrm{virial}}`$ is the virial mass of the Virgo cluster derived by Tully & Shaya (Tully84 (1984)) $`=7.5\times 10^{14}M_{\mathrm{}}R_{\mathrm{Virgo}}/16.8\mathrm{Mpc}`$. Because of fixed infall velocity of the Local Group (LG) into the centre of LSC and because $`H_0`$ was fixed from external considerations the distance to centre of LSC became to be $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$. For further details of the TB-model adopted cf. Paper II. Additional discussion can be found in Paper I, Ekholm & Teerikorpi (Ekholm94 (1994)) and Ekholm (Ekholm96 (1996)).
The observed systemic velocity vs. distance $`R_{\mathrm{gal}}`$ diagrams are presented in Figs. 1-5. In the first four figures galaxies belonging to a $`\mathrm{\Theta }<30\mathrm{°}`$ cone are shown for different angular intervals: galaxies having $`\mathrm{\Theta }<10\mathrm{°}`$ are shown as black bullets, galaxies having $`10\mathrm{°}\mathrm{\Theta }<20\mathrm{°}`$ as grey bullets and galaxies having $`20\mathrm{°}\mathrm{\Theta }<30\mathrm{°}`$ as circles. The TB-curves are given for the mean angular distance, $`\mathrm{\Theta }`$, for each angular interval as thick black curve for $`\mathrm{\Theta }=4.5\mathrm{°}`$, as thick grey curve for $`\mathrm{\Theta }=15.6\mathrm{°}`$ and as thin black curve for $`\mathrm{\Theta }=25.8\mathrm{°}`$. Comparison between the data and the mean predictions were made for different presumed distances to the centre of LSC: $`R_{\mathrm{Virgo}}=16\mathrm{Mpc}`$ (Fig. 1), $`R_{\mathrm{Virgo}}=18\mathrm{Mpc}`$ (Fig. 2), $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ (Fig. 3) and $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ (Fig. 4). We remind that our model is formulated in terms of the relative distance $`d_{\mathrm{gal}}=R_{\mathrm{gal}}/R_{\mathrm{Virgo}}`$. So the TB-curves show different behaviour depending on the normalization.
The thick black line in each figure corresponds to the Hubble law based on $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, $`H_0=67\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, $`H_0=57\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, respectively. The line is drawn through the centre of LSC in order to emphasize our basic assumption that the centre is at rest with respect to the cosmological background. This also allows one to appreciate the infall of the Local Group with an assumed velocity $`V_{\mathrm{LG}}^{\mathrm{in}}=220\mathrm{km}\mathrm{s}^1`$.
Figs. 1 and 2 immediately reveal that the shorter distances are not acceptable because the background galaxies fall far below the expected curves. Correction for any residual Malmquist bias would make situation even worse. Neither is $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$, the distance found favourable in Paper II, totally satisfying. Although the clump of galaxies at $`R_{\mathrm{gal}}32\mathrm{Mpc}`$ and $`V_{\mathrm{sys}}800\mathrm{km}\mathrm{s}^1`$ in Fig. 3 follow the prediction as some other galaxies, the maximum of the velocity amplitude is clearly behind the presumed centre.
This led us to test a longer distance to Virgo. The result is shown in Fig. 4. It is rather remarkable that such a distance gives better fit than the shorter ones. On the other hand $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ together with the adopted cosmological velocity bring about $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. Such a small value has for decades been advocated by Sandage and his collaborators and is within the error bars of our determinations (Theureau et al. Theureau97 (1997); Ekholm et al. 1999b ) as well. It is encouraging that galaxies outside the $`30\mathrm{°}`$ cone follow well the Hubble law for this $`H_0`$. Virgo has only a weak influence on them, and if the Malmquist bias is present these galaxies should predict larger value for $`H_0`$. The dashed line in Fig. 5 is the Hubble law for $`H_0=60\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. It is clearly an upper limit thus giving us a lower limit for the distance to Virgo: $`R_{\mathrm{Virgo}}20\mathrm{Mpc}`$.
## 4 Have we found the true TB signature of Virgo?
So far we have studied the $`V_{\mathrm{sys}}`$ vs. $`R_{\mathrm{gal}}`$ diagram in a simple way by moving the curves for the TB-solution by choosing different normalizing distances to Virgo. The best agreement with the maximum observed amplitude and the curves was found at a rather large distance, namely $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$. Such a long distance leads one to ask whether we have actually found Virgo. We examine this question by comparing our sample given in Table 1 with the sample given by Federspiel et al. (Federspiel98 (1998)) from which they found $`R_{\mathrm{Virgo}}=20.7\mathrm{Mpc}`$. We found 33 galaxies in common when requiring $`\mathrm{\Theta }<6\mathrm{°}`$. We present these galaxies in Fig. 6. For an easy reference each galaxy is assigned a number given also as a superscript after the name in Table 1.
We give each galaxy a symbol following the classification of Federspiel et al. (Federspiel98 (1998)). Following Binggeli et al. (Binggeli93 (1993)) galaxies were divided into subgroup “big A” for galaxies close to M87 (‘A’) and into “B” for galaxies within $`2\stackrel{}{.}4`$ of M49 (‘B’). They also examined whether a galaxy is within the X-ray isophote $`0.444\mathrm{counts}\mathrm{s}^1\mathrm{arcmin}^1`$ based on ROSAT measurements of diffuse X-ray emission of hot gas in the Virgo cluster (Böhringer et al. Bohringer94 (1994)) (‘A,X’, ‘B,X’). Galaxies belonging to subgroup A and within the X-ray contour are labelled in Fig. 6 as bullets and outside the contour with an open circle. Similarly, galaxies in subgroup B are labelled with a filled or open triangle. Federspiel et al. also listed galaxies within the X-ray contour but not classified as members of A or B. The galaxies are marked with a filled square. They also included in their Table 3 some galaxies which fall outside A and B and the X-ray contour (we label them with an open square).
We also give an error estimate for the TF-distance for each galaxy calculated from the $`1\sigma `$ error in the distance modulus:
$$\sigma _\mu =\sqrt{\sigma _B^2+\sigma _{M_p}^2}.$$
(3)
The error in the corrected total B-band magnitude is taken from column (8) in our Table 1 and the intrinsic dispersion of the absolute magnitude $`M`$ for each $`p`$, $`\sigma _{M_p}`$ is estimated to be $`0.3^{\mathrm{mag}}`$.
The straight solid line is the Hubble law for $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ shifted downwards by $`V_{\mathrm{LG}}^{\mathrm{in}}=220\mathrm{km}\mathrm{s}^1`$ in order to make the line go through the centre at $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ which is presumed to be at rest with respect to the cosmological background. The TB-curves are given for $`\mathrm{\Theta }=2\mathrm{°}`$ (thick black curve), $`\mathrm{\Theta }=3\stackrel{}{.}5\mathrm{°}`$ (thick grey curve) and $`\mathrm{\Theta }=5\mathrm{°}`$ (thin black curve), respectively.
To begin with, there are 22 galaxies (67 %) which agree with the TB-solution within $`1\sigma `$ in $`\sigma _\mu `$. Only four galaxies (1,3,19,26; 12 %) <sup>2</sup><sup>2</sup>2Also NGC 4216 (7) should probably be counted to this group, because it differs by $`2\sigma `$ and clearly belongs to the same substructure as the other four disagreeing galaxies. In other words, 15 % of the sample does not agree with the model within $`2\sigma `$. do not agree with the model within $`2\sigma `$. Though we have not reached the traditional 95 % confidence level, the agreement is, at the statistical level found, satisfying enough. Furthermore, we find in the range $`R=24\pm 2\mathrm{Mpc}`$ nine galaxies out of which seven were classified by Federspiel et al. (Federspiel98 (1998)) as ‘A,X’ galaxies hence presumably lying in the very core of Virgo. The remaining two galaxies are ‘A’ galaxies. In the range $`R=16\pm 2\mathrm{Mpc}`$ we find only three ‘A,X’ galaxies and one ‘A’ galaxy. Federspiel et al. (Federspiel98 (1998)) following Guhathakurta et al. (Guhathakurta88 (1988)) listed five galaxies (17, 18, 19, 20 and 26 in Table 1) as HI-deficient. If these galaxies are removed one finds four ‘A,X’ and two ‘A’ galaxies in the range $`R=24\pm 2\mathrm{Mpc}`$, and one ‘A,X’ and one ‘A’ galaxy in the range $`R=16\pm 2\mathrm{Mpc}`$. The numbers are still clearly more favorable for a long distance to Virgo.
Four galaxies in this sample have also distances from the extragalactic $`PL`$-relation (Lanoix Lanoix99 (1999); Lanoix et al. 1999a , 1999b , 1999c ). These galaxies are NGC 4321 (13) with $`R_{PL}=15.00\mathrm{Mpc}`$, NGC 4535 (23) with $`R_{PL}=15.07\mathrm{Mpc}`$, NGC 4548 (24) with $`R_{PL}=15.35\mathrm{Mpc}`$ and NGC 4639 (28) with $`R_{PL}=23.88\mathrm{Mpc}`$. These positions are shown as diamonds in Fig. 6. The mean distance to Virgo using the ‘A,X’ galaxies 13, 21, 24, 27, 28, 29 and 32 (i.e. the HI-deficient galaxies excluded) with TF-distance moduli is $`\mu =31.81`$ or $`R_{\mathrm{Virgo}}=22.98\mathrm{Mpc}`$ and when using the $`PL`$-distance moduli available for the three galaxies (13, 24 and 28) $`\mu =31.60`$ or $`R_{\mathrm{Virgo}}=20.93\mathrm{Mpc}`$. The difference is not large, and in both cases these Virgo core galaxies predict a distance $`R_{\mathrm{Virgo}}>20\mathrm{Mpc}`$.
We find also some other interesting features in Fig. 6. There are three galaxies (15, 16 and 22) which Federspiel et al. (Federspiel98 (1998)) classified as ‘B,X’ and two (11,12) classified as ‘B’. Together they form a clearly distinguishable substructure. It is the region D of Paper I, there interpreted as a tight background concentration. The mean distance for ‘B,X’ galaxies is $`R=30.61\mathrm{Mpc}`$ corresponding to $`\mu =32.43`$. This region is $`0.53^{\mathrm{mag}}`$ more distant than our presumed centre. We find this result satisfying because Federspiel et al. (Federspiel98 (1998)) estimated that the subgroup ‘B’ (region D in Paper I) is, on average, about $`0.46^{\mathrm{mag}}`$ farther distance than subgroup ‘A’. That our sample brings about approximately the correct relative distance between these subgroups lends additional credence to the distance estimation made in the previous section.
The region B of Paper I described as an expanding component is also conspicuously present in Fig. 6 <sup>3</sup><sup>3</sup>3 That such galaxies with negative velocity may be within a small angular distance from the Virgo cluster and still be well in the foreground was explained in Paper I as due to two things: 1) The expansion velocity must decrease away from the massive Virgo, and 2) because of projection effects , the largest negative velocities, belonging to galaxies at small distances from Virgo, are seen close to the Virgo direction.. There is, however, no clear trace of the region C of Paper I (galaxies of high velocities but lying behind the centre; cf. Fig. 8 in Paper I) unless NGC 4568 (25) actually lies at the same distance as NGC 4567 ($`R=34.28\mathrm{Mpc}`$). It should be remembered that NGC 4567/8 is classified as an interacting pair. There are, however, in Fig. 4 many galaxies at larger angular distances around NGC 4567. It is possible that they form the region C. In Paper I region C was divided into two subregions, C1 and C2. C1 was interpreted as the symmetrical counterpart to the region B (these galaxies behind Virgo are expanding away from it) and C2 was considered as a background contamination. Galaxies in region A (galaxies with high velocities lying in front of the centre) were proposed in Paper I to be presently falling into Virgo. As regards regions A and C1 it is now easy to understand that they are not separate regions but reflect the behaviour of the TB-curve: A is on the rising part and C1 on the declining part of the curve in front of the structure.
We conclude that from the expected distance-velocity pattern we have accumulated quite convincing evidence for a claim that the distance to the Virgo cluster is $`R_{\mathrm{Virgo}}=20`$$`24\mathrm{Mpc}`$ or in terms of the distance modulus $`\mu =31.51`$$`31.90`$. $`\mathrm{\Delta }\mu =0.39`$ is within $`1\sigma `$ uncertainty of our TF-sample.
## 5 The velocity field as seen from the centre of LSC
In the first part of this paper we have approached the problem of the dynamical behaviour of LSC in a more or less qualitative manner. We now proceed to present the results in a physically more relevant manner. The main difficulty in the presentation used e.g. in Figs. 14 is that the systemic velocity depends not only on the distance from LG but also on the angular distance from the centre. Basically, for each galaxy there is a unique “S-curve” depending on $`\mathrm{\Theta }`$.
Formally, the $`\mathrm{\Theta }`$-dependence is removed if the velocity-distance law is examined from the origin of the metric instead of from LG, as was done in Sect. 4.5 of Paper I. The velocity as seen from Virgo for a galaxy is solved from:
$$v(d_\mathrm{c})=\pm \frac{V_{\mathrm{obs}}(d_{\mathrm{gal}})V_{\mathrm{Virgo}}^{\mathrm{obs}}\mathrm{cos}\mathrm{\Theta }}{\sqrt{1\mathrm{sin}^2\mathrm{\Theta }/d_\mathrm{c}^2}}.$$
(4)
The relative distance from the centre $`d_\mathrm{c}=R_\mathrm{c}/R_{\mathrm{Virgo}}`$, where $`R_\mathrm{c}`$ is the distance between the galaxy considered and the centre of Virgo, is solved from Eq. 14 of Paper II and the sign is $`()`$ for $`d_{\mathrm{gal}}<\mathrm{cos}\mathrm{\Theta }`$ and $`(+)`$ otherwise.
There are, however, some difficulties involved. We are aware that the calculation of the virgocentric velocity is hampered by some sources of error. Suppose that the cosmological fluid has a perfect radial symmetry about the origin of the TB-metric. Also, the fluid elements do not interact with each other, i.e. each element obeys exactly the equations of motion of the TB-model. It follows that the measured line-of-sight velocity is a genuine projection of the element’s velocity with respect to the origin. It is presumed that the observer has made the adequate corrections for the motions induced by his immediate surroundings (e.g. Sun’s motion with respect to the Galaxy, Galaxy’s motion with respect to the LG).
Now, in practice, $`V_{\mathrm{obs}}`$ is bound to contain also other components than simply the TB-velocity. We may also have mass shells which have travelled through the origin and are presently expanding near it instead being falling in. Such a shell has experienced strong pressures (in fact, a singularity has formed to the origin) i.e. there is no causal connection to the rest of the TB-solution. Also, shells may have crossed. Again singularity has formed and the TB-solution fails (recall that TB-model describes a pressure-free cosmological fluid). Incorrect distance $`R_{\mathrm{gal}}`$ (and the scaling length $`R_{\mathrm{Virgo}}`$) will cause an error in $`v(d_\mathrm{c})`$ even when $`V_{\mathrm{obs}}`$ could be considered as a genuine projection of $`v(d_\mathrm{c})_{\mathrm{TB}}`$.
### 5.1 The two-component mass model
So far we have used a rather simple density model. From hereon we use the “two-component” model of Paper I. In this model one assumes that mass within $`\mathrm{\Theta }=6\mathrm{°}`$ at Virgo distance ($`d_{\mathrm{virial}}=0.105`$) is proportional to the Virgo virial mass and that outside this region the mass is evaluated from the simple density law (Eq. 9 in Paper II):
$$M(d_\mathrm{c})=M(d_\mathrm{c})_\alpha M(d_{\mathrm{virial}})_\alpha +\beta M_{\mathrm{virial}}.$$
(5)
The important quantity is the parameter $`A(R,T_0)`$ (Eq. 6 in Paper II). Following Ekholm (Ekholm96 (1996)) we now proceed to express it in terms of the relative distance “measured” from the origin of the metric $`dd_\mathrm{c}`$ and the deceleration parameter $`q_0`$. In terms of $`d`$ it reads:
$$A(d,T_0)=\sqrt{\frac{GM(d)}{d^3R_{\mathrm{Virgo}}^3}}\times T_0.$$
(6)
Because (cf. Eq. 9 in Paper II)
$$M(d)_\alpha =\frac{q_0H_0^2}{G}d^3R_{\mathrm{Virgo}}^3[1+k^{}d^\alpha ],$$
(7)
we find
$`A(d,T_0)`$ $`=`$ $`H_0T_0\sqrt{q_0}[1+k^{}d^\alpha (d_{\mathrm{virial}}/d)^3(1+k^{}d_{\mathrm{virial}}^\alpha )`$ (8)
$`+(\beta GM_{\mathrm{virial}})/(d^3R_{\mathrm{Virgo}}^3H_0^2)]^{1/2}.`$
Now, using $`M_{\mathrm{virial}}=7.5\times 10^{14}M_{\mathrm{}}R_{\mathrm{Virgo}}/16.8\mathrm{Mpc}`$, $`H_0T_0=C(q_0)`$ (e.g. the function $`C(q_0)=2/3`$ for $`q_0=0.5`$) and $`H_0R_{\mathrm{Virgo}}=V_{\mathrm{Virgo},\mathrm{cosm}}`$, Eq. 8 takes its final form
$`A(d,q_0)`$ $`=`$ $`C(q_0)\sqrt{q_0}[1+k^{}d^\alpha (d_{\mathrm{virial}}/d)^3(1+k^{}d_{\mathrm{virial}}^\alpha )`$ (9)
$`+(\beta \times \mathrm{cst})/(d^3V_{\mathrm{Virgo},\mathrm{cosm}}^2)]^{1/2},`$
where $`\mathrm{cst}=7.5\times 10^{14}M_{\mathrm{}}G/16.8\mathrm{Mpc}=1.92\times 10^5\mathrm{km}^2\mathrm{s}^2`$.
### 5.2 $`v(d_\mathrm{c})`$ vs. $`d_\mathrm{c}`$ diagram for $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$
We show the virgocentric diagram for $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ in the left panel of Fig. 7. The galaxies are now selected in the following manner. From the initial sample we take galaxies having $`0.105<d_\mathrm{c}1.0`$ but make no restriction on $`\mathrm{\Theta }`$. In this way we get a symmetric sample around the centre. Because the angular dependence is no longer relevant, we show the data for different ranges of $`\mathrm{log}V_{\mathrm{max}}`$: black bullets are for $`\mathrm{log}V_{\mathrm{max}}2.4`$, grey bullets for $`\mathrm{log}V_{\mathrm{max}}[2.3,2.4[`$, circles for $`\mathrm{log}V_{\mathrm{max}}[2.2,2.3[`$ and triangles for $`\mathrm{log}V_{\mathrm{max}}[2.1,2.2[`$. The straight line is Hubble law as seen from the centre and the curves (predicted velocity $`v^{}(d_\mathrm{c})`$ vs. $`d_\mathrm{c}`$) correspond to different solutions to the two-component model. We have assumed $`\alpha =2.5`$ and solved the TB-equations with Eq. 9 for $`\beta =0.5`$, $`1.0`$, $`1.5`$ and $`2.0`$ yielding mass excesses $`k^{}=0.701`$, $`0.504`$, $`0.307`$ and $`0.109`$, respectively.
Because the gradient of the $`v^{}(d_\mathrm{c})`$-curve gets quite steep as $`d_\mathrm{c}0`$, it is easier to study the difference between calculated and predicted velocities
$$\mathrm{\Delta }v(d_\mathrm{c})=v(d_\mathrm{c})v^{}(d_\mathrm{c})$$
(10)
as a function of $`d_\mathrm{c}`$. This is shown in the right panel of Fig. 7. The model values $`v^{}(d_\mathrm{c})`$ were based on $`\beta =2.0`$. In this panel we also show the mean $`\mathrm{\Delta }v`$ for each $`\mathrm{log}V_{\mathrm{max}}`$ range. For $`\mathrm{log}V_{\mathrm{max}}2.4`$ ($`N=9`$, $`\mathrm{\Delta }v=579\mathrm{km}\mathrm{s}^1`$) it is given as a black thick line, for $`\mathrm{log}V_{\mathrm{max}}[2.3,2.4[`$ ($`N=26`$, $`\mathrm{\Delta }v=646\mathrm{km}\mathrm{s}^1`$) as a grey thick line, for $`\mathrm{log}V_{\mathrm{max}}[2.2,2.3[`$ ($`N=76`$, $`\mathrm{\Delta }v=359\mathrm{km}\mathrm{s}^1`$) as a dashed line, and for $`\mathrm{log}V_{\mathrm{max}}[2.1,2.2[`$ ($`N=55`$, $`\mathrm{\Delta }v=385\mathrm{km}\mathrm{s}^1`$) as a dotted line. We note that our sample is clearly divided into two subgroups by $`\mathrm{log}V_{\mathrm{max}}=2.3`$. The slower rotators show a better fit to our chosen model. In general, galaxies in this sample have on average higher velocities than the model predicts, possibly due to some residual Malmquist bias (cf. also Figs. 5 – 6 of Paper I). It is, however, clear that the overall TB-pattern is seen in the left panel of Fig. 7 as a general decrease in $`v(d)_\mathrm{c}`$ when one approaches the centre.
### 5.3 Evidence from galaxies with $`PL`$-distances
How do the galaxies with $`PL`$-distances behave in this virgocentric representation? When selected in a similar fashion as above we find 23 galaxies shown in Fig. 8. We saw that $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ was a rather high value for them but now $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ together with $`\alpha =2.5`$ and $`\beta =2.0`$ brings about a remarkable accordance. This is particularly important in the light of the complications mentioned in the introduction to this section. It seems that at least when using high quality distances such as $`PL`$-distances those difficulties do not hamper the diagrams significantly. When this result is compared with the findings of Paper II, the distance estimate given there seems to be more and more acceptable.
There are four galaxies which show anomalous behaviour. NGC 2541 is a distant galaxy as seen from Virgo ($`R_{\mathrm{gal}}=11.59\mathrm{Mpc}`$, $`\mathrm{\Theta }=63.8\mathrm{°}`$, $`V_{\mathrm{obs}}=645\mathrm{km}\mathrm{s}^1`$) and is also close to the tangential point where small errors in distance cause large projection errors in velocity. We tested how much one needs to move this galaxy in order to find the correct predicted velocity. At $`R_{\mathrm{gal}}=13.93\mathrm{Mpc}`$, $`V_{\mathrm{pred}}=645.1\mathrm{km}\mathrm{s}^1`$ and $`v(d_\mathrm{c})=884.7\mathrm{km}\mathrm{s}^1`$ with $`\mathrm{\Delta }v=0.5\mathrm{km}\mathrm{s}^1`$. Note also that even a shift of $`1\mathrm{Mpc}`$ to $`R_{\mathrm{gal}}=12.59\mathrm{Mpc}`$ will yield $`\mathrm{\Delta }v=360.0\mathrm{km}\mathrm{s}^1`$, which is quite acceptable. When NGC 4639 ($`R_{\mathrm{gal}}=23.88\mathrm{Mpc}`$, $`\mathrm{\Theta }=3.0\mathrm{°}`$, $`V_{\mathrm{obs}}=888\mathrm{km}\mathrm{s}^1`$) is moved to $`R_{\mathrm{gal}}=21.0\mathrm{Mpc}`$, one finds $`V_{\mathrm{pred}}=886.8\mathrm{km}\mathrm{s}^1`$ and $`v(d_\mathrm{c})=3419.5\mathrm{km}\mathrm{s}^1`$ with $`\mathrm{\Delta }v=45.6\mathrm{km}\mathrm{s}^1`$. What is interesting in this shift is that in Paper II most of the galaxies tended to support $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ except this galaxy and NGC 4548. Now NGC 4639 fits perfectly. Recently Gibson et al. (Gibson99 (1999)) reanalyzed some old HST measurements finding for NGC 4639: $`\mu =31.564`$ or $`R_{\mathrm{gal}}=20.55\mathrm{Mpc}`$. As regards the two other discordant galaxies (NGC 4414 and NGC 4548) the shift to remove the discrepancy would be too large to be reasonable. At this point we cannot explain their behaviour except by assuming that they are region B galaxies of Paper I (cf. below).
Also, when galaxies with TF-distances were selected according to this normalizing distance we find better concordance with the model than for $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ (cf. Fig. 9). Note also that now only the fastest rotators differ from the rest of the sample: for $`\mathrm{log}V_{\mathrm{max}}2.4`$: $`N=12`$, $`\mathrm{\Delta }v=904\mathrm{km}\mathrm{s}^1`$ (black thick line), for $`\mathrm{log}V_{\mathrm{max}}[2.3,2.4[`$: $`N=23`$, $`\mathrm{\Delta }v=602\mathrm{km}\mathrm{s}^1`$ (grey thick line), for $`\mathrm{log}V_{\mathrm{max}}[2.2,2.3[`$: $`N=65`$, $`\mathrm{\Delta }v=512\mathrm{km}\mathrm{s}^1`$ (dashed line) and for $`\mathrm{log}V_{\mathrm{max}}[2.1,2.2[`$: $`N=49`$, $`\mathrm{\Delta }v=665\mathrm{km}\mathrm{s}^1`$ (dotted line). At relatively large distances from the centre the points in the right panel of Fig. 9 follow on average well a horizontal trend. As one approaches the centre one sees how the velocity difference $`\mathrm{\Delta }v`$ gets larger and larger. This systematic increase explains why the mean values are so high. Note that also the Cepheid galaxies NGC 4414 and NGC 4548 (and NGC 4639 if one accepts the larger distance) show a similar increasing tendency towards the centre.
Because the inward growth of $`\mathrm{\Delta }v`$ appears for both distance indicators one suspects that this behaviour is a real physical phenomenon (we cannot explain it in terms of a large scatter in the TF-relation). Neither can we explain it by a bad choice of model parameters: the effect is much stronger than the variations between different models. A natural explanation is an expanding component (referred to above as region B): galaxies with very high $`\mathrm{\Delta }v`$ are on mass shells which have fallen through the origin in past and have re-emerged as a “second generation” of TB-shells. The very quick decay of the positive velocity residuals supports this picture. The mass of the Virgo cluster is expected to slow down these galaxies quite fast (Sect. 6 in Paper I), so the effect appears at small $`d_\mathrm{c}`$.
## 6 Discussion
We found using the two-component mass model (Eq. 5) and the high quality $`PL`$-distances (Fig. 8) an acceptable fit with parameters $`\alpha =2.5`$ and $`\beta 2.0`$. Our larger TF-sample did not disagree with this model though the scatter for these galaxies is rather large. $`\beta `$ gives the Virgo cluster mass estimate in terms of the virial mass given by Tully & Shaya (Tully84 (1984)). With a distance $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ it is $`M_{\mathrm{TS}}=9.375\times 10^{14}M_{\mathrm{}}`$. By allowing some tolerance ($`\beta =1.5`$$`2.0`$) we get an estimate:
$$M_{\mathrm{Virgo}}=(1.41.875)\times 10^{15}M_{\mathrm{}}$$
(11)
### 6.1 The Virgo cluster mass, $`q_0`$, and behaviour of $`M/L`$
We have confirmed the large value of the mass-luminosity ratio for the Virgo cluster (Tully & Shaya Tully84 (1984); Paper I):
$$(M/L)_{\mathrm{Virgo}}440\beta \times (16.8\mathrm{Mpc}/R_{\mathrm{Virgo}}).$$
(12)
With $`\beta =1`$$`2`$ and $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$, $`(M/L)_{\mathrm{Virgo}}`$ ranges from 350 to 700. Note that some calculations of Paper I for different $`q_0`$ (e.g. Table 2), which were based on $`H_0=70\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`R_{\mathrm{Virgo}}=16.5\mathrm{Mpc}`$, remain valid when $`H_0=55\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`R_{\mathrm{Virgo}}=(70/55)\times 16.5\mathrm{Mpc}=21\mathrm{Mpc}`$. For example, Fig. 7 of Paper I shows that if $`(M/L)_{\mathrm{Virgo}}`$ applies everywhere, rather high values of $`q_0`$ ($`>0.1`$$`0.2`$) are favoured. A very small $`q_0`$, say 0.01, would require that $`M/L`$ outside of Virgo is several times smaller than in Virgo, i.e. the density of dark matter drops much more quickly than the density of luminous matter.
This happens also – though less rapidly – with $`q_0=0.5`$ used in this paper. This is seen from
$$\frac{(M/L)_{\mathrm{sur}}}{(M/L)_{\mathrm{Virgo}}}=\frac{M_{\mathrm{sur}}}{M_{\mathrm{Virgo}}}\times \frac{L_{\mathrm{Virgo}}}{L_{\mathrm{sur}}}.$$
(13)
The surroundings is defined as $`d_\mathrm{c}]0.105,1[`$. The luminosity ratio is $`L_{\mathrm{Virgo}}/L_{\mathrm{sur}}1/4`$ (Tully Tully82 (1982)). The mass ratio is calculated using the two-component mass model (Eq. 5) with the help of Eq. 18 of paper II. For $`M_{\mathrm{Virgo}}=2`$ the parameters needed are $`k^{}=0.109`$, $`\alpha =2.5`$, $`q_0=0.5`$ and $`h_0=0.57`$, which yield $`M(d_\mathrm{c}=1)_\alpha =4.155`$ and $`M(d_\mathrm{c}=0.105)_\alpha =0.137`$. We find $`M_{\mathrm{sur}}=M(d_\mathrm{c}=1)_\alpha M(d_\mathrm{c}=0.105)_\alpha =4.0184`$. Both masses are given in units of the Virgo virial mass. The mass-luminosity ratio becomes $`(M/L)_{\mathrm{sur}}/(M/L)_{\mathrm{Virgo}}0.5`$. When $`M_{\mathrm{Virgo}}=1`$ $`(M/L)_{\mathrm{sur}}/(M/L)_{\mathrm{Virgo}}1.25`$ and when $`M_{\mathrm{Virgo}}=1.5`$ $`(M/L)_{\mathrm{sur}}/(M/L)_{\mathrm{Virgo}}0.75`$.<sup>4</sup><sup>4</sup>4The total mass within $`d_\mathrm{c}=1`$ is 6.018 for $`\beta =2`$, 6.019 for $`\beta =1`$ and 6.020 for $`\beta =1.5`$. The Model 1 of Paper II ($`k^{}=0.606`$ and $`\alpha =2.85`$) gives 6.017 as the total mass. Thus our computational scheme works correctly because the total mass should not depend on how we distribute the matter within our mass shell. This means that with a Virgo mass slightly larger than the virial mass there is a case where the mass-luminosity ratio is constant in and outside Virgo.
How would luminous matter distribute itself? Consider the following simple exercise. Suppose the luminous matter follows a power law $`\rho _{\mathrm{lum}}(r)r^{\alpha _{\mathrm{lum}}}`$ and that the mass ratio is:
$$\frac{_{0.105}^1r^{2\alpha }𝑑r}{_0^{0.105}r^{2\alpha }𝑑r}=\frac{L_{\mathrm{sur}}}{L_{\mathrm{Virgo}}}.$$
(14)
With the luminosity ratio given above one derives for the galaxies $`\alpha _{\mathrm{lum}}2.3`$, indeed smaller than our preferred value of 2.5. Is such a steep value at all reasonable in the light of theoretical work on structure formation?
### 6.2 Comparison with the universal density profile
Tittley & Couchman (Tittley99 (1999)) discussed recently the hierarchical clustering, the universal density profile, and the mass-temperature scaling law of galaxy clusters. Using simulated clusters they studied the dark matter density profile in a Einstein-deSitter universe with $`\mathrm{\Omega }_{DM}=0.9`$, $`\mathrm{\Omega }_{\mathrm{gas}}=0.1`$ and $`\mathrm{\Lambda }=0`$. They assumed $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. Different profiles fitted their simulated data equally well. It is their discontinuous form in the first derivative which interests us:
$$\frac{\rho (r)}{\rho _c}=\{\begin{array}{cc}\delta _\gamma ^{}r^\gamma ^{},\hfill & r<r_s\hfill \\ \delta _\gamma r^\gamma ,\hfill & r>r_s\hfill \end{array}$$
(15)
They connect the overdensities as
$$\delta _\gamma =\frac{r_s^\gamma }{r_s^\gamma ^{}}\delta _\gamma ^{}.$$
(16)
Because the characteristic length $`r_s<R_{200}`$, where $`R_{200}`$ is the radius where the density contrast equals 200, the near field governed by $`\gamma `$ is not important to us. With $`\alpha =2.5`$ and $`\beta =2.0`$ in our model the mass excess $`k^{}=0.109`$. This translates into $`k=(3\alpha )\times k^{}R_{\mathrm{Virgo}}^\alpha /3=36.71`$ in the density law of Paper II: $`\delta (r)=\rho (r)/\rho _0=1+kr^\alpha `$. $`\rho _0`$ is the background density equal to the critical density $`\rho _c`$ when $`q_0=0.5`$. At the defined boundary of the Virgo cluster ($`d=0.105`$ or $`r=2.205\mathrm{Mpc}`$) we have a density excess $`\delta =5`$. For $`\beta =1.5`$, $`k=103.4`$ and $`\delta =14.32`$, and for $`\beta =1.0`$, $`k=168.4`$ and $`\delta =23.33`$. Also, because $`1+kr^\alpha kr^\alpha `$ as $`rr_s`$ comparison between our $`\alpha `$ and the $`\gamma `$ of Tittley & Couchman is acceptable. For hierarchical clustering they find $`\gamma =2.7`$ and for the non-hierarchical case $`\gamma =2.4`$. The density profile fitting dynamical behaviour of the galaxies with $`PL`$-distances is within these limits. Our mass estimate tends to be closer to the maximum values Tittley & Couchman give in their Table 3.
## 7 Summary and conclusions
In this third paper of our series we have extended the discussion of Ekholm et al. (1999a ; Paper II) to the background of Virgo cluster by selecting galaxies with as good distances as possible from the direct B-band magnitude Tully-Fisher (TF) relation. In the following list we summarize our main results:
1. Although having a rather large scatter the TF-galaxies reveal the expected Tolman-Bondi (TB) pattern well. We compared our data with TB-solutions for different distances to the Virgo cluster. It turned out that when $`R_{\mathrm{Virgo}}<20\mathrm{Mpc}`$ the background galaxies fell clearly below the predicted curves. Hence the data does not support such distance scale (cf. Figs. 1 and 2).
2. When we examined the Hubble diagram for galaxies outside the Virgo $`\mathrm{\Theta }=30\mathrm{°}`$ cone (Fig. 5) we noticed that $`H_0=60\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ is a clear upper limit for these galaxies. Together with our preferred cosmological velocity of Virgo ($`1200\mathrm{km}\mathrm{s}^1`$) we concluded that $`R_{\mathrm{Virgo}}=20\mathrm{Mpc}`$ is a lower limit.
3. In both cases any residual Malmquist bias would move the sample galaxies further away and thus make the short distances even less believable.
4. We compared our sample galaxies with $`\mathrm{\Theta }<6\mathrm{°}`$ with the Table 3 of Federspiel et al. (Federspiel98 (1998)) and found 33 galaxies in common. We established a plausible case for $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ corresponding to $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ (cf. Fig. 6). The difference between $`R_{\mathrm{Virgo}}=20\mathrm{Mpc}`$ and $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ is – in terms of the distance moduli – only $`\mathrm{\Delta }\mu =0.39`$, which is within the $`1\sigma `$ scatter of the TF-relation. Due to this scatter it is not possible to resolve the distance to Virgo with higher accuracy. Hence we claim that $`R_{\mathrm{Virgo}}=20`$$`24\mathrm{Mpc}`$.
5. Some of the kinematical features identified in Paper I were revealed also here, in particular the concentration of galaxies in front with very low velocities (interpreted as an expanding component; region B in Paper I) and the tight background concentration (region D in Paper I). The symmetric counterpart of region B (region C1) may actually be part of the primary TB-pattern.
6. The need for a better distance indicator (e.g. the I-band TF-relation) is imminent. As seen e.g. from Fig.9, the scatter in the B-band TF-relation is disturbingly large. It is also necessary to re-examine the calibration of the TF-relation with the new, and better, $`PL`$-distances. It seems that the $`PL`$-distances and the TF-distances from Theureau et al. (Theureau97 (1997)) are not completely consistent. The former tend to be somewhat smaller. This is also seen from Figs. 6 and 8. TF-distances support $`R_{\mathrm{Virgo}}=24\mathrm{Mpc}`$ and $`PL`$-distances $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$. It is, however, worth reminding that our dynamical conclusions are insensitive to the actual distance scale.
7. When we examined the Hubble diagram as it would be seen from the origin of the TB-metric, galaxies with distances from the extragalactic $`PL`$-relation fitted best to a solution with $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ in concordance with Paper II and with Federspiel et al. (Federspiel98 (1998)). We are, however, not yet confident enough to assign any error bars to this value.
8. For $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$ the region D follows well the TB-pattern (cf. Fig. 3) lending some additional credence to this distance. We quite clearly identified this background feature as the subgroup “B” of Federspiel et al. (Federspiel98 (1998)).
9. These high quality galaxies also clearly follow the expected velocity-distance behaviour in the virgocentric frame with much smaller scatter than for galaxies in Paper I or for the TF-galaxies used in this paper. The zero-velocity surface was detected at $`d_\mathrm{c}0.5`$.
10. As in Teerikorpi et al. (Teerikorpi92 (1992); Paper I), the amplitude of the TB-pattern requires that the Virgo cluster mass must be at least its standard virial mass (Tully & Shaya Tully84 (1984)) or more. Our best estimate is $`M_{\mathrm{Virgo}}=(1.5`$$`2)\times M_{\mathrm{virial}}`$, where $`M_{\mathrm{virial}}=9.375\times 10^{14}M_{\mathrm{}}`$ for $`R_{\mathrm{Virgo}}=21\mathrm{Mpc}`$.
11. Our results indicate that the density distribution of luminous matter is shallower than that of the total gravitating matter. The preferred exponent in the density power law, $`\alpha 2.5`$, agrees with the theoretical work on the universal density profile of dark matter clustering (Tittley & Couchman Tittley99 (1999)) in the Einstein-deSitter universe.
###### Acknowledgements.
This work has been partly supported by the Academy of Finland (project 45087: “Galaxy Streams and Structures in the nearby Universe” and project “Cosmology in the Local Galaxy Universe”). We have made use of the Lyon-Meudon Extragalactic Database LEDA and the Extragalactic Cepheid Database. We would like to thank the referee for useful comments. |
warning/0001/gr-qc0001002.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The scalar field theories, in which the global symmetry $`G`$ is spontaneously broken to $`H`$ in such a way that vacuum manifold $`G/H`$ has nontrivial homotopy group $`\pi _3\left(G/H\right)`$, predict the existence of the matter with an equation of state $`\epsilon +3p=0`$ called as the texture matter or k-matter.<sup>2</sup><sup>2</sup>2Also, the matter with such an equation of state naturally appears in string cosmology theories. One can obtain it from the “string-driven” effective equation of state $`\epsilon +dpnq=0`$, where $`d`$ and $`n`$ are respectively the numbers of expanding and internal contracting dimensions, $`p`$ and $`q`$ are respectively the pressure in the expanding space and shrinking dimensions . Indeed, in the “after-string” era the internal dimensions were compactified (therefore, $`n0`$), the number of expanding dimensions became the usual one, $`d=3`$. This texture-dominated era continued till epoches of standard-model particles, quarks and leptons, when the decreasing temperature and density led textures to couple up. Nowadays we probably could observe some tracks of it not only at cosmological scales but also among the fundamental properties of recently observed particles . Of course, the terminology ”texture matter” does not seem to be perfectly apposite because the real $`\sigma `$-model textures are dynamical defects and the equation of state above is not valid in general case. However, in numerous papers such a terminology was fixedly settled (see and references therein) thus we will follow it in present paper as well.<sup>3</sup><sup>3</sup>3The dualities found in some string cosmology models suggest that the matter with this equation of state can be regarded as the matter of low-energy string origin which in some sense is dual to radiation fluid (incoherent radiation). Therefore, the term “dual-radiation matter” seems to be appropriate for it as well. Some known properties of textures say that it is probably another kind of vacuum similar to the de Sitter vacuum $`\epsilon +p=0`$ (the latter is known also as the bubble matter). Let us consider, for instance, the O(4)$``$O(3) textures arising in the scalar fourplet theory described by the action
$$S\left(\stackrel{}{\varphi }\right)=\left[^\mu \stackrel{}{\varphi }_\mu \stackrel{}{\varphi }+\lambda \left(\stackrel{}{\varphi }\stackrel{}{\varphi }\eta ^2\right)^2\right]\sqrt{g}\text{d}^4x$$
in a closed FRW universe ($`0\xi \pi `$)
$$\text{d}s^2=\text{d}t^2a^2\left(t\right)\left[\text{d}\xi ^2+\text{sin}^2\xi \left(\text{d}\theta ^2+\text{sin}^2\theta \text{d}\phi ^2\right)\right].$$
Then the texture solution of winding number one,
$$\stackrel{}{\varphi }=\eta \left[\begin{array}{c}\hfill \mathrm{cos}\phi \mathrm{sin}\theta \mathrm{sin}\xi \\ \hfill \mathrm{sin}\phi \mathrm{sin}\theta \mathrm{sin}\xi \\ \hfill \mathrm{cos}\theta \mathrm{sin}\xi \\ \hfill \mathrm{cos}\xi \end{array}\right],$$
has the following stress-energy tensor,
$$T_\nu ^\mu =\frac{\eta ^2}{2a^2}\text{diag}(3,1,1,1),$$
which evidently satisfies with the above-mentioned equation of state. The zero-zero component of this tensor will be compared in Sec. 3 with a surface case.
The gravitational effects caused by 3D texture matter were intensively studied in many works . The main aim of present paper is to study the 2D fluid of global textures which forms spherically symmetric singular hypersurfaces (surfaces of discontinuities of second kind). These hypersurfaces can be interpreted both as the thin-wall approximation of the layer of bulk matter and as the brane-like objects embedded in spacetime of higher dimensionality. As such, the singular model turns to be simple enough to obtain important and instructive exact results not only when studying classical dynamics but also when considering quantum aspects. With respect to the 3D case this model appears to be the thin-wall approximation, which can elicit main features common for 2D and 3D cases.
The paper is organized as follows. In section 2 we give a comparative description of thermodynamics of 2D and 3D texture matter at finite temperature with respect to each other and with respect to bubble matter and ordinary matter represented by radiation fluid. Section 3 is devoted to classical dynamics of the isentropic singular shells “made” from 2D texture fluid. In section 4 we perform minisuperspace quantization of the singular model with provision for both the through (wormhole-like) and ordinary topology. Conclusions are made in section 5.
## 2 Comparative thermodynamics
Let us consider the thermodynamical properties of texture matter as such and in comparison with those for radiation fluid (quasi-counterpart of texture) and bubble matter $`\epsilon +p=0`$. First of all, we try to answer the question, what is thermodynamical information we can obtain from an equation of state.
The first thermodynamical law says:
$$\text{d}E=T\text{d}Sp\text{d}V.$$
(1)
On the other hand, following the definition of the entropy as a function of volume and temperature, one can write
$$\text{d}S=\frac{S}{T}\text{d}T+\frac{S}{V}\text{d}V.$$
(2)
Comparing these equations, we obtain
$`{\displaystyle \frac{S}{T}}={\displaystyle \frac{1}{T}}{\displaystyle \frac{E}{T}},`$ (3)
$`{\displaystyle \frac{S}{V}}={\displaystyle \frac{1}{T}}\left(p+{\displaystyle \frac{E}{V}}\right).`$ (4)
Then the equality of mixed derivatives yields the expression
$$p+\frac{E}{V}=T\frac{p}{T},$$
(5)
which gives opportunities to obtain internal energy as a function of volume and temperature from an equation of state. Let us introduce the densities of energy and entropy such that
$$E=\epsilon \left(T\right)V,S=s\left(T\right)V,$$
(6)
and consider barotropic matter with linear equation of state (LEOS)
$$p=\eta \epsilon .$$
(7)
Then (5) reads
$$\eta T\frac{\text{d}\epsilon }{\text{d}T}=\left(\eta +1\right)\epsilon ,$$
(8)
and we obtain the energy density
$$\epsilon =\epsilon _0T^{1+1/\eta }.$$
(9)
For instance, for 3D radiation fluid this expression yields the expected Stefan-Boltzmann law describing energy of incoherent radiation with respect to temperature:
$$\epsilon =\alpha _{SB}T^4.$$
The internal energy and pressure are, respectively,
$`E=\epsilon _0T^{1+1/\eta }V,`$ (10)
$`p=\eta \epsilon _0T^{1+1/\eta }.`$ (11)
Further, from (4), (6), (7) and (9) one can see that entropy has to be
$$S=\left(\eta +1\right)\epsilon _0T^{1/\eta }V+S_0.$$
(12)
The above-mentioned special cases of LEOS matter are illustrated in table 1.<sup>4</sup><sup>4</sup>4The features presented there (first of all, the unusual inverse dependence of energy and entropy on temperature which means that increasing of temperature is energetically favorable) can serve also for revealing of free texture matter in present era, e.g., inside superhot objects with appropriate energy density (of order of GUT scale, $`10^{13}`$ TeV). We can observe, e.g., that texture matter cannot approach zero temperature even formally (without the third law of thermodynamics) because its energy diverges; bubble and texture matter have nonzero minimal energy unlike ordinary matter $`\eta >0`$ including ultrarelativistic radiation fluid. It seems to be another argument to the advantage of interpretation of the texture matter as a specific vacuum state similar to the de Sitter one.
## 3 Thin-wall model
Beginning from the classical works formalism of surface layers has been widely described in the literature (see Refs. for details). The three-dimensional singular embeddings appear to be both interesting extended objects as such, and simple (but realistic) models of four-dimensional phenomena. From the viewpoint of the general physics of extended objects the concept “singular hypersurface” has to be the next-order approximation, after the “point particle” one, which takes into account both external, kinetic and dynamical, properties and internal structure (surface pressure, mass density, temperature etc.).
So, one considers the infinitely thin isentropic layer of matter with the surface stress-energy tensor of a perfect fluid in general case (we use the units $`\gamma =c=1`$, where $`\gamma `$ is the gravitational constant)
$$S_{ab}=\sigma u_au_b+p\left(u_au_b+{}_{}{}^{\left(3\right)}g_{ab}^{}\right),$$
where $`\sigma `$ and $`p`$ are the surface mass-energy density and pressure respectively, u is the unit tangent vector, $`{}_{}{}^{\left(3\right)}g_{ab}^{}`$ is the three-metric of the shell’s hypersurface. We suppose that this shell is spherically symmetric, closed, and hence divides the whole manifold into the two regions $`\mathrm{\Sigma }^\pm `$. Also we suppose the metrics of the space-times outside $`\mathrm{\Sigma }^+`$ and inside $`\mathrm{\Sigma }^{}`$ of a spherically symmetric shell to be of the form
$$\text{d}s_\pm ^2=\left[1+\mathrm{\Phi }^\pm \left(r\right)\right]\text{d}t_\pm ^2+\left[1+\mathrm{\Phi }^\pm \left(r\right)\right]^1\text{d}r^2+r^2\text{d}\mathrm{\Omega }^2,$$
(13)
where $`\text{d}\mathrm{\Omega }^2`$ is the metric of the unit two-sphere. Of course, we have some loss of generality but it is enough for further. It is possible to show that if one introduces the proper time $`\tau `$, then the 3-metric of a shell can be written in the form
$${}_{}{}^{\left(3\right)}\text{d}s^2=\text{d}\tau ^2+R^2\text{d}\mathrm{\Omega }^2,$$
(14)
where $`R\left(\tau \right)`$ is a proper radius of a shell. Define a simple jump of the second fundamental forms across a shell as $`\left[K_b^a\right]=K_b^{a+}K_b^a`$, where
$$K_b^{a\pm }=\underset{n\pm 0}{lim}\frac{1}{2}{}_{}{}^{\left(3\right)}g_{}^{ac}\frac{}{n}{}_{}{}^{\left(3\right)}g_{cb}^{},$$
(15)
where $`n`$ is a proper distance in normal direction. The Einstein equations on a shell then yield equations which are the well-known Lichnerowicz-Darmois-Israel junction conditions
$$\left(K_b^a\right)^+\left(K_b^a\right)^{}=4\pi \sigma \left(2u^au_b+\delta _b^a\right),$$
(16)
Besides, an integrability condition of the Einstein equations is the energy conservation law for shell matter. In terms of the proper time it can be written as
$$\text{d}\left(\sigma {}_{}{}^{\left(3\right)}g\right)+p\text{d}\left({}_{}{}^{\left(3\right)}g\right)+{}_{}{}^{\left(3\right)}g\left[T\right]\text{d}\tau =0,$$
(17)
where $`\left[T\right]=\left(T^{\tau n}\right)^+\left(T^{\tau n}\right)^{}`$, $`T^{\tau n}=T_\alpha ^\beta u^\alpha n_\beta `$ is the projection of stress-energy tensors in the $`\mathrm{\Sigma }^\pm `$ space-times on the tangent and normal vectors, $`{}_{}{}^{\left(3\right)}g=\sqrt{det(^{\left(3\right)}g_{ab})}=R^2\mathrm{sin}\theta `$.
We assume that our shell carries no charges on a surface and contains no matter inside itself. If we define $`M`$ to be the total mass-energy of the shell then one can suppose the external and internal spacetimes to be Schwarzschild and Minkowskian respectively:
$$\mathrm{\Phi }^+=\frac{2M}{R},\mathrm{\Phi }^{}=0.$$
(18)
After straightforward computation of extrinsic curvatures the $`\theta \theta `$ component of (16) yields the equation of motion of the perfect fluid neutral hollow shell
$$ϵ_+\sqrt{1+\dot{R}^2\frac{2M}{R}}ϵ_{}\sqrt{1+\dot{R}^2}=\frac{m}{R},$$
(19)
where
$$m=4\pi \sigma R^2$$
(20)
is interpreted as the (effective) rest mass, $`\dot{R}=\text{d}R/\text{d}\tau `$ is a proper velocity of the shell, $`ϵ_+=\text{sign}\left(\sqrt{1+\dot{R}^22M/R}\right)`$, $`ϵ_{}=\text{sign}\left(\sqrt{1+\dot{R}^2}\right)`$. It is well-known that $`ϵ=+1`$ if $`R`$ increases in the outward normal direction to the shell, and $`ϵ=1`$ if $`R`$ decreases. Thus, under the choice $`ϵ_+=ϵ_{}`$ we have an ordinary (black hole type) shell, whereas at $`ϵ_+=ϵ_{}`$ we have the thin-shell traversable wormhole .
Let us consider the conservation law (17). One can obtain that $`\left[T\right]`$ is identically zero for the spacetimes (18). Further, if we assume the 2D texture equation of state of the shell’s matter,
$$\sigma +2p=0,$$
(21)
then, solving the differential equation (17) with respect to $`\sigma `$, we obtain
$$\sigma =\frac{\alpha }{2\pi R},$$
(22)
hence
$$m=2\alpha R,$$
(23)
where $`\alpha `$ is the dimensionless integration constant which can be determined via surface mass density (or pressure) at fixed $`R`$. The surface energy density determined by (22) appears to be the 2D analogue of the cosmological $`T_0^0`$ component from the Sec. 1 if one takes into account the reduction of dimensionality. This is an expected result: from the viewpoint of the 2D observer “living” on the shell it seems for him to be the whole universe with the scale factor $`R`$. Thus, our 2D fluid model indeed not only considers the established trace properties of the texture stress-energy tensor but also restores its components for the surface case. In this connection the integration constant $`\alpha `$ obtains the sense of the (squared) topological charge $`\eta `$. The topological nature of the textures will brightly show itself at the end of of this section when we will study the texture fluid singular layers with the vanishing total gravitational mass-energy.
Equations (19) and (23) together with the choice of the signs $`ϵ_\pm `$ completely determine the motion of the thin-wall texture. In conventional general relativity it is usually supposed that masses are nonnegative. However, keeping in mind possible wormhole and quantum extensions of the theory , we will not restrict ourselves by positive values and consider general case of arbitrary (real) masses. Then forbidden and permitted signs of this values can be determined from table 2. Let us find now the trajectories of 2D textures. Integrating (19) we obtain the transcendental equation for $`R\left(\tau \right)`$
$$\tau /M=J\left(R/M\right)J\left(R_0/M\right),$$
(24)
where
$`J\left(y\right)=\{\begin{array}{cc}\frac{1}{\alpha ^21}\left\{\frac{1}{\alpha }\sqrt{Z_1}+\frac{1}{2\sqrt{1\alpha ^2}}\mathrm{arcsin}Z_2\right\},\hfill & \alpha ^2<1,\hfill \\ \pm \frac{1}{6}\sqrt{4y+1}\left(2y1\right),\hfill & \alpha =\pm 1,\hfill \\ \frac{1}{\alpha ^21}\left\{\frac{1}{\alpha }\sqrt{Z_1}\frac{1}{2\sqrt{\alpha ^21}}\text{arccosh}Z_2\right\},\hfill & \alpha ^2>1,\hfill \end{array}`$ (28)
$`Z_1=\alpha ^2\left(\alpha ^21\right)y^2+\alpha ^2y+1/4,`$
$`Z_2=2\alpha \left(\alpha ^21\right)y+\alpha .`$
Thus, in dependence on the parameter $`\alpha ^2`$ one can distinguish elliptical, parabolical and hyperbolical trajectories. Let us consider below the consistency conditions which yield permitted domains of $`\alpha `$ and $`y=R/M`$ for each from three cases $`\alpha ^2`$.
(a) Hyperbolic trajectories ($`\alpha ^2>1`$).
Following (24) the next two conditions should be satisfied jointly:
$$Z_10,Z_20.$$
(29)
Define
$`y_\pm ={\displaystyle \frac{1}{2\alpha }}{\displaystyle \frac{1}{\alpha \pm 1}},`$ (30)
$`\overline{y}={\displaystyle \frac{1}{2\left(1\alpha ^2\right)}},`$ (31)
and consider the two subcases:
(a.1) $`\alpha <1`$. Then $`y_+<\overline{y}<y_{}<0`$ and inequalities (29) can be reduced respectively to
$$\left\{yy_+\right\}\left\{yy_{}\right\},y\overline{y},$$
that yields
$$yy_+.$$
(32)
(a.2) $`\alpha >1`$. Then $`y_{}<\overline{y}<y_+<0`$ and inequalities (29) can be reduced respectively to
$$\left\{yy_{}\right\}\left\{yy_+\right\},y\overline{y},$$
that yields
$$yy_+.$$
(33)
Thus, inequalities (32) and (33) determine permitted regions $`\{\alpha ,R/M\}`$ for hyperbolical trajectories.
(b) Elliptic trajectories ($`\alpha ^2<1`$).
In the same way as above we can obtain the next restrictions:
$$Z_10,1Z_21,$$
(34)
and consider the two subcases:
(b.1) $`1<\alpha <0`$. Then $`y_+>0`$ and $`y_{}<0`$, and inequalities (34) read
$$y_{}yy_+.$$
(35)
(b.2) $`0<\alpha <1`$. Then $`y_+<0`$ and $`y_{}>0`$, and
$$y_+yy_{}.$$
(36)
(c) Parabolic trajectories ($`\alpha ^2=1`$).
We obtain that $`y`$ should obey
$$y1/4.$$
(37)
The cases (a)-(c) are illustrated in figure 1 which represents dependence $`y=R/M`$ on $`\alpha `$. Note, we did not restrict signs of mass, therefore, table 2 should be kept in mind.
Let us study now equilibrium states of thin-wall textures. Differentiating (19) with respect to $`\tau `$, we obtain
$$\frac{\ddot{R}+M/R^2}{ϵ_+\sqrt{1+\dot{R}^22M/R}}\frac{\ddot{R}}{ϵ_{}\sqrt{1+\dot{R}^2}}=0,$$
(38)
which independently could be obtained from junction conditions (16). Then in equilibrium state $`\dot{R}=\ddot{R}=0`$ we obtain
$$M=0,$$
i. e., the texture fluid in equilibrium has zero total gravitational mass that is already well-known . Another way to show this feature is to generalize (19), (38) by inserting the mass $`M_{}`$ inside the shell, then the external and internal spacetimes turn to be the Schwarzschild ones with masses $`M^+`$ and $`M^{}`$ respectively. Performing the analogical calculations we would obtain that in equilibrium: at $`ϵ_+=ϵ_{}`$ the static masses $`M^+=M^{}`$ and $`\alpha =0`$ (that evidently corresponds to the already decayed shell because $`\alpha `$ is the genuine criterion of existence and distinguishability of the shell) whereas at $`ϵ_+=ϵ_{}`$ the static masses $`M_\pm `$ should vanish but $`\alpha `$ should not, giving the nonzero value for static radius, i.e. again the static texture shell makes no contribution to the total gravitational mass of the system.
In other words, if in the (generalized) equations (19), (38) we even suppose $`M^+=M^{}=0`$ identically then we do not obtain $`\alpha 0`$ with necessity. It illustrates the fact that at some choice of signs $`ϵ_\pm `$ we come to a non-trivial case despite the total masses are zero. Indeed, at $`\dot{R}=\ddot{R}=0`$ we have $`\alpha 0`$ if $`ϵ_+=ϵ_{}`$, i.e., for wormhole shells (as for the ordinary hollow texture-shells, then always $`\alpha _{\text{st}}=0`$, and thus they cannot have equilibrium states). Thus there exists the so-called zeroth traversable texture wormhole (ZTTW): one can see that junction remains to be possible at $`ϵ_+=ϵ_{}`$ and $`\alpha 0`$ (among the rest linear equations of state the texture’s one (21) appears to be unique in this sense). Therefore, ZTTW represents itself the specific vacuum-like topological barrier (characterized only by $`\alpha `$, see eq. (22) and comments after it) between two flat spacetimes<sup>5</sup><sup>5</sup>5It does mean also that in very general case the flat spacetime cannot be regarded as absolutely matter-free one even on the classical level: despite textures are the defects of quantum-field nature and (after)string origin, their condensate, texture matter, is macroscopic. The locally flat spacetime is globally determined at least up to the foliation by nested ZTTW’s walls. which has no observable mass but possesses nontrivial internal structure and inertial external dynamics
$`\sqrt{1+\dot{R}^2}=\left|\alpha \right|,\ddot{R}0,`$
thereby the restriction $`ϵ_{}\alpha >0`$ should be satisfied as it can be readily seen from eqs. (19), (23), (38).
## 4 Minisuperspace quantization
Following the Wheeler-DeWitt’s approach, in quantum cosmology the whole Universe is considered quantum mechanically and is described by a wave function. The minisuperspace approach appears to be the direct application of Wheeler-DeWitt’s quantization procedure for (2+1)-dimensional singular hypersurfaces having own internal three-metric (see and references therein). So, let us consider the minisuperspace model described by the Lagrangian:
$$L=\frac{m\dot{R}^2}{2}\alpha \left(1\alpha ^2\right)R+\alpha M+\frac{M^2}{4\alpha R},$$
(39)
where $`m`$ was defined by (23). If we define
$$U=\alpha \left(1\alpha ^2\right)R\alpha M\frac{M^2}{4\alpha R},$$
then the equation of motion following from this Lagrangian is
$$\frac{\text{d}}{\text{d}\tau }\left(m\dot{R}\right)=\frac{m_R\dot{R}^2}{2}U_R,$$
(40)
where subscript “$`R`$” means derivative with respect to $`R`$. Using time symmetry we can easy decrease order of this differential equation and obtain
$$\dot{R}^2=\frac{2}{\alpha R}\left(HU\right),$$
(41)
where $`H`$ is the integration constant. This equation coincides with double squared (19) at (23) when we suppose $`H=0`$ as a constraint. Thus, our Lagrangian indeed describes dynamics of the thin-wall texture up to the topological wormhole/blackhole division which was described by the signs $`ϵ_\pm `$. However, we always can restore the topology $`ϵ_\pm `$ both at classical (rejecting redundant roots) and quantum (considering appropriate boundary conditions for the corresponding Wheeler-DeWitt equation, see below) levels.
Further, at $`\mathrm{\Pi }=m\dot{R}`$ the (super)Hamiltonian is
$$=\mathrm{\Pi }\dot{R}L=H=0.$$
(42)
The prefix “super” means that in general case $``$ has to be a functional defined on the superspace which is the space of all admissible metrics and accompanying fields. In spherically symmetric case the world sheet of a singular hypersurface is determined by a single function, proper radius $`R\left(\tau \right)`$.
To perform quantization we replace momentum by the operator $`\widehat{\mathrm{\Pi }}=i/R`$ (we assume Planckian units), and (42) yields the Wheeler-DeWitt equation for the wave function $`\mathrm{\Psi }\left(R\right)`$:
$$\mathrm{\Psi }_{RR}+\left[M^2+4M\alpha ^2R4\alpha ^2\left(1\alpha ^2\right)R^2\right]\mathrm{\Psi }=0.$$
(43)
One can see the main advantage of the minisuperspace approach, namely, it does not require any time slicing on the basic manifold.
Further, the important remark should be made now. Last equation can be reduced to that for quantum harmonical oscillator, but not in all cases: the oscillator’s equation is defined on the line $`(\mathrm{},+\mathrm{})`$ whereas in our case the extension of an application domain on the whole axis $`R(\mathrm{},+\mathrm{})`$ seems to be physically ill-grounded in the major cases, therefore we should study the quantum theory on the half-line $`[0,\mathrm{})`$. Strictly speaking, such a situation happens also in quantum field theory then the mathematical procedure known as the Langer modification had been used there . The similar transformation we perform below to obtain the required solution.
In the case $`R[0,+\mathrm{})`$ equation (43) in general cannot be resolved in terms of the parabolical cylinder functions and Hermite polynomials. Fortunately, a solution can be expressed in terms of the functions well-defined on the half-line $`[0,+\mathrm{})`$. To show up this feature let us perform, at first, the following substitution
$$x=Rb/2ax[b/2a,+\mathrm{}),$$
(44)
where $`a=4\alpha ^2\left(1\alpha ^2\right)`$, $`b=4\alpha ^2M`$. Then (43) can be rewritten as
$$\mathrm{\Psi }_{xx}+\frac{\mathrm{\Delta }4a^2x^2}{4a}\mathrm{\Psi }=0,$$
(45)
where $`\mathrm{\Delta }=\left(4\alpha M\right)^2`$. Further, considering (44) it can be seen that at $`a>0`$ ($`\alpha ^2<1`$) the substitution
$$z=\sqrt{a}x^2$$
(46)
has not to be the injective mapping and acts like the baker’s transformation around $`x=0`$ that provides the important property
$$z[0,+\mathrm{}).$$
(47)
Then the transformation
$$\mathrm{\Psi }\left(x\right)=e^{z/2}\sqrt{z}\omega \left(z\right)$$
(48)
rewrites (45) in a form of the confluent hypergeometric equation
$$z\omega _{zz}+\left(c^{}z\right)\omega _za^{}\omega =0,$$
(49)
whose solutions are the regular Kummer functions $`\text{M}(a^{},c^{};z)`$, where
$$a^{}=\frac{3}{4}\frac{\mathrm{\Delta }}{16a^{3/2}},c^{}=\frac{3}{2}.$$
Therefore, the true solutions of (43) at $`R[0,+\mathrm{})`$ are the functions (up to multiplicative constant):
$$\mathrm{\Psi }=e^{z/2}\sqrt{z}\text{M}(\frac{3}{4}\frac{\mathrm{\Delta }}{16a^{3/2}},\frac{3}{2};z),$$
(50)
where
$$z=2\left|\alpha \right|\sqrt{1\alpha ^2}\left[R\frac{M}{2\left(1\alpha ^2\right)}\right]^2.$$
(51)
Further, if we wish to determine bound states we should require $`\mathrm{\Psi }\left(R=+\mathrm{}\right)=0`$; the also required condition $`\mathrm{\Psi }\left(R=0\right)=0`$ (that corresponds, according to aforesaid, to $`\mathrm{\Psi }\left(z=0\right)=0`$) has been already satisfied by the choice of the solution (50) when a one integration constant was used (the second constant always remains to normalize a solution) . In this case
$$a^{}=n,$$
(52)
$`n`$ is a nonnegative integer, and the Kummer function moves to the Laguerre polynomials. From last expression we obtain the mass spectrum of the thin-wall ordinary texture in a bound state:
$$M_n=\pm \sqrt{2\left|\alpha \right|\left(4n+3\right)}\left(1\alpha ^2\right)^{3/4},$$
(53)
which evidently has to be a subset of the oscillator’s spectrum, the Laguerre polynomials are connected with the Hermite ones through the transformation (46). Thus, our procedure has cut out from the oscillator’s eigenfunctions and eigenvalues those which satisfy with the boundary conditions on a half-line.
## 5 Conclusion
In present paper we considered texture matter and singular hypersurfaces made from it. First of all, we studied thermodynamical properties of 2D and 3D texture matter in comparison with radiation fluid and bubble matter. These properties say that textures can be imagined as the specific vacuum state having the congeniality with the already known de Sitter vacuum and duality with radiation fluid . We obtained equations of motion of selfgravitating texture objects, showed that classical motion can be elliptical (finite), parabolical or hyperbolical, thereby permitted and forbidden regions of motion was determined. We showed up that neutral textures in equilibrium have zero total gravitational mass as was expected. Moreover, it was established that there can exist the nontrivial wormhole-textures having vanishing total mass and matching two flat spacetimes. Finally, we considered quantum aspects of the theory by means of Wheeler-DeWitt’s minisuperspace quantization procedure, obtained the exact wave function and spectrum of bound states. |
warning/0001/nucl-th0001034.html | ar5iv | text | # Correlation and Finite Interaction-Range Effects in High-Energy Electron Inclusive Scattering
## 1 Introduction
Experimental data on high-energy electron-nucleus inclusive scattering have been accumulated for a decade . The experiments were mostly done at SLAC, and will be done at the Jefferson Lab. Also a new facility called MUSES (Multi-Use Experimental Storage ring) is planned at RIKEN . Precise data will become available with these facilities and provide us with a way of studying where and how the quark-gluon degrees of freedom may come into sight.
The process that we study is the $`(e,e^{})`$ inclusive reaction (Fig. 1) in the nucleon quasi-elastic region, i.e.,
$$e+Ae^{}+N+(A1)^{}.$$
(1.1)
The incident electron has a four-momentum, $`k_\mu `$ $`=(ϵ,𝐤_\mathrm{e})`$, and the outgoing electron, $`k_\mu ^{}`$ $`=(ϵ^{},𝐤_{}^{}{}_{\mathrm{e}}{}^{})`$. The momentum transfer is $`q_\mu `$ $`=k_\mu k_\mu ^{}`$ $`=(\omega ,𝐪)`$. We are interested in the quasi-elastic region, i.e., $`|𝐪|2[\mathrm{GeV}/\mathrm{c}]`$, and $`\omega 2[\mathrm{GeV}]`$. So the elementary process is considered to be mostly the electron-nucleon elastic scattering, $`e+N`$ $`e+N`$, and the final-state interaction (FSI) of the struck nucleon is what we wish to study here.
In this work we calculate the effect of the FSI, taking into account only the nucleonic degrees of freedom. We use the Glauber approximation for the FSI, including the nuclear correlation. We believe that our treatment is a more systematic approach than the previous works.
The theoretical treatment of this reaction is similar to that we used for the nuclear transparency in the $`(e,e^{}p)`$ reaction at large-momentum transfers , . This process has recently attracted many researchers in connection with a speculated phenomenon of the color transparency . The present formulation allows us to describe the $`(e,e^{})`$ and $`(e,e^{}p)`$ reactions in a unified way.
The nuclear transparency for electron scattering can be defined as the ratio of the $`(e,e^{}p)`$ response to the inclusive $`(e,e^{})`$ response in the quasi-elastic region, and is conveniently formulated by the Green’s function method . We will give here a way of calculating the response functions based on the Glauber approximation for the Green’s function, without taking account of the proton internal structure.
The contents of this paper are as follows: In sec. 2 we explain our formulation. We explain our standpoint about the elementary process, i.e., the electron-nucleon cross section, in subsec. 2.1. We derive a closed form of the inclusive cross section based on the Glauber approximation in subsec. 2.2. The definition of the FSI-function is also given there. The expression of the response function for nuclear matter is derived in subsec. 2.3, and the definition of our “convolution” function is found. We derive the approximate expressions of the FSI-function when we apply zero-range approximations for the nucleon-nucleon potential in subsec. 2.4. In sec. 3 we compare our treatment of the FSI with others. In subsec. 3.1 we discuss the difference and the relation between our formulation and the theory of Gersch, Rodriguez, and Smith , , . A similar comparison with the optical potential formalism , is given in subsec. 3.2. In sec. 4 we show our numerical results and discuss their indications. We show the FSI-function in subsec. 4.1, and the inclusive cross sections in subsec. 4.2. We summarize our results in sec. 5. In the Appendix we show a method for deriving a nucleon-nucleon potential from the experimental scattering amplitude in the eikonal approximation.
## 2 Formulation
In this section we explain our formulation of the electron-nucleus inclusive scattering based on the Green’s function method . We employ the Glauber approximation for the FSI of the struck nucleon.
### 2.1 Choice of the $`eN`$ Cross Section
Let us first explain our standpoint about the elementary process of the $`(e,e^{})`$ inclusive reaction.
We start with the $`(e,e^{})`$ inclusive differential cross section on a target nucleus, $`A`$, which is assumed to be expressed as
$`{\displaystyle \frac{d\sigma _{\mathrm{eA}}}{d\mathrm{\Omega }d\omega }}={\displaystyle \frac{d\sigma _{\mathrm{eN}}}{d\mathrm{\Omega }}}_{\mathrm{el},\mathrm{on}}S(\omega ,𝐪),`$ (2.1)
where the cross section on the r.h.s. is that of on-shell electron-nucleon elastic scattering for the same incident energy and scattering angle, and $`\mathrm{}`$ implies averaging over the spin, and the isospin, of the target nucleons. The difference between the longitudinal and the transverse responses is also neglected.
This factorized form, eq. (2.1), is an approximation, on which some comments are in order: The nucleon struck by the electron is not in free space, but in the nuclear medium both in the initial and in the final states. Such a nucleon is often said to be off-the-mass shell, and its electromagnetic form factors to be used in the calculation of the cross section are generally different from those of the on-shell nucleon. However, we have no reliable way of estimating the differences. In fact, we need a model for the internal structure of the nucleon as well as a model of its interaction with the nuclear medium to study the off-shell effects. We note here that the off-shell effects are essentially dynamical and cannot be obtained by kinematical considerations , . In the following we discuss how off-shell the nucleon will be in the case of inclusive $`(e,e^{})`$ on nuclear matter.
Let us assume that the single nucleon spectrum in the nuclear matter is given by
$`(EV)^2=(m_\mathrm{N}+S)^2+𝐩^2,`$ (2.2)
where $`E`$, $`𝐩`$ and $`m_\mathrm{N}`$ are the energy, the spatial momentum and the mass of the nucleon, respectively, with $`V`$ $`(S)`$ being the single particle vector (scalar) potential. The degree of the nucleon being off-the-mass shell is measured by
$`\delta m^2`$ $``$ $`E^2𝐩^2m_\mathrm{N}^2=2EVV^2+2m_\mathrm{N}S+S^2.`$ (2.3)
The Walecka model gives $`V`$ $`=300[\mathrm{MeV}]`$, $`S`$ $`=350[\mathrm{MeV}]`$ (Case I). We will also consider the two extreme cases, i.e., $`V`$ $`=0[\mathrm{MeV}]`$, $`S`$ $`=50[\mathrm{MeV}]`$ (Case II), and $`V`$ $`=50[\mathrm{MeV}]`$, $`S`$ $`=0[\mathrm{MeV}]`$ (Case III), keeping the sum, $`S+V`$, to be the same. Figure 2 shows $`\delta m^2/m_\mathrm{N}^2`$ as a function of $`E`$ for these three cases. In all the cases, the nucleon in the initial state, where $`Em_\mathrm{N}`$, is not very off-shell, because $`\delta m^2/m_\mathrm{N}^2`$ $`0.1`$. On the other hand, the nucleon in the final state, where $`E`$ $`23m_\mathrm{N}`$, is far off-shell, i.e., $`\delta m^2/m_\mathrm{N}^2`$ $`O(1)`$, in the case (I), while it is not very off-shell in the cases (II) and (III). This exercise tells us that one should be careful in choosing the off-shell kinematics for the inclusive $`(e,e^{})`$ response. What may be important in the high-energy inclusive reaction is the off-shellness of the nucleon in the final state, not in the initial state. This has been stressed previously by, e.g., Uchiyama, et al. , and Ciofi and Simula .
The off-shell cross section proposed so far is that of the half-off-shell , . In the prescription the bound nucleon in the initial state is treated as off-the-mass shell, while the nucleon in the final state is treated on-the-mass shell. This is reasonable for the $`(e,e^{}p)`$ semi-inclusive reaction, but not for the $`(e,e^{})`$ inclusive. Furthermore, for the $`(e,e^{})`$ inclusive process the nuclear matter limit ($`A\mathrm{}`$) can be taken, while for the $`(e,e^{}p)`$ semi-inclusive process that limit cannot be defined, because the nucleon in the final state cannot be free due to the infinitely extended nuclear matter. Therefore, the off-shell cross section, which is presently available, is not suitable for the elementary process of the $`(e,e^{})`$ inclusive reaction.
We thus choose to avoid using such off-shell cross sections, and use the on-shell cross section. This is the basis for our choice of the on-shell kinematics for the initial nucleon. The construction of the models describing the internal structure of the nucleon and its interaction with the nuclear medium is beyond the scope of this work, and we use the simple factorized form with the on-shell form factors.
Therefore the agreement of the numerical results with the experimental data is not necessarily the aim of this work. In the following we concentrate on a consistent treatment of the FSI of the struck nucleon in the Green’s function method with the Glauber approximation.
### 2.2 Green’s Function Method
We now proceed to the derivation of the closed form of the response function, $`S(\omega ,𝐪)`$, in eq. (2.1) in the Green’s function method with the Glauber approximation. The response function is defined by
$`S(\omega ,𝐪)`$ $``$ $`{\displaystyle \underset{X}{}}|X|\widehat{O}(𝐪)|A|^2\delta (\omega E_X)`$ (2.4)
$`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}R(\omega ,𝐪),`$
where
$`R(\omega ,𝐪)A|\widehat{O}^{}(𝐪){\displaystyle \frac{1}{\omega \widehat{H}_A+i\eta }}\widehat{O}(𝐪)|A,(\eta >0).`$ (2.5)
The initial state, $`A`$, is the ground state of the target nucleus, and the sum over the final states, $`X`$, is taken, because we are considering the inclusive process. $`\widehat{H}_A`$ is the full hamiltonian of the target nucleus. The energy of the final state, $`E_X`$, in eq. (2.4) is measured from that of the ground state of the target.
$`\widehat{O}(𝐪)`$ in eq. (2.4) is a hard interaction operator defined by
$`\widehat{O}(𝐪)`$ $``$ $`\mathrm{exp}\{i𝐪𝐫_1\},`$ (2.6)
where $`𝐫_1`$ is the coordinate of the struck nucleon denoted as $`1`$. We neglect the interference effects among different nucleons with which the photon couples.
We decompose the full hamiltonian, $`\widehat{H}_A`$, in eq. (2.4) as
$`\widehat{H}_A=\widehat{K}_1+\widehat{H}_{A1}+\widehat{V}_{1,A1}.`$ (2.7)
Here $`\widehat{K}_1`$ is the kinetic-energy operator for the struck nucleon, $`1`$, and $`\widehat{H}_{A1}`$ is the hamiltonian for the residual nucleus. $`\widehat{V}_{1,A1}`$ is the interaction of the struck nucleon with other nucleons in the target nucleus, which is assumed to be expressed as a sum of the two-body interactions, $`V_{\mathrm{NN}}(𝐫)`$,
$`V_{1,A1}(𝐫_1;𝐫_2,\mathrm{},𝐫_A)={\displaystyle \underset{j=2}{\overset{A}{}}}V_{\mathrm{NN}}(𝐫_1𝐫_j).`$ (2.8)
For numerical calculations we construct the nucleon-nucleon potential, $`V_{\mathrm{NN}}(𝐫)`$, from the phenomenological scattering amplitude describing the observed cross sections. The potential is complex in the energy region where we are interested in. We will leave the details in the Appendix.
Introducing the Green’s function for the nucleon, $`1`$, interacting with the residual nucleus in the fixed-scatterer approximation which is given by
$`G(\omega ;𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)𝐫_1|{\displaystyle \frac{1}{\omega \widehat{K}_1\widehat{V}_{1,A1}+i\eta }}|𝐫_{}^{}{}_{1}{}^{},`$ (2.9)
the response function is given as
$`S(\omega ,𝐪)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}{\displaystyle 𝑑𝐫_1𝑑𝐫_{}^{}{}_{1}{}^{}𝑑𝐫_2\mathrm{}𝑑𝐫_A\mathrm{exp}\{i𝐪(𝐫_{}^{}{}_{1}{}^{}𝐫_1)\}}`$
$`\times \mathrm{\Psi }_A^{}(𝐫_1,𝐫_2,\mathrm{},𝐫_A)\mathrm{\Psi }_A(𝐫_{}^{}{}_{1}{}^{},𝐫_2,\mathrm{},𝐫_A)G(\omega ;𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A),`$
where $`\mathrm{\Psi }(𝐫_1,𝐫_2,\mathrm{},𝐫_A)`$ is the $`A`$-body nuclear wave function. The fixed-scatterer approximation implies that the excitation energy, $`E_n`$, of the nucleus is neglected, and that the nucleon-nucleon interaction is assumed to be local . The struck nucleon is considered to become energetic enough to justify the fixed-scatterer approximation for the FSI due to $`V_{\mathrm{NN}}(𝐫)`$.
The Green’s function is now a one-body Green’s function with a multi-centered potential, and we here introduce the eikonal expression for it ,
$`G(\omega ;𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)`$ (2.11)
$`=`$ $`{\displaystyle \frac{i}{v}}\delta (𝐛_1𝐛_{}^{}{}_{1}{}^{})\theta (z_1z_1^{})\mathrm{exp}\{ip(z_1z_1^{}){\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }{\displaystyle \underset{j=2}{\overset{A}{}}}V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_j)\}`$
$`+(\mathrm{backward}\mathrm{piece}),`$
where $`p`$ is related to the energy loss, $`\omega `$, by
$`\omega +m_\mathrm{N}E_p=\sqrt{p^2+m_\mathrm{N}^2},`$ (2.12)
with $`m_\mathrm{N}`$ being the nucleon mass, and $`v`$ $`=p/E_p`$. We are free to choose the direction of the $`z`$-axis at this stage, but will take it to be in the direction of $`𝐪`$, since the exponential in the integrand of eq. (2.2) is strongly oscillating in this direction for large-$`|𝐪|`$. Here we have also neglected $`E_n`$, because it is small compared to $`\omega `$, which is roughly 2 \[GeV\]. The vectors, $`𝐫_1`$, $`𝐫_{}^{}{}_{1}{}^{}`$, and $`𝐫_{}^{\prime \prime }{}_{1}{}^{}`$, are thus decomposed as,
$`𝐫_1=(𝐛_1,z_1),𝐫_{}^{}{}_{1}{}^{}=(𝐛_{}^{}{}_{1}{}^{},z_1^{}),𝐫_{}^{\prime \prime }{}_{1}{}^{}=(𝐛_1,z_1^{\prime \prime }).`$ (2.13)
The backward piece in eq. (2.11) becomes important only when we discuss the sum rule, and will be omitted in the following. Then the imaginary part of the Green’s function is obtained as
$`\mathrm{Im}G(\omega ;𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A){\displaystyle \frac{1}{2i}}\{G(\omega ;𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)G^{}(\omega ;𝐫_{}^{}{}_{1}{}^{},𝐫_1;𝐫_2,\mathrm{},𝐫_A)\}`$ (2.14)
$`=`$ $`{\displaystyle \frac{1}{v}}\delta (𝐛_1𝐛_{}^{}{}_{1}{}^{})\mathrm{exp}\{ip(z_1z_1^{})\}`$
$`\times [\theta (z_1z_1^{})\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}dz_1^{\prime \prime }{\displaystyle \underset{j=2}{\overset{A}{}}}V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_j)\}`$
$`+\theta (z_1^{}z_1)\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1}^{z_1^{}}}dz_1^{\prime \prime }{\displaystyle \underset{j=2}{\overset{A}{}}}V_{\mathrm{NN}}^{}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_j)\}].`$
Substituting this expression into the r.h.s. of eq. (2.2), we obtain the response function, $`S(\omega ,𝐪)`$, in the Glauber approximation.
Although the response function is given in a closed form, its evaluation contains a multi-dimensional integral over the $`A`$-body density matrix,
$`\rho _A(𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)=\mathrm{\Psi }_A^{}(𝐫_1,𝐫_2,\mathrm{},𝐫_A)\mathrm{\Psi }_A(𝐫_{}^{}{}_{1}{}^{},𝐫_2,\mathrm{},𝐫_A),`$ (2.15)
which is, of course, impossible for nuclear matter ($`A\mathrm{}`$). We previously carried out such a multi-dimensional integral for a light nucleus, <sup>16</sup>O, in a similar expression for the nuclear transparency , and examined various approximations, which could be used for heavier nuclei including nuclear matter. A possible approximation is to choose the $`A`$-body density matrix as
$`\rho _A(𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)=\rho (𝐫_1,𝐫_{}^{}{}_{1}{}^{}){\displaystyle \underset{j=2}{\overset{A}{}}}\rho (𝐫_j)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_j|)g(|𝐫_1𝐫_j|),`$ (2.16)
where $`\rho (𝐫,𝐫^{})`$ is the one-body density matrix with $`\rho (𝐫)`$ $`\rho (𝐫,𝐫^{})`$, and the one-body density is normalized to one, i.e., $`𝑑𝐫\rho (𝐫)`$ $`=1`$. $`g(|𝐫_i𝐫_j|)`$ describes both dynamical and statistical two-body correlation between the nucleon, $`1`$, and nucleon, $`j`$. This approximation consists of including only the two-body correlation and of neglecting the so-called spectator effect , , which amounts to dropping two-body correlations among the nucleons, $`2,\mathrm{},A`$, and has been found accurate enough when the finite interaction-range of $`V_{\mathrm{NN}}(𝐫)`$ is taken into account . With the above $`A`$-body density matrix, the multi-dimensional integral over $`𝐫_2,`$ $`\mathrm{},𝐫_A`$ factorizes, i.e.,
$`{\displaystyle 𝑑𝐫_2\mathrm{}𝑑𝐫_A\mathrm{exp}\{\frac{i}{v}_{z_1^{}}^{z_1}𝑑z_1^{\prime \prime }\underset{j=2}{\overset{A}{}}V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_j)\}\times \rho _A(𝐫_1,𝐫_{}^{}{}_{1}{}^{};𝐫_2,\mathrm{},𝐫_A)}`$ (2.17)
$``$ $`\rho (𝐫_1,𝐫_{}^{}{}_{1}{}^{})`$
$`\times \left({\displaystyle 𝑑𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)\mathrm{exp}\{\frac{i}{v}_{z_1^{}}^{z_1}𝑑z_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\}}\right)^{A1}`$
$`=`$ $`\rho (𝐫_1,𝐫_{}^{}{}_{1}{}^{})\times [{\displaystyle }d𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)`$
$`\times (11+\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}dz_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\})]^{A1}`$
$``$ $`\rho (𝐫_1,𝐫_{}^{}{}_{1}{}^{})\times \mathrm{exp}[(A1){\displaystyle }d𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)`$
$`\times (1\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}dz_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\})],`$
where we have used the relation between the $`A`$-body and the one-body density matrices requiring
$`{\displaystyle 𝑑𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}1,`$ (2.18)
and assumed $`A`$ to be large.
Thus, we obtain the expression for the response function in the following form,
$`S(\omega ,𝐪)`$ $`=`$ $`{\displaystyle \frac{1}{\pi v}}{\displaystyle 𝑑𝐛_1𝑑z_1𝑑z_1^{}\rho (𝐫_1,𝐫_{}^{}{}_{1}{}^{})\mathrm{exp}\{i(pq)(z_1z_1^{})\}}`$
$`\times \left[\theta (z_1z_1^{})\mathrm{exp}\{F(𝐛_1;z_1,z_1^{})\}+\theta (z_1^{}z_1)\mathrm{exp}\{F^{}(𝐛_1;z_1^{},z_1)\}\right],`$
where $`F(𝐛_1;z_1,z_1^{})`$ expresses the effect of the FSI (FSI-function),
$`F(𝐛_1;z_1,z_1^{})`$ $``$ $`(A1){\displaystyle 𝑑𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$ (2.20)
$`\times \left(1\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\}\right),`$
and $`𝐛_{}^{}{}_{1}{}^{}`$ $`=𝐛_1`$ due to the eikonal approximation, eq. (2.11). Note that we have not used the expansion in terms of the nucleon-nucleon potential, because the effect of the interaction with each nucleon cannot be assumed to be small. The expressions, eqs. (2.2) and (2.20), are our main results of this work.
One characteristic point of our formulation appears in this expression. Reflecting the fact that the Green’s function for the struck nucleon, eq (2.9), is off-diagonal with respect to $`𝐫_1`$ and $`𝐫_{}^{}{}_{1}{}^{}`$, the correlation function appears as a product of $`g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)`$$`g(|𝐫_1𝐫_2|)`$. This feature is missing in the formulation based on the optical potential . We will discuss the difference in the subsec. 3.2.
### 2.3 Expressions for Nuclear Matter
For later convenience, we here derive a simplified expression for the response function, eq. (2.2), for nuclear matter. It becomes
$`S(\omega ,𝐪)`$ $`=`$ $`{\displaystyle \frac{V}{\pi v}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑Z\stackrel{~}{W}(Z)\mathrm{exp}\{i(pq)Z\}`$
$`\times [\theta (Z)\mathrm{exp}\{F(Z)\}+\theta (Z)\mathrm{exp}\{F^{}(Z)\}],`$
where $`Zz_1z_1^{}`$, and $`V`$ is the volume of the system. Since we normalize the one-body density to one, $`\rho (𝐫)`$ becomes $`\rho _0`$ $`1/V`$. The FSI-function, eq. (2.20), is expressed as
$`F(Z)`$ $`=`$ $`\rho _{\mathrm{NM}}{\displaystyle 𝑑𝐫_2g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$ (2.22)
$`\times \left(1\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\}\right),`$
where $`\rho _{\mathrm{NM}}`$ $`=A\rho _0`$ is the nuclear matter density, and we use the value, $`\rho _{\mathrm{NM}}`$ $`=A/V`$ $`=0.17[\mathrm{fm}^3]`$. Since both the FSI-function and the one-body density matrix become functions of the difference of the arguments for the uniform nuclear matter, we simply write $`F(Z)`$ for the FSI-function and $`\stackrel{~}{W}(Z)`$ for the one-body density matrix in eq. (2.3),
$`\stackrel{~}{W}(Z)=\stackrel{~}{W}(𝐫𝐫^{})\rho (𝐫,𝐫^{}).`$ (2.23)
Using the fact that $`\stackrel{~}{W}(Z)`$ is an even function, i.e., $`\stackrel{~}{W}(Z)`$ $`=\stackrel{~}{W}(Z)`$, and introducing a dimensionless function, $`\stackrel{~}{w}(Z)`$, defined by
$`\stackrel{~}{w}(Z)V\stackrel{~}{W}(Z)=\rho (𝐫,𝐫^{})/\rho _0,`$ (2.24)
we further reduce the expression in the following way,
$`S(\omega ,𝐪)`$ $`=`$ $`{\displaystyle \frac{V}{\pi v}}{\displaystyle _0^{\mathrm{}}}𝑑Z\stackrel{~}{W}(Z)`$ (2.25)
$`\times [\mathrm{exp}\{i(pq)Z+F(Z)\}+\mathrm{exp}\{i(pq)Z+F^{}(Z)\}]`$
$`=`$ $`{\displaystyle \frac{2}{\pi v}}{\displaystyle _0^{\mathrm{}}}𝑑Z\stackrel{~}{w}(Z)\mathrm{exp}\{\mathrm{Re}F(Z)\}\mathrm{cos}\{(pq)Z+\mathrm{Im}F(Z)\},`$
where we have decomposed $`F(Z)`$ $`\mathrm{Re}F(Z)+i\mathrm{Im}F(Z)`$. This is the expression that we use for the numerical calculations for the inclusive cross sections in subsec. 4.2.
To compare our formulation with the formulation based on the optical potential in subsec. 3.2. it is convenient to define the “convolution” function, $`\zeta (Z)`$. For this purpose we introduce the PWIA response function, $`S_0(\omega ,𝐪)`$, which is defined by
$`S_0(\omega ,𝐪)`$ $``$ $`{\displaystyle \frac{1}{\pi v}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑Z\stackrel{~}{w}(Z)\mathrm{exp}\{i(pq)Z\}`$ (2.26)
$`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑Z\stackrel{~}{S}_0(Z)\mathrm{exp}\{i(pq)Z\},`$
where $`\stackrel{~}{S}_0(Z)`$ $`\stackrel{~}{w}(Z)/(\pi v)`$ is the Fourier transform of $`S_0(\omega ,𝐪)`$. Note that $`\stackrel{~}{S}_0(Z)`$ is actually a function of $`(pq)`$ with $`p`$ related to $`\omega `$ by eq. (2.12), due to our approximation of neglecting the residual nucleus excitation energies (the fixed-scatterer approximation). Using $`\stackrel{~}{S}_0(Z)`$, we can write the response function as
$`S(\omega ,𝐪)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑Z\stackrel{~}{S}_0(Z)\mathrm{exp}\{i(pq)Z\}\zeta (Z),`$ (2.27)
where $`\zeta (Z)`$ is our “convolution” function defined by
$`\zeta (Z)`$ $``$ $`\theta (Z)\mathrm{exp}\{F(Z)\}+\theta (Z)\mathrm{exp}\{F^{}(Z)\}.`$ (2.28)
$`\zeta (Z)`$ includes all the information on the FSI. If there is no FSI, then $`\zeta (Z)`$ $`=1`$, and $`S(\omega ,𝐪)`$ coincides with $`S_0(\omega ,𝐪)`$.
### 2.4 Zero-range Approximations
In this subsection we discuss the zero-range approximations for the nucleon-nucleon potential, $`V_{\mathrm{NN}}(𝐫)`$, for the FSI-function, eq. (2.20), which are often employed in the literatures. Within our knowledge, the approximations have never been examined in the Glauber approach. Here we will study them in two stages within our framework, and reduce the numerical work greatly.
The first one is to make the zero-range approximation only in the $`z`$-direction (ZR1). This approximation amounts to replacing $`V_{\mathrm{NN}}(𝐫)`$ by $`V_{}(𝐛)\delta (z)`$. The $`𝐛`$-dependence of the last factor on the r.h.s. of eq. (2.20) factorizes and we obtain
$`1\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)\}`$ $``$ $`\mathrm{\Gamma }(𝐛_1𝐛_2)\theta (z_1z_2)\theta (z_2z_1^{}),`$ (2.29)
where $`𝐫_2`$ $`=(𝐛_2,z_2)`$. The FSI-function, eq. (2.20), becomes
$`F_{\mathrm{ZR1}}(𝐛_1;z_1,z_1^{})`$ $`=`$ $`(A1){\displaystyle 𝑑𝐛_2_{z_1^{}}^{z_1}𝑑z_2\rho (𝐛_2,z_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$ (2.30)
$`\times \mathrm{\Gamma }(𝐛_1𝐛_2),`$
where $`\mathrm{\Gamma }(𝐛)`$ is the Fourier transform of the nucleon-nucleon scattering amplitude, and the definition is shown in eq. (A.5) in the Appendix. ZR1 should be a good approximation for the finite-range interaction if $`|z_1z_1^{}|`$ is large, because the step function in the $`z`$-direction due to the zero-range approximation will be hard to discriminate from the smoothed one over the range of the interaction for large $`|z_1z_1^{}|`$.
Though the full calculation of the FSI-function for the finite-range interaction is probably difficult, the behavior of the small $`|z_1z_1^{}|`$, where $`|_{z_1^{}}^{z_1}dz_1^{\prime \prime }`$ $`V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)|`$ is small, can be looked into by expanding the l.h.s. of eq. (2.29) in terms of the potential up to the first order. We write the expression (2.20) for $`F_{\mathrm{FR}}(𝐛_1;z_1,z_1^{})`$ <sup>1</sup><sup>1</sup>1FR stands for Finite Range. as
$`F_{\mathrm{FR}}(𝐛_1;z_1,z_1^{})`$ $`=`$ $`{\displaystyle \frac{i}{v}}(A1){\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }{\displaystyle 𝑑𝐫_2\rho (𝐫_2)V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)}`$ (2.31)
$`\times g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|).`$
This expression is useful to see how $`F_{\mathrm{FR}}(𝐛_1;z_1,z_1^{})`$ and that of ZR1 are different in the small $`|z_1z_1^{}|`$ region. The numerical results are shown in subsec. 4.1.
The second one is to make the zero-range approximation in all the direction (ZR2). One cannot simply replace $`V_{\mathrm{NN}}(𝐫)`$ by $`V_0\delta (𝐫)`$, because the appearance of the $`\delta `$-function in the exponential is meaningless, which does not happen for ZR1. We have to be careful for determining the expression of $`\mathrm{\Gamma }(𝐛)`$ in this approximation. From eq. (A.8), $`\mathrm{\Gamma }(𝐛)`$ is written as
$`\mathrm{\Gamma }(𝐛)`$ $`=`$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}{\displaystyle \frac{d^2𝐪}{(2\pi )^2}\mathrm{exp}\{i𝐪𝐛\}f(𝐪)}.`$ (2.32)
The zero-range approximation in the transverse directions implies that $`f(𝐪)`$ is independent of $`𝐪`$. Thus eq. (2.32) becomes
$`\mathrm{\Gamma }(𝐛)`$ $``$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}f(0){\displaystyle \frac{d^2𝐪}{(2\pi )^2}\mathrm{exp}\{i𝐪𝐛\}}`$ (2.33)
$`=`$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}f(0)\delta ^{(2)}(𝐛).`$
Substituting this formula into eq. (2.30), we obtain
$`F_{\mathrm{ZR2}}(𝐛_1;z_1,z_1^{})`$ $`=`$ $`(A1){\displaystyle \frac{2\pi }{i|𝐩|}}f(0){\displaystyle 𝑑𝐫_2\rho (𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$
$`\times \delta ^{(2)}(𝐛_1𝐛_2)\theta (z_1z_2)\theta (z_2z_1^{})`$
$`=`$ $`(A1){\displaystyle \frac{2\pi }{i|𝐩|}}f(0){\displaystyle _{z_1^{}}^{z_1}}𝑑z_2\rho (𝐛_1,z_2)g(|z_1^{}z_2|)g(|z_1z_2|).`$
The two approximations are numerically compared in the form of the FSI-function in subsec. 4.1 and in the inclusive cross section in subsec. 4.2.
## 3 Comparison with other Formulations
In the following two subsections, we compare our formulation based on the Glauber approximation with the theory of Gersch, Rodriguez, and Smith (GRS theory) , and a formulation with the optical potential . For the latter case, we show some numerical estimates in the next section.
### 3.1 GRS Theory
In this subsection, we point out the features of the GRS theory , , and compare them with those of our formulation by using $`R(\omega ,𝐪)`$ defined in eq. (2.5),
$`R(\omega ,𝐪)A|\widehat{O}^{}(𝐪){\displaystyle \frac{1}{\omega \widehat{H}_A+i\eta }}\widehat{O}(𝐪)|A,`$ (3.1)
where $`\widehat{O}(𝐪)`$ $`(=\mathrm{exp}\{i𝐪𝐫_1\})`$ is the hard scattering operator defined in eq.(2.6), and $`\widehat{H}_A|A`$ $`=0`$. The response function is obtained by taking its imaginary part as in eq. (2.4). We decompose the full hamiltonian, $`\widehat{H}_A`$ $`=\widehat{K}_1+\widehat{H}_{A1}+\widehat{V}_{1,A1}`$, in the same way as in eq. (2.7) of sec. 2.2. By expanding $`R(\omega ,𝐪)`$ in $`\widehat{V}_{1,A1}`$ up to the first order, we compare the Glauber theory with the GRS theory.
First, we consider the expansion of $`R(\omega ,𝐪)`$ in the GRS theory. Since the hard scattering operator, $`\widehat{O}(𝐪)`$, is the momentum-shift operator, it operates on the Green’s function in eq. (3.1) to shift the momentum. The kinetic-energy operator of the struck nucleon, $`\widehat{K}_1`$ $`=\widehat{𝐩}_1^2/(2m)`$, is the only term affected in the hamiltonian, $`\widehat{H}_A`$, and is shifted as
$`\widehat{O}^{}(𝐪)\widehat{K}_1\widehat{O}(𝐪)=\widehat{K}_1+{\displaystyle \frac{\widehat{𝐩}_1𝐪}{m}}+{\displaystyle \frac{𝐪^2}{2m}}.`$ (3.2)
Substituting this expression into eq. (3.1), we obtain the expansion in terms of $`\widehat{H}_A`$,
$`R(\omega ,𝐪)`$ $`=`$ $`A|{\displaystyle \frac{1}{\omega \widehat{H}_A\widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}|A`$ (3.3)
$`=`$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}|A`$
$`+A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}\widehat{H}_A{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}|A+\mathrm{}`$
$`=`$ $`R_0(\omega ,𝐪)+R_1(\omega ,𝐪)+\mathrm{}.`$
Using $`\widehat{H}_A|A`$ $`=0`$, we can rewrite $`R_1(\omega ,𝐪)`$ as follows:
$`R_1(\omega ,𝐪)`$ $``$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}\widehat{H}_A{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}|A`$
$`=`$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}[\widehat{H}_A,{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}]|A`$
$`=`$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}[\widehat{V}_{1,A1},{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)}}]|A.`$
Only $`\widehat{V}_{1,A1}`$ does not commute with $`\widehat{𝐩}_1𝐪/m`$.
$`R_0(\omega ,𝐪)`$ and $`R_1(\omega ,𝐪)`$ in eq. (3.3) gives $`F_0(𝐪)`$ and $`F_1(𝐪)`$ of Rinat and Taragin . They apply the GRS theory to the $`(e,e^{})`$ inclusive reaction for the first time. In their work, after they obtain the expression of $`F_0(𝐪)`$ and $`F_1(𝐪)`$, they exponentiate the potential term of $`F_1(𝐪)`$ to obtain their final expression.
Next, we consider the expansion of $`R(\omega ,𝐪)`$ in terms of $`\widehat{V}_{1,A1}`$ in the Glauber theory. According to the usual perturbation theory we write
$`{\displaystyle \frac{1}{\omega \widehat{H}_A}}`$ $`=`$ $`{\displaystyle \frac{1}{\omega \widehat{K}_1\widehat{H}_{A1}}}`$ (3.5)
$`+{\displaystyle \frac{1}{\omega \widehat{K}_1\widehat{H}_{A1}}}\widehat{V}_{1,A1}{\displaystyle \frac{1}{\omega \widehat{K}_1\widehat{H}_{A1}}}+\mathrm{}.`$
We apply the fixed-scatterer approximation and the eikonal approximation to the above expression. By the fixed-scatterer approximation we imply the replacement with $`\widehat{H}_{A1}`$ to $`\overline{E}_{A1}`$, where $`\overline{E}_{A1}`$ is the average value of the excitation energy of the residual nucleus. By the eikonal approximation we imply
$`A|\widehat{O}^{}(𝐪){\displaystyle \frac{1}{\omega \widehat{K}_1\overline{E}_{A1}}}\widehat{O}(𝐪)|A`$ $`=`$ $`A|{\displaystyle \frac{1}{\omega \widehat{K}_1\widehat{𝐩}_1𝐪/m𝐪^2/(2m)\overline{E}_{A1}}}|A`$ (3.6)
$``$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\overline{E}_{A1}}}|A.`$
To obtain the second equality, we have put $`\widehat{K}_1|A`$ $`0`$, because the initial momentum of the struck nucleon is small compared to $`𝐪`$. Corresponding to the expansion of eq. (3.5), $`R(\omega ,𝐪)`$ of eq. (3.1) is expanded as
$`R(\omega ,𝐪)`$ $`=`$ $`\overline{R}_0(\omega ,𝐪)+\overline{R}_1(\omega ,𝐪)+\mathrm{},`$ (3.7)
where $`\overline{R}_0(\omega ,𝐪)`$ and $`\overline{R}_1(\omega ,𝐪)`$ are written as
$`\overline{R}_0(\omega ,𝐪)`$ $``$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\overline{E}_{A1}}}|A,`$ (3.8)
$`\overline{R}_1(\omega ,𝐪)`$ $``$ $`A|{\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\overline{E}_{A1}}}\widehat{V}_{1,A1}`$ (3.9)
$`\times {\displaystyle \frac{1}{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\overline{E}_{A1}}}|A.`$
Now let us discuss the comparison of the GRS theory with the Glauber theory. For up to the first order expansion we can prove the following equation:
$`R_0(\omega ,𝐪)+R_1(\omega ,𝐪)`$ $`=`$ $`\overline{R}_0(\omega ,𝐪)+\overline{R}_1(\omega ,𝐪)+\mathrm{\Delta }_1\overline{R}_0(\omega ,𝐪)+\mathrm{\Delta }_2\overline{R}_0(\omega ,𝐪),`$ (3.10)
where $`\mathrm{\Delta }_1\overline{R}_0(\omega ,𝐪)`$ is the correction to the eikonal approximation defined by
$`\mathrm{\Delta }_1\overline{R}_0(\omega ,𝐪)`$ $`=`$ $`A|{\displaystyle \frac{1}{\{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\}^2}}\widehat{K}_1|A,`$ (3.11)
and $`\mathrm{\Delta }_2\overline{R}_0(\omega ,𝐪)`$ is the correction to the fixed-scatterer approximation defined by
$`\mathrm{\Delta }_2\overline{R}_0(\omega ,𝐪)`$ $`=`$ $`A|{\displaystyle \frac{1}{\{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\}^2}}(\widehat{H}_{A1}\overline{E}_{A1})|A.`$ (3.12)
The proof of eq. (3.10) is as follows: The difference of the zeroth order in $`\widehat{V}_{1,A1}`$ is written as
$`\overline{R}_0(\omega ,𝐪)R_0(\omega ,𝐪)`$ $`=`$ $`A|{\displaystyle \frac{1}{\{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\}^2}}\overline{E}_{A1}|A.`$ (3.13)
Adding the two correction terms, eqs. (3.11) and (3.12), to the zeroth order expression, eq. (3.13), we obtain
$`\overline{R}_0(\omega ,𝐪)R_0(\omega ,𝐪)+\mathrm{\Delta }_1\overline{R}_0(\omega ,𝐪)+\mathrm{\Delta }_2\overline{R}_0(\omega ,𝐪)`$ (3.14)
$`=`$ $`A|{\displaystyle \frac{1}{\{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\}^2}}(\widehat{K}_1+\widehat{H}_{A1})|A`$
$`=`$ $`A|{\displaystyle \frac{1}{\{\omega \widehat{𝐩}_1𝐪/m𝐪^2/(2m)\}^2}}\widehat{V}_{1,A1}|A`$
$`=`$ $`R_1(\omega ,𝐪)\overline{R}_1(\omega ,𝐪),`$
where we have used $`\widehat{H}_A|A`$ $`=0`$. This completes the proof of eq. (3.10).
The relation of the two formulations, eq. (3.10), implies that the expansion up to this order is exact in the GRS theory, while some corrections are needed to equate them in the Glauber theory. Thus one can realize that up to the first order in $`\widehat{V}_{1,A1}`$ the GRS theory includes less approximation than the Glauber theory, though the corrections should be small in the high energy region.
If we go beyond the first order, the situation will change. In the GRS theory it looks difficult to estimate higher order terms, while higher order terms can be summed up in the Glauber theory. This is the crucial difference, because the convergence of the series is not good from our discussion below eq. (2.20) in subsec. 2.2. Since the higher order terms are important, we should sum them up. Furthermore, the fixed-scatterer approximation and the eikonal approximation employed in the Glauber theory are established in high-energy regime. In this sense, we believe that the Glauber theory is superior to the GRS theory for descriptions of the high energy reactions.
### 3.2 A Formulation with Optical Potential
The optical-potential formulation has been used by various authors to treat the effects of the FSI in inclusive scattering (see e.g., Refs. , , , , ). Here we briefly review the treatment of the FSI and of the nuclear correlation in the formulation of Benhar, et al. . The main difference from our formulation will be seen to lie in the treatment of the nuclear correlation.
In the optical-potential formulation the effect of FSI is included in the convolution form. The target nuclear tensor, $`W_{\mu \nu }^A(q,\omega )`$, which corresponds to the response function, $`S(\omega ,𝐪)`$, in our approach, is expressed as
$`W_{\mu \nu }^A(q,\omega )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}F(\omega \omega ^{})W_{\mu \nu ,IA}^A(q,\omega ^{}V(q)),`$ (3.15)
where $`W_{\mu \nu ,\mathrm{IA}}^A(q,\omega )`$ is the target nuclear tensor in the impulse approximation which gives the PWIA cross section corresponding to $`S_0(\omega ,𝐪)`$ in eq. (2.26). $`F(\omega )`$ is the convolution function expressed as
$`F(\omega \omega ^{})`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑te^{i(\omega \omega ^{})t}e^{W(q,t)|t|}`$ (3.16)
$`=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Re}{\displaystyle _0^{\mathrm{}}}𝑑te^{i(\omega \omega ^{})t}e^{W(q,t)t}`$
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑t\mathrm{cos}\{(\omega \omega ^{})t\}e^{W(q,t)t}.`$
$`V(q)iW(q,t)`$ is an “optical potential”. The $`t`$-dependence of the imaginary part is introduced to take account of the initial-state correlation between the struck nucleon and the other nucleons. This is thus not the usual optical potential to be used for the elastic scatterings and its foundation is not clear. The $`\omega `$-dependence of $`V(q)`$ and $`W(q,t)`$ is assumed to be absent. If $`W(q,t)=0`$, then $`F(\omega \omega ^{})`$ $`=\delta (\omega \omega ^{})`$.
The above expression, eq. (3.16), can be directly compared with our convolution function, eq. (2.28), because both of them contain full information of FSI, and appear in the convolution form with the PWIA response function. The numerical comparison will be shown in subsec. 4.1.
Let us write $`z=vt`$ instead of $`t`$ as a variable of $`W(q,z)`$. Using the zero-range approximation for $`V_{\mathrm{NN}}(𝐫)`$, the authors of Ref. give
$`W(q,z)={\displaystyle \frac{\mathrm{}}{2}}\rho _{\mathrm{NM}}v(q)\sigma _{\mathrm{NN}}^{total}(q){\displaystyle \frac{1}{z}}{\displaystyle _0^z}𝑑z^{}g(z^{}).`$ (3.17)
$`g(r)`$ is the pair-distribution function, which corresponds to our correlation function, defined by
$`g(r)\overline{\rho }_{\mathrm{NN}}(\rho _{\mathrm{NM}},r)/\rho _{\mathrm{NM}}.`$ (3.18)
$`\overline{\rho }_{\mathrm{NN}}(\rho _{\mathrm{NM}},r)`$ is the average two-body density in nuclear matter at density $`\rho _{\mathrm{NM}}`$. The average two-body density is defined as
$`\overline{\rho }_{\mathrm{NN}}(𝐫_{12}){\displaystyle \frac{1}{A}}{\displaystyle 𝑑𝐑_{12}\rho _{\mathrm{NN}}(𝐫_1,𝐫_2)},`$ (3.19)
with $`𝐑_{12}(𝐫_1+𝐫_2)/2`$ and $`𝐫_{12}𝐫_1𝐫_2`$. For nuclear matter $`\overline{\rho }_{\mathrm{NN}}(𝐫_{12})`$ is a function of $`|𝐫_{12}|`$, If there is no correlation, i.e., $`g(r)`$ $`=1`$, then
$`W(q,z)={\displaystyle \frac{\mathrm{}}{2}}\rho _{\mathrm{NM}}v(q)\sigma _{\mathrm{NN}}^{total}(q).`$ (3.20)
We cannot directly compare the correlation function defined in eq. (3.18) with our correlation function defined in eq. (2.16). However judging from how they appear in the term of the FSI effect, $`W(q,z)`$ of eq. (3.18) and the FSI-function of eq. (2.20), it would be appropriate to identify their $`g(r)`$ with our $`g^2(r)`$.
## 4 Numerical Results and Discussion
In this section we will show our numerical results and the discussions.
### 4.1 “Convolution” Function
To calculate the convolution functions defined by eq. (2.28), we need a specific form of the correlation function, $`g(r)`$. We take the correlation function of the form
$`g(r)=1c_1e^{r^2/a_1^2},`$ (4.1)
where $`c_1`$ $`=0.84[\mathrm{fm}]`$, and $`a_1`$ $`=0.7[\mathrm{fm}]`$ (Fig. 3), which roughly simulates that of Benhar, et al. (Fig. 6 of Ref. ). In our figure we plot $`g^2(r)`$ instead of $`g(r)`$ in order to compare it easily with their pair-distribution function, because our $`g^2(r)`$ corresponds to their $`g(r)`$ as we discussed at the end of subsec. 3.2. As one can see from those figures, the way how our correlation function approaches unity is slightly different from theirs. The difference in the curvature may cause a nontrivial effect on the cross section. In this case, we need another Gaussian term in eq. (4.1). Since what we would like to see is the difference coming from different manipulations for the FSI and the correlation, we believe that our choice of the correlation function causes no big problem.
One more thing which we should comment on is the normalization of the wave function. For the case of a finite nucleus, we should be very careful for the normalization including the effect of the nuclear correlation . Fortunately the correction should be the order of $`1/A`$, which can be neglected for nuclear matter.
By using eq. (4.1), the approximate expressions for $`F(𝐛_1;z_1,z_1^{})`$, eqs. (2.30) and (2.4), can be calculated explicitly for nuclear matter. The expression for ZR1, eq. (2.30), where the zero-range approximation is applied only in the $`z`$-direction, becomes
$`F_{\mathrm{ZR1}}(Z)`$ $``$ $`F_{\mathrm{ZR1}}(𝐛_1;z_1,z_1^{})`$
$`=`$ $`\rho _{\mathrm{NM}}{\displaystyle 𝑑𝐛_2_{z_1^{}}^{z_1}𝑑z_2g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)\mathrm{\Gamma }(𝐛_1𝐛_2)}`$
$`=`$ $`\rho _{\mathrm{NM}}{\displaystyle \frac{f(\mathrm{𝟎})}{2i|𝐩|\gamma }}[4\gamma \pi Z{\displaystyle \frac{c_1\pi }{1/a_1^2+1/(4\gamma )}}a_1\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{a_1^2}})\}`$
$`+{\displaystyle \frac{c_1^2\pi }{2/a_1^2+1/(4\gamma )}}{\displaystyle \frac{a_1}{\sqrt{2}}}\mathrm{exp}({\displaystyle \frac{Z^2}{2a_1^2}})\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{2a_1^2}})\}],`$
where $`Z=z_1z_1^{}`$ and $`(A1)\rho _0`$ $`\rho _{\mathrm{NM}}`$. Here $`\mathrm{\Gamma }(z,p)`$ is the incomplete Gamma function defined by
$`\mathrm{\Gamma }(z,p)={\displaystyle _p^{\mathrm{}}}𝑑te^tt^{z1},(\mathrm{Re}z>0).`$ (4.3)
A useful formula related to the incomplete Gamma function is
$`{\displaystyle _a^{\mathrm{}}}𝑑xe^{x^2}={\displaystyle \frac{1}{2}}\mathrm{\Gamma }({\displaystyle \frac{1}{2}},a^2).`$ (4.4)
The expression for ZR2 where the zero-range approximation is applied in the whole direction, eq. (2.4), becomes
$`F_{\mathrm{ZR2}}(Z)`$ $``$ $`F_{\mathrm{ZR2}}(𝐛_1;z_1,z_1^{})`$ (4.5)
$`=`$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}f(\mathrm{𝟎})\rho _{\mathrm{NM}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_2g(|z_1^{}z_2|)g(|z_1z_2|).`$
$`=`$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}f(\mathrm{𝟎})\rho _{\mathrm{NM}}[Zc_1a_1\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{a_1^2}})\}`$
$`+c_1^2{\displaystyle \frac{a_1}{\sqrt{2}}}\mathrm{exp}({\displaystyle \frac{Z^2}{2a_1^2}})\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{2a_1^2}})\}],`$
Using the Gaussian series of the potential, eq. (A.12), in the Appendix, we write the expression for FR, eq. (2.31), as
$`F_{\mathrm{FR}}(Z)`$ $``$ $`F_{\mathrm{FR}}(𝐛_1;z_1,z_1^{})`$
$`=`$ $`{\displaystyle \frac{i}{v}}\rho _{\mathrm{NM}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }{\displaystyle 𝑑𝐫_2V_{\mathrm{NN}}(𝐫_{}^{\prime \prime }{}_{1}{}^{}𝐫_2)g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$
$`=`$ $`{\displaystyle \frac{i}{v}}\rho _{\mathrm{NM}}\pi ^{3/2}{\displaystyle \underset{n=1}{\overset{N_0}{}}}V_n[\{\left({\displaystyle \frac{n}{4\gamma }}\right)^{3/2}+c_1^2\left({\displaystyle \frac{2}{a_1^2}}\right)^{3/2}\mathrm{exp}({\displaystyle \frac{Z^2}{2a_1^2}})\}Z`$
$`c_1{\displaystyle \frac{\{(n/(4\gamma ))(1/a_1^2)\}^{1/2}}{n/(4\gamma )+1/a_1^2}}\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{(n/(4\gamma ))(1/a_1^2)}{n/(4\gamma )+1/a_1^2}}Z^2)\}].`$
Note that this expression is valid only for small-$`Z`$ $`(=z_1z_1^{})`$ region. From the discussion at the end of the Appendix, the expansion should be valid for $`|Z|`$ $`<0.4`$ \[fm\] for the parameterization which we use below.
To clarify the relation of ZR1 with FR we introduce one more expression, ZR1’. This has no physical implication, but is meaningful for obtaining some insight on the relation between ZR1 and FR. In order to obtain ZR1’, we apply the expansion in terms of the potential to $`\mathrm{\Gamma }(𝐛)`$ in eq. (4.1),
$`\mathrm{\Gamma }(𝐛)`$ $``$ $`1\mathrm{exp}\{{\displaystyle \frac{i}{v}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}V_{\mathrm{NN}}(𝐛,z^{})\}`$ (4.7)
$``$ $`{\displaystyle \frac{i}{v}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}V_{\mathrm{NN}}(𝐛,z^{}).`$
This approximation is similar to that of FR, but the integration range is different. Substituting this expression for eq. (4.1), we obtain,
$`F_{\mathrm{ZR1}^{}}(Z)`$ $``$ $`F_{\mathrm{ZR1}^{}}(𝐛_1;z_1,z_1^{})`$
$`=`$ $`{\displaystyle \frac{i}{v}}\rho _{\mathrm{NM}}\sqrt{\pi }{\displaystyle 𝑑𝐛_2_{z_1^{}}^{z_1}𝑑z_2g(|𝐫_{}^{}{}_{1}{}^{}𝐫_2|)g(|𝐫_1𝐫_2|)}`$
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}dz_1^{\prime \prime }V_{\mathrm{NN}}(𝐫_1^{\prime \prime }𝐫_2)`$
$`=`$ $`{\displaystyle \frac{i}{v}}\rho _{\mathrm{NM}}\sqrt{\pi }{\displaystyle \underset{n=1}{\overset{N_0}{}}}V_n\sqrt{{\displaystyle \frac{4\gamma }{n}}}[{\displaystyle \frac{4\gamma }{n}}\pi Z{\displaystyle \frac{c_1\pi }{1/a_1^2+n/(4\gamma )}}a_1\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{a_1^2}})\}`$
$`+{\displaystyle \frac{c_1^2\pi }{2/a_1^2+n/(4\gamma )}}{\displaystyle \frac{a_1}{\sqrt{2}}}\mathrm{exp}({\displaystyle \frac{Z^2}{2a_1^2}})\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{Z^2}{2a_1^2}})\}].`$
The first term of this expression is the leading term for large-$`Z`$. This term is the same as that of FR of eq. (4.1). Therefore, ZR1’ and FR should show the same large-$`Z`$ behavior.
The numerical results for the convolution functions are shown in Figs. 4 a)-c). In Figs. 4a) and b) we plot
$`\mathrm{Re}\zeta (Z)=\mathrm{exp}\{\mathrm{Re}F(Z)\}\mathrm{cos}\left(\mathrm{Im}F(Z)\right),`$ (4.9)
because this form appears in eq. (2.25) for $`p`$ $`=q`$, and it carries full information on FSI. In Fig. 4a) the results with the nuclear correlation are shown. The important finding here is that the slope of FR at $`Z`$ $`=0`$ is the same as that of ZR1. This implies that the ZR1 can be a good approximation of the finite-range interaction case even for small-$`Z`$. Since we consider ZR1 as a good approximation of FR for large-$`Z`$, we expect it to be so for the whole range of $`Z`$. ZR2 shows a similar behavior to ZR1 in large-$`Z`$ region, but in small-$`Z`$ region ZR2 behaves differently from ZR1. Thus ZR2 cannot be expected to be generally applicable. ZR1’ shows the same behavior as FR in the large-$`Z`$ region, and slightly different in the small-$`Z`$. This is what we expected for the small-$`Z`$ expansion and thus justifies our expectation that ZR1 is a good approximation for the whole region.
In Fig. 4b) the results without the nuclear correlation are shown. ZR1 and ZR2 degenerate as one sees from eqs. (4.1) and (4.5), and they show different behavior from FR even in small-$`Z`$. From those results, we realize that the nuclear correlation makes the zero-range approximations better. The nuclear correlation plays an important role in this context.
In order to compare our formulation with that of the optical potential , we make the Fourier transform of our convolution function, $`\zeta (Z)`$, because their results were shown as a function of $`\omega `$,
$`\stackrel{~}{\zeta }(\omega )`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑te^{i\omega t}\zeta (t)`$ (4.10)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}𝑑t\left(\mathrm{exp}\{+i\omega t+F(t)\}+\mathrm{exp}\{i\omega t+F^{}(t)\}\right)`$
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑t\mathrm{exp}\{\mathrm{Re}F(t)\}\mathrm{cos}\left(\omega t+\mathrm{Im}F(t)\right).`$
Here we use $`vt`$ $`=Z`$.
The numerical results are shown in Fig. 4c). In the figure we plot the Fourier transform of ZR1, ZR2, and ZR2 without the real part of $`V_{\mathrm{NN}}(r)`$. This figure should be compared with Fig. 4 of Ref. . The characteristic feature of our results of ZR1 and ZR2 is the positive slope at $`\omega `$ $`=0`$. This behavior comes from the appearance of $`\mathrm{Im}F(t)`$ in the argument of cosine of eq. (4.10), and is the effect of the real part of $`V_{\mathrm{NN}}(r)`$. To confirm this statement, we plot ZR2 without the real part of $`V_{\mathrm{NN}}(r)`$. The curve turns out to be monotonically decreasing.
For reference, we plot $`\stackrel{~}{\zeta }_{\mathrm{OP}}(\omega )`$ defined by
$`\stackrel{~}{\zeta }_{\mathrm{OP}}(\omega )`$ $``$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑t\mathrm{cos}(\omega t)e^{W(q,t)t}.`$ (4.11)
which corresponds to the case of $`\mathrm{Im}F(t)`$ $`=0`$ in eq. (4.10), and is the same form as eq. (3.16), i.e., the correlated Glauber formulation of Ref. . We use eq. (3.17) for $`W(q,t)`$ with our parameter set, and $`g(r)`$ in eq. (3.17) is replaced our correlation function, $`\{g(r)\}^2`$ of eq. (4.1). Therefore, the following expression of $`W(q,t)`$ is not completely the same as that of Ref. :
$`W(q,t)`$ $``$ $`{\displaystyle \frac{1}{2}}\rho _{\mathrm{NM}}v\sigma _{\mathrm{NN}}^{total}{\displaystyle \frac{1}{t}}{\displaystyle _0^t}𝑑t^{}\{g(vt)\}^2`$ (4.12)
$`=`$ $`{\displaystyle \frac{1}{2}}\rho _{\mathrm{NM}}v\sigma _{\mathrm{NN}}^{total}[1c_1{\displaystyle \frac{a_1}{vt}}\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{v^2t^2}{a_1^2}})\}`$
$`+c_1^2{\displaystyle \frac{a_1}{2\sqrt{2}vt}}\{\sqrt{\pi }\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{2v^2t^2}{a_1^2}})\}].`$
The $`q`$-dependence of $`W(q,t)`$ is implicitly included in $`v`$ and $`\sigma _{\mathrm{NN}}`$. As one can see from Fig. 4c), $`\stackrel{~}{\zeta }_{\mathrm{OP}}(\omega )`$ shows a quite similar behavior to the result of $`\stackrel{~}{\zeta }(\omega )`$ of ZR2 without Re $`V_{\mathrm{NN}}`$ except for the magnitude.
### 4.2 Inclusive Cross Section
We numerically calculate the cross sections of the $`(e,e^{})`$ inclusive scattering off nuclear matter. Since we are interested in the treatment of the FSI and the nuclear correlation, other information, such as the one-body density matrix, is taken as an input.
We determine the shape of the one-body density matrix, $`\rho (𝐫,𝐫^{})`$, from numerical results of the momentum distribution of a nucleon, $`W(𝐤)`$. The one-body density matrix is related to the momentum distribution as
$`\rho (𝐫,𝐫^{})\rho _0{\displaystyle \frac{d𝐤}{(2\pi )^3}e^{i𝐤(𝐫𝐫^{})}W(𝐤)},`$ (4.13)
where $`\rho _0`$ $`=1/V`$.
For the momentum distribution of the nuclear matter in eq. (4.13), we use the numerical results of Refs. . They use a realistic nuclear force to obtain the results. We just use their numbers by $`\chi `$-square fitting in terms of the following series:
$`W(𝐤)`$ $``$ $`(2\pi )^3n_0(k)`$
$`=`$ $`(2\pi )^3\left(\theta (k_Fk)\{\alpha +{\displaystyle \underset{j=1}{\overset{N_1}{}}}n_{g,j}e^{jk^2/k_F^2}\}+\theta (kk_F){\displaystyle \underset{j=1}{\overset{N_2}{}}}n_{e,j}e^{jk/k_0}\right),`$
where $`\alpha `$ $`=7.156\times 10^2`$, $`k_F`$ $`=1.33[\mathrm{fm}^1]`$, and $`k_0`$ $`=0.588[\mathrm{fm}^1]`$. $`N_1`$ $`=15`$, and $`N_2`$ $`=10`$. Since the bare data are given by $`n_0(k)`$ in Refs. , we show the relation of $`n_0(k)`$ with our $`W(k)`$ in eq. (4.2). We impose the condition that the Gaussian series and the exponential series are to be connected continuously at $`k`$ $`=k_F`$ when we calculate the coefficients. The normalization is
$`{\displaystyle \frac{d𝐤}{(2\pi )^3}W(𝐤)}=1.`$ (4.15)
The resulting momentum distribution is shown in Fig. 5. For comparison, we plot the momentum distribution of the Fermi gas, $`W(𝐤)`$ $`=W_0\theta (k_Fk)`$. The tail shows a reflection of the nuclear correlation.
The numerical result for the one-body density matrix obtained from eq. (4.13) is shown in Fig. 6. For comparison, we plot that of the Fermi-gas model, which takes the form ,
$`\rho (𝐫,𝐫^{})=\rho _0{\displaystyle \frac{3j_1(k_FR)}{k_FR}},`$ (4.16)
where $`R`$ $`|𝐫𝐫^{}|`$, and $`j_1(x)`$ $`=(\mathrm{sin}xx\mathrm{cos}x)/x^2`$ is the spherical Bessel function.
Putting the above ingredients into eq. (2.25), we obtain the electron inclusive cross sections. The numerical results are shown in Figs. 7 a)-c). For the electron-nucleon cross section we use the so-called Rosenbluth cross section , and average over the proton and the neutron, i.e.,
$`{\displaystyle \frac{d\sigma _{\mathrm{eN}}}{d\mathrm{\Omega }}}_{\mathrm{el},\mathrm{on}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{d\sigma _{\mathrm{ep}}}{d\mathrm{\Omega }}}+{\displaystyle \frac{d\sigma _{\mathrm{en}}}{d\mathrm{\Omega }}}\right).`$ (4.17)
The comparison of the calculated cross sections, ZR1 and ZR2, is shown in Fig. 7a). For comparison, we plot the results of PWIA and the case of no correlation. The PWIA result directly reflects the momentum distribution of a nucleon in the nuclear matter. Inclusion of only the FSI, i.e., $`g(r)`$ $`=1`$, broadens the cross section, but adding the nuclear correlation (ZR1 and ZR2) reduces the broadening. That is, the nuclear correlation reduces the effect of the FSI. Those are the general features of our results, and were also observed in Ref. . By proceeding the zero-range approximation (ZR1 $``$ ZR2), the cross section comes lower in the low-$`\omega `$ region.
In Fig. 7b) we plot the cross sections for the cases of ZR1, ZR2, and no correlation with or without the real part of $`V_{\mathrm{NN}}(r)`$ to see the effect of the real part. As one can easily see from the figure, the effect is small for all cases. From our discussion at the end of the previous subsection, the real part affects the behavior of the FSI-function (see Fig. 4c)), but it does not affect the cross sections. We have thus confirmed the validity of neglecting the real part, which has often been presumed in the literature.
In Fig. 7c) we plot our numerical result of the cross section of ZR1 (times 0.4) with the experimental data. From the discussion in subsec. 4.1, the ZR1 includes the major part of the finite-range effect, but the cross section of ZR1 still overestimates the experimental cross section for low-$`\omega `$ region, although this is a consistent treatment for the inclusive process. This would be due to our use of the on-shell $`eN`$ cross section.
As we mentioned in the beginning of subsec. 2.1, we should use an off-shell $`eN`$ cross section for the elementary process since the struck nucleon can become far off-the-mass-shell in the final state. Actually putting a binding effect by shifting the nucleon mass to the smaller one shifts the cross section to large-$`\omega `$. Even this simple manipulation can make the agreement of the numerical results with the experimental data better, but this is only a prescription. Since one cannot determine the off-shell cross section without a dynamical model describing the internal structure of the nucleon and its interaction with the nuclear medium, we avoid to get involved in the off-shell problems in this work. The problems are interesting and important, but we leave the study for our future work.
In high-$`\omega `$ region, the experimental data highly exceeds our results, because the inelastic channel of the $`eN`$ cross section takes part in the data. To describe the cross section in such a high-$`\omega `$ region is out of our scope of this work.
## 5 Summary and Conclusion
We have given here a method to deal with the FSI in the high-energy $`(e,e^{})`$ inclusive scattering. Since we are interested in the quasi-elastic region, the relevant degrees of freedom are the nucleonic ones. We have applied the Green’s function method, and used the Glauber approximation for the FSI. An advantage of this formulation is that the finite-range effect of the nucleon-nucleon interaction and the nuclear correlation are included in a systematical way.
Though the method allows us to express the FSI effect in a closed form, the full calculation still requires a large amount of numerical works. We have thus examined two kinds of zero-range approximations for the nucleon-nucleon interaction, which greatly simplify the integrals involved in the closed form. The one called ZR1, which uses the zero-range approximation only in the longitudinal direction, has been found to be accurate enough for the actual calculations.
The method has then been compared with the treatments of the FSI in the other approaches to the response functions such as that of Gersch, Rodrigues and Smith, which has been applied to the $`(e,e^{})`$ process in Ref. , and the optical-potential approach used in Ref. .
Our principal objective is to propose a unified framework of calculating the inclusive $`(e,e^{})`$ and the semi-inclusive $`(e,e^{}p)`$ responses with the particular emphasis on the treatment of the FSI. The results of calculation in the present formalism (ZR1) for the $`(e,e^{})`$ cross section off nuclear matter show strong effects of the FSI especially in the low energy transfer ($`\omega `$) region qualitatively similar to those observed in Ref. . The calculated cross section overestimates the experimental one by a factor of about 2 in the peak region and by larger factors in the low-$`\omega `$ region. This would be partly due to our simple choice of on-shell kinematics for the $`eN`$ elastic cross section. We have not discussed the difficult problem of constructing the off-shell cross section, which would be required if we were to make a serious comparison with the observed $`(e,e^{})`$ response. Instead, we have simply pointed out that, in the inclusive processes, the final nucleon is much more off-shell than the initial one. Further studies on the in-medium cross section including the off-shell kinematics are necessary to draw definitive conclusions on the FSI in connection with the color transparency.
We should critically comment here the approximations which we used in this work. We have relied on several approximations, whose validation needs further and extensive work. One approximation is that the total disregard of $`E_n`$ of the residual nucleus, which we have used it to obtain the response function, eq. (2.2). We believe it is good though we have no quantitative basis for it. The other approximation is to our choice of the form, eq. (2.16), for the $`A`$-body density matrix, which is claimed to represent a reasonable approximation for the calculation of the nuclear transparency. However it might be inadequate for the calculation of the low-energy-transfer region of the inclusive cross section. These approximations need to be investigated carefully for their quantitative validity.
Acknowledgements
A.K. would like to express his gratitude to Profs. S.C.Pieper and V.R.Pandharipande for kindly giving their numerical tables of the momentum distributions to him. He is supported by the Special Postdoctoral Researcher Program at RIKEN.
This work is supported by the Grant-in-Aid for Scientific Research of Monbusho(C-08640355), and by the U.S. DOE at CSUN (DE-FG03-87ER40347) and by the U.S. NSF at Caltech (PHY-9722428 and PHY-9420470).
APPENDIX
## Appendix A Nucleon-Nucleon Potential
In this appendix, we explain how we construct the NN potential phenomenologically. This is an extension of the method introduced in Ref. .
We deal with a two-body scattering in free space. $`𝐩`$ is a momentum of the incident proton, $`𝐩^{}`$ is that of the outgoing proton, and the momentum transfer $`𝐪`$ $`=𝐩𝐩^{}`$. Please do not confuse them with the notation of the main part of this paper. In the following we focus on a kinematic region where $`𝐩𝐪`$ $`0`$.
The scattering amplitude is defined by
$`f(𝐪)`$ $``$ $`{\displaystyle \frac{i|𝐩|}{2\pi }}{\displaystyle d^2b\mathrm{\Gamma }(𝐛)\mathrm{exp}\{i(𝐩𝐩^{})𝐛\}}`$ (A.1)
$``$ $`f(0)\mathrm{exp}\{\gamma 𝐪^2\},`$ (A.2)
where $`𝐛`$ is the impact parameter defined by
$`𝐫=𝐛+z{\displaystyle \frac{𝐩}{|𝐩|}},`$ (A.3)
and
$`f(0)={\displaystyle \frac{i+c_0}{4\pi }}\sigma _{\mathrm{NN}}^{total}|𝐩|.`$ (A.4)
The last expression, eq. (A.2), is a phenomenological fit for the forward scattering amplitude. $`\mathrm{\Gamma }(𝐛)`$ in eq. (A.1) is the scattering amplitude in the coordinate space, and it has the following eikonal expression
$`\mathrm{\Gamma }(𝐛)`$ $``$ $`1\mathrm{exp}\{i\chi (𝐛)\}.`$ (A.5)
$`\chi (𝐛)`$ is the profile function defined by
$`\chi (𝐛)`$ $``$ $`{\displaystyle \frac{1}{v}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}V_{\mathrm{NN}}(𝐛,z^{})`$ (A.6)
$`=`$ $`{\displaystyle \frac{1}{i}}\mathrm{log}(1\mathrm{\Gamma }(𝐛)).`$ (A.7)
where $`V_{\mathrm{NN}}(𝐫)`$ is the nucleon-nucleon potential which we would like to obtain here.
With eq. (A.1), we analytically obtain the expression of $`\mathrm{\Gamma }(𝐛)`$ for the phenomenological fit of $`f(𝐪)`$, eq. (A.2), by the Fourier transformation,
$`\mathrm{\Gamma }(𝐛)`$ $`=`$ $`{\displaystyle \frac{2\pi }{i|𝐩|}}{\displaystyle \frac{d^2𝐪}{(2\pi )^2}e^{i𝐪𝐛}f(𝐪)}`$ (A.8)
$``$ $`\mathrm{\Gamma }(0)e^{𝐛^2/4\gamma },`$
where
$`\mathrm{\Gamma }(0)={\displaystyle \frac{1}{2i|𝐩|\gamma }}f(0).`$ (A.9)
Following the original paper of Glauber , we obtain the potential by using the Abel integral equation through eq. (A.6),
$`V_{\mathrm{NN}}(r)`$ $`=`$ $`{\displaystyle \frac{v}{\pi }}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}{\displaystyle _r^{\mathrm{}}}b𝑑b{\displaystyle \frac{\chi (b)}{\sqrt{b^2r^2}}}`$ (A.10)
$`=`$ $`{\displaystyle \frac{v}{\pi }}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}{\displaystyle _r^{\mathrm{}}}b𝑑b{\displaystyle \frac{\mathrm{log}(1\mathrm{\Gamma }(𝐛))}{i\sqrt{b^2r^2}}}`$
$`=`$ $`{\displaystyle \frac{v}{\pi }}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}{\displaystyle _0^{\mathrm{}}}𝑑y{\displaystyle \frac{1}{i}}\mathrm{log}(1\mathrm{\Gamma }(y^2+r^2))`$
$`=`$ $`{\displaystyle \frac{v}{2\gamma \pi i}}{\displaystyle _0^{\mathrm{}}}𝑑y{\displaystyle \frac{\mathrm{\Gamma }(0)e^{(y^2+r^2)/4\gamma }}{1\mathrm{\Gamma }(0)e^{(y^2+r^2)/4\gamma }}},`$
where $`y=\sqrt{b^2r^2}`$. We numerically calculate eq. (A.10), and obtain the potential. We use the following numbers at $`|𝐩|`$ $`=2.0`$ \[GeV/c\] :
$`\sigma _{\mathrm{NN}}^{total}=43.8[\mathrm{mb}],c_0=0.14,\gamma =3.37\times 10^6[(\mathrm{MeV})^2],v=0.905.`$ (A.11)
The numerical results are shown in Fig. 8.
For calculations of the response function, it is convenient to express the nucleon-nucleon potential in terms of the Gaussian series,
$`V_{\mathrm{NN}}^{app}(r)={\displaystyle \underset{j=1}{\overset{N_0}{}}}V_je^{jr^2/4\gamma }.`$ (A.12)
We determine the coefficients by $`\chi ^2`$ fit with the number of terms, $`N_0`$ $`50`$.
With the above parameterization, eq. (A.11), the absolute value of the following integral,
$`{\displaystyle \frac{i}{v}}{\displaystyle _{z_1^{}}^{z_1}}𝑑z_1^{\prime \prime }V_{\mathrm{NN}}(𝐛_1,z_1^{\prime \prime })`$
$`=`$ $`{\displaystyle \frac{i}{v}}{\displaystyle \underset{n=1}{\overset{N_0}{}}}V_n\sqrt{{\displaystyle \frac{4\gamma }{n}}}\mathrm{exp}\{{\displaystyle \frac{n}{4\gamma }}𝐛_1^2\}{\displaystyle \frac{1}{2}}\{\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{n}{4\gamma }}z_1^2)\mathrm{\Gamma }({\displaystyle \frac{1}{2}},{\displaystyle \frac{n}{4\gamma }}z_1^2)\},`$
stays less than unity for $`|z_1z_1^{}|`$ $`<0.4[\mathrm{fm}]`$ with $`𝐛_1`$ $`=\mathrm{𝟎}`$ and $`z_1`$ $`=0`$.
Figure Captions
Figure 1
Kinematic notations in the $`(e,e^{})`$ inclusive reaction.
Figure 2
Comparison of the off-shellness of the nucleon in the medium. $`\delta m^2/m_\mathrm{N}^2`$ is plotted as a function of $`E/m_\mathrm{N}`$. The solid line is the case (I), where $`V`$ $`=300[\mathrm{MeV}]`$, $`S`$ $`=350[\mathrm{MeV}]`$. The dashed line is the case (II), where $`V`$ $`=0[\mathrm{MeV}]`$, $`S`$ $`=50[\mathrm{MeV}]`$. The dash-dotted line is the case (III), where $`V`$ $`=50[\mathrm{MeV}]`$, $`S`$ $`=0[\mathrm{MeV}]`$. The sum, $`S+V`$, is kept to be the same.
Figure 3
Nuclear correlation function, $`g^2(r)`$, of eq. (4.1) as a function of $`r`$.
Figure 4 a)
Comparison of the two zero-range approximations for the nuclear matter with the nuclear correlation. Real part of our “convolution” function, $`\zeta (Z)`$, of eq. (2.28) is plotted as a function of $`Z`$ $`(=z_1z_1^{})`$. The solid curve is ZR1, eq. (4.1), and the dashed curve is ZR2, eq. (4.5). For reference, we plot ZR1’, eq. (4.1), by the dot-dashed curve, and FR, eq. (4.1), for a limited range of $`Z`$, by the dotted curve.
Figure 4 b)
Comparison of the two zero-range approximations for the nuclear matter without the nuclear correlation. The two are indistinguishable in the figure. The meaning of the curves is the same as Fig.4 a).
Figure 4 c)
Comparison of the two zero-range approximations for the nuclear matter with the nuclear correlation. Fourier transform of our “convolution” function, $`\stackrel{~}{\zeta }(Z)`$, of eq. (4.10) is plotted as a function of $`\omega `$. The solid curve is ZR1, eq. (4.1), and the dashed curve is ZR2, eq. (4.5). For reference, we plot ZR2 without the real part of $`V_{\mathrm{NN}}(r)`$ by the dot-dashed curve, and the one based on the optical potential, $`\stackrel{~}{\zeta }_{\mathrm{OP}}(\omega )`$ of eq. (4.11), by the dotted curve.
Figure 5
Momentum distributions of a nucleon in the nuclear matter as a function of $`k[\mathrm{fm}^1]`$. $`W(k)`$ $`(2\pi )^3n_0(k)`$. The solid curve is the result of the fitting, eq. (4.2), and the dotted line is the case of the Fermi gas. The crosses are the data points from Refs. , .
Figure 6
One-body density matrix of the nuclear matter as a function of $`|𝐫𝐫^{}|`$. The solid curve is the case including the nuclear correlation. The dashed curve is that of the Fermi gas.
Figure 7 a)
The cross sections of the inclusive scattering off nuclear matter. The dotted curve is the case of PWIA. The dash-dotted curve is the case of only the FSI (no correlation effect). The solid curve and the dashed curve are the full calculation including both the FSI and the nuclear correlation with zero-range approximations. The solid curve is ZR1, the zero-range approximation in $`z`$-direction only. The dashed curve is ZR2, the zero-range approximation in the whole direction.
Figure 7 b)
Comparison of the cross sections for the cases of ZR1, ZR2, and no correlation with (the solid curve) or without the real part of $`V_{\mathrm{NN}}(r)`$ (the dashed curve).
Figure 7 c)
Comparison of the numerical result of the cross section and the experimental data . The solid curve is the case of ZR1 $`\times 0.4`$.
Figure 8
Nucleon-nucleon potential, eq. (A.10), as a function of $`r`$ for the case of $`|𝐩|`$ $`=2.0`$ \[GeV/c\], and $`\gamma `$ $`=3.37[(\mathrm{GeV}/\mathrm{c})^2]`$. The dashed curve is the real part, and the solid curve is the imaginary part. |
warning/0001/cond-mat0001164.html | ar5iv | text | # Adhesion-induced lateral phase separation in membranes
## I Introduction
Adhesion of membranes and vesicles has attracted considerable experimental and theoretical interest because of its prime importance to many bio-cellular processes . Theoretical treatments of membranes composed of single component lipid bilayers have revealed that generic interactions such as van der Waals, electrostatic or hydration interactions govern the adhesive properties of interacting membranes. It is also worthwhile to mention that related phenomena are found in unbinding transition of nearly flat membranes or adhesion of vesicles to surfaces .
In addition to general non-specific interactions mentioned above, it is known from the works of Bell and coworkers as well as others , that highly specific molecular interactions play an essential role in biological adhesion. This interaction acts between complementary pairs of proteins such as ligand and receptor, or antibody and antigen. Well studied example of such coupled systems is the biotin-avidin complex. The avidin molecule has four biotin binding sites, two on each side, and forms a five-molecules biotin-avidin-biotin complex. The resulting specific interaction is highly local and short-ranged. Measurements by surface force apparatus or atomic force microscopy have shown that the force required to break a biotin-avidin bond is about 170pN. In related experiments measuring chemical equilibrium constants , it was found that the biotin-avidin binding energy is about 30–35 $`k_BT`$ which is larger than thermal fluctuations. Other coupled systems are those of selectins and their sugar ligands where the bond is much weaker, of the order of 5$`k_BT`$ .
More recently several models taking into account thermal fluctuations in membrane adhesion have been proposed. Zuckerman and Bruinsma used a statistical mechanics model which is mapped onto a two-dimensional Coulomb plasma with attractive interactions. They predicted an enhancement of the membrane adhesion due to thermal fluctuations. In another work, Lipowsky considered the adhesion of lipid membranes which includes anchored stickers, i.e., anchored molecules with adhesive segments . It was shown that flexible membranes can adhere if the sticker concentration exceeds a certain threshold. If the multi-component membranes, including lipids and sticker molecules, undergo a phase separation, the adhesion is dominated by the sticker-rich domains. Further studies in this direction using mean-field theory and Monte Carlo simulations obtained a phase separation which is driven both by attractive intra-membrane sticker interactions and fluctuation-induced interactions between stickers.
The problem of multi-component membrane adhesion is intimately related to that of formation of domains (a lateral phase separation). This has been observed by several experiments. For example, the biotin-avidin interaction occurring during vesicle-vesicle adhesion was investigated by a micropipette technique . The adhesion between one avidin-coated vesicle and a second biotinylated vesicle is followed by an accumulation of biotin-avidin complex in the contact zone. This accumulation of cross-bridges between the two vesicles is found to be a diffusion-controlled process.
Adhesion-induced phase separation has been observed by Albersdörfer et al. and results from the interplay between long-range repulsive and short-range attractive interactions . The membrane includes repeller molecules in the form of lipopolymers (modified DOPE lipid with a polyethyleneoxide headgroup), mimicking glycocalix in real biological systems. The other component is a receptor molecule in the form of biotinylated lipids (DOPE-X-biotin). This lipopolymer is responsible for longer-range repulsive interaction, while the short-range attractive interaction is introduced by adding streptavidin to the extra-cellular solution. The streptavidin acts as a connector between the biotinylated lipids on the two membranes. A technique of reflection interference contrast microscopy was used to observe domain formation on a vesicle adhering to a membrane supported on a solid substrate. The lateral phase separation on both membranes leads to the formation of domains of tight adhesion separated by domains of loose adhesion .
In a related work, adhesion between cationic vesicles and anionic supported membranes revealed that electrostatic interactions induce lateral charge segregation on the membrane . This phase separation leads to patches of tight inter-membrane contact and decoupled “blisters”. Furthermore, adhesion of membranes including self-recognizing homophilic molecules and lipopolymers has been investigated . It was found that the initial weak adhesion is followed by slower aggregation into tightly bound domains coexisting with domains of weak adhesion. The result has been interpreted in terms of a double-well inter-membrane interaction potential due to the presence of the lipopolymers. Let us emphasize that in all the above mentioned experiments, it was reported that adhesion molecules aggregate spontaneously and form domains of tight adhesion.
It is generally believed that multi-component biomembranes in physiological conditions are close to their critical point, and membrane functions are partially governed through phase separation processes. Moreover, concentration fluctuations in the vicinity of the critical point may affect biophysical properties of membranes and can be of importance in regulating membrane processes in a robust way. Recently this conjecture was supported by an experiment of an insoluble Langmuir monolayer at the air/water interface . The monolayer was prepared in two different steps. The first mimics the composition of the inner leaflet of a cell biomembrane, while the second mimics the outer leaflet. In both cases, by using fluorescence microscopy technique, it was found that the Langmuir monolayer is close to its corresponding critical point of demixing.
So far, the interplay between lateral phase separation and membrane adhesion has not been considered theoretically in detail except in Refs. . The work in Refs. deal only with the specific case of oppositely charged membranes. In this paper we provide a general phenomenological approach for the adhesion of multicomponent membranes. Using a mean-field theory, we investigate how the lateral phase separation within the membrane is affected by the adhesion of membranes. Like in Refs. , we consider adhesion mediated by sticker molecules. Sticker molecules are polymers or macromolecules anchored to one membrane and interacting with the other membrane by another sticky part of the molecule. They can form bridges between two adjacent membranes (so called trans-interaction) , and play an essential role in the adhesion of cell membranes in biological systems.
We distinguish three types of adhesion depending on the structure of bridges as represented in Fig. 1. (i) “Bolaform-sticker” adhesion where each bridge molecule consists of a single sticker having two sticky ends (Fig. 1(a)). One sticker end is anchored to one membrane while the other end is adhering directly to the second membrane. (ii) “Homophilic-sticker” adhesion where the bridges are formed by two stickers of the same type (Fig. 1(b)). Each sticker is anchored on one of the membranes, while their free ends bind together to form the bridge. (iii) “Lock-and-key” adhesion where the bridges consist of two different stickers forming a ligand-receptor type bond (Fig. 1(c)). This case represents an asymmetric adhesion due to the lack of symmetry between the ligand and receptor. In the present work, we mainly discuss the symmetric bolaform-sticker adhesion (case (i) above) using a model where the equilibrium spacing between two membranes is coupled to the local concentration of stickers. Even in the latter symmetric case, a certain asymmetry can be obtained by controlling separately the sticker chemical potentials on the two membranes. An important consequence of our model is that the lateral phase separation is enhanced.
This paper is organized as follows. In the next section, we explain our phenomenological model of bolaform-sticker adhesion. The mean-field phase diagrams are given in Sec. III. The inter-membrane distance between two coexisting domains is calculated in Sec. IV. Finally discussion is provided in Sec. V where the other types of adhesion mentioned above are considered.
## II Bolaform-sticker adhesion
In this section, we treat the case where the adhesion is mediated by a single type of sticker molecules which are anchored irreversibly to one membrane and stick to the other membrane by another sticky part of the molecules as in Fig. 1(a). The anchor segments consist of a hydrophobic segment and penetrate into the hydrophobic interior of the lipid bilayer. The sticky segments, on the other hand, adhere directly to another membrane having some potential of sticking . As mentioned above, we call this a “bolaform-sticker”. Consider two interacting membranes labeled by $`i=1,2`$ consisting of lipid molecules and bolaform-stickers as schematically shown in Fig. 2. Let the sticker concentration in each membrane be denoted by $`\psi _i(𝒓)`$, where $`𝒓=(x,y)`$ is a two-dimensional planar vector and $`0\psi _i(𝒓)1`$. Note that the average concentrations of stickers on the two membranes, $`\psi _1`$ and $`\psi _2`$, do not have to be the same.
When the adhesion molecules are very flexible, they can bend back to form arches on a single membrane. In order to avoid such a situation in experiments and in the model, the bending rigidity of the sticker molecules should be sufficiently large. Hence we assume that for stiff enough stickers all the bonds are inter-membrane ones connecting the two separate membranes as considered in Refs. .
The interaction between two stickers on the same membrane is called cis-interaction and can be repulsive or attractive. Here we discuss the case in which this interaction is attractive. Then, below a certain critical temperature, the multi-component membrane undergoes a first-order phase transition and stickers form lateral domains. As shown in Fig. 3, a sticker-poor phase coexists with a sticker-rich phase in the two-phase region of the phase diagram. The sticker critical concentration $`\psi _c`$ and the critical temperature $`T_c`$ are assumed to be the same for the two planar membranes. We define the concentration difference $`\varphi _i(𝒓)`$ for each of the membranes with respect to the critical concentration by
$$\varphi _i(𝒓)=\psi _i(𝒓)\psi _c(i=1,2),$$
(1)
where $`\varphi _i`$ varies between $`1\varphi _i1`$.
The total free energy of the two coupled membranes can be written as a sum of several terms detailed below. The first contribution describes the lateral phase separation of each membrane. Motivated by recent experiments on Langmuir monolayers demonstrating that the inner and outer leaflets of biomembranes are close to their critical point, we employ a phenomenological Ginzburg-Landau free energy which is an expansion in powers of the order parameters $`\{\varphi _i\}`$. Hence we have
$$F_1=\frac{1}{2}\underset{i=1,2}{}\text{d}^2𝒓\left[\frac{1}{2}c(\varphi _i)^2+\frac{1}{2}t\varphi _i^2+\frac{1}{4}\varphi _i^4\mu _i\varphi _i\right].$$
(2)
This expansion for the free energy can be justified close to a critical point where the $`\varphi _i`$’s are small enough. The parameter $`c`$ representing the line tension acting at the domain boundary, and the reduced temperature $`t=(TT_c)/T_c`$ are taken to be the same for the two membranes. On the other hand, the chemical potential $`\mu _i`$, coupled to the membrane sticker concentration $`\varphi _i`$, can differ between the two membranes since the sticker concentrations on the two membranes do not have to be the same. We recall that each bolaform-sticker is modeled with one of its ends anchored irreversibly to one membrane, while the second sticky end is attracted to the second membrane. The sticker concentration is associated with the anchored end of the stickers. The normalization factor 1/2 in (2) is introduced in order to write down the free energy per single membrane. The coefficient of the fourth order term can generally be set as a positive constant without loss of generality.
It is convenient to introduce the following new variables for the average and the difference between the two concentrations:
$$\varphi _+=\frac{\varphi _2+\varphi _1}{2}=\frac{\psi _2+\psi _1}{2}\psi _c,\varphi _{}=\frac{\varphi _2\varphi _1}{2}=\frac{\psi _2\psi _1}{2},$$
(3)
where $`1\varphi _+1`$ and $`1\varphi _{}1`$. In terms of these new variables, (2) can be written as
$`F_1=`$ $`{\displaystyle }\text{d}^2𝒓[{\displaystyle \frac{1}{2}}c[(\varphi _+)^2+(\varphi _{})^2]+{\displaystyle \frac{1}{2}}t(\varphi _+^2+\varphi _{}^2)`$ (5)
$`+{\displaystyle \frac{1}{4}}(\varphi _+^4+6\varphi _+^2\varphi _{}^2+\varphi _{}^4)\mu _+\varphi _+\mu _{}\varphi _{}],`$
where
$$\mu _+=\frac{\mu _2+\mu _1}{2},\mu _{}=\frac{\mu _2\mu _1}{2}.$$
(6)
The chemical potential $`\mu _{}`$ associated with the order parameter $`\varphi _{}`$ is non-zero when the symmetry between the two membranes is explicitly broken. Namely, the two interacting membranes have different average concentrations of stickers.
Next we consider the out-of-plane deformation energy of the two membranes. As depicted in Fig. 2, the membrane shape is parameterized by their heights $`\mathrm{}_1(𝒓)`$, $`\mathrm{}_2(𝒓^{})`$, above the $`x`$-$`y`$ reference plane. Working in the Monge representation it is implicitly assumed that the membranes remain flat on average and have no overhangs. This approach can be also useful to treat adhesion of vesicles in their contact zone. When the vesicle is large enough, it will be roughly flat close to the contact region, and the entire vesicle can be thought of as a reservoir for the stickers. Returning to the deformation energy, it can be written as the sum of the bending energy and the surface tension of each of the two membranes separately, as well as the interacting potential energy between them :
$`F_2=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1,2}{}}{\displaystyle \text{d}^2𝒓\left[\frac{1}{2}\kappa (^2\mathrm{}_i)^2+\frac{1}{2}\sigma (\mathrm{}_i)^2\right]}`$ (8)
$`+{\displaystyle \text{d}^2𝒓v(\mathrm{}_1\mathrm{}_2;\psi _1,\psi _2)},`$
where $`\kappa `$ is the bending rigidity, $`\sigma `$ is the mechanical surface tension acting on the membranes, and $`v`$ is the potential energy per unit area representing the inter-membrane interactions. For simplicity, $`\kappa `$ and $`\sigma `$ are assumed to be equal for the two membranes and do not vary as a function of the sticker concentration $`\psi _i`$. The potential $`v(\mathrm{}_1\mathrm{}_2;\psi _1,\psi _2)`$ can be generally assumed to be a function of the local relative height coordinate $`\mathrm{}_1\mathrm{}_2`$ and the sticker concentration $`\psi _i`$. The former assumption is the so-called Derjaguin approximation . The dependence on the sticker concentration $`\psi _i`$ will be considered later.
We now make a change of variables and transform to the center-of-mass and relative coordinates given, respectively, by
$$L=\frac{\mathrm{}_2+\mathrm{}_1}{2},\mathrm{}=\frac{\mathrm{}_2\mathrm{}_1}{2}.$$
(9)
Only terms which depend on $`\mathrm{}`$ can be considered in the case where the center-of-mass is stationary, hence $`L`$ is a constant of motion. Then (8) can be written as
$$F_2=\text{d}^2𝒓\left[\frac{1}{2}\kappa (^2\mathrm{})^2+\frac{1}{2}\sigma (\mathrm{})^2+v(\mathrm{};\psi _1,\psi _2)\right].$$
(10)
In the high-temperature phase, the stickers are homogeneously distributed, and each of the membrane is in a one-phase region on the phase diagram. We assume that even in the absence of sticker molecules, such membranes are bound to each other due to the balance between the short-range repulsive (e.g., hydration interaction) and longer-range attractive interactions (e.g., van der Waals interaction). Hence, we do not consider the interesting problem of the unbinding transition . Although the membranes are always bound together, their equilibrium distance $`\mathrm{}`$ depends on the sticker concentration. Let us consider the potential $`v(\mathrm{};\psi _1,\psi _2)`$ in (10) for $`\varphi _+=0`$. Note that $`\varphi _+=0`$ means that $`(\psi _1+\psi _2)/2=\psi _c`$, namely, the average sticker concentration on the two membranes is at its critical value $`\psi _c`$. The inter-membrane potential $`v(\mathrm{};\varphi _+=0)`$ is assumed to have a single minimum at a certain inter-membrane distance $`\mathrm{}=\mathrm{}_0`$ for $`\varphi _+=0`$. This gives the equilibrium distance between the two bound membranes. The deviation of the inter-membrane distance from $`\mathrm{}_0`$ is defined by the dimensionless quantity $`\delta (𝒓)`$ given by
$$\delta (𝒓)=\frac{\mathrm{}(𝒓)\mathrm{}_0}{\mathrm{}_0}.$$
(11)
For small deviations from the minimum of the potential, $`v(\mathrm{};\varphi _+=0)`$ can be expanded to second order. This is known as the harmonic approximation and gives
$`v(\mathrm{};\varphi _+=0)`$ $``$ $`v(\mathrm{}_0)+{\displaystyle \frac{1}{2}}v^{\prime \prime }(\mathrm{}_0)(\mathrm{}\mathrm{}_0)^2`$ (12)
$`=`$ $`v(\mathrm{}_0)+{\displaystyle \frac{1}{2}}V\delta ^2,`$ (13)
where $`v^{\prime \prime }(\mathrm{}_0)`$ is the second derivative of $`v`$ with respect to $`\mathrm{}`$ evaluated at $`\mathrm{}=\mathrm{}_0`$, $`Vv^{\prime \prime }(\mathrm{}_0)\mathrm{}_0^2`$, and $`v^{}(\mathrm{})=0`$ at $`\mathrm{}=\mathrm{}_0`$ . Using (11) and (13), (10) can be written as
$$F_2\text{d}^2𝒓\left[\frac{1}{2}K(^2\delta )^2+\frac{1}{2}\mathrm{\Sigma }(\delta )^2+\frac{1}{2}V\delta ^2\right],$$
(14)
with $`K\kappa \mathrm{}_0^2`$ and $`\mathrm{\Sigma }\sigma \mathrm{}_0^2`$. This is the expression of the deformation energy within the harmonic approximation and it served as a starting point to many calculations on membrane adhesion .
Now we will include the effect of the adhesion on the phase separation and suggest a lowest-order coupling between the composition $`\varphi _i(𝒓)`$ and the inter-membrane distance $`\delta (𝒓)`$. When the membranes are quenched into a two-phase region of the phase diagram a sticker-poor phase coexists with a sticker-rich phase. As shown in Fig. 4, this can lead to different inter-membrane distance for the different membrane domains. Since the sticky segments of the bridges adhere directly onto the two membranes, the coupling is proportional to the sum of the local sticker concentrations of the two membranes. This can be phenomenologically represented by the following coupling term:
$$F_3=\frac{\alpha }{2\mathrm{}_0}\text{d}^2𝒓(\psi _1+\psi _2)\mathrm{}=\alpha \text{d}^2𝒓\varphi _+\delta +\mathrm{},$$
(15)
where the coupling constant $`\alpha `$ is positive preferring smaller separation $`\delta <0`$ in regions where the average concentration $`\varphi _+`$ is positive (or $`\psi _1+\psi _2>2\psi _c`$). In the last expression of (15), we have neglected the linear terms in $`\varphi _+`$ and $`\delta `$, which merely shift the chemical potential or minimum of the potential, respectively. Depending on the value of $`\varphi _+`$, this coupling term not only introduces a shift of the minimum of the potential $`v`$ but also changes the minimum value of the potential. Note also that (15) is symmetric with respect to the exchange of the two membranes $`12`$.
The above linear coupling energy can also be understood in the following way. Let us first consider a single flexible membrane with sticker molecules adhering to a flat substrate. Suppose $`v_r(\mathrm{})`$ and $`v_f(\mathrm{})`$ are the potentials for sticker-rich and sticker-free membrane, respectively. Following the same discussion as in (13), each of the potential is parabolic around a different separation: $`v_ra(\mathrm{}\mathrm{}_r)^2`$ and $`v_fa(\mathrm{}\mathrm{}_f)^2`$. The effective potential can be obtained by a linear combination of these two potentials, i.e., $`(1\psi )v_f(\mathrm{})+\psi v_r(\mathrm{})`$. By expanding $`v_f(\mathrm{})`$ and $`v_r(\mathrm{})`$, we get a coupling term which is proportional to $`\psi \mathrm{}`$. In the case of adhesion between two membranes, we add the contributions from both of the membranes and obtain the coupling energy as given in (15). The same argument can be repeated for any arbitrary adhesion potentials, $`v_f(\mathrm{})`$ and $`v_r(\mathrm{})`$, provided each of them has a single well-defined minimum at some distance $`\mathrm{}`$.
The total free energy considered in our model is the sum of (5), (14), and (15):
$$F=F_1+F_2+F_3.$$
(16)
Here it is convenient to convert to Fourier space. The Fourier transform of any function $`f(𝒓)`$ is defined as
$$\stackrel{~}{f}(𝒒)=\text{d}^2𝒓f(𝒓)e^{i𝒒𝒓},$$
(17)
where $`𝒒`$ is the two-dimensional in-plane wavevector. The total free energy can be expressed as
$`F`$ $`=`$ $`{\displaystyle \frac{1}{A}}{\displaystyle \underset{𝒒}{}}[{\displaystyle \frac{1}{2}}(t+cq^2)(|\stackrel{~}{\varphi }_+(𝒒)|^2+|\stackrel{~}{\varphi }_{}(𝒒)|^2)`$ (21)
$`+{\displaystyle \frac{1}{2}}(V+\mathrm{\Sigma }q^2+Kq^4)|\stackrel{~}{\delta }(𝒒)|^2+\alpha \stackrel{~}{\varphi }_+(𝒒)\stackrel{~}{\delta }(𝒒)]`$
$`+{\displaystyle }\text{d}^2𝒓[{\displaystyle \frac{1}{4}}(\varphi _+^4+6\varphi _+^2\varphi _{}^2+\varphi _{}^4)`$
$`\mu _+\varphi _+\mu _{}\varphi _{}],`$
where $`A`$ is the area of the membranes projected on the $`x`$-$`y`$ plane. For convenience the free energy (21) is written as a combination of real space and Fourier space terms.
Within the mean-field level, the free energy in terms of $`\varphi _+`$ and $`\varphi _{}`$ is obtained by functionally minimizing $`F`$ with respect to $`\stackrel{~}{\delta }(𝒒)`$. Then we find
$$\stackrel{~}{\delta }(𝒒)=\frac{\alpha \stackrel{~}{\varphi }_+(𝒒)}{V+\mathrm{\Sigma }q^2+Kq^4}.$$
(22)
Hence the inter-membrane distance $`\delta =(\mathrm{}\mathrm{}_0)/\mathrm{}_0`$ is fully determined by the value of $`\varphi _+`$. By inserting (22) into (21), the resulting free energy depends only on $`\varphi _+`$ and $`\varphi _{}`$, and becomes
$`F`$ $`=`$ $`{\displaystyle \frac{1}{A}}{\displaystyle \underset{𝒒}{}}[{\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Gamma }}_+(𝒒)|\stackrel{~}{\varphi }_+(𝒒)|^2+{\displaystyle \frac{1}{2}}(t+cq^2)|\stackrel{~}{\varphi }_{}(𝒒)|^2]`$ (25)
$`+{\displaystyle }\text{d}^2𝒓[{\displaystyle \frac{1}{4}}(\varphi _+^4+6\varphi _+^2\varphi _{}^2+\varphi _{}^4)`$
$`\mu _+\varphi _+\mu _{}\varphi _{}],`$
where
$$\stackrel{~}{\mathrm{\Gamma }}_+(𝒒)=t+cq^2\frac{\alpha ^2}{V+\mathrm{\Sigma }q^2+Kq^4}.$$
(26)
If we expand the last term in (26) for small $`q`$, we obtain
$$\stackrel{~}{\mathrm{\Gamma }}_+(𝒒)(t\gamma )+\left(c+\frac{\alpha ^2\mathrm{\Sigma }}{V^2}\right)q^2,$$
(27)
with
$$\gamma \frac{\alpha ^2}{V}.$$
(28)
The parameter $`\gamma `$ is an important parameter characterizing the coupling strength. The first two terms in (27) implies an upward shift of the transition temperature, as will be discussed in detail in the next section. We also find that the presence of the coupling ($`\alpha 0`$) increases the line tension $`c`$ provided the mechanical surface tension $`\mathrm{\Sigma }`$ is non-zero; $`cc+\alpha ^2\mathrm{\Sigma }/V^2`$.
## III Phase Diagrams
In this Section, we calculate the mean-field phase diagrams for bolaform-sticker adhesion using the free energy explained in the previous section. In order to study the bulk properties of the system, we set $`𝒒=0`$ and study the homogeneous solutions, $`\varphi _i`$’s and $`\delta `$ being constants. From (22), the inter-membrane distance which minimizes the free energy is given by
$$\delta =\frac{\alpha \varphi _+}{V}.$$
(29)
Since $`\alpha `$ is positive, $`\delta `$ is negative (smaller inter-membrane distance) for positive $`\varphi _+`$, and $`\delta `$ is positive (larger inter-membrane distance) for negative $`\varphi _+`$. By substituting back this $`\delta `$ into the free energy $`f`$ per unit area for homogeneous (constant) $`\varphi _+`$ and $`\varphi _{}`$, we obtain
$`f=`$ $`{\displaystyle \frac{1}{2}}\left(t\gamma \right)\varphi _+^2+{\displaystyle \frac{1}{2}}t\varphi _{}^2`$ (31)
$`+{\displaystyle \frac{1}{4}}(\varphi _+^4+6\varphi _+^2\varphi _{}^2+\varphi _{}^4)\mu _+\varphi _+\mu _{}\varphi _{},`$
where $`\gamma `$ is defined in (28). Notice that $`\gamma `$ is never negative and vanishes only when $`\alpha =0`$. Therefore, when $`\mu _+=\mu _{}=0`$, the field $`\varphi _+`$ will order before $`\varphi _{}`$, and the phase with $`\varphi _+0`$ and $`\varphi _{}=0`$ is expected . Although the phase behavior of this free energy can be examined in general, we concentrate here only on two particular cuts in the parameter space, i.e., $`\mu _{}=0`$ and $`\mu _+=0`$. In these cases one can clearly see the effect of adhesion on the lateral phase separation.
### A The case $`\mu _{}=0`$
When $`\mu _{}=0`$ the two membranes have the same chemical potential $`\mu _1=\mu _2`$. Since the chemical potential $`\mu _+`$ is coupled to $`\varphi _+`$, $`f`$ can be minimized first with respect to $`\varphi _{}`$. A “symmetric phase” is obtained for $`t+3\varphi _+^2>0`$ with
$$\varphi _{}=0,$$
(32)
where the two membranes have the same concentrations, $`\varphi _1=\varphi _2`$. Likewise, two “asymmetric phases” are obtained for $`t+3\varphi _+^2<0`$ with
$$\varphi _{}=\pm \sqrt{t3\varphi _+^2}.$$
(33)
In the asymmetric phase the two membranes have different concentrations $`\varphi _1\varphi _2`$. After inserting these expressions into (31) with $`\mu _{}=0`$, the free energy becomes
$$f_1=\{\begin{array}{ccc}\frac{1}{2}\left(t\gamma \right)\varphi _+^2+\frac{1}{4}\varphi _+^4\mu _+\varphi _+\hfill & & \text{for }t+3\varphi _+^2>0\hfill \\ \frac{1}{4}t^2\frac{1}{2}\left(2t+\gamma \right)\varphi _+^22\varphi _+^4\mu _+\varphi _+\hfill & & \text{for }t+3\varphi _+^2<0\hfill \end{array}.$$
(34)
Notice that this free energy is continuous at $`t=3\varphi _+^2`$. This free energy $`f_1`$ can now be minimized with respect to $`\varphi _+`$. The resulting equation of state is written as
$$\mu _+=\{\begin{array}{ccc}\left(t\gamma \right)\varphi _++\varphi _+^3\hfill & & \text{for }t+3\varphi _+^2>0\hfill \\ \left(2t+\gamma \right)\varphi _+8\varphi _+^3\hfill & & \text{for }t+3\varphi _+^2<0\hfill \end{array}.$$
(35)
The phase diagram can now be calculated and the two-phase region is obtained by the Maxwell construction. The phase diagram for $`\mu _{}=0`$ is illustrated in Fig. 5.
For $`t>0`$, only the symmetric phase with $`\varphi _{}=0`$ can appear since $`t+3\varphi _+^2>0`$. Two symmetric phases with different $`\varphi _+`$ can coexist when $`t<\gamma `$. The coexistence curve is simply given by
$$\varphi _+=\pm \sqrt{t+\gamma },$$
(36)
and the associated critical point is located at
$$(t,\varphi _+,\mu _+)_c=(\gamma ,0,0).$$
(37)
We stress that the critical temperature is increased from $`t_c=0`$ to $`t_c=\gamma =\alpha ^2/V`$ due to the coupling between the composition $`\varphi _+`$ and the inter-membrane distance $`\delta `$ as given in (15). In other words, the phase separation is enhanced by the adhesion of membranes. As presented in Fig. 4, the two coexisting values of $`\varphi _+`$ given by (36) lead to different inter-membrane distances $`\delta `$ according to (29). Since $`\alpha >0`$, $`\delta `$ is negative ($`\mathrm{}<\mathrm{}_0`$) in the sticker-rich domain, and this phase is called the “tight phase” (T). On the other hand, $`\delta `$ is positive ($`\mathrm{}>\mathrm{}_0`$) in the sticker-poor domain and this phase is called the “loose phase” (L). However, for each of the coexisting tight and loose phases, $`\varphi _{}=0`$, which means that the sticker concentration is the same in the two membranes, $`\varphi _1=\varphi _2`$.
For $`\gamma /2<t<0`$, the asymmetric phase with $`\varphi _{}0`$ is always unstable, and the tight and loose phases coexist according to (29) and (36). For $`t<\gamma /2`$, the asymmetric phase can be locally stable but it is only metastable. Namely, its free energy is higher than that of the symmetric phase. Hence the coexistence between the tight and loose phases given by (29) and (36) preempts the asymmetric phase. The limit of metastability of the asymmetric phase is obtained by calculating the second derivative of the second equation of (34) with respect to $`\varphi _+`$. This leads to
$$\varphi _+=\pm \sqrt{\frac{2t\gamma }{24}},$$
(38)
which is also shown as a dotted line inside the L+T coexisting region of Fig. 5(a).
In summary, for $`\mu _{}=0`$, the asymmetric phase $`\varphi _{}0`$ does not exist as a stable phase for any temperature. At most it is metastable and occurs within the L+T coexistence region. The tight and loose phases coexist for $`t<\gamma =\alpha ^2/V`$ according to (29) and (36).
### B The case $`\mu _+=0`$
Next we consider the case of $`\mu _+=0`$ but with $`\mu _{}0`$. This means that the chemical potentials of the two membranes have the same magnitude but opposite sign, i.e., $`\mu _1=\mu _2`$. This is a special case of the more general situation where the symmetry between the two membranes is explicitly broken. Now $`f`$ in (31) can be minimized with respect to $`\varphi _+`$ first. As long as $`t+3\varphi _{}^2>\gamma `$, the only solution is
$$\varphi _+=0.$$
(39)
This is called the “middle phase” (M) where the inter-membrane distance is exactly $`\mathrm{}_0`$ (or $`\delta =0`$). Again note that $`\varphi _+=0`$ means that $`\psi _1+\psi _2=2\psi _c`$. For $`t+3\varphi _{}^2<\gamma `$, we have the tight (or loose) phase with
$$\varphi _+=\pm \sqrt{t+\gamma 3\varphi _{}^2},$$
(40)
where $`\mathrm{}`$ deviates from $`\mathrm{}_0`$ (or $`\delta 0`$) according to (29). Since $`\mu _+=0`$, both the tight and the loose phases are energetically degenerated and they coexist. By substituting $`\varphi _+`$ back into (31) with $`\mu _+=0`$, the free energy becomes
$$f_2=\{\begin{array}{ccc}\frac{1}{2}t\varphi _{}^2+\frac{1}{4}\varphi _{}^4\mu _{}\varphi _{}\hfill & & \text{for }t+3\varphi _{}^2>\gamma \hfill \\ \frac{1}{4}\left(t\gamma \right)^2+\frac{1}{2}\left(2t+3\gamma \right)\varphi _{}^22\varphi _{}^4\mu _{}\varphi _{}\hfill & & \text{for }t+3\varphi _{}^2<\gamma \hfill \end{array}.$$
(41)
After minimizing with respect to $`\varphi _{}`$, the equation of state is given as
$$\mu _{}=\{\begin{array}{ccc}t\varphi _{}+\varphi _{}^3\hfill & & \text{for }t+3\varphi _{}^2>\gamma \hfill \\ \left(2t+3\gamma \right)\varphi _{}8\varphi _{}^3\hfill & & \text{for }t+3\varphi _{}^2<\gamma \hfill \end{array}.$$
(42)
The calculated phase diagrams for $`\mu _+=0`$ are shown in Fig. 6. The phase diagram is symmetric about $`\varphi _{}=0`$ and $`\mu _{}=0`$ as a consequence of the $`\varphi _+^2\varphi _{}^2`$ coupling term, and lack of any odd terms in $`\varphi _+`$ in the free energy. For $`t>\gamma `$ there is a one-phase region of the middle phase with $`\varphi _+=0`$ since $`t+3\varphi _{}^2>\gamma `$. For $`5\gamma /6<t<\gamma `$, the system undergoes a second-order phase transition between the middle phase ($`\varphi _+=0`$) and the tight (or loose) phase ($`\varphi _+0`$). The analytical expressions of the second-order phase transition lines are
$$\varphi _{}=\pm \sqrt{\frac{t+\gamma }{3}},$$
(43)
and
$$\mu _{}=\pm \frac{2t+\gamma }{3}\sqrt{\frac{t+\gamma }{3}},$$
(44)
in Fig. 6, respectively. On the second-order phase transition line, $`\varphi _+`$ goes continuously to zero.
For $`t<5\gamma /6`$, the transition changes to first order. This has been numerically determined by the Maxwell construction. The point which connects the first- and second-order phase transition lines is a tricritical point . In our model, it is located at
$$(t,\varphi _{},\mu _{})_{tcp}=(\frac{5}{6}\gamma ,\pm \frac{1}{3\sqrt{2}}\gamma ^{1/2},\pm \frac{4\sqrt{2}}{27}\gamma ^{3/2}).$$
(45)
The first-order phase transition corresponds to the coexistence of the middle phase with $`\varphi _+=0`$ and the tight (or loose) phase with $`\varphi _+0`$. The obtained two-phase coexistence region is indicated by “M+T” in Fig. 6(a). Within the present Ginzburg-Landau expansion, the tight phase persists even if we go to low temperatures. Because of the degeneracy between tight and loose phases, the first-order line near the tricritical point actually corresponds to coexistence of three phases: tight, loose, and middle phases.
In continuation to the discussion of the previous subsection ($`\mu _{}=0`$), we see that the phase separation is also enhanced for the $`\mu _+=0`$ parameter space. It occurs at higher temperatures, since the tricritical temperature $`t_{tcp}=5\gamma /6=5\alpha ^2/6V`$ is positive for $`\alpha 0`$. It is important to notice that in the middle phase with $`\varphi _+=0`$, the inter-membrane distance is $`\mathrm{}_0`$ since $`\delta =0`$. On the other hand, in the tight (loose) phase with $`\varphi _+>0`$ ($`\varphi _+<0`$), according to (29) and (36), $`\mathrm{}<\mathrm{}_0`$ ($`\mathrm{}>\mathrm{}_0`$). The coexisting membrane domains between tight and middle phases, or between loose and middle phases is schematically represented in Fig. 7.
We end this section by commenting on the general case when both $`\mu _+`$ and $`\mu _{}`$ are non-zero. When $`\mu _+`$ becomes non-zero, the degeneracy between the tight and the loose phases is lifted. In such a case, instead of the three-phase coexistence for $`\mu _+=0`$, there is a coexistence between either tight and middle phases, or between loose and middle phases as shown in Fig. 7(a) and (b), respectively. Notice that the tricritical point exists only when $`\mu _+=0`$. In a more general phase diagram drawn in the ($`t,\mu _+,\mu _{}`$) space, three second-order lines meet at the tricritical point. In the three-dimensional parameter space, these second-order lines lie on the perimeter of two-phase coexistence planes between either tight and middle phases (T+M) or between loose and middle phases (L+M).
## IV Non-Monotonous Membrane Profile
One of our assumptions was that the inter-membranes potential $`v(\mathrm{};\varphi _+=0)`$ has a single minimum at $`\mathrm{}=\mathrm{}_0`$ when $`\varphi _+=0`$. In the absence of thermal fluctuations, two homogeneous membranes are bound with inter-membrane distance $`\mathrm{}_0`$ for $`\varphi _+=0`$. In this section, we calculate the profile of the inter-membrane distance between two membranes which are quenched below the phase separation temperature.
We first expand the potential $`v(\mathrm{};\varphi _+=0)`$ up to the fourth order terms in $`\delta `$;
$$v(\mathrm{};\varphi _+=0)v(\mathrm{}_0)+\frac{1}{2}V\delta ^2+\frac{1}{4}U\delta ^4,$$
(46)
where $`Vv^{\prime \prime }(\mathrm{}_0)\mathrm{}_0^2`$ as before and $`Uv^{(4)}(\mathrm{}_0)\mathrm{}_0^4>0`$. Both $`V`$ and $`U`$ are positive constants because of the convexity of $`v`$ at its minimum. Suppose that each of the membranes is in its high-temperature phase ($`t>0`$). Then the fourth-order $`\varphi _i`$ terms in the Ginzburg-Landau expansion (5) can be neglected since the second-order terms are positive. The resulting free energy with $`\mu _+=\mu _{}=0`$ is
$`F`$ $`=`$ $`{\displaystyle }\text{d}^2𝒓[{\displaystyle \frac{1}{2}}c[(\varphi _+)^2+(\varphi _{})^2]`$ (49)
$`+{\displaystyle \frac{1}{2}}t(\varphi _+^2+\varphi _{}^2)+\alpha \varphi _+\delta `$
$`+{\displaystyle \frac{1}{2}}K(^2\delta )^2+{\displaystyle \frac{1}{2}}\mathrm{\Sigma }(\delta )^2+{\displaystyle \frac{1}{2}}V\delta ^2+{\displaystyle \frac{1}{4}}U\delta ^4]`$
$`=`$ $`{\displaystyle \frac{1}{A}}{\displaystyle \underset{𝒒}{}}[{\displaystyle \frac{1}{2}}(t+cq^2)(|\stackrel{~}{\varphi }_+(𝒒)|^2+|\stackrel{~}{\varphi }_{}(𝒒)|^2)`$ (52)
$`+\alpha \stackrel{~}{\varphi }_+(𝒒)\stackrel{~}{\delta }(𝒒)`$
$`+{\displaystyle \frac{1}{2}}(V+\mathrm{\Sigma }q^2+Kq^4)|\stackrel{~}{\delta }(𝒒)|^2]+{\displaystyle }\text{d}^2𝒓{\displaystyle \frac{1}{4}}U\delta ^4.`$
We now minimize $`F`$ with respect to the concentrations $`\stackrel{~}{\varphi }_+(𝒒)`$ and $`\stackrel{~}{\varphi }_{}(𝒒)`$ and obtain
$$\stackrel{~}{\varphi }_+(𝒒)=\frac{\alpha }{t+cq^2}\stackrel{~}{\delta }(𝒒),\stackrel{~}{\varphi }_{}(𝒒)=0.$$
(53)
By inserting these equations into (52) and expanding for small $`q`$, the free energy can be written as
$$F=\text{d}^2𝒓[\frac{1}{2}K_\mathrm{e}(^2\delta )^2+\frac{1}{2}\mathrm{\Sigma }_\mathrm{e}(\delta )^2+\frac{1}{2}V_\mathrm{e}\delta ^2+\frac{1}{4}U_\mathrm{e}\delta ^4],$$
(54)
with
$`K_\mathrm{e}K{\displaystyle \frac{\alpha ^2c^2}{t^3}},`$ $`\mathrm{\Sigma }_\mathrm{e}\mathrm{\Sigma }+{\displaystyle \frac{\alpha ^2c}{t^2}},`$ (55)
$`V_\mathrm{e}V{\displaystyle \frac{\alpha ^2}{t}},`$ $`U_\mathrm{e}U.`$ (56)
We see that for $`t>0`$ the coupling always increases the mechanical tension $`\mathrm{\Sigma }_\mathrm{e}>\mathrm{\Sigma }`$, but reduces the rigidity $`K_\mathrm{e}<K`$ and the potential strength $`V_\mathrm{e}<V`$.
Let us consider the strong coupling case when $`V_\mathrm{e}<0`$ but still having $`K_\mathrm{e}>0`$, namely,
$$V<\frac{\alpha ^2}{t}<K\left(\frac{t}{c}\right)^2.$$
(57)
For $`t>0`$, although no phase separation occurs in the absence of coupling ($`\alpha =0`$), it occurs for non-zero $`\alpha `$. The minimum free energy of the membranes is given by solving the Euler-Lagrange equation obtained by minimizing (54) with respect to the inter-membrane distance $`\delta `$:
$$K_\mathrm{e}^4\delta \mathrm{\Sigma }_\mathrm{e}^2\delta +V_\mathrm{e}\delta +U_\mathrm{e}\delta ^3=0.$$
(58)
The two uniform (bulk) solutions of (58) are
$$\delta _0=\pm \sqrt{V_\mathrm{e}/U_\mathrm{e}}.$$
(59)
We assume a one-dimensional profile $`\delta (x)`$ describing the inter-membrane distance along the $`x`$-direction. A typical profile determined by a numerical solution of the Euler-Lagrange equation (58) using a relaxational method is shown in Fig. 8. It is convenient to rescale the variables $`\delta `$ and $`x`$ as $`\zeta =U_\mathrm{e}^{1/3}\delta `$ and $`u=K_\mathrm{e}^{1/4}U_\mathrm{e}^{1/12}x`$, respectively, yielding the following one-dimensional profile equation:
$$\frac{d^4\zeta }{du^4}\left(\frac{\mathrm{\Sigma }_\mathrm{e}}{K_\mathrm{e}^{1/2}U_\mathrm{e}^{1/6}}\right)\frac{d^2\zeta }{du^2}+\left(\frac{V_\mathrm{e}}{U_\mathrm{e}^{1/3}}\right)\zeta +\zeta ^3=0.$$
(60)
Only two independent combinations of the four parameters $`K_\mathrm{e}`$ $`\mathrm{\Sigma }_\mathrm{e}`$, $`V_\mathrm{e}`$ and $`U_\mathrm{e}`$ exist. In Fig. 8 they are set to be $`\mathrm{\Sigma }_\mathrm{e}/(K_\mathrm{e}^{1/2}U_\mathrm{e}^{1/6})=0.1`$ and $`V_\mathrm{e}/U_\mathrm{e}^{1/3}=1`$, respectively. The profile has a large slope at the interface $`x=u=0`$, but relaxes to the bulk values $`\pm \delta _0`$ at $`x=\pm \mathrm{}`$ in a non-monotonic fashion with two symmetric overshoots, having a height greater than $`\delta _0`$. These overshoots are suppressed by increasing $`\mathrm{\Sigma }_\mathrm{e}`$ or by increasing the coupling strength $`\alpha `$. The maximum value of $`\delta `$ at the overshoot scales as $`|V_\mathrm{e}|^{1/2}`$ as can be seen from (59). The overshoot of the profile is followed by a damped oscillation which minimizes the curvature energy. This behavior is similar to the nonlinear response of membranes to local pinning sites or membranes adhering to a geometrically structured substrate and is a result of the 4th order derivative in the profile equation. The oscillatory decay has been also predicted for the membrane profile between two inclusions such as proteins .
The configuration of the phase separated membranes corresponding to the above inter-membrane distance $`\delta `$ is schematically represented in Fig. 9. In the case of the adhesion of a single flexible membrane onto a supported membrane, the supported membrane cannot have any shape fluctuations. Therefore, the inter-membrane distance profile calculated in this section can be regarded as a distance of the flexible membrane from the substrate with respect to its equilibrium distance $`\mathrm{}_0`$.
## V Discussion
### A Main Findings
In this paper, the interplay between adhesion and lateral phase separation of multicomponent membranes is investigated. We consider the “bolaform-sticker” adhesion where adhesive bridges are formed by a single sticker having two sticky segments and adhere directly onto the two membranes, as shown in Fig. 1(a). We proposed a phenomenological free energy consisting of three parts: (i) the free energy describing the lateral phase separation of stickers on each membrane (see (2)); (ii) the deformation energy of the two membranes, which is the sum of the bending energy, the surface tension, and the potential energy (see (8)); and, (iii) the coupling energy between the inter-membrane distance and the average concentration of stickers on both membranes (see (15)). The difference of the chemical potentials between the two membranes is also taken into account because the sticker concentrations do not have to be the same.
We calculate the phase diagrams describing the bulk properties for two particular choices of the chemical potentials, i.e., $`\mu _{}=0(\mu _1=\mu _2)`$ and $`\mu _+=0(\mu _1=\mu _2)`$. In the case of $`\mu _{}=0`$, the critical temperature increases depending on the coupling strength and the potential strength (see (37)). Hence the lateral phase separation is enhanced due to the adhesion. This is one of the main consequences of our model. When the phase separation takes place, the inter-membrane distance is smaller for the domains rich in the sticker molecules (“tight phase”), and larger for the domains poor in the stickers (“loose phase”). In the case of $`\mu _+=0`$, our model exhibits a tricritical behavior. The upward shift of the tricritical temperature also indicates the enhancement of the lateral phase separation.
We find that the line tension for the lateral phase separation increases because of the coupling effect as long as the mechanical surface tension is non-zero. We have also calculated the inter-membrane distance profile between the two membranes which are quenched below their phase separation temperature. Because the membrane shape is governed by the bending rigidity, the inter-membrane distance profile relaxes to the bulk values in a non-monotonic way with two symmetric overshoots.
### B Membrane Adhesion on Solid Surfaces and Supported Membranes
So far, we have mainly discussed the adhesion of two membranes. Our model also applies to the case where a single flexible membrane with sticker molecules adheres to a flat substrate or a supported membrane . Let us discuss these two cases separately. For a flat substrate without any supported membrane on it, the contributions from the second membrane (say $`i=2`$) can be dropped from the model. The coupling term (15) simply reduces to $`\alpha \varphi _1\delta `$ because the stickers are assumed to adhere directly to the substrate. The second case is that of a supported membrane with sticker molecules. Unlike the case of two fluctuating membranes discussed in Sec. 8, the supported membrane does not have any shape fluctuations. However, even in such situations, there is an enhancement of the phase separation due to the coupling effect and the upward shift of the critical temperature is given by (37). Another related situation is the case where a membrane is composed of two different lipids and the membrane is put close to a flat substrate. If the two lipids feel different hydration force and prefer different distances from the substrate, the phase separation between the two components will be enhanced by the adhesion for the same reason described in this paper.
### C Relation to Other Models
There exists an analogy between the phase behavior of our membrane system with that of metamagnets (magnets which undergo fist-order phase transitions in an increasing magnetic field) or <sup>3</sup>He-<sup>4</sup>He mixtures described by the BEG (Blume-Emery-Griffith) spin-one model . Moreover, the phase diagrams for $`\mu _+=0`$ are analogous to those describing the phase separation of two-component mixtures in fluid bilayers which also exhibits tricritical behavior and other related amphiphilic systems . However, in the former case of two-component bilayers the concentration difference between the two leaflets of the membrane is linearly coupled to the curvature of the bilayer and the difference in the chemical potential is not taken into account.
In our paper, we did not address the problem of the unbinding transition. We rather assumed that the membranes are always bound together, even in the absence of any sticker molecules. This assumption is partially motivated by the experimental study of Ref. where suspended membrane (part of the giant vesicle) was claimed to be bound to the supporting membrane even in the absence of sticker molecules. In this case, the inter-membrane distance $`\mathrm{}`$ stays finite and it is permissible to expand the free energy around the minimum. Hence the phase separation consists of loosely and tightly bound patches. The interplay between unbinding transition and phase separation of multi-component membranes has been considered in other theoretical works . The adhesion there is only brought about by sticker molecules, and the phase separation is induced both by attractive interactions and fluctuation-induced interactions between the stickers. Although their model treats a different aspect of the more general problem, the fluctuation effect yields similar consequences compared to ours. We assumed that the $`cis`$-interaction in (2) is attractive, and tracing over the inter-membrane distance $`\delta `$ yields a term proportional to $`\gamma \varphi _+^2`$. Since this term does not depend on the sign of $`\varphi _+`$, it has a similar effect as fluctuations although our treatment is restricted to the mean-field level.
It is worthwhile to comment here the difference between the present study and that of Ref. . In their paper, it is found that the adhesion between the membranes including homophilic recognition molecules and repeller molecules is controlled by lateral phase separation. The multiple competing states of adhesion is attributed to the double-well inter-membrane interaction potential generated by the competition of two forces; attraction between homophilic molecules and the repulsion between repeller molecules. By changing the repeller concentration, the double-minimum potential causes the first-order transition between a state with inter-membrane spacing set by the thickness of the repeller molecules to a state with a spacing set by the bare potential (van der Waals plus hydration interactions). In our work, the effect of repeller molecules is not taken into account and the minimum of the potential depends on the sticker concentration through the coupling term (15). When the stickers are phase separated and two different values of the sticker concentration coexist, the inter-membrane potential has double-minimum. However the physical origin of this double-minimum potential is different from that in Ref. because it is not due to the presence of repeller molecules.
### D Other Types of Sticker Molecules
As mentioned in the introduction, “homophilic-sticker” adhesion occurs when the adhesive bridges are formed by two stickers of the same type bound together by their two sticky segments (see Fig. 1(b)). Suppose $`\psi _i`$ ($`i=1,2`$) denotes the sticker concentration on each membrane. Then the inter-membrane distance depends on the product of each sticker concentration expressing the probability to have two stickers – one on each membrane – at the same position. Using (1) this coupling term can be written in terms of $`\varphi _i`$ as
$$\psi _1\psi _2=\varphi _1\varphi _2+\psi _c(\varphi _1+\varphi _2)+\psi _c^2.$$
(61)
An interesting remark can be made for homophilic-sticker adhesion. The resulting phase separation within each membrane leads to three different values for the inter-membrane distance. The inter-membrane distance between domains rich in stickers on both membranes (rich-rich), as well as between rich-poor domains, and poor-poor domains can be different . Notice that these three different inter-membrane distances correspond to the tight, middle, and loose phases in our model.
A third case is that of “lock-and-key” adhesion due to the formation of chemical bonds between lock-and-key types of stickers, e.g., ligands and receptors (see Fig. 1(c)). Suppose that both types of stickers are distributed on the two membranes and $`\psi _i`$ now represents the local concentration, say, of the lock molecules. First let us assume that the membranes are saturated with sticker molecules (no lipid). Then, $`1\psi _i(𝒓)`$ represents the concentration of key molecules. Since domains rich in lock (key) molecules on one membrane adhere with domains rich in key (lock) molecules on the other membrane, the coupling tern in the free energy $`F_3`$ will be a coupling between the inter-membrane distance and
$`\psi _1(1\psi _2)+(1\psi _1)\psi _2`$ (62)
$`=`$ $`2\varphi _1\varphi _2+(12\psi _c)(\varphi _1+\varphi _2)+2\psi _c(1\psi _c).`$ (63)
This term is symmetric with respect to the exchange of two membranes. Both in (61) and (63), we see that the lowest order term in the concentration (except the constant term) is proportional to $`\varphi _+`$. If there is a linear coupling between the inter-membrane distance and $`\varphi _+`$ in these cases, we expect an upward shift in the transition temperature and the phase separation will be enhanced as argued above. Due to the presence of higher order terms in (61) and (63), however, the phase behavior will be more complex.
Let us now take into account the presence of lipids in the lock-and-key adhesion. If a single type of sticker is present on each membrane, namely, lock molecules on membrane 1 and key molecules on membrane 2, we can regard $`\psi _1`$ and $`\psi _2`$ as the concentrations of lock and key molecules embedded in the lipid membrane, respectively. Then $`1\psi _1`$ and $`1\psi _2`$ describe the concentration of the second component (lipid) on each membrane, respectively. In this case, the inter-membrane distance depends on $`\psi _1\psi _2`$ as in (61). When both lock and key molecules are present on both membranes, one has to start with a three-component mixture for each of the membranes. The generic lattice model to study the behavior of ternary membranes of monolayers is the BEG spin-one model . Here one has to include the coupling between the two membranes. If we denote the concentration of lock and key molecules on each membrane as $`\psi _i^L`$ and $`\psi _i^K`$ (and hence the concentration of the dilution lipid is $`1\psi _i^L\psi _i^K`$), the inter-membrane distance now depends on $`\psi _1^L\psi _2^K+\psi _1^K\psi _2^L`$ which is similar to (63). More detailed calculations for the homophilic stickers and lock-and-key stickers and their influence on membrane adhesion are left for future studies.
###### Acknowledgements.
We have greatly benefited from the discussions and correspondence with R. Netz, J. Rädler, E. Sackmann, and T. Weikl. SK would like to thank the Ministry of Education, Science and Culture, Japan for providing financial support during his visit to Israel. Support from the exchange program between the Japan Society for the Promotion of Science (JSPS) and the Israel Ministry of Science and Technology is also gratefully acknowledged. DA acknowledges partial support from the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities, Centers of Excellence Program and the US-Israel Binational Science Foundation (BSF) under grant number 98-00429.
Fig.1 Komura and Andelman
Fig.2 Komura and Andelman
Fig.3 Komura and Andelman
Fig.4 Komura and Andelman
Fig.5(a) Komura and Andelman
Fig.5(b) Komura and Andelman
Fig.6(a) Komura and Andelman
Fig.6(b) Komura and Andelman
Fig.7(a) Komura and Andelman
Fig.7(b) Komura and Andelman
Fig.8 Komura and Andelman
Fig.9 Komura and Andelman |
warning/0001/astro-ph0001318.html | ar5iv | text | # Review of Big Bang Nucleosynthesis and Primordial Abundances
## 1 Introduction
There are now four main observations which validate the Big Bang theory: the expansion of the universe, the Planck spectrum of the Cosmic Microwave Background (CMB), the density fluctuations seen in the slight CMB anisotropy and in the local galaxy distribution, and BBN. Together, they show that the universe began hot and dense .
BBN occurs at the earliest times at which we have a detailed understanding of physical processes. It makes predictions which are relatively precise (10% – 0.1%), and which have been verified with a variety of data. It is critically important that the standard theory (SBBN) predicts the abundances of several light nuclei (H, D, <sup>3</sup>He <sup>4</sup>He and <sup>7</sup>Li ) as a function of a single cosmological parameter, the baryon to photon ratio, $`\eta n_b/n_\gamma `$ . The ratio of any two primordial abundances should give $`\eta `$, and the measurement of the other three tests the theory.
The abundances of all the light elements have been measured in a number of terrestrial and astrophysical environments. Although it has often been hard to decide when these abundances are close to primordial, it has been clear for decades (e.g. , ) that there is general agreement with the BBN predictions for all the light nuclei. The main development in recent years has been the increased accuracy of measurement. In 1995 a factor of three range in the baryon density was considered $`\mathrm{\Omega }_b=0.0070.024`$. The low end of this range allowed no significant dark baryonic matter. Now the new D/H measurements towards quasars give $`\mathrm{\Omega }_b=0.019\pm 0.0024`$ (95%) – a 13% error, and there have been improved measurements of the other nuclei.
### 1.1 Other Reviews
Many reviews of BBN have been published recently: e.g. , , , , , , and , some of which are lengthy: e.g. , and . All modern cosmology texts contain a summary. Several recent books contain the proceedings of meetings on this topic: , , and . The 1999 meeting of the International Astronomical Union (Symposium 198 in Natal, Brazil) was on this topic, as are many reviews in upcoming special volumes of Physics Reports and New Astronomy, both in honor of the major contributions by David N. Schramm.
## 2 Physics of BBN
Excellent summaries are given in most books on cosmology e.g.: , , , , , and most of the reviews listed above, including , and .
### 2.1 Historical Development
The historical development of BBN is reviewed by , , , and . Here we mention a few of the main events.
The search for the origin of the elements lead to the modern Big Bang theory in the early 1950s. The expansion of the universe was widely accepted when Lemaitre suggested that the universe began in an explosion of a dense unstable “primeval atom”. By 1938 it was well established that the abundances of the elements were similar in different astronomical locations, and hence potentially of cosmological significance. Gamow , asked whether nuclear reactions in the early universe might explain the abundances of the elements. This was the first examination of the physics of a dense expanding early universe, beyond the mathematical description of general relativity, and over the next few years this work developed into the modern big bang theory. Early models started with pure neutrons, and gave final abundances which depended on the unknown the density during BBN. Fermi & Turkevich showed that the lack of stable nuclei with mass 5 and 8 prevents significant production of nuclei more massive than <sup>7</sup>Li, leaving <sup>4</sup>He as the most abundant nucleus after H. Starting instead with all possible species, Hayashi first calculated the neutron to proton (n/p) ratio during BBN, and Alpher realized that radiation would dominate the expansion. By 1953 the basic physics of BBN was in place. This work lead directly to the prediction of the CMB (e.g. Olive 1999b ), it explained the origin of D, and gave abundance predictions for <sup>4</sup>He similar to those obtained today with more accurate cross-sections.
The predicted abundances have changed little in recent years, following earlier work by Peebles (1964) , Hoyle & Tayler (1964) , and Wagoner, Fowler & Hoyle (1967) . The accuracy of the theory calculations have been improving, and they remain much more accurate than the measurements. For example, the fraction of the mass of all baryons which is <sup>4</sup>He, $`Y_p`$, is predicted to within $`\delta `$$`Y_p`$ $`<\pm 0.0002`$ . In a recent update, Burles et al. uses Monte-Carlo realizations of reaction rates to find that the previous estimates of the uncertainties in the abundances for a given $`\eta `$ were a factor of two too large.
## 3 Key Physical Processes
### 3.1 Baryogenesis
The baryon to photon ratio $`\eta `$ is determined during baryogenesis , , . It is not known when baryogenesis occurred. Sakharov noted that three conditions are required: different interactions for matter and anti-matter (CP violation), interactions which change the baryon number, and departure from thermodynamic equilibrium. This last condition may be satisfied in a first order phase transition, the GUT transition at $`10^{35}`$ s, or perhaps the electroweak transition at $`10^{11}`$ s. If baryogenesis occurred at the electroweak scale, then future measurements may lead to predictions for $`\eta `$, but if, alternatively, baryogenesis is at the GUT or inflation scale, it will be very hard to predict $`\eta `$ (J. Ellis personal communication).
The matter/anti-matter asymmetry of the universe (the $`\eta `$ value) is attracting discussion in the popular science press because of the inauguration of major experiments to study CP violation in B mesons (, ; Economist, May 8 1999, 85-87).
### 3.2 The main physical processes in BBN
At early times, weak reactions keep the n/p ratio close to the equilibrium Boltzmann ratio. As the temperature, T, drops, n/p decreases. The n/p ratio is fixed (“frozen in”) at a value of about 1/6 after the weak reaction rate is slower than the expansion rate. This is at about 1 second, when $`T1`$MeV. The starting reaction n+p $`{}_{}{}^{}D+\gamma `$ makes D. At that time photodissociation of D is rapid because of the high entropy (low $`\eta `$) and this prevents significant abundances of nuclei until, at 100 sec., the temperature has dropped to 0.1 MeV, well below the binding energies of the light nuclei. About 20% of free neutrons decay prior to being incorporated into nuclei. The <sup>4</sup>He abundance is then given approximately by assuming that all remaining neutrons are incorporated into <sup>4</sup>He.
The change in the abundances over time for one $`\eta `$ value is shown in Figure 1, while the dependence of the final abundances on $`\eta `$ is shown in Figure 2, together with some recent measurements.
In general, abundances are given by two cosmological parameters, the expansion rate and $`\eta `$. Comparison with the strength of the weak reactions gives the n/p ratio, which determines $`Y_p`$. $`Y_p`$ is relatively independent of $`\eta `$ because n/p depends on weak reactions between nucleons and leptons (not pairs of nucleons), and temperature. If $`\eta `$ is larger, nucleosynthesis starts earlier, more nucleons end up in <sup>4</sup>He, and $`Y_p`$ increases slightly. D and <sup>3</sup>He decrease simultaneously in compensation. Two channels contribute to the abundance of <sup>7</sup>Li in the $`\eta `$ range of interest, giving the same <sup>7</sup>Li for two values of $`\eta `$.
## 4 Measurement of Primordial Abundances
The goal is to measure the primordial abundance ratios of the light nuclei made in BBN. We normally measure the ratios of the abundances of two nuclei in the same gas, one of which is typically H, because it is the easiest to measure.
The two main difficulties are the accuracy of the measurement and departures from primordial abundances. The state of the art today (1 $`\sigma `$) is about 3% for $`Y_p`$, 10% for D/H and 8% for <sup>7</sup>Li, for each object observed. These are random errors. The systematic errors are hard to estimate, usually unreliable, and potentially much larger.
By the earliest time at which we can observe objects, redshifts $`z6`$, we find heavy elements from stars in most gas. Although we expect that large volumes of the intergalactic medium (IGM) remain primordial today , we do not know how to obtain accurate abundances in this gas. Hence we must consider possible modifications of abundances. This is best done in gas with the lowest abundances of heavy elements, since this gas should have the least deviations caused by stars.
The nuclei D, <sup>3</sup>He, <sup>6</sup>Li and <sup>7</sup>Li are all fragile and readily burned inside stars at relatively low temperatures of a few $`10^6`$ K. They may appear depleted in the atmosphere of a star because the gas in the star has been above the critical temperature, and they will be depleted in the gas returned to the interstellar medium (ISM). Nuclei <sup>3</sup>He, <sup>7</sup>Li and especially <sup>4</sup>He are also made in stars.
### 4.1 From Observed to Primordial Abundances
Even when heavy element abundances are low, it is difficult to prove that abundances are primordial. Arguments include the following.
Helium is observed in the ionized gas surrounding luminous young stars (H II regions), where O abundances are 0.02 to 0.2 times those in the sun. The <sup>4</sup>He mass fraction $`Y`$ in different galaxies is plotted as a function of the abundance of O or N. The small change in $`Y`$ with O or N is the clearest evidence that the $`Y`$ is almost entirely primordial (e.g. Fig 2). Regression gives the predicted $`Y_p`$ for zero O or N . The extrapolation is a small extension beyond the observed range, and the deduced primordial $`Y_p`$ is within the range of $`Y`$ values for individual H II regions. The extrapolation should be robust , but some algorithms are sensitive to the few galaxies with the lowest metal abundances, which is dangerous because at least one of these values was underestimated by Olive, Skillman & Steigman .
For deuterium we use a similar argument. The observations are made in gas with two distinct metal abundances. The quasar absorbers have from 0.01 to 0.001 of the solar C/H, while the ISM and pre-solar observations are near solar. Since D/H towards quasars is twice that in the ISM, 50% of the D is destroyed when abundances rise to near the solar level, and less than 1% of D is expected to be destroyed in the quasar absorbers, much less than the random errors in individual measurements of D/H. Since there are no other known processes which destroy or make significant D (e.g., ), we should be observing primordial D/H in the quasar absorbers.
Lithium is more problematic. Stars with a variety of low heavy element abundances (0.03 – 0.0003 of solar) show very similar abundances of <sup>7</sup>Li ( Fig 3), which should be close to the primordial value. Some use the observed values in these “Spite plateau” stars as the BBN abundance, because of the small scatter and lack of variation with the abundances of other elements, but three factors should be considered. First, the detection of <sup>6</sup>Li in two of these stars suggests that both <sup>6</sup>Li and some <sup>7</sup>Li was been created prior to the formation of these stars. Second, the possible increase in the abundance of <sup>7</sup>Li with the iron abundance also indicates that the <sup>7</sup>Li of the plateau stars is not primordial. If both the iron and the enhancement in the <sup>7</sup>Li have the same origin we could extrapolate back to zero metals , as for <sup>4</sup>He, but the enhanced <sup>7</sup>Li may come from cosmic ray interactions in the ISM, which makes extrapolation less reliable. Third, the amount of depletion is hard to estimate. Rotationally induced mixing has a small effect because there is little scatter on the Spite plateau, but other mechanisms may have depleted <sup>7</sup>Li. In particular, gravitational settling should have occurred, and left less <sup>7</sup>Li in the hotter plateau stars, but this is not seen, and we do not know why. More on this later.
The primordial abundance of <sup>3</sup>He is the hardest to estimate, because stars are expected to both make and destroy this isotope, and there are no measurements in gas with abundances well below the solar value.
### 4.2 Key observational Requirements
By way of introduction to the data, we list some of the key goals of ongoing measurements of the primordial abundances.
* <sup>4</sup>He: High accuracy, robust measurement in a few places with the lowest metal abundances.
* <sup>3</sup>He: Measurement in gas with much lower metal abundances, or an understanding of stellar production and destruction and the results of all stars integrated over the history of the Galaxy (Galactic chemical evolution).
* D: The discovery of more quasar absorption systems with minimal H contamination.
* <sup>7</sup>Li: Observations which determine the amount of depletion in halo stars, or which avoid this problem. Measurement of <sup>6</sup>Li, Be and B to help estimate production prior to halo star formation, and subsequent depletion.
Since we are now obtaining “precision” measurements, it now seems best to make a few measurements with the highest possible accuracy and controls, in places with the least stellar processing, rather than multiple measurements of lower accuracy. We will now discuss observations of each of the nuclei, and especially D, in more detail.
## 5 Deuterium in quasar spectra
The abundance of deuterium (D or <sup>2</sup>H) is the most sensitive measure of the baryon density . No known processes make significant D, because it is so fragile (, , and ). Gas ejected by stars should contain zero D, but substantial H, thus D/H decreases over time as more stars evolve and die.
We can measure the primordial abundance in quasar spectra. The measurement is direct and accurate, and with one exception, simple. The complication is that the absorption by D is often contaminated or completely obscured by the absorption from H, and even in the rare cases when contamination is small, superb spectra are required to distinguish D from H.
Prior to the first detection of D in quasar spectra , D/H was measured in the ISM and the solar system. The primordial abundance is larger, because D has been destroyed in stars. Though generally considered a factor of a few, some papers considered a factor of ten destruction . At that time, most measurements of <sup>4</sup>He gave low abundances, which predict a high primordial D/H, which would need to be depleted by a large factor to reach ISM values .
Reeves, Audouze, Fowler & Schramm noted that the measurement of primordial D/H could provide an excellent estimate of the cosmological baryon density, and they used the ISM <sup>3</sup>He +D to concluded, with great caution, that primordial D/H was plausibly $`7\pm 3`$ $`\times 10^5`$.
Adams suggested that it might be possible to measure primordial D/H towards low metallicity absorption line systems in the spectra of high redshift quasars. This gas is in the outer regions of galaxies or in the IGM, and it is not connected to the quasars. The importance of such measurements was well known in the field since late 1970s , but the task proved too difficult for 4-m class telescopes (, , ). The high SNR QSO spectra obtained with the HIRES echelle spectrograph on the W.M. Keck 10-m telescope provided the breakthrough.
There are now three known absorption systems in which D/H is low: first, D/H $`=3.24\pm 0.3\times 10^5`$ in the $`z_{abs}=3.572`$ Lyman limit absorption system (LLS) towards quasar 1937–1009 , ; second, D/H = $`4.0_{0.6}^{+0.8}\times 10^5`$ in the $`z_{abs}=2.504`$ LLS towards quasar 1009+2956 , and third, D/H $`<6.7`$ $`\times 10^5`$ towards quasar 0130–4021 . This last case is the simplest found yet, and seems especially secure because the entire Lyman series is well fit by a single velocity component. The velocity of this component and its column density are well determined because many of its Lyman lines are unsaturated. Its Ly$`\alpha `$ line is simple and symmetric, and can be fit using the H parameters determined by the other Lyman series lines, with no additional adjustments for the Ly$`\alpha `$ absorption line. There is barely enough absorption at the expected position of D to allow low values of D/H, and there appears to be no possibility of high D/H. Indeed, the spectra of all three QSOs are inconsistent with high D/H.
There remains uncertainty over a case at $`z_{abs}=0.701`$ towards quasar 1718+4807, because we lack spectra of the Lyman series lines which are needed to determine the velocity distribution of the Hydrogen, and these spectra are of unusually low signal to noise, with about 200 times fewer photons per kms<sup>-1</sup> than those from Keck. Webb et al. , assumed a single hydrogen component and found D/H = $`25\pm 5\times 10^5`$, the best case for high D/H. Levshakov et al. allow for non-Gaussian velocities and find D/H $`4.4\times 10^5`$, while Tytler et al. find $`8\times 10^5<`$ D/H $`<57\times 10^5`$ (95%) for a single Gaussian component, or D/H as low as zero if there are two hydrogen components, which is not unlikely. This quasar is then also consistent with low D/H.
Recently Molaro et al. claimed that D/H might be low in an absorber at $`z=3.514`$ towards quasar APM 08279+5255, though they noted that higher D/H was also possible. Only one H I line, Ly$`\alpha `$, was used to estimate the hydrogen column density $`N_{HI}`$ and we know that in such cases the column density can be highly uncertain. Their Figure 1 (panels a and b) shows that there is a tiny difference between D/H = 1.5 $`\times 10^5`$ and 21 $`\times 10^5`$, and it is clear that much lower D is also acceptable because there can be H additional contamination in the D region of the spectrum. Levshakov et al. show that $`N_{HI}`$= 15.7 (too low to show D) gives an excellent fit to these spectra, and they argue that this is a more realistic result because the metal abundances and temperatures are then normal, rather than being anomalously low with the high $`N_{HI}`$ preferred by Molaro et al.
The first to publish a D/H estimate using high signal to noise spectra from the Keck telescope with the HIRES spectrograph were Songaila et al. , who reported an upper limit of D/H $`<25\times 10^5`$ in the $`z_{abs}=3.32`$ Lyman limit system (LLS) towards quasar 0014+813. Using different spectra, Carswell et al. reported $`<60\times 10^5`$ in the same object, and they found no reason to think that the deuterium abundance might be as high as their limit. Improved spectra support the early conclusions: D/H $`<35\times 10^5`$ for this quasar. High D/H is allowed, but is highly unlikely because the absorption near D is at the wrong velocity, by $`17\pm 2`$ km s<sup>-1</sup>, it is too wide, and it does not have the expected distribution of absorption in velocity, which is given by the H absorption. Instead this absorption is readily explained entirely by H (D/H $`0`$) at a different redshift.
Very few LLS have a velocity structure simple enough to show deuterium. Absorption by H usually absorbs most of the quasar flux near where the D line is expected, and hence we obtain no information of the column density of D. In these extremely common cases, very high D/H is allowed, but only because we have essentially no information.
All quasar spectra are consistent with low primordial D/H ratio, D/H $`3.4\times 10^5`$. Two quasars (1937–1009 & 1009+2956) are inconsistent with D/H $`5\times 10^5`$, and the third (0130–4021) is inconsistent with D/H $`6.7\times 10^5`$. Hence D/H is low in these three places. Several quasars allow high D/H, but in all cases this can be explained by contamination by H, which we discuss more below, because this is the key topic of controversy.
### 5.1 ISM D/H
Observations of D in the ISM are reviewed by Lemoine et al. . The first measurement in the ISM, D/H $`=1.4\pm 0.2`$ $`\times 10^5`$, using Lyman absorption lines observed with the Copernicus satellite , have been confirmed with superior HST spectra. A major program by Linsky et al. , has given a secure value for local ISM ($`<20`$ pc) D/H = $`1.6\pm 0.1`$ $`\times 10^5`$.
Some measurements have indicated variation, and especially low D/H, in the local and more distant ISM towards a few stars , . Vidal-Madjar & Gry concluded that the different lines of sight gave different D/H, but those early data may have been inadequate to quantify complex velocity structure . Variation is expected, but at a low level, from different amounts of stellar processing and infall of IGM gas, which leaves differing D/H if the gas is not mixed in a large volume.
Lemoine et al. suggested variation of D/H towards G191-B2B, while Vidal-Madjar et al. described the variation as real, however new STIS spectra do not confirm this, and give the usual D/H value. The STIS spectra show a simpler velocity structure, and a lower flux at the D velocity, perhaps because of difficulties with the background subtraction in the GHRS spectra.
H$`\stackrel{´}{\mathrm{e}}`$brand et al. report the possibility of low D/H $`<1.6`$ $`\times 10^5`$ towards Sirius A, B.
The only other instance of unusually low D/H from recent data is D/H $`=0.74_{0.13}^{+0.19}`$ $`\times 10^5`$ (90%) towards the star $`\delta `$ Ori . We would much like to see improved data on this star, because a new instrument was used, the signal to noise is very low, and the velocity distribution of the D had to be taken from the N I line, rather than from the H I.
Possible variations in D/H in the local ISM have no obvious connections to the D/H towards quasars, where the absorbing clouds are 100 times larger, in the outer halos young of galaxies rather in the dense disk, and the influence of stars should be slight because heavy element abundances are 100 to 1000 times smaller.
Chengalur, Braun & Burton report D/H $`=3.9\pm 1.0`$ $`\times 10^5`$ from the marginal detection of radio emission from the hyper-fine transition of D at 327 MHz (92 cm). This observation was of the ISM in the direction of the Galactic anti-center, where the molecular column density is low, so that most D should be atomic. The D/H is higher than in the local ISM, and similar to the primordial value, as expected, because there has been little stellar processing in this direction.
Deuterium has been detected in molecules in the ISM. Some of these results are considered less secure because of fractionation and in low density regions, HD is more readily destroyed by ultraviolet radiation, because its abundance is too low to provide self shielding, making HD/H<sub>2</sub> smaller than D/H.
However, Wright et al. deduce D/H = $`1.0\pm 0.3\times 10^5`$ from the first detection of the 112 $`\mu `$m pure rotation line of HD outside the solar system, towards the dense warm molecular clouds in the Orion bar, where most D is expected to be in HD, so that D/H $``$ HD/H<sub>2</sub>. This D/H is low, but not significantly lower than in the local ISM, especially because the H<sub>2</sub> column density was hard to measure.
Lubowich et al. , report D/H = $`0.2\pm 1`$ $`\times 10^5`$ from DCN in the Sgr A molecular cloud near the Galactic center, later revised to 0.3 $`\times 10^5`$ (private communication 1999). This detection has two important implications. First, there must be a source of D, because all of the gas here should have been inside at least one star, leaving no detectable D. Nucleosynthesis is ruled out because this would enhance the Li and B abundances by orders of magnitude, contrary to observations. Infall of less processed gas seems likely. Second, the low D/H in the Galactic center implies that there is no major source of D, otherwise D/H could be very high. However, this is not completely secure, since we could imagine a fortuitous cancellation between creation and destruction of D.
We eagerly anticipate a dramatic improvement in the data on the ISM in the coming years. The FUSE satellite, launched in 1999, will measure the D and H Lyman lines towards thousands of stars and a few quasars, while SOFIA (2002) and FIRST (2007) will measure HD in dense molecular clouds. The new GMAT radio telescope should allow secure detection of D 82 cm emission from the outer Galaxy, while the Square Kilometer Array Interferometer would be able to image this D emission in the outer regions of nearby galaxies; regions with low metal abundances. These data should give the relationship between metal abundance and D/H, and especially determine the fluctuations of D/H at a given metal abundance which will better determine Galactic chemical evolution, and, we expect, allow an accurate prediction of primordial D/H independent of the QSO observations.
### 5.2 Solar System D/H
The D/H in the ISM from which the solar system formed 4.6 Gyr ago can be deduced from the D in the solar system today, since there should be no change in D/H, except in the sun.
Measurement in the atmosphere of Jupiter will give the pre-solar D/H provided (1) most of Jupiter’s mass was accreted directly from the gas phase, and not from icy planetessimals, which, like comets today, have excess D/H by fractionation, and (2) the unknown mechanisms which deplete He in Jupiter’s atmosphere do not depend on mass. Mahaffy et al. find D/H = $`2.6\pm 0.7`$ $`\times 10^5`$ from the Galileo probe mass spectrometer. Feuchtgruber et al. used infrared spectra of the pure rotational lines of HD at 37.7 $`\mu `$m to measure D/H $`=5.5_{1.5}^{+3.5}`$$`\times 10^5`$ in Uranus and $`6.5_{1.5}^{+2.5}`$$`\times 10^5`$ in Neptune, which are both sensibly higher because these planets are known to be primarily composed of ices which have excess D/H.
The pre-solar D/H can also be deduced indirectly from the present solar wind, assuming that the pre-solar D was converted into <sup>3</sup>He. The present <sup>3</sup>He/ <sup>4</sup>He ratio is measured and corrected for (1) changes in <sup>3</sup>He/H and <sup>4</sup>He/H because of burning in the sun, (2) the changes in isotope ratios in the chromosphere and corona, and (3) the <sup>3</sup>He present in the pre-solar gas. Geiss & Gloeckler reported D/H = $`2.1\pm 0.5`$ $`\times 10^5`$, later revised to $`1.94\pm 0.36`$ $`\times 10^5`$ .
The present ISM D/H $`=1.6\pm 0.1`$$`\times 10^5`$ is lower, as expected, and consistent with Galactic chemical evolution models, which we now mention.
### 5.3 Galactic Chemical Evolution of D
Numerical models are constructed to follow the evolution of the abundances of the elements in the ISM of our Galaxy.
The main parameters of the model include the yields of different stars, the distribution of stellar masses, the star formation rate, and the infall and outflow of gas. These parameters are adjusted to fit many different data. These Galactic chemical evolution models are especially useful to compare abundances at different epochs, for example, D/H today, in the ISM when the solar system formed, and primordially.
In an analysis of a variety of different models, Tosi et al. concluded that the destruction of D in our Galaxy was at most a factor of a few, consistent with low but not high primordial D. They find that all models, which are consistent with all Galactic data, destroy D in the ISM today by less than a factor of three. Such chemical evolution will destroy an insignificant amount of D when metal abundances are as low as seen in the quasar absorbers.
Others have designed models which do destroy more D , , , , for example, by cycling most gas through low mass stars and removing the metals made by the accompanying high mass stars from the Galaxy. These models were designed to reduce high primordial D/H, expected from the low $`Y_p`$ values prevalent at that time, to the low ISM values. Tosi et al. describe the generic difficulties with these models. To destroy 90% of the D, 90% of the gas must have been processed in and ejected from stars. These stars would then release more metals than are seen. If the gas is removed (e.g. expelled from the galaxy) to hide the metals, then the ratio of the mass in gas to that in remnants is would be lower than observed. Infall of primordial gas does not help, because this brings in excess D. These models also fail to deplete the D in quasar absorbers, because the stars which deplete the D, by ejecting gas without D, also eject carbon. The low abundance of carbon in the absorbers limits the destruction of D to $`<`$1% .
### 5.4 Questions About D/H
Here we review some common questions about D/H in quasar spectra.
#### 5.4.1 Why is saturation of absorption lines important?
Wampler suggested that the low D/H values might be inaccurate because in some cases the H absorption lines have zero flux in their cores; they are saturated. Songaila, Wampler & Cowie suggested that this well known problem might lead to errors in the H column density, but later work, using better data and more detailed analyses has shown that these concerns were not significant, and that the initial result was reliable.
Neutral deuterium (D I) is detected in Lyman series absorption lines, which are adjacent to the H I lines. The separation of 82 km s<sup>-1</sup> is easily resolved in high resolution spectra, but it is not enough to move D out of the absorption by the H. The Lyman series lines lie between 1216Å and 912Å, and can be observed from the ground at redshifts $`>2.5`$.
Ideally, many (in the best cases $`>20`$) Lyman lines are observed, to help determine the column density ($`N_{HI}`$, measured in H I atoms per cm<sup>-2</sup> along the line of sight) and velocity width (b values, $`b=\sqrt{2}\sigma `$, measured in km s<sup>-1</sup>) of the H. But in some cases only Ly$`\alpha `$ has been observed (Q1718+4807, APM 08279+5255), and these give highly uncertain D/H, or no useful information.
The column densities of H and D are estimated from the precise shapes of their absorption lines in the spectra. For H, the main difficulties are the accuracy of the column density and the measurement of the distribution in velocity of this H. For D the main problem is contamination by H, which we discuss below.
It is well known that column densities are harder to measure when absorption lines become saturated. The amount of absorption increases linearly with the column density as long as only a small fraction of the photons at the line central wavelength are absorbed. Lines saturate when most photons are absorbed. The amount of absorption then increases with the log of the column density.
Wampler has suggested that D/H values could be 3 – 4 times higher in Q1937–1009 than measured by Tytler, Fan & Burles . He argued that saturation of the H Lyman series lines could allow lower $`N_{HI}`$. This would lead to residual flux in the Lyman continuum, which would contradict the data, but Wampler suggested that the background subtraction might have been faulty, which was not a known problem with HIRES.
Tytler & Burles explained why Wampler’s general concerns were not applicable to the existing data on Q1937–1009. Thirteen Lyman series lines were observed and used to obtain the $`N_{HI}`$. The cross section for absorption (oscillator strength) decreases by 2000 from the Ly$`\alpha `$ to the Ly-19 line. This means that the lines vary significantly in shape, and this is readily seen in spectra with high resolution and high signal to noise. The background subtraction looked excellent because the line cores were near zero flux, as expected.
Songaila, Wampler & Cowie measured the residual flux in the Lyman continuum of the D/H absorber in Q1937–1009. They found a lower $`N_{HI}`$ and hence a higher D/H. Burles & Tytler presented a more detailed analysis of better data, and found a lower $`N_{HI}`$, consistent with that obtained from the fitting of Lyman series lines. They explained that Songaila, Wampler and Cowie had underestimated $`N_{HI}`$ because they used poor estimates of the continuum level and the flux in the Lyman continuum.
In summary, saturation does make the estimation of $`N_{HI}`$ harder. Column densities of H might be unreliable in data with low spectral resolution, or low signal to noise, and when only a few Lyman lines are observed. The above studies show that it is not a problem with the data available on Q1937–1009, Q1009+2956, Q0014+8118 and Q0130-4021. For the first two quasars, we obtain the same answer by two independent methods, and for the last three the higher order Lyman lines are not saturated.
Saturation is avoided in absorbers with lower $`N_{HI}`$, but then the D lines are weaker, and contamination by H lines becomes the dominant problem.
#### 5.4.2 Hidden Velocity Structure
To obtain D/H we need to estimate the column densities of D and H. Column densities depend on velocity distributions, and when lines are saturated, it is hard to deduce these velocity distributions. Similar line profiles are made when the velocity dispersion is increased to compensate for a decrease in the column density. We mentioned above that this degeneracy is broken when we observe lines along the Lyman series. For Q1937–1009, which has the most saturated H lines of the quasars under discussion, Burles & Tytler showed that the D/H did not change for arbitrary velocity structures, constrained only by the spectra. The same conclusion was obtained for Q1009+2956 . The favorable results for these two quasars do not mean that we will always be able to break the degeneracy. That must be determined for each absorption system.
There are two reasons why hidden velocity structure is not expected to be a major problem. First, we are concerned about hidden components which have high columns and low enough velocity dispersions that they hide inside the wider lines from lower column gas. Such gas would be seen in other lines which are not saturated: the D lines and the metal lines from ions with similar (low) ionization. Second, we search for D in absorbers with the simplest velocity distributions. They tend to have both narrow overall velocity widths and low temperatures, which makes it much harder to hide unseen components. Typically, the main component accounts for all of the absorption in the higher order Lyman lines, and these lines are too narrow for significant hidden absorption.
### 5.5 Correlated Velocity Structure: Mesoturbulence
In a series of papers, Levshakov, Kegel & Takahara , , have demonstrated a viable alternative model for the velocity distribution.
In most papers, absorption lines are modelled by Voigt profiles. The line width is the sum of the thermal broadening, turbulent broadening, and the instrumental resolution, each of which is assumed to be Gaussian. When an absorption line is more complex than a single Voigt, gas centered at other velocities is added to the model. As the signal to noise increases, we typically see that more velocity components are required to fit the absorption. Each component has its own physical parameters: central velocity, velocity dispersion (rms of thermal and turbulent broadening), ionization, column densities and elemental abundances. Prior to its use with quasars, this fitting method was developed for the ISM, where it represents gas in spatially separate clouds.
Levshakov and co-workers have proposed a different type of model, the mesoturbulent model, in which the gas velocities are correlated, and the column density per unit velocity is varied to fit the absorption line profiles. They assume that the absorption comes from a single region in space, and they calculate the distribution of the gas density down the line of sight. To simplify the calculations, in early Reverse Monte-Carlo models, they assumed that the gas temperature and density were constant along the line of sight, which is not appropriate if there are separate discrete clouds of gas with differing physical conditions.
The effects of mesoturbulence on the D/H absorbers towards Q1937–1009 , Q1009+2956 and Q1718+4807 were examined in detail using this early model. In the first paper they allowed the $`N_{HI}`$ to vary far from the observed value ($`N_{HI}`$$`=7.27\times 10^{17}`$ ), and consequently they found a variety of $`N_{HI}`$, but when the $`N_{HI}`$ is held within range, the D/H is 3.3$`\times 10^5`$, exactly the same as with the usual model . For the second quasar, the D/H obtained is again similar to that obtained in the usual way. The results are the same as with the usual model in part because the H and D line widths are dominated by thermal and not turbulent motions, and for these two quasars the total $`N_{HI}`$ is not affected, because it is measured from the Lyman continuum absorption, which does not depend on velocity.
Recently they have developed a new model called MCI , appropriate for absorption systems which sample different densities. They now use H I and metal ions to solve for two random fields which vary independently along the line of sight: the gas density and the peculiar velocities. This model allows the temperature, ionization and density to all vary along the line of sight.
The mesoturbulent model of Levshakov et al. and the microturbulent Voigt model give the same column densities and other parameters when one of the following conditions apply: 1) The line of sight through the absorbing gas traverses many correlation lengths. 2) If each velocity in a spectrum corresponds to gas at a unique spatial coordinate. 3) The absorbing regions are nearly homogeneous, with at most small fluctuations in density or peculiar velocities, or equivalently, thermal broadening larger than the turbulent broadening.
The Voigt model could give the wrong result when two or more regions along the line of sight, with differing physical conditions, give absorption at the same velocity. A remarkable and unexpected example of this was reported by Kirkman & Tytler who found a Lyman limit system which comprised five main velocity components. Each component showed both C IV and O VI absorption at about the same velocity, but in each of the five components, the O VI had a larger velocity dispersion, and hence came from different gas than the C IV. While this LLS is much more complex than those in which we can see D, this type of velocity structure could be common.
All authors other than Levshakov and collaborators use standard Voigt fitting methods to determine column densities, for several reasons. The Voigt method was used, with no well known problems, for many decades to analyze absorption in the ISM, and the ISM is well modeled by discrete clouds separated in space. The Levshakov et al. methods are more complex. In early implementations, Levshokov et al. made assumptions which are not suitable for all absorbers. The current methods require weeks of computer time, and in many cases the two methods have given the same results.
We conclude that, when we have sufficient data, velocity structure is not a problem for the absorbers like those now used for D/H.
#### 5.5.1 Was the primordial D high but depleted in the absorbers?
The idea here is that the average BBN D/H was high, and it has been depleted in the three absorbers which show low D. There are two options: local depletion in some regions of the universe, and uniformly global depletion. We conclude that there is no known way to deplete D locally, and global depletion seems unlikely.
First we list seven observations which together rule out local depletion, including that suggested by Rugers & Hogan .
1. We note that D/H is also low in our Galaxy, and that Galactic chemical evolution accounts for the difference from the low primordial D. Hence we know of four places where D is low and consistent with a single initial value.
2. If the BBN D/H was high, let us say ten times larger at 34$`\times 10^5`$, then the depletion in all four, widely separated in space, must be by a similar factor: Q1937–1009: $`0.90\pm 0.02`$; Q1009+2956: $`0.88\pm 0.02`$; Q0130-4021: $`>0.80`$; local ISM in our Galaxy: 0.86 – 0.93, where for the Galaxy alone we assume that Galactic chemical evolution reduced the initial D/H by a factor of 1.5 – 3 .
3. The quasar absorption systems are large – a few kpc along the line of sight , far larger than can be influenced by a single star or supernovae. The gas today in the local ISM is a mixture of gas which was also distributed over a similar large volume prior to Galaxy formation.
4. The abundance of the metals in the quasar cases are very low; too low for significant ($`>1`$%) destruction of D in stars .
5. The quasar absorbers are observed at high redshifts, when the universe is too young for low mass stars ($`<2`$ solar masses) to have evolved to a stage where they eject copious amounts of gas.
6. The quasar absorbers are observed at about the time when old stars in the halo of our Galaxy were forming. These stars may have formed out of gas like that seen in the quasar spectra, but with high density. We expect that much of the gas seen in absorption is in the outer halo regions of young galaxies, and that some of it was later incorporated into galaxies and halo stars.
7. The ratio of the abundances of Si/C in the quasar absorbers is similar to that in old stars in the halo of our Galaxy. This abundance ratio is understood as the result of normal chemical evolution.
Global destruction of D prior to $`z=3`$, or in the early universe, remains a possibility, but it seems contrived.
Gnedin & Ostriker discuss photons from early black holes. Sigl et al. show that this mechanism creates 10 times more <sup>3</sup>He than observed, and Jedamzik & Fuller find the density of gamma ray sources is improbably high.
Holtmann, Kawasaki & Moroi , showed that particles which decay just after BBN might create photons which could photodissociate D. With very particular parameters, the other nuclei are not changed, and it is possible to get a D/H which is lower than from SBBN with the same $`\mathrm{\Omega }_b`$. Hence low D and low $`Y_p`$ can be concordant. An exception is <sup>6</sup>Li which is produced with <sup>6</sup>Li/ H $`10^{12}`$, which is about the level observed in two halo stars. There is no conflict with the usual conclusion that most <sup>6</sup>Li is made by Galactic cosmic rays prior to star formation, because the observed <sup>6</sup>Li has been depleted by an uncertain amount. This scenario has two difficulties: Burles (private communication) notes that there would be a conflict with the $`\mathrm{\Omega }_b`$ measured in other ways, and it seems unlikely that the hypothetical particle has exactly the required parameters to change some abundances slightly, within the range of measurement uncertainty, but not catastrophically.
Most conclude that there are no likely ways to destroy or make significant D.
#### 5.5.2 Could the D/H which we observe be too high?
The answer to this question from Kirshner is, that the D/H could be slightly lower than we measure, but not by a large amount. We discuss two possibilities: measurement problems and biased sampling of the universe.
First we consider whether the D/H in the quasar absorbers could be less than observed. This can readily happen if the D is contaminated by H, but a large reduction in D/H is unlikely because the D line widths match those expected in Q1937–1009 and Q1009+2956. We do not know how the ISM D/H values could be too high, and Galactic chemical evolution requires primordial D/H to be larger than that in the ISM, and similar to the low value from quasars. Hence it is unlikely that the D/H is much below the observed value.
Second, we consider whether the absorbers seen in the quasar spectra are representative. The absorbers are biased in three ways: they represent regions of the universe with well above (100 – 1000 times) the average gas density at $`z=3`$, and amongst such high density regions, which are observed as Lyman Limit absorption systems, they have relatively low temperatures ($`2\times 10^4`$ K), and simple quiescent velocity structures. The last two factors are necessary to prevent the H absorption from covering up that from D, while the high density follows from the high density of neutral H which is needed to give detectable neutral D. It is likely that the gas in the absorbers at $`z=3`$ has by today fallen into a Galaxy, though this is not required because some gas will be heated as galaxies form, preventing infall. The low temperatures and quiescent velocities argue against violent astrophysical events, and there are no reasons to think that the absorbers are any less representative than, say, the gas which made up our Galaxy.
We should also consider whether the quasar absorbers might be unrepresentative because of inhomogeneous BBN. In this scenario regions with above average density will have below average D/H, but the evolution of density fluctuations could be such that the low density regions fill more volume , , so that they are more likely to dominate the observed universe today. In that scenario the $`\mathrm{\Omega }_b`$ derived from the D/H would be below the universal average, and the observed (low) value of D/H would be “high” compared to expectation for SBBN with the same $`\mathrm{\Omega }_b`$. This scenario will be tested when we have observations of many more quasars.
measurements of D/H towards QSOs.
#### 5.5.3 Is there spatial variation in D/H towards quasars?
It seems highly likely that the D is low in the three quasars which show low D, and we discussed above why it is hard to imagine how this D could have been depleted or created since BBN. Hence we conclude that the low D/H is primordial.
Are there other places where D is high? All quasar spectra are consistent with a single low D/H value. The cases which are also consistent with high D are readily explained by the expected H contamination. We now explain why we have enough data to show that high D must be rare, if it occurs at all.
High D should be much easier to find than low D. Since we have not found any examples which are as convincing as those of low D, high D must be very rare. If D were ten times the low value, the D line would be ten times stronger for a given $`N_{HI}`$, and could be seen in spectra with ten times lower signal to noise, or 100 times fewer photons recorded per Å. If such high D/H were common, it would have been seen many times in the high resolution, but low signal to noise, spectra taken in the 1980’s, when the community was well aware of the importance of D/H. High D would also have been seen frequently in the spectra of about 100 quasars taken with the HIRES spectrograph on the Keck telescope. In these spectra, which have relatively high signal to noise, high D could be detected in absorption systems which have 0.1 of the $`N_{HI}`$ needed to detect low D. Such absorbers are about 40 – 60 times more common than those needed to show low D/H, and hence we should have found tens of excellent examples.
#### 5.5.4 Why is there lingering uncertainty over D?
Today it is widely agreed that D is low towards a few quasars. There remains uncertainty over whether there are also cases of high D, for the following reasons:
* measurements have been made in few places;
* contamination of D by H looks very similar to D, and resembles high D;
* both the low $`Y_p`$ values reported during the last 25 years, and the <sup>7</sup>Li abundance in Spite plateau halo stars, with no correction for depletion, imply low $`\mathrm{\Omega }_b`$, low $`\eta `$, and high D/H for SBBN; and
* the first claims were for high D.
In most cases, the apparent conflicts over D/H values concern whether the absorption near the expected position of D is mostly D or mostly H. Steigman and all observational papers discussed this contamination of D by H.
Carswell et al. noted that contamination was likely in Q0014+813 and hence the D/H could be well below the upper limit. Songalia et al. stated: “because in any single instance we can not rule out the possibility of a chance H contamination at exactly the D offset, this result \[the high D/H\] should be considered as an upper limit until further observations of other systems are made.” Burles et al. showed that Q0014+813 is strongly contaminated, does not give a useful D/H limit. For Q1718+4807 we and Levshokov, Kegel & Takahara have argued that contamination is again likely.
There are many reasons why contamination is extremely common:
* H absorption looks just like that from D,
* H is 30,000 times more common,
* spectra of about 50 quasars are needed to find one example of relatively uncontaminated D,
* high signal to noise spectra are needed to determine if we are seeing H or D, and
* these spectra should cover all of the Lyman series and metal lines, because we need all possible information.
When H contaminates D, the resulting D/H will be too high.
It is essential to distinguish between upper limits and measurements. There are only two measurements (Q1937–1009 and Q1009+2956). They are measurements because we were able to show that the D absorption line has the expected width for D. All other cases are upper limits, and there is no observational reason why the D/H should be at the value of the limit. In many cases, all of the D can be H, and hence and D/H $`=0`$ is an equally good conclusion from the data.
Only about 2% of QSOs at $`z3`$ have one absorption systems simple enough to show D. All the rest give no useful information on D/H. Typically, they do not have enough H to show D, or there is no flux left at the position of D. In such cases the spectra are consistent with high, or very high, D/H, but it is incorrect to conclude that D/H could be high in $`98`$% of abosrption systems because these systems are not suitable to rule out high D/H. Rather, we should concentrate on the few systems which could rule out both high and low D/H.
We will continue to find cases like Q1718+4807 which are consistent with both low and high D/H. As we examine more QSOs we will find some cases of contamination which look exactly like D, even in the best spectra, by chance. But by that time we will have enough data to understand the statistics of contamination. We will know the distribution function of the contaminating columns and velocities, which we do not know today because the D/H absorbers are a rare and special subset of all Lyman limit absorbers. When absorbers are contaminated we will find a different D/H in each case, because the $`N_{HI}`$, velocity and width of the contaminating H are random variables. But we will be able to predict the frequency of seeing each type of contamination. If there is a single primordial D/H then we should find many quasars which all show this value, with a tail of others showing apparently more D/H, because of contamination. We will be able to predict this tail, or alternatively, to correct individual D/H for the likely level of contamination. When we attempted to correct for contamination in the past , , , we used the statistics of H I in the Ly$`\alpha `$ forest because we do not have equivalent data about the H I near to the special LLS which are simple enough to show D/H. Such data will accumulate at about the same rate as do measurements of D/H, since we can look for fake D which is shifted to the red (not blue) side of the H I.
There are large differences in the reliability and credibility of different claimed measurements of D/H in quasar spectra, and hence much is missed if all measurements are treated equally. It also takes time for the community to criticize and absorb the new results. Early claims of high D/H , in Q0014+8118 are still cited in a few recent papers, after later measurements with better data, have shown that this quasar gives no useful information, and that the high D/H came from a “spike” in the data which was unfortunately an artifact of the data reduction.
In summary, the lack of high quality spectra, which complicates assessment of contamination by H, is the main reasons why there remains uncertainty over whether some absorbers contain high D.
#### 5.5.5 Why we believe that the D/H is Primordial
Here we review why we believe that the low D/H is primordial. These arguments are best made without reference to the other nuclei made in BBN, because we wish to use the abundances of these nuclei to test SBBN theory.
* D/H is known to be low in four widely separated locations: towards three quasars, and in the ISM of our Galaxy.
* The extraction of D/H from quasar spectra is extremely direct, except for corrections for contamination by H, which make D/H look too large.
* Since contamination is common, all data are consistent with low D/H, and no data require high D/H.
* High D/H is rare, or non-existent, because it should be easy to see in many existing spectra, but we have no secure examples.
* The low D/H in the quasars, pre-solar system and in the ISM today are all consistent with Galactic chemical evolution.
* The quasar absorption systems are large – many kpc across, as was the initial volume of gas which collapsed to make our Galaxy.
* The abundance of the metals in the quasar cases are very low, and much too low for significant ($`>1`$%) destruction of D in stars.
* The quasar absorbers are observed at high redshifts, when the universe is too young for low mass stars to have evolved to a stage where they eject copious amounts of gas.
* The ratio of the abundances of Si/C in the absorbers is normal for old stars in the halo of our galaxy, indicating that these elements were made in normal stars.
* In the quasar absorbers, the temperatures and velocities are low, which argues against violent events immediately prior to the absorption.
* If BBN D/H were high, the hypothetical destruction of D would have to reduce D/H by similar large amounts in all four places.
* The above observations make local destruction of D unlikely.
* There are no known processes which can make or destroy significant D.
* Global destruction of D by photodissociation in the early universe requires very specific properties for a hypothetical particle, and is limited by other measures of $`\mathrm{\Omega }_b`$.
#### 5.5.6 Conclusions from D/H from quasars
Most agree that D is providing the most accurate $`\eta `$ value , although some have one remaining objection, that there might also be quasar absorbers which show high values of D/H , .
The D/H from our group (Burles & Tytler , , ), together with over 50 years of theoretical work and laboratory measurements of reaction rates, leads to the following values for cosmological parameters (unlike most errors quoted in this review, which are the usual $`1\sigma `$ values, the following are quoted with 95% confidence intervals):
* D/H = $`3.4\pm 0.5\times 10^5`$ (measured in quasar spectra)
* $`\eta `$ $`=5.1\pm 0.5\times 10^{10}`$ (from BBN and D/H)
* $`Y_p=0.246\pm 0.0014`$ (from BBN and D/H)
* $`{}_{}{}^{7}\mathrm{Li}/\mathrm{H}=3.5_{0.9}^{+1.1}\times 10^{10}`$ (from BBN and D/H)
* 411 photons cm<sup>-3</sup> (from the CMB temperature)
* $`\rho _b=3.6\pm 0.4\times 10^{31}\mathrm{gcm}^3`$ (from CMB and $`\eta `$ )
* $`\mathrm{\Omega }_bh^2=0.019\pm 0.0024`$ (from the critical density $`\rho _c`$)
* $`N_\nu <3.20`$ (from BBN, D/H and $`Y_p`$ data).
If we accept that D/H is the most accurate measure of $`\eta `$, then observations of the other elements have two main roles. First, they show that the BBN framework is approximately correct. Second, the differences between the observed and predicted primordial abundances teach us about subsequent astrophysical processes. Recent measurements of <sup>4</sup>He agree with the predictions. It appears that some <sup>7</sup>Li has been destroyed in halo stars , and <sup>3</sup>He is both created and destroyed in stars.
## 6 Helium
The high abundance of <sup>4</sup>He allows accurate measurements in many locations. However, <sup>4</sup>He is also produced by stars, and since such high accuracy is required, the primordial abundance is best measured in locations with the least amounts of stellar production. High accuracy is desired, since D/H predicts $`Y_p`$ to within 0.0014 ($`\delta `$$`Y_p`$/$`Y_p`$= 0.006, 95% confidence), which is well beyond the typical accuracy of astronomical abundance determinations. In the local ISM, the amount of <sup>4</sup>He from stars is about $`Y=0.010.04`$; much less than $`Y_p`$, but ten times the desired accuracy for $`Y_p`$.
Helium has been seen in the intergalactic medium, where Carbon abundances are $`<0.01`$ solar, and possibly zero in much of the volume. Strong absorption is seen from the He II Ly$`\alpha `$ line at 304Å in the redshifted spectra of quasars , however it is difficult to obtain an abundance from these measurements, because nearly all He is He III which is unobservable, and we do not know the ratio He II/He to within an order of magnitude. However, the strength of the He II absorption does mean that there is abundant He in the intergalactic gas , which has very low metal abundances, which is consistent with BBN, and probably not with a stellar origin for the <sup>4</sup>He.
The best estimates of the primordial abundance of He are from ionized gas surrounding hot young stars (H II regions) in small galaxies. The two galaxies with the lowest abundances have 1/55 and 1/43 of the solar abundance. The <sup>4</sup>He and H abundances come from the strengths of the emission lines which are excited by photons from near by hot stars.
Values for $`Y_p`$ from these extragalactic H II regions have been reported with small errors for more than 25 years, e.g.:
* $`Y_p`$$`=0.216\pm 0.02`$
* $`Y_p`$$`=0.230\pm 0.004`$
* $`Y_p`$$`=0.234\pm 0.008`$
* $`Y_p`$$`=0.236\pm 0.005`$
* $`Y_p`$$`=0.228\pm 0.005`$
* $`Y_p`$$`=0.234\pm 0.002\pm 0.005`$ . (random, and systematic errors)
* $`Y_p`$$`=0.246\pm 0.0014`$ (95% prediction from low D/H and SBBN).
These values are lower than the value now predicted by low quasar D/H and they appear incompatible, because of the small errors. However Skillman et al. argued that errors could be much larger than quoted, allowing $`Y_p`$$`<0.252`$, and Pagel (and private communication 1994) agreed this was possible.
The measurement of $`Y_p`$ involves three steps. Emission line flux ratios must be measured to high accuracy, which requires good detector linearity and flux calibration, and corrections for reddening and stellar He I absorption. These fluxes must be converted to an abundance, which requires correction for collisional ionization and neutral He. Correction for unseen neutral He depends on the spectral senegy distribution adopted for the ionizing radiation and might change $`Y_p`$ by 1 – 2 percent. Then the primordial abundance must be deduced from the $`Y`$ values in different galaxies.
Izotov, Thuan & Lipovetsky , have been pursuing a major observational program to improve the determination of $`Y_p`$. They have found many more low metallicity galaxies and have been reporting consistently higher $`Y_p`$ values, most recently in their clear and persuasive paper :
* $`Y_p`$$`=0.244\pm 0.002`$ from regression with O/H and
* $`Y_p`$$`=0.245\pm 0.001`$ from regression with N/H.
The four main reasons why these values are higher are as follows, in order of importance , , , (Skillman, and Thuan personal communication 1998).
1. When stellar He I absorption lines underlying He emission lines are not recognized, the derived $`Y_p`$ is too low. This is a important for IZw18 which has the lowest metallicity and hence great weight in the derivation of $`Y_p`$, and perhaps for many other galaxies.
2. The emission line fluxes must be corrected for collisional excitation from the metastable level. At low abundances, which correlate with high temperatures, these corrections can be several percent. The amount of correction depends on the density. There are no robust ways to measure these densities, and differing methods, used by different groups, give systematically different results. Izotov and Thuan solve for the He II density, while Olive, Skillman and Steigman use an electron density from the S II lines.
3. Izotov & Thuan have spectra which show weaker lines, and they use the five brightest He lines, while Olive et al. usually use only HeI 6678.
4. Izotov & Thuan , correct for fluorescent enhancement, which increases the $`Y`$ values from for a few galaxies.
For these reasons Izotov & Thuan obtain higher $`Y`$ values for individual galaxies which have also been observed by Olive, Skillman & Steigman , and Izotov & Thuan find a shallower slope for the regression to zero metal abundance (see Fig 6). And most importantly, using higher quality Keck telescope spectra, they obtain high $`Y_p`$= $`0.2452\pm 0.0015`$ (random errors), from the two galaxies with the lowest metal abundances .
These measurement difficulties, combined with the recent improvements, lead most to conclude that the $`Y_p`$ is in accord with the SBBN. The Izotov & Thuan values are very close to the low D/H predictions, while the lower $`Y_p`$ quoted by Olive , $`0.238\pm 0.002\pm 0.005`$, is also consistent when the systematic error is used.
It is clear that the systematic errors associated with the $`Y_p`$ estimates have often been underestimated in the past, and we propose that this is still the case, since two methods of analyzing the same Helium line fluxes give results which differ by more than the quoted systematic errors. While the Izotov & Thuan method has advantages, we do not know why the method used by Olive, Skillman & Steigman should give incorrect answers. Hence the systematic error should be larger than the differences in the results: 0.007 using the most recent values, or 0.011 using earlier results.
## 7 <sup>3</sup>He
The primordial abundance of <sup>3</sup>He has not been measured. This is most unfortunate, since it is nearly as sensitive as D to the baryon density during BBN. <sup>3</sup>He is harder to measure than D because the difference in wavelength of <sup>3</sup>He and <sup>4</sup>He lines is smaller than for D, and the Lyman series lines of He II, main absorption lines of He in the IGM, are in the far ultraviolet at 228 – 304Å which is hard to observe because of absorption in the Lyman continuum of H I at $`<912`$Å.
Rood, Steigman & Tinsley argued that it was unlikely that <sup>3</sup>He could be used to supplement cosmological information from D because low mass stars should make a lot of <sup>3</sup>He, increasing the current ISM value to well above that in the pre-solar system ISM, and in potential conflict with observations at that time. This conflict has been confirmed. Measurements do show enhanced <sup>3</sup>He in Planetary nebulae, as expected from the production in the associated low mass stars, but this is not reflected in the ISM as a whole. The pre-solar and current <sup>3</sup>He abundances are similar , in contradiction with expectation , , for unknown reasons.
It was suggested (, see review by Hata et al. ) that the uncertainty over the amount of destruction of D could be circumvented using the sum of the abundances of D + <sup>3</sup>He, since the destroyed D should become <sup>3</sup>He, and <sup>3</sup>He is relatively hard to destroy. The primordial D + <sup>3</sup>He should then be $``$ the same sum observed today, as more <sup>3</sup>He is made in stars over time. However, there are two problems with this scenario. First, the <sup>3</sup>He should increase over time, which it does not, implying that some stars destroy <sup>3</sup>He, and second, the <sup>3</sup>He abundance should be about constant in the ISM today, which it appeared not to be in early data . Hence, just prior to the measurement of D in quasars, most concluded that D + <sup>3</sup>He in the Galaxy does not provide secure cosmological information , and summaries by , .
Balser et al. report on a 14 year program to measure <sup>3</sup>He in the Galactic H II regions. Using models for the gas density structure, they find an average <sup>3</sup>He/H $`=1.6\pm 0.5`$ $`\times 10^5`$ for a sub-sample of seven simple nebulae. No variation is seen with Oxygen abundance over a factor of ten, and there is little scatter . This value may represent the average in the ISM today, but it is not known how to use this to obtain primordial abundances.
These measurements are relevant to stellar nucleosynthesis and Galactic chemical evolution, and are consistent with a cosmological origin for the <sup>3</sup>He, but we suggest that gas with much lower metal abundances will need to be observed to derive a secure primordial abundance for <sup>3</sup>He.
## 8 Lithium
Lithium is observed in the solar system, the atmospheres of a wide variety of stars and in the ISM. Arnould & Forestini review light nuclei abundances in a variety of stars and related stellar and interstellar processes, while halo stars are reviewed by , , and .
Old halo stars which formed from gas which had low iron abundances show approximately constant <sup>7</sup>Li$`/H1.6\times 10^{10}`$ and little variation with iron abundance or surface temperature from 5600 – 6300 K. The lack of variation amongst these “Spite plateau” stars (references in ) shows that their <sup>7</sup>Li is close to primordial.
Since the halo stars formed about ten times more <sup>7</sup>Li has been produced in the inner Galaxy. Abundances of <sup>7</sup>Li/H $`10^9`$ are common, although some stars show more, presumably because they make <sup>7</sup>Li. Stars typically destroy <sup>7</sup>Li when they evolve, accounting for the low abundances, $`<10^{11}`$, in evolved stars. Stars with deeper convection zones, such as halo stars with lower surface temperatures, show less <sup>7</sup>Li, because they have burnt it in their interiors.
Here, and in the next section on <sup>6</sup>Li, we will the following topics:
* measurement of current surface abundances on the Spite plateau,
* change in <sup>7</sup>Li with iron abundance,
* creation of <sup>7</sup>Li and <sup>6</sup>Li after BBN and prior to halo star formation,
* depletion of these nuclei in the halo stars,
* stars with differing <sup>7</sup>Li, and
* gravitational settling.
The recent homogeneous data on 22 halo stars with a narrow range of temperature on the “Spite plateau” have very small random errors and show that most (not all) stars with similar surface properties have the same <sup>7</sup>Li/H . Earlier data showed more scatter, which some considered real (references in ), and hence evidence of depletion.
The Ryan, Norris, & Beers sample shows a clear increase of <sup>7</sup>Li with iron abundance, as had been found earlier. This trend appears to be real, because the data and stellar atmosphere models used to derive the <sup>7</sup>Li abundance do not depend on metallicity. But it was not found by Bonifacio & Molaro , perhaps because of larger scatter in temperatures and iron abundance. This trend is not understood, and there are several possible explanations. It may have been established in the gas from which the stars formed, perhaps from cosmic rays in the ISM, or from AGB stars. Alternatively, we speculate that it might instead relate to depletion of the <sup>7</sup>Liin the stars. In either case, the BBN <sup>7</sup>Li will be different from that observed: smaller if the <sup>7</sup>Li was created prior to the star formation, and higher if the trend is connected to destruction in the stars. More on this below.
Creation of <sup>7</sup>Li in the ISM by cosmic ray spallation prior to the formation of the halo stars is limited to 10 – 20% because Be would also be enhanced by this process , .
A clear summary of arguments for and against significant depletion is given by Cayrel . There are two main reasons why depletion is believed to be small: the negligible dispersion in <sup>7</sup>Li for most halo stars on the plateau, and the presence of <sup>6</sup>Li. The main arguments for depletion are that it is expected, it clearly occurs in some stars, some halo stars on the plateau show differing abundances, and star in the globular cluster M92 which should have similar ages, composition and structure, show a factor of two range in <sup>7</sup>Li.
Different depletion mechanisms include mixing induced by rotation or gravity waves, mass loss in stellar winds and gravitational settling. Some models predict either variation from star to star, or trends with temperature, which are not seen for the stars on the plateau. For example, the rotationally induced mixing model implies that stars with different angular momentum histories will today show different <sup>7</sup>Li. Ryan, Norris & Beers find that the small scatter in their data, especially after the removal of the correlation with the iron abundance, limits the mean depletion in these models to $`<30`$%, much less than the factor of two needed to make <sup>7</sup>Li agree exactly with the predicted abundance from low D/H.
Some stars which should lie on the plateau have very low <sup>7</sup>Li, while others show a range of abundances (see ref. in ). Differences are also seen between halo field stars and stars in the globular cluster M92 , , which show a factor of two spread in <sup>7</sup>Li. These observations are not understood.
Gravitational settling (diffusion) of heavier elements reduces the <sup>7</sup>Li in the atmospheres of stars. However, the depletion should be most in the hottest (highest mass) stars, which is not seen, and not understood. Vauclair & Charbonnel proposed that small stellar winds might be balancing the settling. Vauclair & Charbonnel noted that the peak abundances inside the stars are independent of both mass and iron abundance. Normal stellar models predict that these peak abundances will not be seen in the stellar atmospheres, because convection does not reach this far down into the stars. However they point out that if some mechanism does mix gas from the <sup>7</sup>Li peak zone into the bottom of the convection zone then the stars on the plateau would have similar abundances as observed. Assuming that the observed abundances are those from the peaks inside the stars, they find that the initial abundance in the stars was <sup>7</sup>Li/H $`=2.2\pm 0.6\times 10^{10}`$, without free parameters, which is still below but statistically consistent with the prediction from low D of $`3.5_{0.9}^{+1.1}`$ $`\times 10^{10}`$.
### 8.1 Primordial <sup>7</sup>Li
Ryan, Norris and Beers conclude <sup>7</sup>Li/H $`10^{10}`$, with small random errors and three sources of systematic error, each up to a factor of 1.3, from the effective temperatures, stellar atmospheres and enhancement prior to star formation. Bonifacio & Molaro found <sup>7</sup>Li/H $`=1.73\pm 0.05\pm 0.2`$ $`\times 10^{10}`$. These abundances are both below the value of $`3.5_{0.9}^{+1.1}\times 10^{10}`$ (95%) from BBN and our D/H, but unlike , we feel they are not inconsistent given the quoted systematic errors, the lack of understanding of depletion, and the variation amongst similar stars. We do not know how to estimate the systematic errors connected with these issues. Given the comparative simplicity of D/H, we prefer to use it and SBBN, and we stick with our earlier suggestion that <sup>7</sup>Li in the Spite plateau halo stars is depleted by about a factor of two. Most, but not all agree that this is reasonable. Depletion by much larger factors, which was discussed a few years back, is now our of favor because of improved models. Improved modelling of rotational mixing, has lead to better fits to high metal abundance (population I) stars, which can be applied to halo (population II) stars, while the initial rotation rates of the halo stars may be lower than was assumed (Deliyannis private communication).
In summary, both the data and theory tells us that the <sup>7</sup>Li on the Spite plateau is not exactly the primordial value. The correction is probably small, less than a factor of two, but we do not yet know its value.
If we are to attain a primordial <sup>7</sup>Li abundance we must either (1) understand why its abundance varies from star to star, and learn to make quantitative predictions of the level of depletion, or (2) make measurements in relatively unprocessed gas.
We are optimistic that primordial <sup>7</sup>Li will be measured to high precision. Compared to D and He, the observations are simple: 15 – 20 mÅ absorption lines in relatively empty spectra of often bright stars (V=11). The best data have small errors. We anticipate that further studies will determine the amount of <sup>7</sup>Li produced prior to the formation of the stars, and the subsequent depletion in these stars. The possible increase in <sup>7</sup>Li with iron abundance is a clue, as are the <sup>6</sup>Li, Be and B abundances in the same stars.
### 8.2 <sup>6</sup>Li
The primordial <sup>6</sup>Li abundance has not been observed, but <sup>6</sup>Li/H has been measured in two stars on the Spite plateau. The abundance is well below that expected from SBBN, but <sup>7</sup>Li is used to help determine the primordial BBN <sup>6</sup>Li abundance in two ways. First, the presence of <sup>6</sup>Li limits the amount of destruction of <sup>7</sup>Li, because <sup>6</sup>Li is more fragile than <sup>7</sup>Li. Second, if the observed <sup>6</sup>Li was made prior to the formation of the stars, then some, perhaps much <sup>7</sup>Li , may have been made by similar processes. The first point is often presented as evidence that the <sup>7</sup>Li on the Spite plateau is close to primordial (e.g. less that a factor of two depletion, according to ), but the second point is cause for caution.
<sup>6</sup>Li has been detected in only two stars on the Spite plateau, because the absorption line at 6707.97Å is weak and fully blended with <sup>7</sup>Li at 6707.81Å. This is a difficult observation. The <sup>6</sup>Li makes the absorption line slightly asymmetric, and this is detected using models of the line broadening, which are tested on other absorption lines which are expected to have similar profiles because they arise in the same layers of the stellar atmosphere. Following the impressive first detection by and , and Cayrel et al. report <sup>6</sup>Li/<sup>7</sup>Li$`=0.052\pm 0.019`$ in HD84937, while Smith, Lambert & Nissen report <sup>6</sup>Li/<sup>7</sup>Li$`=0.06\pm 0.03`$ in BD+26 3578. It is not known whether these detections are representative of halo stars on the Spite plateau. Most assume that they are, but they could be above normal, perhaps by a lot; Smith et al. report <sup>6</sup>Li/<sup>7</sup>Li$`=0.00\pm 0.03`$ for six other stars.
The SBBN makes <sup>6</sup>Li/H $`10^{13.9}`$ , , using the $`\eta `$ from D/H, which is 500 times less than the measured abundance of $`7\times 10^{12}`$ in the two halo stars. The SBBN isotope ratio is <sup>6</sup>Li/<sup>7</sup>Li= 3 $`\times 10^5`$, a factor of 2000 less than observed in these two stars. This is not considered a contradiction with SBBN, because <sup>6</sup>Li, and some <sup>7</sup>Li at the same time, can be made elsewhere.
The <sup>6</sup>Li is usually assumed to have been present in the gas when the stars formed, but it could be made later, e.g. when cosmic rays strike the star or in stellar flares . Production by cosmic rays in the ISM prior to the star formation is most favored . With this assumption, the effects on <sup>7</sup>Li can be calculated in two steps. First, determine the ratio of <sup>6</sup>Li/ <sup>7</sup>Li in the production process (the production ratio). Second, correct for the depletion of <sup>6</sup>Li in the stars to determine the initial abundance of <sup>6</sup>Li. The amount of <sup>7</sup>Li produced along with the initial <sup>6</sup>Li is then specified.
Cosmic rays in the early ISM could have made <sup>6</sup>Li and some <sup>7</sup>Li prior to the formation of the Spite plateau halo stars. The production ratio depends on the reaction and energies (e.g. ). Two reactions of cosmic rays in the ISM are considered to produce <sup>6</sup>Li. Smith, Lambert & Nissen find that <sup>6</sup>Li/Be ratios imply that most <sup>6</sup>Li was made in $`\alpha \alpha `$ fusion reactions, rather than in spallation (e.g. O + p $``$ <sup>6</sup>Li) which is favored by and . The production ratio is <sup>6</sup>Li/ –<sup>7</sup>Li $`2`$ for the $`\alpha \alpha `$ reaction.
Standard stellar models predict that much of the initial <sup>6</sup>Li will have been destroyed in the stars. The more that was destroyed, the more <sup>6</sup>Li and non-BBN <sup>7</sup>Li should have been in the initial gas to give the observed abundances. Depending on the destruction mechanism, the destruction of <sup>6</sup>Li may also destroy <sup>7</sup>Li, but this is usually ignored.
When we choose the amount of depletion of <sup>6</sup>Li, we fix the amount present when the stars formed. If the <sup>6</sup>Li has been depleted by a large factor, $`100`$, then the stars would have begun with <sup>6</sup>Li/<sup>7</sup>Li similar to the production ratio, and essentially all of the <sup>7</sup>Li would be non-primordial , which is an unusual conclusion!
Ryan, Norris and Beers assume that 50% of the <sup>6</sup>Li and none of the <sup>7</sup>Li was destroyed, and use a production ratio of 1.5 to conclude that the BBN <sup>7</sup>Li was 0.84 of that now in the stars. Since nearly all observations of Li are made at low resolution, the <sup>6</sup>Li and <sup>7</sup>Li lines are not resolved, they correct for the <sup>6</sup>Li . If the two stars with observed <sup>6</sup>Li are normal, then the BBN <sup>7</sup>Li is about 79% of the observed Li absorption.
Many other papers discuss this topic. Olive & Fields give a summary. Cayrel et al. use models for the formation of Li, Be and B and calculate the expected abundance of <sup>6</sup>Li when the star formed, and find that the observed abundance implies little depletion of <sup>6</sup>Li, and a <sup>7</sup>Li depletion of less than 25%. Vangiono-Flam et al. also argue that <sup>6</sup>Li is not much depleted, and find that its BBN abundance, extrapolated back to before the production by spallation, is compatible with a BBN abundance of $`3\times 10^{13}`$$`5.6\times 10^{14}`$.
All eagerly await the measurement of <sup>6</sup>Li, together with Beryllium, in more stars.
## 9 Beryllium
The primordial abundance of Beryllium has not been observed. The production in SBBN is <sup>9</sup>Be/H $`<10^{17}`$ , , orders of magnitude below the observed level. Inhomogeneous BBN allows much higher abundances, possibly approaching detection .
Be is observed. It is created in the ISM when cosmic rays strike C, N and O nuclei, and it is destroyed in stars. It is difficult to use Be to constrain the cosmic ray production of Li because the production ratio is highly model dependent . Beryllium is observed in the atmospheres of halo stars, including those on the Spite plateau. Boesgaard et al. have found that Be increases with Iron, and that Be increases 8 times faster than Oxygen, a rate consistent with cosmic ray creation. There is some evidence for a spread in Be as a given Fe/H, but no sign of a primordial plateau, down to Be/H $`=10^{13.5}`$.
## 10 Are the different nuclei concordant or is there a crisis?
Nearly everyone believes that the primordial abundances are consistent with BBN (e.g. , , , , , ), but there are many lingering questions about the measurements. The reader will readily detect the two attitudes described by Audouze : “optimistic”, and “agnostic and perhaps heretical” in many papers. Each of us tends to adopt differing attitudes for each nucleus and astrophysical processes. This review favors D/H because it is simple and familiar.
Steigman noted that there was “a hint of an emerging crisis” because the <sup>4</sup>He abundances appeared to be lower than expected using the $`\eta `$ from the other nuclei, but he recommended much more careful study of the uncertainty in BBN predictions, chemical evolution, and observational uncertainties including systematic effects. Hata et al. and Steigman stated that “there is a conflict”, referring to the differences in $`\eta `$ implied by low D and low $`Y_p`$ values.
Whether or not there is a crisis depends on the confidence assigned to the answers to three questions:
* Is primordial D/H low everywhere, or are there also a some places with high values?
* Is $`Y_p`$ low, high, or uncertain?
* Has <sup>7</sup>Li in halo stars been depleted by a factor of two?
Some combinations of answers are not consistent with SBBN. Recent data make low D/H seem secure in three quasars plus the ISM, hence the issue is whether there are also other places with high primordial D. Low D/H is compatible with high $`Y_p`$ and depleted <sup>7</sup>Li, but not with low $`Y_p`$ or undepleted <sup>7</sup>Li. High D is compatible with low $`Y_p`$ and undepleted <sup>7</sup>Li, but it is incompatible with the three sites which show low D/H and with Galactic chemical evolution. A factor of ten D depletion would be required in all four places. Low $`Y_p`$ is compatible with undepleted <sup>7</sup>Li and high D, but is incompatible with the low D.
A good case has been made for high $`Y_p`$, explanations have been given why earlier results gave lower values, and the uncertainty appears to be larger than quoted. Hence D and $`Y_p`$ are in agreement.
The <sup>7</sup>Li observed in stars on the Spite plateau is lower than values consistent with low D. Depletion might provide an explanation, but the amount of depletion and the dominant mechanism are not known. The lack of scatter implies little depletion, less than expected, which , conclude is not sufficient to match low D. Bonifacio & Molaro find a higher <sup>7</sup>Li, but still below the level required to match low D without depletion.
## 11 Non-standard BBN
The many different forms of non-standard BBN have been reviewed by Coles & Lucchin and Jedamzik . Much work has been devoted to inhomogeneous baryon distributions during BBN, additional relativistic particles, decaying particles, large neutrino chemical potentials (e.g. ), sterile neutrinos (e.g. ), magnetic fields (e.g. ), anti-matter domains (e.g. ), and alternative theories of gravity (e.g. ).
### 11.1 Inhomogeneous BBN
Following early discussion of inhomogeneous BBN (IBBN) by Epstein & Lattimer and Hogan , many detailed studies of different types of inhomoegeneity have been published. Malaney & Mathews and Kainulainen, Kurki-Suonio, & Sihvola give reviews.
IBBN has been discussed to allow larger $`\mathrm{\Omega }_b`$ than standard BBN, to allow differing values of D/H in the universe, and to reconcile low $`Y_p`$ with low D/H values.
One exciting goal of this work was to determine whether inhomogeneity could give the observed abundances with $`\mathrm{\Omega }_b`$ much larger than the usual value, and perhaps large enough to account for all gravitating matter, without the need for non-baryonic dark matter (e.g. , , ). The best upper limit on $`\mathrm{\Omega }_b`$ comes from the lowest observed D/H, which until recently was in the ISM. In standard BBN, a higher $`\mathrm{\Omega }_b`$ is ruled out because BBN would make less than the observed ISM D/H, and no other way to make D is known. In IBBN the D/H in the ISM comes from low density regions, allowing a higher average density. The current observations, with some exceptions, fit SBBN well, and hence IBBN allows only a slight increase in $`\mathrm{\Omega }_b`$.
Inhomogenieties can be imagined over a wide range of distance scales. The smallest scales, $`<10^5`$ pc, mix prior to BBN, leaving homogeneous SBBN. Small scales mix during BBN. Intermediate scales which mix after BBN give abundances which are constant in space today, but the abundances are different from SBBN with the same $`\mathrm{\Omega }_b`$. Extra D would be made in regions with low density during BBN, giving enhanced D/H everywhere today. Large scales ($`>1`$ kpc) may have avoided mixing, and could give different D/H in different locations today. The near isotropy of the CMB limits inhomogenieties to $`<1`$ Mpc.
Jedamzik & Fuller found it difficult to match observed abundances of <sup>7</sup>Li with large scale primordial isocurvature baryon number fluctuations. Most overly dense regions of the universe with masses greater than the local baryon Jeans mass would have to collapse (to prevent observation of the <sup>7</sup>Li which is overproduced) and smaller scale fluctuations would have to be absent or suppressed. Gnedin, Ostriker & Rees and Copi, Olive & Schramm reached similar conclusions. Copi, Olive & Schramm also showed that large scale ($`>>`$1Mpc) isocurvature perturbations conflict with the smoothness of the CMB, but do not rule out inhomogeneity .
Kainulainen, Kurki-Suonio, & Sihvola review IBBN. The $`\mathrm{\Omega }_b`$ can be higher than in SBBN provided the distance scale of the baryon inhomogeneity is near to optimal to maximize neutron diffusion effects. The distance scale expected for inhomogeneities arising in the electroweak transition are too small ($`10^6`$ to $`10^3`$pc today) to have major effects, although not below the accuracy of BBN abundance calculations. QCD inhomogeneities are not so limited. However, a low D/H $`<5`$ $`\times 10^5`$ still requires $`Y_p`$$`>0.240`$ even in IBBN, which helps reconcile low D/H and low $`Y_p`$ measurements, especially when we accept that the errors on $`Y_p`$ are larger than quoted.
Rehm & Jedamzik studied BBN in the presence of anti-matter domains. Annihilation is preferentially on neutrons, and in a limiting case the resulting universe is without light nuclei, in violation of the measured abundances. With small amount of anti-matter, both the low $`Y_p`$ and low D/H measurements are matched.
Early IBBN results looked promising. Today it appears that the scales are too small to have major effects, and measurements of primordial abundances, especially upper limits on <sup>7</sup>Li, with modest depletion ($`<`$ factor of two), are usually used to give limits on the inhomogeneity, rather than to argue that inhomogeneity helps explain discordant data or allows different conclusions about $`\mathrm{\Omega }_b`$.
### 11.2 The number of Relativistic Particles and their Decays
The main idea here is that the <sup>4</sup>He abundance depends on the number of relativistic particles during BBN. Extra particles, such as neutrinos or supersymmetric particles, which are relativistic during BBN, lead to faster expansion, larger n/p and a larger $`Y_p`$.
Steigman, Schramm & Gunn calculated that BBN limited the number of families to $`N_\nu <5`$ to match the <sup>4</sup>He abundance. The range allowed by SBBN and laboratory measurements have both narrowed over the years and agree well today , . A recent update gives $`N_\nu <3.20`$ (95%) from SBBN, although a larger range is obtained if a wider variety of measured abundances are accepted . March-Russell et al. note that additional relativistic degrees of freedom are allowed if there is a large compensating asymmetry in the electron neutrino number.
Shi, Fuller & Abazajian follow the time evolution of a lepton number asymmetry arising from active – sterile neutrino transformations during BBN. For $`\nu _e`$ mixing with $`\nu _s`$, $`Y_p`$ was allowed to change from –1% to +9%, while $`\nu _\tau `$ or $`\nu _\mu `$ mixing with $`\nu _s`$ allowed –2% to +5%. Hence the $`Y_p`$ predicted by low D/H in SBBN could be lowered to 0.241, which is between the high and low measurements.
For many years past, observations suggested that $`Y_p`$ was smaller than expected for the low D/H in the ISM and now QSOs. It is hard to make $`Y_p`$ lower, since this requires fewer, not more, particles than in SBBN. Holtmann et al. , proposed decays of neutrinos, but this is nearly ruled out by the Kamiokande results on atmospheric neutrinos .
Lindley found that massive particles decaying into photons must have lifetimes in excess of a few thousand seconds, to avoid the destruction of BBN D. Audouze, Lindley & Silk noted that such radiative decays could photodisintegrate <sup>4</sup>He and make D and <sup>3</sup>He, removing the upper bound on $`\mathrm{\Omega }_b`$. Other references are given by Holtmann et al. , who discuss weakly interacting massive (100 GeV) particles which decay of order $`10^6`$ s after BBN. The authors give limits on the abundance and lifetimes of gravitinos and neutralinos, for a wide range of light nuclei primordial abundances.
Kohri & Yokoyama give limits on the mass fraction in primordial black holes with masses $`10^83\times 10^{10}`$ g which evaporate during BBN and change the abundances.
L$`\stackrel{´}{\mathrm{o}}`$pez-Su$`\stackrel{´}{\mathrm{a}}`$rez & Canal combine inhomogeneous nucleosynthesis and particles which decay at a late time to reassess the limits on $`\mathrm{\Omega }_b`$. They find parameters which allow $`\mathrm{\Omega }_b`$$`<0.130.18h_{70}^2`$ ($`h_{70}^2`$ is the Hubble constant in units of 70 km s<sup>-1</sup> Mpc<sup>-1</sup>). Such high $`\mathrm{\Omega }_b`$ might appear to remove the need for non-baryonic dark matter, but there would then be conflicts with other measures of $`\mathrm{\Omega }_b`$, especially the baryon fraction in clusters of galaxies, if all those baryons were observable today.
## 12 Cosmological Baryon Density
The measurement of the baryon density is now a highly active area of research. In the coming years, we anticipate that higher accuracy measurements of the baryon density, from the CMB, clusters of galaxies, and the Ly$`\alpha `$ forest, will give a new rigorous test of BBN . This test can be viewed from two directions. First, we can use the baryon density to fix the last free parameter in BBN, and second, we can compare the different baryon density measurements, which should be identical if SBBN is correct, and all baryons are counted in the measurements made at later times.
In addition to BBN, the baryon density is measured in four ways: in the IGM, in clusters of galaxies, using simulations of galaxy formation, and directly from the CMB. All agree with the value from SBBN using low D/H, but today they are each about an order of magnitude less accurate.
### 12.1 $`\mathrm{\Omega }_b`$ from the IGM Lyman-$`\alpha `$ forest absorption
The gas in the IGM is observed through H I Ly$`\alpha `$ absorption in the spectra of all QSOs. Gunn & Peterson discussed how redshift produces continuous absorption in the ultraviolet spectra of QSOs. Density fluctuations in the IGM trun this continuous absorption into the Ly$`\alpha `$ forest absorption lines. The IGM fills the volume of space, and at redshifts $`z>1`$ it contains most of the baryons.
The baryon density is estimated from the total amount of H I absorption, correcting for density fluctuations which change the ionization. The gas is photoionized, recombination times are faster in the denser gas, and hence this gas shows more H I absorption per unit gas. Using the observed ionizing radiation from QSOs, we have a lower limit on the ionizing flux, and hence a lower limit on the ionization of the gas. If the gas is more ionized than this, then we have underestimated the baryon density in the IGM.
Three different groups obtained similar results , , : $`\mathrm{\Omega }_b`$$`>0.035h_{70}^2`$. This seems to be a secure lower limit, but not if the IGM is less ionized than assumed, because there is more neutral gas in high density regions, and these were missing from simulations which lack resolution.
We do not have similar measurements at lower redshifts, because the space based data are not yet good enough, and the universe has expanded sufficiently that simulations are either too small in volume or lack resolution. Cen & Ostriker have shown that by today, structure formation may have heated most local baryons to temperatures of $`10^510^7`$K, which are extremely hard to detect , .
### 12.2 Clusters of Galaxies
Clusters of galaxies provide an estimate of the baryon density because most of the gas which they contain is hot and hence visible. The baryons in gas were heated up to 8 keV through fast collisions as the clusters assembled. The mass of gas in a cluster can be estimated from the observed X-ray emission, or from the scattering of CMB photons in the Sunyaev-Zel’dovich (SZ) effect. Other baryons in stars, stellar remnants and cool gas contribute about 6% to the total baryon mass.
The cosmological baryon density is obtained from the ratio of the baryonic mass to the total gravitating mass . Numerical simulations show that the value of this ratio in the clusters will be similar to the cosmological average, because the clusters are so large and massive, but slightly smaller, because shock heating makes baryons more extended than dark matter , . The total mass of a cluster, $`M_t`$, can be estimated from the velocity dispersion of the galaxies, from the X-ray emission, or from the weak lensing of background galaxies. We then use $`\mathrm{\Omega }_b`$/ $`\mathrm{\Omega }_m`$$`M_b/M_t`$. The baryon fraction in clusters in the last factor is about $`0.10h_{70}^1`$ (SZ effect: ), or $`0.050.13h_{70}^{3/2}`$ (X-ray: ), or $`0.11h_{70}^{3/2}`$ (X-ray: , ). Using $`\mathrm{\Omega }_m`$$`=0.3\pm 0.2`$ from a variety of methods , we get $`\mathrm{\Omega }_b`$$`0.03`$, with factor of two errors. These $`\mathrm{\Omega }_b`$ estimates are lower limits, since there might be additional unobserved baryons.
### 12.3 Local Dark Baryonic Matter
The baryon density estimated in the Ly$`\alpha `$ forest at $`z3`$ and in local clusters of galaxies are both similar to the that from SBBN using low D/H. This implies that there is little dark baryonic matter in the universe . This result seems conceptually secure, since there is little opportunity to remove baryons from the IGM at $`z<3`$ or to hide them in dense objects without making stars which we would see , and the clusters are believed to be representative of the contents of the universe as a whole today. However, the numerical estimates involved are not yet accurate enough to rule out a significant density (e.g. 0.5 $`\mathrm{\Omega }_b`$) of baryonic MACHOS.
### 12.4 Simulations of the formation of Galaxies
Ostriker (private communication) notes that the $`\mathrm{\Omega }_b`$ can be constrained to a factor of two of that derived from SBBN using low D/H by the requirement that these baryons make galaxies. Semi-analytic models can also address the distribution of baryons in temperature and the total required to make observed structures (Frenk & Baugh, personal communication).
### 12.5 CMB
The baryon density can be obtained from the amplitude of the fluctuations on the sky of the temperature of the CMB. The baryons in the IGM at $`z1300`$ scattered the CMB photons. The amplitude of the fluctuations is a measure of $`\mathrm{\Omega }_b`$h<sup>2</sup>, and other parameters. Published data favor large $`\mathrm{\Omega }_b`$, with large errors, however dramatic improvements are imminent, and future constraints may approach or exceed the accuracy of $`\mathrm{\Omega }_b`$ from SBBN , .
## 13 The Achievements of BBN
Standard Big Bang Nucleosynthesis (SBBN) is a major success because the theory is well understood, close connections have developed between theory and observation, and observations are becoming more reliable.
The early attempts to include physics in the mathematical model of the expanding universe lead to an understanding of the creation of the elements and the development of standard big bang theory, including the predictions of the CMB.
The general success of SBBN is based on the robustness of the theory, and the resulting predictions of the abundances of the light nuclei. The abundances of <sup>4</sup>He, <sup>7</sup>Li and D can be explained with a single value for the free parameter $`\eta `$, and the implied $`\mathrm{\Omega }_b`$ agrees with other estimates.
This agreement is used to limit physics beyond that in SBBN, including alternative theories of gravity, inhomogeneous baryon density, extra particles which were relativistic during BBN, and decays of particles after BBN. After decades owere f detailed study, no compelling major departures from SBBN have been found, and few departures are allowed.
Using SBBN predictions and measured abundances, we obtain the best estimates for the cosmological parameters $`\eta `$ and $`\mathrm{\Omega }_b`$.
The abundances of D, <sup>4</sup>He and <sup>7</sup>Li have all been measured in gas where there has been little stellar processing. In all three cases, the observed abundance are near to the primordial value remaining after SBBN. The D/H measured toward QSOs has the advantage of simplicity: D is not made after BBN, there are no known ways to destroy D in the QSO absorbers, and D/H can be extracted directly from the ultraviolet spectra, without corrections. There are now three cases of low D/H which seem secure. There remains the possibility that D/H is high in other absorbers seen towards other QSOs, but such high D must be very rare because no secure cases have been found, yet they should be an order of magnitude easier to find than the examples which show low D.
We use low D/H as the best estimator of $`\eta `$ and the baryon density. SBBN then gives predictions of the abundance of the other light nuclei. These predictions suggest that $`Y_p`$ is high, as suggested by Izotov, Thuan and collaborators. Low D also implies that <sup>7</sup>Li has been depleted by about a factor of two in the halo stars on the Spite plateau, which is more than some expect.
The high $`\mathrm{\Omega }_b`$ from SBBN plus low D/H is enough to account for about 1/8th of the gravitating matter. Hence the remaining dark matter is not baryonic, a result which was established decades ago using SBBN and D/H in the ISM.
The near coincidence in the mass densities of baryons and non-baryonic dark matter is perhaps explained if the dark matter is a supersymmetric neutralino .
At redshifts $`z3`$ the baryons are present and observed in IGM with an abundance similar to $`\mathrm{\Omega }_b`$. Hence there was no dark, or missing baryonic matter at that time. Today the same is true in clusters of galaxies. Outside clusters the baryons are mostly unseen, and they may be hard to observe if they have been heated to $`10^5`$$`10^7`$K by structure formation.
The number of free parameters in BBN has been decreasing over the years: Fermi & Terkovich gave nuclear reaction rates, the half-life of the neutron was measured, and then the number of families of neutrinos was measured. In standard BBN we are now left with one parameter, the baryon density, which is today measured with D/H using SBBN. When, in the next few years, this parameter is also measured, SBBN will have no free parameters. When free parameters can be adjusted to obtain consistency with the data, it is hard to tell if a hypothesis is correct. The agreement between SBBN theory and measurement has grown stronger over the decades, as more parameters were constrained by independent measurements, and abundance measurements improved. This is the most convincing evidence that BBN happened and has been understood.
This work was funded in part by grant G-NASA/NAG5-3237 and by NSF grants AST-9420443 and AST-9900842. We are grateful to Steve Vogt, the PI for the Keck HIRES instrument which enabled our work on D/H. Scott Burles and Kim Nollet kindly provided the figures for this paper. It is a pleasure to thank Scott Burles, Constantine Deliyannis, Carlos Frenk, George Fuller, Yuri Izotov, David Kirkman, Hannu Kurki-Suonio, Sergei Levshakov, Keith Olive, Jerry Ostriker, Evan Skillman, Gary Steigman and Trinh Xuan Thuan for suggestions and many helpful and enjoyable discussions. We thank the organizers of this meeting, Lars Bergstrom, Per Carlson & Claes Fransson for their gracious hospitality. |
warning/0001/math0001073.html | ar5iv | text | # On the Hydrodynamic Equilibrium of a Rod in a Lattice Fluid
## 1 Introduction.
Studying the coupled dynamics of granular matter of different shapes and sizes is of great interest for a range of phenomena. One example is the size segregation of particles as a result of vibrations. A typical realization is a can with nuts; upon shaking the larger (often taken to be Brazil) nuts rise to the top. Because of its wide interest the phenomenon has been considered and reconsidered and while some of the aspects are well-understood not everything has stopped surprising. If one asks for an analysis starting at the microscopic level the situation is not so satisfactory and even simple models have escaped serious mathematical handling (cf. , for further references).
In this paper we consider such a microscopic — albeit stochastic — dynamics for the motion of a large particle or rod in a lattice fluid composed of monomers. The problem of the present paper is however not quite similar to the canonical Brazil nuts scenario as we are interested in the equilibrium dynamics. In fact, as we will see, on the time scale of the motion of the rod, the monomers are in equilibrium for a reversible density profile. The rod then finds its hydrodynamic equilibrium at a vertical height where the density of the fluid is about equal to its own density. Going beyond equilibrium conditions, e.g. starting from a homogeneous density for the lattice fluid, gives rise to additional mathematical problems that we will only be touching at the end of the paper (see Section 4, Remark 2), and which will be the subject of future work.
The result of this paper can be classified under the heading: how to obtain a Markovian reduced dynamics? This problem is of course a very common one in nonequilibrium statistical mechanics where one considers the system composed of various types of degrees of freedom. The dynamics is globally defined in which the various degrees of freedom are coupled. In some circumstances and under some limit procedures one then expects that some degrees of freedom of the system effectively decouple giving rise to an autonomous (in many cases, Markovian) dynamics for a subset of degrees of freedom. In our case, it is the shaking, the speeding up of the monomer dynamics in the horizontal direction, that does the job. In this way, between any two moves of the polymer, the monomer configuration has the time to relax to its reversible measure and the polymer always sees the fluid in equilibrium.
Our main result is a mathematically rigorous proof of this dynamical decoupling between the motion of the rod and the monomer fluid when the monomer dynamics is (infinitely) speeded up (at least) in the horizontal direction (orthogonal to the motion of the polymer). In that limit of excessive horizontal shaking the reduced dynamics of the polymer becomes that of a random walker with rates directly given in terms of the equilibrium fluid density. When $`N`$ (the length of the polymer) is sufficiently big (depending on the rates for jumping up or down) the polymer finds its most probable height around its mean position of order $`\mathrm{log}N`$ with a variance of order 1. In the next section we describe the model and the result. The third section is devoted to the proofs. The final section contains an open problem and some additional remarks.
## 2 Model and Results.
### 2.1 Model.
#### 2.1.1 Configuration.
For convenience we put the system on the square lattice $`\text{ }\text{}^2`$. A point $`i=(x,y)`$ of the lattice has a ‘vertical’ coordinate $`y`$ and a ‘horizontal’ coordinate $`x`$. We also write $`i=(i_1,i_2)`$ if, in the notation of the coordinates, we want to remember the site $`i`$.
The system contains a rigid polymer (large particle, rod) whose position at time $`t`$ is denoted by $`Y_t`$. For simplicity we allow the rod to move only vertically. The horizontal coordinate is fixed (at 0) and $`Y_t`$ takes values in $`\text{ }\text{}`$ (thought of as the ‘vertical’ axis). The same results would hold if the polymer also jumps horizontally at rate 1. The polymer occupies $`N\{2,3,\mathrm{}\}`$ lattice sites. If the polymer has position $`Y_t=y`$, then it occupies the region
$$A_N(y)=\{(0,y),(1,y),\mathrm{},(N1,y)\}.$$
This region is forbidden for the monomers (fluid particles). The monomer configuration is denoted by $`\eta \{0,1\}^{\text{ }\text{}^2}`$ and we use $`\eta _t`$ to denote the random field of monomers at time $`t`$. We have that $`\eta _t(i)=0`$ if there is no monomer at site $`i`$ at time $`t`$; $`\eta _t(i)=1`$ if there is a monomer at site $`i`$ at time $`t`$. The dynamics will always be subject to the restriction (exclusion) that $`\eta _t(i)=0`$ for $`iA_N(Y_t)`$ (the rod acts as an obstacle for the fluid motion). The full configuration space is denoted by $`\mathrm{\Omega }=\{0,1\}^{\text{ }\text{}^2}\times \text{ }\text{}`$.
#### 2.1.2 Dynamics.
We now define the coupled dynamics for the polymer-monomers system. All motion is via jumping to vacant sites. There are the horizontal jumps of the monomers (which we take symmetric and at rate $`\gamma _1`$), the vertical jumps of the monomers (asymmetric at rate $`\gamma _2`$) and the vertical jumps of the rod (asymmetric at rate 1). The asymmetry in the vertical direction models the presence of an external (e.g. gravitational) field acting on fluid matter and polymer but can in general be different for monomers and polymer. Increasing the rates $`\gamma _1`$ and $`\gamma _2`$ can be used to simulate the greater mobility of the smaller particles upon shaking. We are most interested in the case where $`\gamma _21`$ and $`\gamma _11`$ (horizontal shaking).
Fig.1: A polymer between monomers
Fig.2: Rates of jumping for monomer and polymer where $`\gamma _2=\gamma `$ and $`\gamma _1=1`$.
Here comes the formal definition of the generators of these motions. Let $`f`$ be a local function on $`\mathrm{\Omega }`$ (i.e., a function that depends on the configuration in a finite region of $`\text{ }\text{}^2`$). The first part of the generator represents horizontal monomer-jumping:
$$_\mathrm{h}(\eta ,y)=\frac{1}{2}\underset{ij:i_2=j_2}{}I[ijA_N(y)=\mathrm{}][f(\eta ^{i,j},y)f(\eta ,y)]$$
(2.1)
where $`\eta ^{i,j}(k)=\eta (i)`$ if $`k=j,=\eta (j)`$ if $`k=i`$ and $`=\eta (k)`$ otherwise; the summation is over nearest neighbor pairs $`ij`$ with the same vertical coordinate ($`i_2=j_2`$). The notation $`I[]`$ will always stand for the indicator function.
Second comes the vertical monomer-jumping: with $`p<q`$,
$`_\mathrm{v}f(\eta ,y)`$ $`=`$ $`{\displaystyle \underset{i}{}}\{p\eta (i)(1\eta (i_1,i_2+1))I[(i_1,i_2+1)A_N(y)]`$ (2.2)
$`\times [f(\eta ^{i,(i_1,i_2+1)},y)f(\eta ,y)]`$
$`+q\eta (i)(1\eta (i_1,i_21))I[(i_1,i_21)A_N(y)]`$
$`\times [f(\eta ^{i,(i_1,i_21)},y)f(\eta ,y)]\}.`$
Finally, there is the polymer-jumping: with $`a,b\text{I}\text{R}^+`$,
$`_{\mathrm{poly}}f(\eta ,y)`$ $`=`$ $`aI[\eta (i)=0,iA_N(y+1)][f(\eta ,y+1)f(\eta ,y)]`$ (2.3)
$`+bI[\eta (i)=0,iA_N(y1)][f(\eta ,y1)f(\eta ,y)].`$
We will then choose $`p/q,a/b<1`$ to represent an external field in the vertical direction driving all particles, big and small, downward. E.g. in the case of a gravitational field, we could have $`p/q=\mathrm{exp}(mg/kT),a/b=\mathrm{exp}(Mg/kT)`$, where $`m,M`$ denotes the mass of a monomer, resp. polymer.
The formal generator $``$ of the full dynamics consists of three pieces:
$$=\gamma _1_\mathrm{h}+\gamma _2_\mathrm{v}+_{\mathrm{poly}},$$
(2.4)
where $`\gamma _1,\gamma _2>0`$ are additional parameters governing the rates of the monomer-jumping. Notice that $`_{\mathrm{mono}}=\gamma _1_\mathrm{h}+\gamma _2_\mathrm{v}`$ works on the configuration of monomers only (for fixed rod position), while $`_{\mathrm{poly}}`$ works on the polymer configuration (for fixed monomers). The only interaction is by excluded volume. The generator (2.4) can be rewritten in the form:
$$f(\eta ,y)=_\eta f_\eta (y)+_yf_y(\eta )$$
(2.5)
where $`f_\eta ()=f(\eta ,)`$ and $`f_y()=f(,y)`$.
#### 2.1.3 Initial distribution and extra remarks.
At time 0 (starting time) we put the polymer at the origin: $`Y_{t=0}=0`$. Then fix a real parameter $`\kappa `$ and distribute the monomers independently with density
$$\rho (i)=\frac{\kappa (p/q)^{i_2}}{1+\kappa (p/q)^{i_2}}$$
(2.6)
varying in the vertical direction (constant in the horizontal direction), conditioned on $`\eta (i)=0,iA_N(0)`$.
More precisely, we let $`\nu _\rho `$ denote the product measure on $`\{0,1\}^{\text{ }\text{}^2}`$ with density
$$\nu _\rho (\eta (i))=\rho (i),$$
(2.7)
defined by (2.6). This measure is reversible for each of the monomer generators process without polymer —i.e. the generators defined by (2.1) and (2.2) but without the indicator functions prohibiting jumps. The proof of this fact is a simple computation. The one-dimensional analogue is well known, see .
For any given $`y\text{ }\text{}`$, we write
$$\nu _\rho ^y=\nu _\rho (|\eta (i)=0,iA_N(y)),$$
(2.8)
At time 0, we put the distribution $`\mu _0`$ on $`\mathrm{\Omega }`$ defined by
$$\mu _0(d\eta ,y)=\delta _{y,0}\nu _\rho ^y(d\eta ),$$
(2.9)
where $`\delta _{y,0}`$ stands for the Kronecker-delta.
From the initial condition described above and the dynamics defined via (2.4) the process $`(\eta _t,Y_t)`$ is generated. The measure at time $`t0`$ is denoted by $`\mu _t`$. Of course this depends on the choice of parameters $`p,q,a,b,\gamma _1`$ and $`\gamma _2`$ and we will sometimes make this explicit in the notation.
A useful way to imagine the process is by associating two exponential clocks (at rate $`a`$ respectively $`b`$) to the polymer: one clock gives rise to the trial times for the polymer to jump up, the other indicates the trials for the polymer to jump down. If, just before the trial time $`\tau `$, say for jumping up, there are no monomers right above the polymer ($`\eta _\tau ^{}(i)=0,iA_N(Y_\tau ^{}+1)`$), then the jump is performed and at time $`\tau `$ the polymer is at height $`Y_\tau =Y_\tau ^{}+1`$, otherwise it stays where it was. Between the trial times of the polymer, only the monomers move. The dynamics for the monomers for a fixed position of the polymer (say at $`y`$) is generated by
$$_{\mathrm{mono}}^yf(\eta )=\gamma _1_\mathrm{h}^y+\gamma _2_\mathrm{v}^y$$
(2.10)
which can be read off from (2.1) and (2.2). The associated semigroup is denoted by $`S^y(t)`$. Now, the important thing where the ‘equilibrium’ in the title of this paper refers to, is that $`\nu _\rho ^y`$ is a reversible measure for $`S^y(t)`$. This will be proven as Lemma 3.1 in Section 3.
### 2.2 Results.
#### 2.2.1 Limiting random walk.
In the limit $`\gamma _1+\mathrm{}`$ the motion of the rod will decouple from the monomer dynamics. It will be a random walk. We first introduce this limiting rod motion.
For $`a,b\text{I}\text{R}^+`$ consider the continuous time random walk on $`\text{ }\text{}`$ with generator
$$^{\mathrm{RW}}f(y)=a[1\rho (y+1)]^N[f(y+1)f(y)]+b[1\rho (y1)]^N[f(y1)f(y)],$$
(2.11)
where the density profile $`\rho `$ is obtained from (2.6). Remark that, in the notation of (2.5),
$$^{\mathrm{RW}}f(y)=\nu _\rho ^y(d\eta )_\eta ^yf(y)$$
(2.12)
$`^{\mathrm{RW}}`$ generates a continuous time random walk $`Y_t^{\mathrm{RW}}`$ which we start at $`Y_{t=0}^{\mathrm{RW}}=0`$ and with rate for moving one step upward $`a[1\rho (y+1)]^N`$ and rate moving one step downward equal to $`b[1\rho (y1)]^N`$. We fix the initial state to be $`0`$ for the sake of definiteness. Our results hold for any other initial (deterministic or random) state.
###### Proposition 2.1
If $`a/b>(p/q)^N`$, then the random walk with generator (2.11) defined above has a unique reversible probability measure $`m`$ on $`\text{ }\text{}`$, which is given by
$$m(y)=\frac{1}{Z}\frac{(a/b)^y}{\left(1+(p/q)^y\right)^N},$$
(2.13)
where $`Z`$ is a normalizing constant. In particular, the random walk is positive recurrent.
Proof: Reversibility of $`m(y)`$ is immediate, and the condition $`a/b>(p/q)^N`$ guarantees that $`m(y)`$ can be normalized (i.e. $`Z<\mathrm{}`$). Positive recurrence follows immediately from the existence of a reversible probability measure.
Remark that the condition $`a/b>(p/q)^N`$ in the case of a gravitational field just means $`M/N<m`$, i.e. the density of the polymer is smaller than the density of the monomer-fluid. It is thus very natural that in this case the polymer will drift up and will float at a height where the fluid density is proportional to $`1/N`$, see (2.20) and .
In order to study some global properties of the limiting random walk, in particular its behavior for large $`N`$, we replace the discrete distribution $`m(y)`$ on $`\text{ }\text{}`$ by a continuous distribution:
$$m(dx):=\frac{\mathrm{exp}(\alpha x)}{(1+\mathrm{exp}(\beta x))^N}\frac{1}{Z(\alpha ,\beta ,N)}dx.$$
(2.14)
Here
$$Z(\alpha ,\beta ,N)=_{\mathrm{}}^{\mathrm{}}𝑑x\frac{\mathrm{exp}(\alpha x)}{(1+\mathrm{exp}(\beta x))^N}=\frac{1}{\beta }\frac{\mathrm{\Gamma }(\alpha /\beta )\mathrm{\Gamma }(N\alpha /\beta )}{\mathrm{\Gamma }(L)},$$
(2.15)
and $`e^\alpha =a/b`$, $`e^\beta =p/q`$. From (2.15) we can calculate the cumulants of the continuous distribution $`m(dx)`$: in particular
$$xm(dx)=\frac{1}{\beta }\left(\psi (\frac{\alpha }{\beta })\psi (L\frac{\alpha }{\beta })\right),$$
(2.16)
where $`\psi (x)=\mathrm{\Gamma }^{}(x)/\mathrm{\Gamma }(x)`$. Using the asymtotic expansion
$$\psi (z)=\mathrm{log}z\frac{1}{2z}\frac{1}{12z^2}+\mathrm{}$$
(2.17)
we obtain
$$xm(dx)=\frac{1}{\beta }\mathrm{log}N+O(1),\text{as}N\mathrm{}$$
(2.18)
and all higher order cumulants are of order 1 as $`N`$ tends to infinity. The modus of $`m`$ (the position where $`m(x)`$ reaches its maximum) is
$`\text{Mo}(m)`$ $`=`$ $`\left(\mathrm{log}(p/q)\right)^1\mathrm{log}\left({\displaystyle \frac{\mathrm{log}(p/q)^N}{\mathrm{log}(a/b)}}1\right)`$ (2.19)
$``$ $`{\displaystyle \frac{1}{\beta }}\mathrm{log}N,\text{ as}N\mathrm{}.`$
#### 2.2.2 Main result.
Our main result states
###### Theorem 2.1
Let $`0p<q<\mathrm{}`$ and $`a,b\text{I}\text{R}^+`$ and consider the joint monomer-polymer process with generator (2.4). For any finite time-interval $`K`$, the marginal law of the polymer motion $`(Y_t^{\gamma _1}:tK)`$ converges, as $`\gamma _1\mathrm{}`$, to the law of the random walk $`(Y_t^{\mathrm{RW}}:tK)`$ defined by (2.11).
#### 2.2.3 Discussion.
Since for $`a/b>(p/q)^N`$ the limiting motion is an ergodic random walk in a countable state space, the process starting from any initial distribution will converge to the (unique) invariant measure. Hence, by (2.18), the polymer will rise from the zero level to a level at height proportional to $`\frac{1}{\beta }\mathrm{log}N`$. If it starts in equilibrium, then it will perform a random walk around this position. This is exactly what we would expect from general hydrodynamics, see . After all, the fluid density at height $`\frac{1}{\beta }\mathrm{log}N`$ is precisely, cf (2.6):
$$\rho (\frac{1}{\beta }\mathrm{log}N)=\frac{\kappa /N}{1+\kappa /N}\frac{\kappa }{N}$$
(2.20)
confirming Archimedes’ law in this model of granular matter.
## 3 Proofs.
### 3.1 Outline of proof
In this section we state the main steps of the proof of Theorem 2.1. The reader may use this section as a guideline to the next section. The main idea of the proof is that in the limit $`\gamma _1\mathrm{}`$ the monomers are moving very fast in the horizontal directions and thus can reach equilibrium in the time between two successive jumps of the polymer. Therefore the rate at which the polymer jumps, which is a function of the whole monomer configuration, can be replaced by the expectation of that rate in the equilibrium distribution of the monomer configuration.
As a first step (Lemma 3.1) we identify the reversible equilibrium measure for the monomers for fixed position of the polymer. This is (by reversibility) the original reversible measure of the monomer gas without polymer, conditioned on having no monomers on the lattice sites occupied by the fixed polymer.
In a second step (Lemma 3.2-Proposition 3.1) we prove that in the limit $`\gamma _1\mathrm{}`$ any time dependent expectation of a function $`f(Y_t,\eta _t)`$ of both polymer position $`Y_t`$ and monomer gas configuration $`\eta _t`$ can be replaced by the expectation of a new function depending only on the polymer position, and obtained from $`f`$ by integrating out the $`\eta `$ variables over the equilibrium measure. The main ingredients in the proof of that statement are
1. Discrepancies in the asymmetric exclusion process move as “second class particles” which are a kind of random walkers. When $`\gamma _1\mathrm{}`$, this “random walker” diffuses away very quickly.
2. The distribution of the monomers at any jumping time of the polymer is absolutely continuous w.r.t. the monomer-fixed polymer equilibrium measure.
In the first two steps we obtain convergence of the distribution of the polymer position $`Y_t^{\gamma _1}`$ to the distribution of the random walk $`Y_t`$. To finish our proof, we still have to prove that the whole process $`\{Y_t^{\gamma _1}:t0\}`$ converges to the whole process $`\{Y_t:t0\}`$ (i.e. the distributions on trajectories converge). This final step is made by first proving that any limiting process is Markovian and next that there exists a limiting process (tightness).
### 3.2 Proof of Theorem 2.1
We start this section with an easy lemma on reversible Markov processes.
###### Lemma 3.1
Let $`\{\eta _t:t0\}`$ be a Markov process on $`\mathrm{\Omega }`$ with generator $``$ and let $`\mu `$ be a reversible measure for $``$. Suppose $`A\mathrm{\Omega }`$ such that $`\mu (A)>0`$ and such that $`1_A`$ is in the domain of the generator. Consider the process with generator
$$_Af=1_A(1_Af)(1_A1_A)f$$
(3.21)
That is, $`_A`$ corresponds to a process with “forbidden region” $`A^c`$ (i.e., jumps from $`A`$ to $`A^c`$ are suppressed). Then the measure $`\mu _A:=\mu (|A)`$ is reversible for $`_A`$.
Proof: Because the second term in the right hand side of (3.21) is just multiplication with the function $`1_A(1_A)`$, it suffices to show that $`\stackrel{~}{}_Af:=1_A(1_Af)`$ defines a symmetric operator on $`L^2(A,\mu _A)`$. Let $`f,g`$ be in the domain of $`\stackrel{~}{}_A`$. Since $`d\mu _A=(1/\mu (A))1_Ad\mu `$, we get, using the symmetry of $``$ in $`L^2(\mu )`$:
$`{\displaystyle g(\stackrel{~}{}_Af)𝑑\mu _A}`$ $`=`$ $`{\displaystyle \frac{1}{\mu (A)}}{\displaystyle 1_Ag\stackrel{~}{}_Af𝑑\mu }`$ (3.22)
$`=`$ $`{\displaystyle \frac{1}{\mu (A)}}{\displaystyle (g1_A)1_Af𝑑\mu }`$
$`=`$ $`{\displaystyle \stackrel{~}{}_Agf𝑑\mu _A}.`$
Note that reversibility is crucial in the proof of this lemma. Indeed if $`\mu `$ is only stationary, then we cannot conclude in general that $`\mu _A`$ will be stationary for the process with generator $`_A`$. Indeed, one easily computes
$$_Af𝑑\mu _A=\frac{1}{\mu (A)}(1_A^{}1_A1_A1_A)f𝑑\mu ,$$
(3.23)
i.e., $`\mu _A`$ will be stationary iff $`1_A^{}1_A1_A1_A=0`$ $`\mu `$-a.s. Since the profile measures are reversible for the exclusion process of the monomers without polymer, we can apply lemma 3.1 for $`\mu =\nu _\rho `$, $`A=\{\eta \{0,1\}^{\text{ }\text{}^2}:_{zA(y,N)}\eta (z)0\}`$, i.e., those monomer configurations which are excluded when the polymer is at vertical position $`y`$. This yields:
###### Corollary 3.1
For fixed polymer position at $`y\text{ }\text{}`$ the measure $`\nu _\rho ^y`$ is reversible for the monomer dynamics with semigroup $`S^y(t)`$.
###### Lemma 3.2
Fix $`y\text{ }\text{}`$. Let $`f`$ be a local function on $`\{0,1\}^{\text{ }\text{}^2}`$ which only depends on the monomer configuration in the layers at height $`y+1`$ and $`y1`$. Suppose that
$$\nu _\rho ^y(f)=0.$$
Then, for any $`t>0`$,
$$\underset{\gamma _1+\mathrm{}}{lim}S^y(t)f_{L^2(\nu _\rho ^y)}=0.$$
Proof: Abbreviate $`\mu :=\nu _\rho ^y`$ and consider the case $`f_y(\eta )=1_A(\eta )\mu (A)`$ for a set $`A`$ in the space of configurations depending only on a finite number of coordinates in labels $`y1`$ and $`y+1`$. The extension to general local $`f`$ is straightforward. Denote
$$D_A:=\{x\text{ }\text{}^2:1_A(\eta )1_A(\eta ^x)\text{for some}\eta \},$$
(3.24)
the dependence set of $`A`$. By reversibility:
$`{\displaystyle (S^y(t)f_y)^2𝑑\mu }`$ $`=`$ $`{\displaystyle (S^y(2t)f_y)f_y𝑑\mu }.`$ (3.25)
$`=`$ $`\mu (A)\left(\text{I}\text{E}_{\mu (|A)}^y(1_A(\eta _t))\text{I}\text{E}_\mu ^y(1_A(\eta _t))\right).`$
To compute the difference of the expectations in the above expression we realize the processes with initial configurations $`\eta `$ and $`\zeta `$ in the same probability space (coupling).
To construct this coupling we first associate two Poisson clocks to each site of $`\text{ }\text{}`$ with parameters $`\gamma _2p`$ and $`\gamma _2q`$ respectively and use them to decide the times of the vertical attempted jumps. A jump from $`(i_1,i_2)`$ to $`(i_1,i_2+1)`$ is performed at time $`t`$ if an event of the Poisson process of rate $`p`$ occurs at that time, a particle is present at $`(i_1,i_2)`$ and no particle is present at $`(i_1,i_2+1)`$ at time $`t`$. Similarly, a jump from $`(i_1,i_2)`$ to $`(i_1,i_21)`$ is performed at time $`t`$ if an event of the Poisson process of rate $`q`$ occurs at that time, a particle is present at $`(i_1,i_2)`$ and no particle is present at $`(i_1,i_21)`$ at time $`t`$. Jumps either to or from sites occupied by the rod are suppressed. This takes care of the vertical jumps. See Ferrari (1992) for details of this construction. For the horizontal jumps we associate Poisson clocks with rate $`\gamma _1`$ to pairs of horizontal nearest-neighbor sites. When the clock associated with sites $`(i_1,i_2)`$ and $`(i_1+1,i_2)`$ rings, the contents of those sites are interchanged. Also here, if at least one of the sites is occupied by the rod, the jump is suppressed. The horizontal motion is also called *stirring* process. See Arratia (1986) for details of this construction. More rigourosly, let $`(N_t(i,j):i=(i_1,i_2)\text{ }\text{},j=(i_1+1,i_2))`$, $`(N_t^+(i):i\text{ }\text{})`$ and $`(N_t^{}(i):i\text{ }\text{})`$ three independent families of independent Poisson processes of rates $`\gamma _1/2`$, $`p\gamma _2`$ and $`q\gamma _2`$, respectively —the Poisson clocks. Use the notation $`\mathrm{d}N_t()=1`$ if there is an event of the Poisson process $`(N_t())`$ at time $`t`$, otherwise it is zero. The motion is defined by
$`\mathrm{d}f(\eta _t)`$ $`=`$ $`{\displaystyle \underset{ij:i_2=j_2}{}}\mathrm{d}N_t(i,j)I[ijA_N(y)=\mathrm{}][f(\eta ^{i,j},y)f(\eta ,y)]`$ (3.26)
$`+{\displaystyle \underset{i}{}}\{\mathrm{d}N_t^+(i)\eta (i)(1\eta (i_1,i_2+1))I[(i_1,i_2+1)A_N(y)]`$
$`\times [f(\eta ^{i,(i_1,i_2+1)},y)f(\eta ,y)]`$
$`+\mathrm{d}N_t^{}(i)\eta (i)(1\eta (i_1,i_21))I[(i_1,i_21)A_N(y)]`$
$`\times [f(\eta ^{i,(i_1,i_21)},y)f(\eta ,y)]\}.`$
Standard arguments, see for instance Durrett (1993) show that (3.26) defines a process $`\eta _t=\mathrm{\Phi }(\eta _0;N[0,t])`$, with initial configuration $`\eta _0`$, where $`\mathrm{\Phi }`$ is the function induced by (3.26) and $`N[0,t]:=(N_s(,),N_s^+(),N_s^{}():\mathrm{\hspace{0.17em}0}st)`$; furthermore it is immediate to see that $`\eta _t`$ has generator $`_y`$. Given two initial configurations $`\eta `$ and $`\zeta `$, the coupling of their evolutions is constructed using the same Poisson processes: define
$$(\eta _t,\zeta _t):=(\mathrm{\Phi }(\eta ;N[0,t]),\mathrm{\Phi }(\zeta ;N[0,t])).$$
Let $`\text{I}\text{E}_{(\eta ,\zeta )}^y`$ denote expectation in the coupling starting with $`(\eta ,\zeta )`$. We need also to couple the initial configurations. Let $`\stackrel{~}{\mu }_A`$ be the law of a pair of configurations $`(\eta ,\zeta )`$ with marginal distributions $`\mu _A`$ and $`\mu `$ and such that $`\eta (x)=\zeta (x)`$ for all $`x\text{ }\text{}^2D_A`$. It is possible to construct a measure with these properties because $`\mu `$ is a product measure. We then have
$$\text{I}\text{E}_{\mu (|A)}^y(1_A(\eta _t))\text{I}\text{E}_\mu ^y(1_A(\eta _t))=\stackrel{~}{\mu }_A(d(\eta ,\zeta ))\text{I}\text{E}_{(\eta ,\zeta )}[1_A(\eta _t)1_A(\zeta _t)]$$
(3.27)
The number of initial discrepancies is finite, that is, $`_xI(\eta (x)\zeta (x))|D_A|<\mathrm{}`$. At each site $`x`$ of $`\text{ }\text{}`$ we have one of three possibilities: $`(\eta (x)\zeta (x))=0`$, no discrepancies; $`(\eta (x)\zeta (x))^+>0`$, positive discrepancies; or $`(\eta (x)\zeta (x))^{}>0`$, negative discrepancies. Following the evolution of the particles and the discrepancies we notice that if a positive discrepancy jumps over a negative one, then both discrepancies collide, giving place to a coupled particle and a hole; if a coupled particle attempts to jump to a discrepancy, the jumps occur and then the discrepancy must jump to the site previously occupied by the coupled particle. These two behaviors only occur when vertical jumps are involved. In the horizontal jumps, discrepancies and coupled particles just interchange positions according to the Poisson horizontal (stirring) clocks.
We say that there is a *first class particles* at site $`i`$ at time $`t`$ when $`\xi _t(i)=\eta _t(i)\zeta _t(i)=1`$, a *positive second class particles* when $`(\eta \zeta )_t(i)=\eta _t(i)\zeta _t(i)=1`$ and a *negative second class particles* when $`(\zeta \eta )_t(i)=\zeta _t(i)\eta _t(i)=1`$. The first class particles occupy initially those sites $`i`$ occupied by both $`\eta `$ and $`\zeta `$. Locally in time, the motion of the first class particles is the one given by generator $`^y`$ but superposed to it there is a pure birth process of first class particles: with rate
$$p(\eta \zeta )_t(i_1,i_21)(\zeta \eta )_t(i_1,i_2)$$
the second class particles at $`(i_1,i_21)`$ and $`(i_1,i_2)`$ annihilate each other and a first class particle appears at $`(i_1,i_2)`$ and an empty site appears at $`(i_1,i_21)`$. Similarly, at rate
$$q(\eta \zeta )_t(i_1,i_2+1)(\zeta \eta )_t(i_1,i_2)$$
the second class particles at $`(i_1,i_2+1)`$ and $`(i_1,i_2)`$ annihilate each other and a first class particle appears at $`(i_1,i_2)`$ and an empty site appears at $`(i_1,i_21)`$.
The marginal distribution of a second class particle between two vertical jumps (or between a jump and an annihilation) corresponds to the law of a nearest neighbor symmetric random walk —with reflection at the rod when at level $`y`$— in the horizontal direction. In the vertical direction the motion is not Markovian —it depends on the configuration of the first and second class particles at the instants of attempted jumps— and either there is an annihilation as described above or the second class particles just change horizontal line. For instance, at time $`t`$, jumps of a $`(\eta \zeta )`$ second class particle from site $`(i_1,i_2)`$ to site $`(i_1,i_2+1)`$ occur with rate
$$p(\eta \zeta )_t(i_1,i_2)(1\xi _t(i_1,i_2+1))+q(\eta \zeta )_t(i_1,i_2)\xi _t(i_1,i_2+1)$$
and similarly for the other cases. The first term corresponds to the jump over an empty site and the second one to the interchange of positions with a first class particle.
This coupling has the property
$$\text{I}\text{P}_{(\eta ,\zeta )}^y\left\{\underset{x}{}I(\eta (x)\zeta (x))\underset{x}{}I(\eta _t(x)\zeta _t(x))\right\}=1,$$
(3.28)
i.e., the number of discrepancies cannot increase.
Since by construction the discrepancies between $`\eta `$ and $`\zeta `$ are all located at $`D_A`$, we have the estimate
$$\text{I}\text{E}_{(\eta ,\zeta )}^y[1_A(\eta _t)1_A(\zeta _t)]\underset{iD_A}{}\underset{zD_A}{}\text{I}\text{P}(X^i(t)=z),$$
(3.29)
where $`X^i(t)`$ is the position of a second class particle initially at $`i`$. If at site $`i`$ there were no discrepancy we use the convention $`X_t^i\text{ }\text{}^2`$ (and hence $`z`$, for all $`zD_A`$). If particles $`i`$ and $`j`$ were discrepancies of different sign and collided before time $`t`$, we also set $`X_t^i,X_t^j\text{ }\text{}^2`$. The process $`X_t^i`$ has rate $`\gamma _1`$ to move symmetrically in the horizontal direction. If the rod were not present, we could dominate $`\text{I}\text{P}(X^i(t)=y)`$ by $`\text{I}\text{P}((X^i(t))_1=y_1)`$, where $`(X^i(t))_1`$ is the first coordinate of the walk. Since without the rod the first coordinate makes just a symmetric random walk at rate $`\gamma _1`$, that probability would be dominated by $`\gamma _1^{1/2}`$ times a constant. But with the rod we have to work a bit more. The process $`X^i(t)`$ has rate at most $`\gamma _2(p+q)`$ to move in the vertical direction. This implies that the time elapsed between the last vertical jump and $`t`$ is dominated by the minimum between an exponential time of rate $`\gamma _2(p+q)`$ and $`t`$. With this in hand it is not difficult to prove that also in this case $`\text{I}\text{P}(X^i(t)=y)`$ is bounded above by $`\gamma _1^{1/2}`$ times a(nother) constant. Here we use that $`\gamma _2`$ remains bounded when $`\gamma _1`$ goes to infinity. We conclude that for any pair $`i`$, $`y`$ in $`D_A`$:
$$\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{P}(X^i(t)=y)=\mathrm{\hspace{0.17em}0}$$
(3.30)
Therefore we conclude, combining (3.25), (3.27 ), (3.29 ) and (3.30) and the fact that $`D_A`$ is a finite set:
$`\underset{\gamma _1\mathrm{}}{lim}{\displaystyle 𝑑\mu f_yS^y(2t)f_y}=\mathrm{\hspace{0.33em}0}.`$ (3.31)
Remark: We postpone until Section 4, Remark 3, an alternative more general proof of Lemma 3.2 which works equally well for a broader class of exclusion dynamics (e.g. with speed change) provided the projection of the invariant measure on horizontal layers is ergodic for the horizontal dynamics.
We now prove an intermediate result which is important for the proof of Theorem 2.1.
###### Proposition 3.1
Let $`f_y`$ be a function depending only on the configuration values at the $`N`$ sites of $`A_N(y1)`$ or $`A_N(y+1)`$. Then we have for all $`t>0`$:
$$\underset{\gamma _1\mathrm{}}{lim\; sup}\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\left(f_{Y_t}(\eta _t)f_{Y_t}(\eta )\nu _\rho ^{Y_t}(d\eta )\right)=0.$$
(3.32)
Proof: We first want to condition on a sequence $`\underset{¯}{T}^ϵ:=(T_1^{ϵ_1},\mathrm{},T_n^{ϵ_n})`$ of marked trial jumps before $`t`$. Here $`ϵ\{1,+1\}`$ is the mark of the jump: $`+1`$ for up, $`1`$ for down. Next we consider $`\alpha _1,\mathrm{},\alpha _n\{0,1\}`$ with interpretation $`\alpha _i=1`$ if $`i`$-th marked trial jump succeeds, $`\alpha _i=0`$ if not. Given $`(T_1^{ϵ_1},\mathrm{},T_n^{ϵ_n})`$ and $`\alpha :=(\alpha _1,\mathrm{},\alpha _n)`$, we define
$$Y_k^\alpha =\underset{j=1}{\overset{k}{}}ϵ_j\alpha _j.$$
(3.33)
This corresponds to the position of the polymer at time $`T_k^{ϵ_k}`$, given succeeded and failed jumps $`(\alpha _1,\mathrm{},\alpha _k)`$. Finally we denote by $`V_p^{\alpha ,ϵ}:=V_{\underset{¯}{T}^ϵ,\alpha _p}^\alpha (\eta )`$ the event that the polymer in $`Y_{p1}^\alpha `$ can (for $`\alpha _p=1`$) or cannot (for $`\alpha _p=0`$) perform the jump to $`Y_{p1}^\alpha +ϵ_p`$. With this notation, we can write
$`\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}(f_{Y_t}(\eta _t){\displaystyle }f(\eta )\nu _\rho ^{Y_t}(d\eta )|T_1^{ϵ_1},\mathrm{},T_n^{ϵ_n};T_n^{ϵ_n}<t<T_{n+1}^{ϵ_{n+1}})`$
$`=`$ $`{\displaystyle \underset{\alpha \{0,1\}^{\{1,\mathrm{},n\}}}{}}\text{I}\text{P}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}(\alpha )`$
$`\times {\displaystyle }d\mu _{\alpha ,Y_n^\alpha }^{(\gamma _1,\gamma _2)}S_{(\gamma _1,\gamma _2)}^{Y_n^\alpha }(tT_n^{ϵ_n})(f_{Y_n^\alpha }{\displaystyle }f_{Y_n^\alpha }(\eta )\nu _\rho ^{Y_n^\alpha }(d\eta )).`$
Here $`\mu _{\alpha ,Y_n^\alpha }^{(\gamma _1,\gamma _2)}`$ denotes the monomer distribution at time $`s=(T_n^{ϵ_n})^+`$, given the successes $`(\alpha _1,\mathrm{},\alpha _n)`$, and $`\text{I}\text{P}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}(\alpha )`$ denotes the probability of the sequence of succeeded and failed jumps prescribed by $`\alpha `$ at the times $`\underset{¯}{T}^ϵ`$. The crucial thing to realize at this point is that the probability measure $`\mu _{\alpha ,Y_n^\alpha }^{(\gamma _1,\gamma _2)}`$ is absolutely continuous with respect to the conditioned Bernoulli measure $`\nu _\rho ^{Y_n^\alpha }`$. In Lemma (3.3) below we shall give a uniform bound on the density
$$\mathrm{\Psi }_{\alpha ,n}^{(\gamma _1,\gamma _2)}:=\frac{d\mu _{\alpha ,Y_n^\alpha }^{(\gamma _1,\gamma _2)}}{d\nu _\rho ^{Y_n^\alpha }}.$$
(3.35)
By dominated convergence, the proof of the proposition is reduced to showing that for any $`\alpha \{0,1\}^{\{1,\mathrm{},n\}}`$ and any $`\delta >0`$:
$$\underset{\gamma _1\mathrm{}}{lim}𝑑\mu _{\alpha ,Y_n^\alpha }^{(\gamma _1,\gamma _2)}S_{(\gamma _1,\gamma _2)}^{Y_n^\alpha }(\delta )\left(f_{Y_n^\alpha }f(\eta )\nu _\rho ^{Y_n^\alpha }(d\eta )\right)=0.$$
(3.36)
The expression inside the limit in the left hand side of (3.36) is bounded by
$$\mathrm{\Psi }_{\alpha ,n}^{(\gamma _1,\gamma _2)}_{\mathrm{}}S_{(\gamma _1,\gamma _2)}^{Y_n^\alpha }(\delta )\left(ff(\eta )\nu _\rho ^{Y_n^\alpha }(d\eta )\right)_{L^2\left(\nu _\rho ^{Y_n^\alpha }\right)}.$$
(3.37)
Therefore, (3.36) is a consequence of Lemma 3.2 and the following estimate on the density $`\mathrm{\Psi }_{\alpha ,n}^{(\gamma _1,\gamma _2)}`$.
###### Lemma 3.3
Put $`c(\rho ,x):=[\rho (x+1)\rho (x1)(1\rho (x))]^N`$. For any $`\alpha \{0,1\}^{\text{I}\text{N}}`$ and for any $`n\text{I}\text{N}`$, we have the estimate:
$$\underset{\gamma _1\mathrm{}}{lim\; sup}\mathrm{\Psi }_{\alpha ,n}^{(\gamma _1,\gamma _2)}\underset{p=0}{\overset{n1}{}}c(\rho ,Y_p^\alpha )$$
(3.38)
Proof: We fix $`\alpha `$ and proceed by induction in $`n`$. First put $`n=1`$. By stationarity of $`\nu _\rho ^0`$ under the evolution $`S_{(\gamma _1,\gamma _2)}^0`$, we have
$$\mu _{\alpha _1,Y_1^{\alpha _1}}^{(\gamma _1,\gamma _2)}=\nu _\rho ^0[|V_1^\alpha ].$$
(3.39)
First consider $`\alpha _1=1`$, i.e., the jump succeeds. Denote $`V(x)`$ the event that the set $`A_N(x)`$ contains no monomers. Then we can write:
$$f(\eta )\nu _\rho ^0[d\eta |V_1^{\alpha _1}]=𝑑\nu _\rho ^{Y_1^{\alpha _1}}[fI(V(0))]\frac{\nu _\rho (V(Y_1^{\alpha _1}))}{\nu _\rho (V(0))}.$$
(3.40)
Hence, we conclude
$$\mathrm{\Psi }_{\alpha _1}^{(\gamma _1,\gamma _2)}=I(V(0))\frac{\nu _\rho (V(Y_1^{\alpha _1}))}{\nu _\rho (V(0))}.$$
(3.41)
And we can estimate
$$\mathrm{\Psi }_{\alpha _1}^{(\gamma _1,\gamma _2)}\frac{1}{\nu _\rho (V(0))}c(\rho ,0).$$
(3.42)
Next consider $`\alpha _1=0`$, i.e., the jump fails (and thus $`Y_1^{\alpha _1}=0`$). We write
$$f(\eta )\nu _\rho ^0[d\eta |V_1^{\alpha _1}]=\nu _\rho ^{Y_1^{\alpha _1}}(d\eta )[fI(V_1^{\alpha _1}(\eta ))]\frac{1}{\nu _\rho ^{Y_1^{\alpha _1}}(V_1^{\alpha _1})}.$$
(3.43)
Hence,
$$\mathrm{\Psi }_{\alpha _1}^{(\gamma _1,\gamma _2)}=\frac{I[V_1^{\alpha _1}]}{\nu _\rho ^0(V_1^{\alpha _1})}$$
(3.44)
So also in that case we have the estimate
$$\mathrm{\Psi }_{\alpha _1}^{(\gamma _1,\gamma _2)}\frac{1}{\nu _\rho ^0(V_1^{\alpha _1})}c(\rho ,0).$$
(3.45)
This proves the claim for $`n=1`$. Suppose the claim is true for $`n=1,\mathrm{},p1`$. Put $`\alpha _p=1`$, the case $`\alpha _p=0`$ can be treated analogously. In order to simplify the notation, we make some further abbreviations:
1. $`\mu _{\alpha ,Y_p^\alpha }^{(\gamma _1,\gamma _2)}:=\mu _p`$.
2. $`\nu _\rho ^{Y_p^{\alpha _1,\mathrm{},\alpha _p}}:=\nu _\rho ^p`$
3. $`\mathrm{\Psi }_{\alpha _1,\mathrm{},\alpha _p}^{(\gamma _1,\gamma _2)}:=\mathrm{\Psi }_p^{(\gamma _1,\gamma _2)}`$
4. $`S_{(\gamma _1,\gamma _2)}^{Y_p^{\alpha _1,\mathrm{},\alpha _p}}(t):=S_p^{(\gamma _1,\gamma _2)}(t)`$
5. $`T_p^{ϵ_p}T_{p1}^{ϵ_{p1}}:=\tau _p`$
6. $`V_p^{\alpha _1,\mathrm{},\alpha _p}:=V_p`$
We compute $`\mathrm{\Psi }_p^{(\gamma _1,\gamma _2)}`$:
$`\mu _p(f)`$ $`=`$ $`\left(\mu _{p1}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)\right)[f|V_p]`$ (3.46)
$`=`$ $`{\displaystyle \frac{𝑑\mu _{p1}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(f1_{V_p})}{𝑑\mu _{p1}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(1_{V_p})}}`$
$`=`$ $`{\displaystyle \frac{𝑑\nu _\rho ^{p1}\left(S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})f1_{V_p}\right)}{𝑑\nu _\rho ^{p1}\left(S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})1_{V_p}\right)}}`$
$`=`$ $`{\displaystyle \frac{𝑑\nu _\rho ^p\left(1_{V_{p1}}fS_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})\right)}{𝑑\nu _\rho ^p\left(1_{V_{p1}}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})\right)}},`$
where in the third step we used reversibility of $`\nu _\rho ^{p1}`$. From (3.46) we read off the density:
$$\mathrm{\Psi }_p^{(\gamma _1,\gamma _2)}=\frac{1_{V_{p1}}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})}{𝑑\nu _\rho ^p\left(1_{V_{p1}}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)}\right)}.$$
(3.47)
We first estimate the nominator of the rhs of (3.47):
$`{\displaystyle 𝑑\nu _\rho ^p1_{V_{p1}}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})}`$ (3.48)
$`=`$ $`{\displaystyle 𝑑\nu _\rho ^{p1}\left(1_{V_p}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)})\right)\frac{\nu _\rho (V_{p1})}{\nu _\rho (V_p)}}`$
$`=`$ $`{\displaystyle 𝑑\nu _\rho ^{p1}\left(S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)(1_{V_p})\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)}\right)\frac{\nu _\rho (V_{p1})}{\nu _\rho (V_p)}}`$
$``$ $`\nu _\rho (V_{p1})`$
$``$ $`{\displaystyle \frac{\nu _\rho (V_{p1})}{\nu _\rho (V_p)}}\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)}_{\mathrm{}}S_{p1}^{(\gamma _1,\gamma _2)}(\tau _p)[1_{V_p}\nu _\rho ^{p1}(V_p)]_{L^2(\nu _\rho ^{p1})}`$
$``$ $`{\displaystyle \frac{1}{c(\rho ,Y_{p1}^\alpha )}}o(\gamma _1),`$
where $`o(\gamma _1)`$ tends to zero as $`\gamma _1\mathrm{}`$ by Lemma 3.2. By the induction hypothesis, we obtain from (3.47), (3.48):
$`\underset{\gamma _1\mathrm{}}{lim\; sup}\mathrm{\Psi }_p^{(\gamma _1,\gamma _2)}_{\mathrm{}}`$ $``$ $`\underset{\gamma _1\mathrm{}}{lim\; sup}\mathrm{\Psi }_{p1}^{(\gamma _1,\gamma _2)}_{\mathrm{}}c(\rho ,Y_{p1}^\alpha )`$ (3.49)
$``$ $`{\displaystyle \underset{k=0}{\overset{p1}{}}}c(\rho ,Y_k^\alpha ).`$
This finishes the proof of Lemma 3.3 and Proposition 3.1. As a first application we obtain convergence of the one-point marginales of the processes $`\{Y_t^{(\gamma _1,\gamma _2)}:t0\}`$. For $`f:\text{ }\text{}\text{I}\text{R}`$ a bounded function, we have, using the notation of (2.5).
$$\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\left(f(Y_t)f(Y_0)_0^t𝑑s(_{\eta _s}f)(Y_s)\right)=0.$$
(3.50)
By Proposition 3.1 we obtain in the limit $`\gamma _1\mathrm{}`$:
$`\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\left(f(Y_t)f(Y_0){\displaystyle _0^t}𝑑s[(_\eta )\nu _\rho ^{Y_s}(d\eta )]f(Y_s)\right)`$ (3.51)
$`=`$ $`\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\left(f(Y_t)f(Y_0){\displaystyle _0^t}𝑑s(^{\mathrm{RW}}f)(Y_s)\right)=0.`$
This implies in particular that
$$\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\nu _\rho ^x\times \delta _x}^{(\gamma _1,\gamma _2)}f(Y_t)=\text{I}\text{E}_x^{\mathrm{RW}}f(Y_t^{\mathrm{RW}}).$$
(3.52)
In order to prove that the processes $`\{Y_t^{(\gamma _1,\gamma _2)}:t0\}`$ converge weakly in the Skorohod space of trajectories to the random walk $`\{Y_t^{\mathrm{RW}}:t0\}`$, i.e. the content of Theorem 2.1, it suffices to show that the process $`\{Y_t^{(\gamma _1,\gamma _2)}:t0\}`$ is asymptotically Markovian. Indeed, then it is uniquely determined by its single time distributions which are those of the random walk $`\{Y_t^{\mathrm{RW}}:t0\}`$. More precisely it is sufficient to prove the following lemma:
###### Lemma 3.4
Let $`\{_t:t0\}`$ denote the $`\sigma `$-field generated by $`\{(\eta _s,Y_s):0st\}`$. We have
$$\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}|\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}(f(Y_t)|_s)\underset{y}{}p_{ts}^{\mathrm{RW}}(Y_s,y)f(y)|=0$$
(3.53)
Proof: By the Markov property of the process $`\{(\eta _t,Y_t):t0\}`$,
$`\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\left(f(Y_t)|_s\right)`$ $`=`$ $`\text{I}\text{E}_{\eta _s\times \delta _{Y_s}}^{(\gamma _1,\gamma _2)}(f(Y_{ts}))`$ (3.54)
$`=`$ $`f(Y_s)+\text{I}\text{E}_{\eta _s\times \delta _{Y_s}}^{(\gamma _1,\gamma _2)}{\displaystyle _0^{ts}}_{\eta _r}f(Y_r)𝑑r.`$
Therefore, it suffices to show that
$$\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\eta _s\times \delta _{Y_s}}^{(\gamma _1,\gamma _2)}\left(_0^{ts}_{\eta _r}f(Y_r)𝑑r_0^{ts}𝑑r_\eta f(Y_r)\nu _\rho ^{Y_r}(d\eta )\right)=0$$
(3.55)
Since the trial jumps of the polymer are on the event times of a Poisson process with rate independent of $`(\gamma _1,\gamma _2)`$, we can write
$$_0^{ts}_{\eta _r}f(Y_r)𝑑r=_0^{ts}\frac{1}{ϵ}_r^{r+ϵ}_{\eta _r^{}}f(Y_r^{})𝑑r^{}+o(ϵ),$$
(3.56)
where $`o(ϵ)`$ goes to zero in $`L^2(\text{I}\text{P}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)})`$, uniformly in $`(\gamma _1,\gamma _2)`$, when $`ϵ`$ tends to zero. Therefore, it is sufficient to show that
$$\underset{\gamma _1\mathrm{}}{lim}\left(\text{I}\text{E}_{\nu _\rho ^0\times \delta _0}^{(\gamma _1,\gamma _2)}\text{I}\text{E}_{\eta _s\times \delta _{Y_s}}^{(\gamma _1,\gamma _2)}\left|\frac{1}{ϵ}_0^ϵf_{Y_r}(\eta _r)𝑑r\nu _\rho ^{Y_s}(d\eta )f_{Y_s}(\eta )\right|\right)=0$$
(3.57)
Following the same strategy as in the proof of Proposition 3.1, i.e., by estimates on the density of the monomer distribution with respect to the appropriate conditioned Bernoulli measure, this reduces to showing that for any $`ϵ>0`$, for any $`y\text{ }\text{}`$ and for $`f_y`$ depending on layer $`y+1`$ or $`y1`$:
$$\underset{\gamma _1\mathrm{}}{lim}\text{I}\text{E}_{\nu _\rho ^y}^{(\gamma _1,\gamma _2),y}\left(\frac{1}{ϵ}_0^ϵ𝑑sf_y(\eta _s)\nu _\rho ^y(d\eta )f_y(\eta )\right)^2=0.$$
(3.58)
Putting $`\stackrel{~}{f}_y:=f_y\nu _\rho ^y(f)`$, the expression inside the limit in the left hand side of (3.58) can be rewritten as
$`{\displaystyle \nu _\rho ^y(d\eta )\frac{1}{ϵ^2}_0^ϵ𝑑s_0^ϵ𝑑r\stackrel{~}{f}_yS_{(\gamma _1,\gamma _2)}^y(|rs|)\stackrel{~}{f}_y}`$ (3.59)
$``$ $`{\displaystyle \frac{1}{ϵ^2}}{\displaystyle _0^ϵ}𝑑s{\displaystyle _0^ϵ}𝑑r\stackrel{~}{f}_y_{L^2}S_{(\gamma _1,\gamma _2)}^y(|rs|)\stackrel{~}{f}_y_{L^2}.`$
Hence we obtain (3.58) as an application of Lemma 3.2. Arrived at this point, we know that any weak limit point of the processes $`\{Y_t^{(\gamma _1,\gamma _2)}:t0\}`$ equals in distribution the random walk $`\{Y_t^{\mathrm{RW}}:t0\}`$. Hence, to finish the proof of Theorem 2.1, it is sufficient to see that such a weak limit point actually exists. This is an easy task:
###### Lemma 3.5
The sequence of processes $`\{Y_t^{(\gamma _1,\gamma _2)}:t[0,T],\}_{(\gamma _1,\gamma _2)}`$ is tight.
Proof: Since the number of jumps the polymer makes in \[0,T\] is bounded by a mean one Poisson process, we have
$$\text{I}\text{P}(\underset{0tT}{sup}|Y_s^{(\gamma _1,\gamma _2)}|M)\frac{2T}{M},$$
(3.60)
and also
$$\underset{\delta 0}{lim}\text{I}\text{P}(\underset{s,t[0,T],|st|\delta }{sup}|Y_s^{(\gamma _1,\gamma _2)}Y_t^{(\gamma _1,\gamma _2)}|>ϵ)=0.$$
(3.61)
This proves tightness (cf. Theorem 1.3 p.51 of ).
## 4 Additional remarks
Remark 1: What happens when the system is out of equilibrium? For instance, start the monomers in a homogeneous product measure. When the density is constant and equal to $`\rho [0,1]`$ (no $`p,q,i_2`$ dependence in (2.6), the measure $`\nu _\rho ^y`$ is no longer invariant for the monomer dynamics $`S^y(t)`$ at fixed rod position $`y`$. However in the limit $`\gamma _1\mathrm{}`$ the polymer will perform a continuous time random walk with rates $`a(1\rho )^N`$ and $`b(1\rho )^N`$ for up and down jumps respectively. Significant corrections in the case $`\gamma _1<\mathrm{}`$ can be expected, cf. .
Another problem is obtained if we start the monomers from a sharp density profile. That is, above the polymer the fluid density is constant $`\rho _1`$ and under the polymer the density is also homogeneous equal to $`\rho _2`$. In this case the vertical density will follow a discrete space noiseless Burgers equation: $`\rho (i,t)[0,1]`$, $`t\text{I}\text{R}`$, $`i\text{ }\text{}`$
$`{\displaystyle \frac{\rho (i,t)}{t}}`$ $`=`$ $`p\rho (i,t)(1\rho (i+1,t))q\rho (i,t)(1\rho (i1,t))`$
$`+p\rho (i1,t)(1\rho (i,t))+q\rho (i+1,t)(1\rho (i,t))`$
with initial condition $`\rho (i,0)=\rho _2I(i0)+\rho _1I(i>0)`$. The limiting motion of the rod will be a non-homogeneous (in time) Markov process described by
$`{\displaystyle \frac{d\text{I}\text{E}(f(Y_t)|_t)}{dt}}`$ $`=`$ $`a[1\rho (Y_t+1,t)]^N[f(Y_t+1)f(Y_t)]`$
$`+b[1\rho (Y_t1,t)]^N[f(Y_t1)f(Y_t)].`$
where $`_t`$ is the sigma field generated by $`\{Y_s:st\}`$. These results can be obtained with the techniques we used to prove Theorem 2.1 and will be the content of a future publication, cf. .
Remark 2: One may wonder how general the results are. As an illustration of this we consider the following somewhat abstract modification of Lemma 3.2. Suppose that $`\mu `$ is a reversible measure on $`\{0,1\}^{\text{ }\text{}^2}`$ both for a monomer dynamics with generator $`_1`$ and one with generator $`_2`$. As an example, we could keep in mind the case where $`_{12}=_1+_2`$ is a Kawasaki dynamics (exclusion process with speed change) at finite temperature with $`_1`$ generating the horizontal and $`_2`$ generating the vertical jumps; $`\mu `$ is the corresponding Gibbs measure. The measure $`\mu `$ is then also reversible for $`_{12}^\gamma =\gamma _1+_2`$. Now we insert the polymer and we fix it at some position $`y\text{ }\text{}`$. The dynamics of the monomers is now conditioned on having no monomers in the excluded volume $`A_N(y)`$: $`\eta _t(i)=0,iA_N(y),t0`$. The new generator is $`_{12}^{y,\gamma }=\gamma _1^y+_2^y`$ obtained by setting all of the original rates equal to zero for all updating that would create a monomer in the region $`A_N(y)`$ (the direct analogue of what was done in (2.1) and (2.2)). It follows then from Lemma 3.1 that $`\mu ^y=\mu (|\eta (i)=0,iA_N(y))`$ is reversible for $`_{12}^y`$. We finally denote by $`\mu _x^y`$ the restriction of $`\mu ^y`$ to the layer at height $`x`$ (i.e., the set $`\{i\text{ }\text{}^2,i_2=x\}`$). This measure is reversible for $`_1^y`$. We have the following result:
###### Proposition 4.1
Denote by $`S_y^\gamma (t)`$ the semigroup with generator $`_{12}^{y,\gamma }`$. Assume that for all $`x`$ $`\mu _x^y`$ is ergodic for $`_1^y`$. Let $`f_x`$ be a function in $`L^2(\mu ^y)`$ with dependence set on layer $`xy`$. We have:
$$\underset{\gamma \mathrm{}}{lim}S_y^\gamma (t)f_x𝑑\mu _x^yf_x_{L^2(\mu ^y)}=0.$$
(4.62)
Proof: By ergodicity $`_1^y`$ has simple eigenvalue $`0`$ with corresponding eigenspace the constant functions. Hence by the spectral theorem,
$$(𝑑\mu _x^yf_x)^2=\text{I}\text{E}_{f_x,f_x}^{_1^y}(\{0\}),$$
(4.63)
where $`\text{I}\text{E}_{f_x,f_x}^{_1^y}`$ denotes the spectral measure of the selfadjoint operator $`_1^y`$. Therefore, we have to show that if $`f_x`$ is a function on layer $`x`$ such that
$$\text{I}\text{E}_{f_x,f_x}^{_1^y}(\{0\})=0,$$
(4.64)
then
$$\underset{\gamma \mathrm{}}{lim}S_y^\gamma (t)_{L^2(\mu ^y)}=0.$$
(4.65)
For every $`\phi `$ in the domain of $`_{12}^{y,\gamma }`$,
$$\underset{\gamma \mathrm{}}{lim}\frac{1}{\gamma }(_{12}^{y,\gamma }\phi )=_1^y\phi .$$
(4.66)
Hence the spectral measures $`\text{I}\text{E}_{f_x,f_x}^{\frac{1}{\gamma }_{12}^{y,\gamma }}`$ converges weakly to the spectral measure $`\text{I}\text{E}_{f_x,f_x}^{_1^y}`$. Therefore, we can estimate
$`S_y^\gamma (t)f_x_{L^2(\mu ^y)}^2`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}e^{\gamma t\lambda }\text{I}\text{E}_{f_x,f_x}^{\frac{1}{\gamma }_{12}^{y,\gamma }}(d\lambda )`$ (4.67)
$``$ $`{\displaystyle _0^\delta }e^{\gamma t\lambda }\text{I}\text{E}_{f_x,f_x}^{\frac{1}{\gamma }_{12}^{y,\gamma }}(d\lambda )+e^{\gamma t\delta }f_x_{L^2(\mu ^y)}`$
$``$ $`\text{I}\text{E}_{f_x,f_x}^{\frac{1}{\gamma }_{12}^{y,\gamma }}([0,\delta ])+e^{\gamma t\delta }f_x_{L^2(\mu ^y)}.`$
Letting $`\gamma `$ tend to infinity, and then $`\delta `$ to zero, using (4.64), we obtain (4.65).
Remark 3: Proposition 4.1 is general but has a strong hypothesis: the ergodicity of the one-layer horizontal dynamics. This is known only in a few cases, in particular in the symmetric simple exclusion process we treated in Lemma 3.2. It is however expected to be true at least for high temperature Kawasaki dynamics. |
warning/0001/quant-ph0001013.html | ar5iv | text | # Untitled Document
EFFECTS OF DICKE SUPERRADIANCE IN THE
CONTEXT OF THE ONE-ATOM MASER
N. Nayak<sup>a</sup>, A. S. Majumdar<sup>b</sup>, and V. Bartzis<sup>c</sup>
<sup>a&b</sup>S. N. Bose National Centre for Basic Sciences,
Block-JD, Sector-3, Salt Lake City, Calcutta-700091, India
<sup>c</sup>General Department of Physics, Chemistry and Material Technology,
Technological Educational Institutions of Athens, Egaleo 12210, Greece
Abstract
We consider a micromaser model to study the influence of Dicke superradiance in the context of the one-atom maser. The model involves a microwave cavity into which two-level Rydberg atoms are pumped in pairs. We consider a random pump mechanism which allows the presence of at most one pair of atoms in the cavity at any time. We analyze the differences between the present system, called the Dicke micromaser, and an equivalently pumped conventional one-atom micromaser. These differences are attributed to the Dicke cooperativity in the two-atom system. We also show that the two-atom Dicke micromaser is equivalent to a one-atom cascade two-photon micromaser. With the introduction of a one-photon detuning, the present theory further describes a true two-photon micromaser. We discuss in detail the role of one-photon detuning in the mechanism of a one-atom two-photon micromaser. This leads us to point out that the two-atom cavity dynamics can be verified by a proper scaling of the results from an equivalent one-atom two-photon micromaser.
PACS Nos : 42.50.Dv, 42.52.+x, 32.80.-t
1. Introduction
The first experimental observation of maser action with just one atom in the so-called micromaser renewed interest in the subject extensively as it opened possibilities of observing cavity-QED effects which are purely of quantum mechanical origin. Indeed, one of the results of the much discussed Jaynes-Cummings interactions , the quantum revival of the oscillations in the atomic population, has been observed in the micromaser device . The micromaser device consists of a high-Q superconducting cavity cooled down to sub-Kelvin temperatures into which Rydberg atoms in the upper of the two masing levels are pumped randomly but one at a time. A clever velocity selector decides a fixed flight time $`\tau `$ through the cavity for every atom. The repetition time $`\overline{t}_R`$, $`\tau `$ and $`t_{cav}`$, the duration in which the cavity is empty of atom, satisfy $`\overline{t}_R=\tau +t_{cav}`$ with $`\tau t_{cav}`$. This is the basis of the one-atom maser theories \[4-6\] with the further assumption of strictly single-atom events.
However, the above arrangements in the experiments cannot eliminate overlap of flight times through the cavity of successive atoms completely, which of course, is less than $`1\%`$. This has been one of the obstacles in realizing certain predictions of one-atom maser theories such as number states of the cavity radiation field. Another limiting factor is the dissipative mechanism acompanying the coherent atom-field interaction . The results in show that the effects of the dissipative mechanism can be minimized by increasing the $`Q`$ factor of the cavity and by decreasing its temperature. But the pump mechanism has not been able to eliminate traces of two atom events in one-atom maser dynamics. Hence it has become necessary to explain the effect of such two-atom events on the photon statistics of the cavity field. Various models incorporating multiple atomic events have been proposed \[7–13\]. The basic difference between such models and one-atom maser theory \[4-6\] is the Dicke atomic cooperative effect acompanying the multiple atomic events. Thus, in order to understand the limitations of maser theories based on strictly single atom events \[4-6\], it is necessary to quantify the effect of atomic cooperation. For this purpose we propose the following model.
We consider a fixed number of atoms entering the cavity with a constant flight time through it, but with random time gaps between successive such events. The cooperative nature of the interaction of the atoms with the cavity field shall be fully reflected in the photon statistics. In this paper we consider atoms being pumped in pairs into the cavity. We limit the number of atoms to two just for the sake of simplicity which can be easily generalized to a larger number. Both the atoms are taken to be in their upper states when they enter the cavity. The arrival of the pair at the cavity is assumed to be Poissonian. We thus have $`\overline{t}_R=\tau +\overline{t}_{cav}`$ with $`t_R=1/\overline{R}`$. We assume $`\tau `$ to be a fixed duration for every pair of atoms, and hence, $`t_{cav}`$ is random. $`\overline{R}`$ is now a Poisson average of R pairs of atoms pumped into the cavity in one second. Due to cooperative interaction of the pair of atoms with the cavity field, we will find its dynamics with $`\overline{R}`$ pairs of atoms pumped into the cavity per second in the present case being different from the one-atom maser pumped with $`2\overline{R}`$ atoms into the cavity per second. This is due to the superradiance effects which will be manifested in the steady-state photon statistics. For these reasons, we call the present system a “Dicke micromaser”.
It is rather a difficult task to achieve such a model experimentally. However, we will show in the following that the present model is closely related to a one-atom two-photon micromaser \[16-18\]. This is a key result of this paper. Two-photon micromaser action with one atom has been experimentally demonstrated . Thus the model we are considering to show the atomic cooperative effects is not beyond experimental verification.
The organization of the paper is as follows. In the section 2, we present the microscopic Hamiltonian and the Dicke states involved in the micromaser action. The equations of motion are derived in section 3. We derive and discuss the photon statistics in section 4. The present system is compared with the two-photon micromasers in section 5. We conclude the paper in section 6.
2. The Model
A pair of Rydberg atoms in the upper of their two levels $`a`$ and $`b`$ enter the cavity at $`t=0`$. We assume that the atomic transition frequency is in resonance with the cavity mode frequency $`\omega `$. Both the atoms take the time $`\tau `$ to travel across the cavity. During this time, the evolution of the system comprising of the two atoms and the single eigen mode of the microwave cavity is governed by the Hamiltonian
$$H=H_0+H_I$$
$`(1a)`$
where
$$H_0=\omega (S_1^Z+S_2^Z+a^{}a)$$
$`(1b)`$
as we consider resonance condition between atomic transition frequencies and the cavity frequency. $`H_I`$ is the Jaynes-Cummings Hamiltonian .
$$H_I=g\underset{i=1}{\overset{2}{}}(S_i^+a+S_i^{}a^{})$$
$`(1c)`$
We assume that the atom-field coupling constant $`g`$ is the same for each atom. The atomic operators obey the relations $`[S_i^+,S_j^{}]=2S_i^Z\delta _{ij}`$ and $`a(a^{})`$ is the photon annihilation (creation) operator obeying the commutation relation $`[a,a^{}]=1`$. In the frame rotating at $`H_0`$, we have $`H=H_I`$.
It is convenient to represent the two-atom system by a Dicke system as both the atoms interact with the same cavity field simultaneously. In doing so, we however, neglect the dipole-dipole interaction as the microwave cavity dimensions are much large compared to particle wavelengths. Defining the collective operators $`S^\pm =_{i=1}^2S_i^\pm `$, Eq.(1c) takes the form
$$H_I=g(S^+a+S^{}a^{})$$
$`(2)`$
which is clearly invariant by atomic permutation. Both the atoms are in their respective upper states at $`t=0`$ and hence, the state representing the atomic system $`|a,a`$ is also invariant under atomic permutation. So it remains in a symmetrical state at all times. These are known as Dicke states with maximum cooperative number $`J=1`$ (we have a two-atom system) and these states are isomorphous to angular momentum states with principal angular momentum $`J=1`$. These states, represented by $`|J,M`$, are generated by
$$S^{}|J,M=\sqrt{(J+M)(JM+1)}|J,M1$$
and
$`(3)`$
$$S^+|J,M=\sqrt{(J+M+1)(JM)}|J,M+1$$
The level $`|J,M`$ is nondegenerate and has atomic energy $`M\omega `$. As $`JMJ`$, this has $`2J+1`$ states, and in the present case, the number of states is three. Since $`J`$ is a constant here, we represent $`|J,M|M`$.
Thus, the pair of two-level atoms can be represented by a three-level system of the Dicke states $`|M=1`$, $`|M=0`$, and $`|M=1`$, having energy $`\omega `$, $`0`$, and $`\omega `$ respectively. This allows us to write the Hamiltonian in Eq.(2) in the matrix form
$$H_I=\left(\begin{array}{ccc}0& g\sqrt{2(n1)}& 0\\ g\sqrt{2(n1)}& 0& g\sqrt{2n}\\ 0& g\sqrt{2n}& 0\end{array}\right)$$
$`(4)`$
where $`n`$ represents the photon number of the cavity field. Eq.(4) resembles the case for a two-photon process in a three-level system . Following the method in , we can obtain the eigenvalues and eigenstates of the interaction of the Dicke states with the cavity radiation field. The eigenvalues
$$\lambda _0=0,$$
$$\lambda _+=g\sqrt{2(2n1)},$$
$`(5)`$
$$\lambda _{}=g\sqrt{2(2n1)},$$
and its states are, respectively,
$$|0,n=x_1^{(n)}|1,n2+y_1^{(n)}|0,n1+z_1^{(n)}|1,n,$$
$$|+,n=x_2^{(n)}|1,n2+y_2^{(n)}|0,n1+z_2^{(n)}|1,n,$$
$`(6)`$
$$|,n=x_3^{(n)}|1,n2+y_3^{(n)}|0,n1+z_3^{(n)}|1,n,$$
where $`|1,n`$, $`|0,n`$ and $`|1,n`$ are the composite Dicke and field states, and
$$x_1^{(n)}=\sqrt{\frac{n}{2n1}},y_1^{(n)}=0,z_1^{(n)}=\sqrt{\frac{n1}{2n1}},$$
$`(7a)`$
$$x_2^{(n)}=\sqrt{\frac{n1}{4n2}},y_2^{(n)}=\frac{1}{\sqrt{2}},z_2^{(n)}=\sqrt{\frac{n}{4n2}},$$
$`(7b)`$
$$x_3^{(n)}=\sqrt{\frac{n1}{4n2}},y_3^{(n)}=\frac{l}{\sqrt{2}},z_3^{(n)}=\sqrt{\frac{n}{4n2}}.$$
$`(7c)`$
The states $`|+,n`$, $`|,n`$ and $`|0,n`$ can be called as dressed states of the interaction of the Dicke states with the cavity field.
3. Derivation of the equations of motion
Using the dressed states in Eq.(6) as the basic states, the equation of motion
$$\dot{\rho }=i[H,\rho ]$$
$`(8)`$
can be derived easily . $`\rho `$ represents the composite atom-field system having the initial conditions
$$\rho _{i,j}^{(n)}(t=0)i,n|\rho |j,n=x_i^{(n)}x_j^{(n)}P_{n2}$$
$`(9)`$
where $`P_n`$ is the photon distribution function of the cavity radiation field. In writing Eq.(6), we have neglected the influences of the atomic as well as the cavity reservoirs on the dynamics. These effects can play a crucial role in the one-atom micromaser . However, these influences have been shown to be negligible for a cavity having very high $`Q`$ and at very low temperatures and a low atomic pump rate. Eq.(8) represents such a system.
The time-dependent solutions for the density matrix elements can be written as
$$\rho _{i,j}^{(n)}(t)=P_{n2}x_i^{(n)}x_j^{(n)}e^{i(\lambda _i^{(n)}\lambda _j^{(n)})t},$$
$`(10)`$
$$i,j0,+,.$$
By inverting the matrix which expresses the dressed states as a linear superposition of the Dicke states, we can obtain the Dicke states $`|i,n`$, $`i=1,0,1`$, in terms of the dressed states. Using these relations and Eq.(10), we obtain the density matrix elements in the Dicke state basis. We follow these steps to get the $`P_n`$ at $`t=\tau `$, which is given by the trace
$$P_n(\tau )=\underset{M=1}{\overset{+1}{}}M,n|\rho |M,n.$$
$`(11)`$
We thus have
$$P_n(\tau )=\mathrm{\Theta }_1^{(n)}(\tau )P_n+\mathrm{\Theta }_2^{(n)}(\tau )P_{n1}+\mathrm{\Theta }_3^{(n)}(\tau )P_{n2}$$
$`(12)`$
where
$$\mathrm{\Theta }_1^{(n)}(\tau )=\frac{(n+2)}{(4n+3)}(x_1^{(n+2)})^2+\frac{(n+1)}{(4n+6)}\left((x_2^{(n+2)})^2+(x_3^{(n+2)})^2\right)$$
$$\frac{2\sqrt{(n+2)(n+1)}}{\sqrt{(4n+3)(4n+6)}}(x_1^{(n+2)}x_2^{(n+2)}Re.(e^{i\lambda _+^{(n+2)}\tau })+x_1^{(n+2)}x_3^{(n+2)}Re.(e^{i\lambda _{}^{(n+2)}\tau }))$$
$$+\frac{2(n+1)}{(4n+6)}x_2^{(n+2)}x_3^{(n+2)}Re.(e^{i(\lambda _+^{(n+2)}\lambda _{}^{(n+2)})\tau }),$$
$`(13a)`$
$$\mathrm{\Theta }_2^{(n)}(\tau )=\frac{1}{2}\left((x_2^{(n+1)})^2+(x_3^{(n+1)})^2\right)x_2^{(n+1)}x_3^{(n+1)}Re.(e^{i(\lambda _+^{(n+1)}\lambda _{}^{(n+1)})\tau }),$$
$`(13b)`$
$$\mathrm{\Theta }_3^{(n)}(\tau )=\frac{n1}{2n1}(x_1^{(n)})^2+\frac{n}{4n2}\left((x_2^{(n)})^2+(x_3^{(n)})^2\right)$$
$$+\frac{\sqrt{2n(n1)}}{(2n1)}(x_1^{(n)}x_2^{(n)}Re.(e^{i\lambda _+^{(n)}\tau })+x_1^{(n)}x_3^{(n)}Re.(e^{i\lambda _{}^{(n)}\tau }))$$
$$+\frac{2n}{4n2}x_2^{(n)}x_3^{(n)}Re.(e^{i(\lambda _+^{(n)}\lambda _{}^{(n)})\tau }).$$
$`(13c)`$
The change in $`P_n`$ at $`t=\tau `$ is then $`\delta P_n=P_n(\tau )P_n`$ where $`P_n`$ is the photon distribution function at the time of a pair of atoms entering the cavity. For a time $`\mathrm{}t`$ such that $`\tau \mathrm{}tt_p`$, where $`t_p`$ is the cavity photon lifetime, we have
$$\mathrm{}P_n=\delta P_nR\mathrm{}t$$
$`(14)`$
where $`R`$ is the number of pairs of atoms passing through the cavity in one second
The coarse-grained time derivative due to gain from the atomic interaction is given by
$$\frac{dP_n}{dt}|_{gain}=R\left((\mathrm{\Theta }_1^{(n)}1)P_n+\mathrm{\Theta }_2^{(n)}P_{n1}+\mathrm{\Theta }_3^{(n)}P_{n2}\right).$$
$`(15)`$
Eq.(15) represents the dynamics during time $`\tau `$. During the time lapse $`t_{cav}`$, between the flights of two successive pairs of atoms, the cavity field interacts with its reservoir, and its equation of motion is given by
$$\dot{P}_n|_{loss}=2(n+1)(\overline{n}_{th}+1)\kappa P_{n+1}2\kappa (n+\overline{n}_{th}+2n\overline{n}_{th})P_n+2n\kappa \overline{n}_{th}P_{n1}$$
$`(16)`$
where $`\kappa =(2t_p)^1`$ is the cavity bandwidth, and $`\overline{n}_{th}`$ is the thermal photon present in the cavity. Under the coarse-graining assumptions, the complete equation of motion combines Eqs.(15) and (16) additively. This assumption has been seen to be valid for a random input of atoms into the cavity, which is the case we are studying in this paper. We thus have,
$$\dot{P}_n=\dot{P}_n|_{gain}+\dot{P}_n|_{loss}$$
$`(17)`$
4. Steady-state photon statistics
The equation of motion for $`P_n`$, given by Eq.(17), involves $`P_{n2}`$, and this makes it difficult to get a steady-state solution by the method followed in case of the one-atom micromaser . However, following Risken , Eq.(17) can be recast in a tri-diagonal matrix equation involving the two-component vector $`T_n=\left(\begin{array}{c}P_{2n}\\ P_{2n+1}\end{array}\right)`$, which provides an expression for $`T_n`$ in the form of matrix continued fractions. Eq.(17) directly gives an expression for $`P_1`$ in terms of $`P_0`$, and $`P_0`$ can be determined from the normalization condition. This completes the determination of $`P_n`$. But numerical evaluations of $`P_n`$ go out of control due to extremely slow convergence of the matrix continued fractions.
Hence, we follow the method given below, which we find computationally efficient. First, we set $`n`$ at a value $`n=n_{max}`$, say, and write all the $`P_n`$ in a vector
$$\stackrel{}{X}=\{P_0,P_1,P_2,\mathrm{}\mathrm{}\mathrm{}.,P_{n_{max}}\}^T$$
$`(18)`$
This vector obviously has $`n_{max}+1`$ elements. We get $`n_{max}`$ equations involving
$`P_0,P_1,P_2,\mathrm{}\mathrm{}.,P_{n_{max}}`$ from Eq.(17). With the normalization condition
$$\underset{n=0}{\overset{n_{max}}{}}P_n=1$$
$`(19)`$
we get the equation
$$M\stackrel{}{X}=\stackrel{}{I}$$
$`(20)`$
where $`M`$ is a $`(n_{max}+1)\times (n_{max}+1)`$ matrix and $`\stackrel{}{I}`$ is another vector having again $`n_{max}+1`$ elements
$$\stackrel{}{I}=\{0,0,0,\mathrm{}\mathrm{}\mathrm{}.,1\}^T$$
$`(21)`$
The solution
$$\stackrel{}{X}=M^1\stackrel{}{I}$$
$`(22)`$
gives all $`P_n`$ for $`0nn_{max}`$.
We analyze the photon statistics of the cavity field by numerically evaluating its first and second moments, that is,
$$<n>=\underset{n=0}{\overset{n_{max}}{}}nP_n$$
$`(23)`$
proportional to the intensity of the cavity field and its normalized variance
$$v=\sqrt{\frac{<n^2><n>^2}{<n>}}$$
$`(24)`$
$`v=1`$ for a coherent state field, and thus, $`v<1`$ indicates the nonclassical nature of the cavity field. The process is repeated for a higher value of $`n_{max}`$ and the corresponding $`<n>`$ and $`v`$ are compared with those obtained from the earlier value of $`n_{max}`$. If they agree to a desired accuracy, the process is stopped, and we have the numerical values of the photon distribution function $`P_n`$, as well as the average $`<n>`$ and variance $`v`$.
The results are displayed in Figs.1-3. For the convenience of describing the micromaser action, we define the pump parameter
$$D=\sqrt{N}g\tau $$
$`(25)`$
where $`N=\overline{R}/2\kappa `$ is the number of pairs of atoms that pass through the cavity in a photon lifetime $`(2\kappa )^1`$. The patterns in the variation of $`<n>`$ and $`v`$ as $`D(\tau )`$ changes, can be understood from the dynamics of a collection of two-level atoms interacting with the radiation field of a single mode cavity with each atom being coupled to the cavity field by the Jaynes-Cummings interaction . For a two-atom case with initial condition $`J=1`$, it has been shown in that the dynamics is controlled by a d.c. term with a prefactor of the order of $`P_n/n`$ and terms oscillating at the first and second harmonic of the Rabi frequancy $`2g\sqrt{n+3/2}`$ of the interaction Hamiltonian in Eq.(1c) and having prefactors proportional to $`P_n`$ and $`P_n/n`$ respectively. In the single-atom case (Jaynes-Cummings model), the well known dynamics is controlled by only one term having the prefactor $`P_n`$ and oscillating at the Rabi frequency $`2g\sqrt{n+1}`$ of the Jaynes-Cummings interaction . These differences make the present results different, in general, from the one-atom micromaser action. Figs.1-3 show the differences clearly where we compare the results for the present Dicke micromaser for $`\overline{N}=100`$ with the one-atom micromaser action \[4-6\] with a pump rate $`\overline{N}=200`$. The numerical values of the two pump rates make the total number of atoms that pass through the two cavities having the same $`Q`$ equal, and thus, justifies the comparison between the two systems. Fig.2a shows that the superradiance nature of the interaction in Eq.2 makes the threshold values of $`D`$ lower compared to the one-atom micromaser. The differences between the two systems become predominant for higher values of $`D`$, that is, for longer interaction times. This is evidently due to the cooperative interaction becoming more important for longer interaction times. Thus, we notice in Fig.1 that $`<n>`$ is in general higher compared to the one-atom micromaser for longer $`\tau `$.
The variance in the cavity field is in general different in the two systems. It is generally very high near threshold, and hence, the sharp peaks in $`v`$ at threshold appear at different values of $`D`$ in the two systems as depicted in Fig.2b. We also notice in Fig.2b that the cavity field gets nonclassical properties ($`v<1`$) for different ranges of $`D`$ in the two systems. The photon distribution function also has different shapes as shown in Figs.3, which clearly shows that the nature of the field depends clearly on the nature of the pump. For example, in Fig.3a, the field in the one-atom case is sub-Poissonian, while the field in the two-atom case is super-Poissonian. The field characteristics are vice-versa in Fig.3c. However, the $`P_n`$ in the present case do not indicate any existence of the so-called trapped states \[4-6\]. The Rabi frequencies in the two-atom Dicke system, represented by the eigen values $`\lambda _0`$, $`\lambda _+`$ and $`\lambda _{}`$ do not provide clear conditions for the trapped states as one gets in the case of the one-atom micromaser . There the conditions are essentially zeros of the function $`Sin\sqrt{n+1}x`$, $`x`$ being the dimensionless atom-field interaction time. The trapped states, are however, washed out by occasional presence of two atoms in the cavity of the one-atom maser . Thus the present analysis gives an easy and transparent understanding of the disappearance of trapped states in multi-atom events.
5. Comparison with one-atom two-photon micromaser
The Dicke micromaser studied in the present paper has similarities with the dynamics of a one-atom two-photon micromaser involving pumping of atoms, individually into a microwave cavity where the atoms make a two-photon transition from the upper to lower level via an intermediate level. This is because the Dicke atomic system with $`J=1`$ is equivalent to a one-atom three-level system with the middle level having equal frequency of separation $`\omega `$ from the upper and the lower levels. It may be noted here that $`\omega `$ is the transition frequency of the individual two-level atoms in the Dicke system. If the cavity eigenmode is in resonance with the two degenerate atomic transitions, and in addition, if they have dipole moments of equal strength, then the present theory can be employed to explain the dynamics of the two-photon cascade micromaser. In other words, the results in Figs.1-3 can represent a one-atom two-photon cascade micromaser if we set the coupling constants of the two degenerate transitions $`g_1=g\sqrt{2}`$. The factor $`\sqrt{2}`$ in the Dicke system represents its cooperative nature .
We need not stop at this comparison. A Dicke micromaser with pumping of three atoms at a time (all in their upper masing levels) will be equivalent to a degenerate five-photon cascade micromaser, and so on.
However, a true two-photon process involves non-resonant one-photon transitions. If $`\omega _1`$ and $`\omega _2`$ are two atomic transitions, and $`\nu `$ is the cavity mode frequency, then a true two-photon process should have large detunings $`\mathrm{}_1=\omega _1\nu `$ and $`\mathrm{}_2=\nu \omega _2`$, that is, $`\mathrm{}_1`$, $`\mathrm{}_2g_1`$. A two-photon resonance means $`\mathrm{}_1+\mathrm{}_2=0`$, or $`\omega _1+\omega _2=2\nu `$. In the cascade two-photon micromaser discussed above, we have $`\mathrm{}_1=\mathrm{}_2=0`$. An extensive comparison of cavity-QED of cascade two-photon with true two-photon processes can be found in where it has been shown that the two dynamics are, in general, different.
The present approach can also be applied to such a two-photon micromaser if we can accomodate the one-photon detuning in the derivation of the photon distribution function. This amounts to just setting the matrix elements $`[H_I]_{1,1}=[H_I]_{3,3}=\mathrm{}`$ in Eq.4. For simplicity, we set $`g_1=g\sqrt{2}`$ here also. We present the photon statistics for the two photon micromaser for various values of $`\mathrm{}`$ in Figs.4. We find that the operational characteristics are strongly dependent on $`\mathrm{}`$. We also notice that as $`\mathrm{}`$ increases, the theory recovers the results for a two-photon micromaser derived by using an effective two-level Hamiltonian, a derivation and discussion of which can be found in .
For smaller values of $`\mathrm{}`$, we notice wriggles in the variation of $`<n>`$ versus $`D`$ in Fig.4a. These wriggles begin at a value of $`D(\tau )`$ and persist as it is further increased. The value of $`D`$ at which the wriggles begin to appear increases with $`\mathrm{}`$. In other words, the smooth variations in $`<n>`$ as in the case of an effective two-level system , occur for values of $`D`$ having an upper bound which increases with $`\mathrm{}`$. Such characteristics have been discussed in detail in the context of the two-photon Jaynes-Cummings model in . In addition to the above results, we also notice that the thresholds shift to higher values of $`D`$ as $`\mathrm{}`$ increases, and also, the range of $`D`$ between two successive thresholds also increases with the one-photon detuning. The cavity field usually has sub-Poissonian photon statistics $`(v<1)`$ between successive thresholds, and thus with a larger detuning, this model provides sub-Poissonian fields for wider range of $`D`$ as shown in Fig.4b.
6. Conclusion
We have presented a theory for a two-atom Dicke micromaser and have brought out the role played by the Dicke cooperativity on the micromaser action. Though at first glance it may look that the Dicke micromaser results mimic that for an equivalently pumped one-atom maser, the two systems are in general different. At places, however, the two systems have similar trends, an example of which can be seen in Fig.3b. We further show that the micromaser dynamics involving the three-level Dicke atomic system is formally equivalent to a one-atom cascade two-photon micromaser. Interestingly, we note that the two one-photon coupling constants in the two-photon cascade transition has to scale $`\sqrt{2}`$ times of the one-atom coupling constants in the Dicke system for the two micromaser actions to have identical results. As mentioned above, this factor of $`\sqrt{2}`$ originates from the cooperative nature of the Dicke system. With the introduction of a one-photon detuning in the two-photon process, the present approach also describes a one-atom two-photon micromaser. We have discussed in detail the role of the one-photon detuning in a two-photon micromaser action. Thus the Dicke superradiance effects can be quantitatively evaluated by scaling, as mentioned above, the results from the one-atom two-photon micromaser .
Thus the Dicke micromaser model discussed in this paper is not beyond the scope of experiments for verifying its results. The present paper gives the clue to verify results from micromaser actions involving pump mechanisms like the one studied here, which are rather difficult to achieve experimentally. This suggests a way to experimentally demonstrate the two-atom cavity-QED results in a one-atom two-photon micromaser . It may be recalled that the one-atom micromaser demonstrated the phenomenon of quantum revival . The same techniques may be employed in the case of the one-atom two-photon micromaser and a proper scaling of its results should demonstrate the two-atom cavity-QED results .
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FIGURE-CAPTIONS
Fig.1
The cavity field intensity, proportional to $`<n>`$, versus the pump parameter $`D`$. $`\overline{n}_{th}=0.1`$. $`\overline{N}=100`$ and $`200`$ for the Dicke micromaser (a) and one-atom micromaser (b) respectively. This makes the total number of atoms that pass through the cavity in the two systems equal. The curve (b) is shifted upwards by $`150`$ for clarity.
Fig.2a
Variation in $`<n>`$ with respect to $`D`$ for shorter interaction time. The parameters are same as in Fig.1. The full and broken curves are for the Dicke micromaser and one-atom micromaser respectively.
Fig.2b
Variation of $`v`$ with $`D`$. The other parameters are same as in Fig.2a.
Fig.3a
Photon distribution function for $`D=25`$. The other parameters are same as in Figs.2. $`P(n)`$ is sub-Poissonion $`(v=0.60761)`$ in the case of one-atom micromaser (broken curve), and is super-Poissonian $`(v=1.27596)`$ in the case of Dicke micromaser.
Fig.3b
$`P(n)`$ for $`D=50`$. The other parameters are same as in Fig.3a. For both the cases, $`P(n)`$ is super-Poissonian at this value of $`D`$. $`v=1.09944`$ (broken curve) and $`v=1.15871`$ (full curve).
Fig.3c
$`D=400`$. $`P(n)`$ is sub-Poissonian $`(v=0.22911)`$ in the Dicke micromaser system (full curve). One-atom micromaser (broken curve) provides a super-Poissonian $`(v=1.05726)`$ in this case.
Fig.4a
$`<n>`$ versus $`D`$ for a two-photon micromaser for three different values of the one-photon detuning $`\mathrm{}`$. $`\overline{n}_{th}=0.1`$ and $`\overline{N}=100`$. $`\mathrm{}=100`$ (curve (a)), $`\mathrm{}=150`$ (curve (b)), and $`\mathrm{}=300`$ (curve (c)). For clarity, the curves (b) and (c) are shifted upwards by $`40`$ and $`80`$ respectively.
Fig.4b
Variation of $`v`$ with $`D`$ for a two-photon micromaser. The other parameters are same as in Fig.4a. Here, the curves (b) and (c) are shifted by $`4`$ and $`8`$ respectively, for clarity. |
warning/0001/cond-mat0001162.html | ar5iv | text | # Anomalous Hopping Exponents of Ultrathin Films of Metals
## I INTRODUCTION
In highly disordered materials, electrical conduction occurs by the hopping of electrons between localized sites. This results in a thermally activated electrical resistance of the form:
$$R(T)=R_0\mathrm{exp}\left(\frac{T_0}{T}\right)^x$$
(1)
where $`T`$ is temperature, and $`R_0`$, $`T_0`$ and $`x`$ are constants which depend on the disorder, the details of the interactions and the dimensionality of the system. Simple activated hopping over a constant barrier results in the Arrhenius form with $`x=1`$. For noninteracting electrons, when the average hopping distance depends on temperature due to the compromise between hopping to sites which are close in energy, but farther away, Mott variable range hopping is expected, with $`x=1/(d+1)`$, where $`d`$ is the dimension. Efros and Shklovskii (ES) showed that including Coulomb interactions between electrons results in a soft gap in the density of states at the Fermi energy, which changes the variable range hopping exponent to $`x=1/2`$ in all dimensions .
Hopping conduction has been investigated in a wide variety of materials, such as doped semiconductors , semiconducting heterostructures , amorphous metals , magnetic materials and superconductors . Both the Mott and the ES forms of variable range hopping have been observed, as well as a crossover between the two regimes . It should be emphasized, however, that it is often hard to distinguish between Mott and ES hopping, particularly in experiments in which the resistance changes only by one or two orders of magnitude. The unambiguous identification of the Mott or ES hopping can be further complicated by factors which are usually neglected, such as the granularity of the system, possible temperature dependence of $`R_0`$, or correlations between electron hops.
While investigating the transport properties of ultrathin quench-condensed films over the course of many years, we have often found that the resistance of the thinnest films was thermally activated with $`x0.75`$. A similar hopping exponent has been reported by other authors , but has rarely been discussed per se. Since there is no theory which predicts this value of the hopping exponent, its origin has been left an open question. Here we report a detailed study of the temperature dependence of the resistance in very disordered films of four different materials: Ag, Bi, Pb and Pd. The films were grown in separate runs over several years, in different cryostats and on different substrates, yet they all show the same thermally activated resistance with an almost identical exponent. A careful analysis of the data points to a new conduction mechanism in this regime, or perhaps calls for a modification of the conventional picture. We compare our results with those of other experiments and available theoretical calculations, and suggest that the model that may possibly explain the anomalous hopping exponent is the collective variable range hopping model of Fisher et al. . Developed to describe vortices in superconductors, this approach has not been considered before in the context of charge transport in disordered electronic systems.
In Section II we survey the various models that have been considered in the discussion of transport in disordered films. We also exhibit the mapping of the model for collective vortex hopping onto the problem of charge transport in disordered systems and estimate the value of the hopping exponent. Experimental details of film growth and resistance measurements are given in Sec. III. In Sec. IV, we analyze the temperature dependence of the film resistance using several different methods to show that the exponent obtained is really a property of the system, rather than a consequence of an improper fit. The results are discussed and compared with other experimental and theoretical work in Sec. V.
## II SURVEY OF HOPPING MODELS
In recent years, extensions of the basic variable range hopping model to include percolation effects and correlations between electron hops have been developed. Deutscher et al. proposed a hopping mechanism which leads to a thermally activated resistance with $`x`$ close to $`1/2`$ without considering Coulomb interactions. The mechanism was based on the superlocalization property of wavefunctions on incipient percolation clusters , and may be relevant for atomically disordered systems as well as for granular percolative structures. The detailed microstructure of ultrathin quench-condensed films is still not known, and although these films are usually considered to be homogeneous, they may actually contain small grains or clusters. It is then possible that the electrons are restricted to move on a sublattice which is fractal over some range of length scales, and that their wavefunctions decay faster then exponentially with distance. Based on this assumption, a hopping conductivity law has been derived near the percolation threshold, which has the form of Eq 1 with $`x=3/7`$. This is experimentally almost indistinguishable from the Efros-Shklovskii law with $`x=1/2`$ if only the temperature dependence of the resistance is studied, but can be identified through the behavior of the parameter $`T_0`$ and the nature of the crossover to the conventional Mott regime .
Generally speaking, there is no reason to assume that the prefactor $`R_0`$ in Eq. 1 is independent of temperature. Van Keuls et al. studied the resistivity in a gated $`GaAs/Al_xGa_{1x}As`$ heterostructures as a function of temperature, electron density and magnetic field. Assuming the prefactor of the form:
$$R_0=bT^m$$
(2)
where $`b`$ and $`m`$ are constants, these workers fit Eq. 1 to their data with $`x=1/3`$ in low magnetic fields and with $`x=1/2`$ in high magnetic fields. The same crossover was observed as a function of electron density and temperature, and it scaled with the separation between the electron layer and the nearby screening gate, as predicted by Aleiner and Shklovskii . In addition to introducing a temperature dependent prefactor, this experiment also raised the issue of the importance of correlations between the electron hops.
In an electron glass, where the screening length is long and the interactions long-ranged, electron hopping may be correlated . Excitations can leave the system far from equilibrium and relaxation occurs through the rearrangement of charge. The energy of a single electron hop may then be significantly reduced by the motion of the surrounding charges. At sufficiently low temperatures, such collective hopping might be the dominant conduction mechanism.
In the analysis of their data, Van Keuls et al. assumed that the number of configurations of occupied states reached by the correlated hopping of a number of electrons is proportional to the single-particle density of states. In that case, the qualitative behavior of the resistance remains unchanged, and the effects of correlations enter through the constants which determine the parameter $`T_0`$ in different regimes.
Yet another issue which can be relevant in extremely thin films is the possibility that the electrons might interact logarithmically rather than as 1/r. As shown by Keldysh , the range of the logarithmic interaction is given by:
$$r_{\mathrm{log}}=\frac{\kappa }{\kappa _s+1}d$$
(3)
where $`\kappa `$ and $`\kappa _s`$ are the dielectric constants of the film and the substrate, respectively, and $`d`$ is the film thickness. One possible consequence of the logarithmic electron-electron interactions is that the Coulomb gap in the density of states (linear in 2D) may change to an exponential form . This leads to a modified variable range hopping law with a temperature dependent exponent.
Alternatively, the behavior of the logarithmically interacting electrons might be similar to that of vortices in 2D superconductors, which are known to interact logarithmically. Collective variable range hopping of vortices in disordered thin-film superconductors was studied by Fisher et al. . Including the effect of correlations, these authors found that multivortex hopping results in a lower energy than single vortex hopping. They also suggested that such multiparticle hopping might dominate single particle hopping even in the case of inverse power law interactions.
This approach may then be mapped onto a disordered two-dimensional system of charges. Following the arguments of Fisher et al. , the energy $`U(r)`$ of a multi-particle excitation of length $`r`$, can be estimated to be:
$$U(r)K(\frac{l}{r})^{1/2}$$
(4)
where $`l`$ is the distance between charges, and $`K`$ is the bare single-particle excitation energy. The latter is equal to $`e^2/\kappa a`$, where $`a`$ is the localization length and $`\kappa `$ is the dielectric constant. The simultaneous hopping of many charges may then result in a lower energy than the hop of a single charge.
The electrical resistance is a product of the probability for an electron to tunnel a distance $`r`$, $`\mathrm{\Gamma }_r`$, and the probability for an excitation with the energy $`U(r)`$ to occur, $`\mathrm{\Gamma }_U`$:
$$R1/\mathrm{\Gamma }_r\mathrm{\Gamma }_U$$
(5)
A lower bound for the multihop rate can be estimated as follows: if all of the $`(r/l)^2`$ electrons in the $`r`$ by $`r`$ region hop a distance comparable to the spacing $`l`$, then the rate should be proportional to the single-hop rate , $`exp(l/a)`$, raised to the power of the number of electrons, resulting in:
$$\mathrm{\Gamma }_re^{(l/a)(r/l)^2}$$
(6)
The probability for an excitation of energy $`U(r)`$ is proportional to $`exp(U(r)/T)`$, or using Eq. 4:
$$\mathrm{\Gamma }_Ue^{(K/T)(l/r)^{1/2}}$$
(7)
The minimum resistance is obtained when the hopping distance $`r_{hop}`$ is:
$$r_{hop}^2al(\frac{T_0}{T})^{4/5}$$
(8)
where
$$T_0=K(l/a)^{1/4}$$
(9)
Substituting Eqs. 6 through 9 back into Eq. 1 results in a hopping form such as that of Eq. 1 with $`x=4/5`$.
In the other limit, the minimum number of electrons participating in a collective hop could be taken as $`(r/d)`$, which leads to $`x=2/3`$. The collective variable range hopping mechanism may therefore result in a resistance of the form of Eq. 1, where the range of exponents is $`2/3<x<4/5`$, depending on the fraction of electrons participating in the process.
## III EXPERIMENTAL METHODS
The temperature dependence of the resistance has been studied in ultrathin quench-condensed films of Ag, Bi, Pd and Pb . The films were deposited on liquid helium cooled substrates and resistance measurements were performed in situ at temperatures down to $`0.15K`$. Ultra-high vacuum conditions and temperatures below $`20K`$ were sustained throughout each run, in order to avoid contamination or crystallization. The substrates were glazed alumina (for Bi and Pd films) or $`SrTiO_3`$ (100) (for Ag and Pb films). The $`SrTiO_3`$ (100) substrates were $`0.75mm`$ thick and had a $`100nm`$ thick Au gate on the back. Such a gate can be used to study the response of the film to a perpendicular electric field and was used to establish glass-like behavior in the most resistive films .
Films were deposited in thickness increments between $`0.05`$ and $`0.5\AA `$ on top of a thin germanium layer $`(510\AA )`$. (The Pd films were the exception to this as they were deposited directly onto glazed alumina substrates where they became connected at monolayer coverage.) Films grown on amorphous Ge are believed to be homogeneous, since they become connected at an average thickness of about one monolayer . The thicknesses of the films studied ranged from $`5\AA `$ up to $`15\AA `$. These nominal values of the film thickness are calculated from the deposition rate, which was measured using a quartz crystal monitor placed in the vicinity of the substrate. The first low-temperature scanning tunneling microscopy studies of the morphology of films grown in a similar manner indicate that the thinnest films may indeed be homogenous, while the thicker ones contain small clusters .
Resistance measurements were carried out using a standard dc four-probe technique. Very low bias currents $`(<10nA)`$ were used to avoid Joule heating of the sample and to make sure that the voltage across the sample was a linear function of the applied current. When measuring very resistive films ($`10^410^8`$ $`\mathrm{\Omega })`$, because of the long time constants of the circuit, it was necessary to monitor the voltage as a function of time after the current was changed, and allow adequate time for the voltage to stabilize. To avoid the voltage offset errors due to thermal EMFs, both polarities of the current were used to determine the resistance.
The resistance of a series of Ag films was also studied in a magnetic field. Magnetic fields of up to 20kG (12 kG) were applied in direction parallel (perpendicular) to the plane of the substrate using a split-coil superconducting magnet. In a regime where the anomalous hopping exponent is observed, the resistance was found to be independent of magnetic field.
## IV RESULTS AND ANALYSIS
The temperature dependence of the sheet resistance (resistance per square) for series of Ag, Bi, Pb and Pd films is shown in Fig. 1. The logarithm of the sheet resistance, plotted as a function of $`T^{0.75}`$ follows a straight line for each film, indicating that the resistance is thermally activated with $`x0.75`$. Using other values for $`x`$, such as $`1`$, $`1/2`$ or $`1/3`$ yielded considerably larger deviations from the data.
Since the prefactor in Eq. 1 is generally expected to be temperature-dependent $`(m0`$ in Eq. 2$`)`$, we attempted to fit the data using different combinations of $`m`$ and $`x`$. As shown in Fig. 2, using values of $`m`$ greater than zero actually increased the error of the fit. The maximum deviation in the fit of the combinations of Eqs. 1 and 2 to the data became much larger than the noise in $`R`$ as $`m`$ was increased. Furthermore if values of $`m`$ were chosen to force either Mott or ES hopping exponents of $`x=1/3`$ or $`1/2`$, respectively, the quality of fits as measured by chi squared would be significantly worse than that with $`m=0`$, in contrast with the findings of Van Keuls et al. .
Assuming that the hopping exponent is $`x0.75`$, the activation energy $`T_0`$ can be extracted from the fit to Eq. 1. The values of the parameter $`T_0`$ for different films of all four materials are shown in Fig. 3 as a function of $`R_{14K}`$, which is the sheet resistance measured at $`14K`$. This quantity is inversely proportional to the film thickness, so by using $`R_{14K}`$ instead of the thickness, one can avoid systematic errors in the nominal thicknesses of the films of different materials. It is apparent in Fig. 3 that $`T_0`$ changes greatly as $`R_{14K}`$ (and therefore also the film thickness) changes, ranging from around $`100K`$ for the thinnest films, to around $`10K`$ for the thickest films. The same qualitative and quantitative behavior was found for all four materials.
A more direct method of determining the exponent $`x`$ (which is exact under the condition $`m<<\left(T_0/T\right)^x`$) has been developed by Zabrodskii and Zinov’eva . The method is based on defining the function $`w=d(\mathrm{log}R)/d(\mathrm{log}T)`$. If $`R`$ is given by Eq. 1, then $`\mathrm{log}wx\mathrm{log}T`$. By plotting $`\mathrm{log}w`$ as a function of $`\mathrm{log}T`$, $`x`$ can be easily extracted from the slope of the resulting straight line. The benefit of the Zabrodskii-Zinov’eva approach is the simplicity of fitting a line rather than a complicated function with up to four adjustable parameters($`b,`$ $`m`$, $`T_{0\text{ }}`$and $`x`$). Once it has been determined that $`m=0`$, this method eliminates the danger of finding a local minimum instead of the best fit. The results of determining $`x`$ this way are shown in Fig. 4 . For very resistive films, plotting $`\mathrm{log}w`$ vs. $`\mathrm{log}T`$ indeed yielded straight lines. The values of $`x`$ varied slightly between different materials, from $`x0.7`$ for Ag to $`x0.8`$ for Bi. Remarkably, the value of $`x`$ did not change between different films in the same series over a significant range of sheet resistances, as shown in Fig. 5.
Even though $`R_{14K}`$ and $`T_0`$ change from film to film, as more material is added to increase the average thickness and decrease the sheet resistance of the film, $`x`$ stays constant over three orders of magnitude in $`R_{14K}`$. For thicker films (smaller $`R_{14K}`$), $`x`$ drops rather abruptly to a value between $`1/3`$ and $`1/2`$. In this regime, the Zabrodskii-Zinov’eva plots no longer produce straight lines, indicating that the hopping exponent changes as a function of temperature. Further increase of the film thickness leads to another crossover to a weakly localized regime where the temperature dependence of the resistance is logarithmic (not shown in Fig. 5).
## V DISCUSSION
An activated temperature dependence of the resistance with an anomalous hopping exponent $`x0.75`$ has been observed in disordered films of four different materials, grown on different substrates and measured in different cryostats. This strongly suggests that the exponent $`x0.75`$ is a general property of ultrathin films of metals in the very strongly localized regime.
The same exponent has been reported by Adkins and Astrakharchik in ultrathin quench-condensed films of Bi with a Ge underlayer. In that experiment, the temperature dependence of the resistance changed to simply activated (with $`x=1`$) when the Coulomb interaction was screened in the presence of a nearby metallic gate. This behavior was ascribed to the fixed range hopping of dipoles in screened films, but no details were given on the origin of $`x0.75`$ in unscreened films. The authors suggest that the films may be in the crossover regime between the variable range hopping and the fixed range hopping regime. Such a crossover can occur when the optimal hopping distance $`r_{hop}`$ becomes comparable with the localization length $`\xi `$.
In our experiment, $`x`$ stays constant over several orders of magnitude in sheet resistance, and then drops abruptly as the sheet resistance decreases further. If our films were merely at the crossover between $`x1/2`$ (or $`1/3`$) and $`x1`$, the change in $`x`$ would be expected to be gradual. The observed constancy of $`x`$ implies that a consistent mechanism may be governing the conduction in this regime, a mechanism different from Mott or Efros-Shklovskii variable range hopping which is usually observed in less resistive films.
Furthermore, it was not possible to obtain a satisfactory fit to the data using a temperature-dependent prefactor, as in the work of Van Keuls et al. It is interesting that these workers obtain $`x0.75`$ in all magnetic fields if $`m`$ is taken to be zero. However, the activation energies obtained from such fits are reported to be unacceptably small .
There are several other mechanisms which may be relevant in a very disordered 2D system. For example, Dai et al. observed an exponent $`x1`$ in Si:B, which changed to $`x1/2`$ when a magnetic field was applied. They suggested that the $`x1`$ was due to the exchange interaction between the electron spins, which is destroyed in a magnetic field. It must be noted that the Ag did not show any magnetoresistance up to the highest field available, 20 kG, so the exchange interaction is most likely not the origin of the anomalous hopping exponent observed in these films.
If we allow a possibility that our films are granular on a very small scale (which we cannot unequivocally rule out), than we must consider the superlocalization mechanism of Deutscher et al. as a possible candidate to explain our data. Without Coulomb interactions, this model predicts $`x0.43,`$ which obviously cannot account for our results. Including the Coulomb interactions may lead to a higher exponent, as proposed by van der Putten et al . These authors studied the hopping conductivity of percolating carbon-black-polymer compounds and found $`x0.66,`$ which is much closer to our result, although still too low. They interpret their results as evidence of superlocalization on a fractal network with Coulomb dominated hopping. The activation energies were found to be independent of the electron concentration, as predicted by Deutscher et al. In contrast, the activation energies found in our experiment depend strongly on the film thickness, which is closely related to the electron concentration. Furthermore, if the Coulomb interactions were screened, one might expect the exponent to decrease towards $`0.43`$, rather than to increase towards $`1`$, as observed in screened Bi films by Adkins and Astrakharchik .
Another possibility is that the anomalous exponent is a consequence of the exponential gap in the density of states, which can arise if the electrons interact logarithmically . In that case, the hopping exponent would be something close to, but smaller than $`1`$ at higher temperatures, and then cross over smoothly to $`1/2`$ at low temperatures. Forcing a fit of Eq. 1 to the data would result in an exponent which changes continuously with temperature. In the less resistive regime where we observe a temperature-dependent exponent, a closer inspection shows the opposite trend: the exponent is close to $`1/2`$ at higher temperatures, and increases with decreasing temperature. On the other hand, we cannot rule out the possibility that we might observe a smooth crossover to $`1/2`$ in the most resistive films if we could measure at much lower temperatures.
Finally, we consider the collective variable range hopping mechanism, proposed in the context of vortices in disordered superconductors by Fisher et al. . The mapping of this model onto a 2D electron system may actually be exact, if the electron-electron interactions are logarithmic over relevant length scales, but the authors suggest that collective hopping may dominate over the single-particle hopping even in the case of a conventional Coulomb interaction. The range of exponents predicted by the collective hopping model is $`2/3<x<4/5`$, depending on the ratio of electrons which participate in the process. The exponent found in our experiment, as well as the exponents found by other groups , are well within that range. The activation energies are expected to depend on the localization length, which in turn depends on the film thickness, as shown in Fig. 3.
In conclusion, we have addressed the issue of the anomalous hopping exponent $`x0.75`$ observed in ultrathin films of metals and related 2D systems. We argue that this hopping exponent is a general property of very strongly disordered systems, rather than a result of an improper fit or a signature of some sort of a crossover behavior. The usual models of hopping conduction do not explain this result. Our data can be explained by a collective variable range hopping mechanism, but our work by no means provides a proof of such a mechanism. Future experimental and theoretical studies will be needed to shed more light on this matter.
We gratefully acknowledge useful discussions with Boris Shklovskii and Leonid Glazman. This work was supported in part by the National Science Foundation under Grant No. NSF/DMR-987681.
Present Address: Department of Applied Physics, Technical University of Delft, the Netherlands.
Present Address: Center for Integrated Systems, Stanford University, Stanford, CA, USA.
Present Address: Seagate Technology, Bloomington, MN, USA. |
warning/0001/gr-qc0001069.html | ar5iv | text | # Regular coordinate systems for Schwarzschild and other spherical spacetimes
## I Introduction
The difficulties of the Schwarzschild coordinates $`(t,r,\theta ,\varphi )`$ at the event horizon of a nonrotating black hole provide a vivid illustration of the fact that in general relativity, the meaning of the coordinates is not independent of the metric tensor $`g_{\alpha \beta }`$. The Schwarzschild spacetime, whose metric is given by (we use geometrized units, so that $`c=G=1`$)
$`ds^2`$ $`=`$ $`fdt^2+f^1dr^2+r^2d\mathrm{\Omega }^2,`$ (1)
$`f`$ $`=`$ $`12M/r,`$ (3)
where $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$, indeed gives one of the simplest example of the failure of coordinates which have an obvious interpretation in one region of the spacetime (the region for which $`r2M`$), but not in another (the region for which $`r2M`$). Understanding this failure of the “standard” coordinate system is one of the most interesting challenges in the study of general relativity. Overcoming this obstacle is one of the most rewarding experiences associated with learning the theory.
Most textbooks on general relativity discuss the continuation of the Schwarzschild solution across the event horizon either via the Kruskal-Szekeres (KS) coordinates, or via the Eddington-Finkelstein (EF) coordinates; both coordinate systems produce a metric that is manifestly regular at $`r=2M`$. The main purpose of this paper is to show that useful alternatives exist. One of them, the Painlevé-Gullstand (PG) coordinates, are especially simple and attractive, and we will consider them in detail. We will also generalize them into a one-parameter family of coordinate systems, and show that the EF and PG coordinates are members of this family.
In a pedagogical context, the KS coordinates come with several drawbacks. First, the explicit construction of the KS coordinates is relatively complicated, and must be carried out in a fairly long series of steps. Second, the fact that $`r`$ is only implicitly defined in terms of the KS coordinates makes working with them rather difficult. Third, the manifold covered by the KS coordinates, with its two copies of each surface $`r=\text{constant}`$, is unnecessarily large for most practical applications; while the extension across the event horizon is desirable, the presence of another asymptotic region (for which $`r2M`$) often is not. While the KS coordinates are not to be dismissed out of hand — they do play an irreplaceable role in black-hole physics, and they should never be left out of a solid education in general relativity — we would advocate, for pedagogical purposes and as a first approach to this topic, the construction of simpler coordinate systems for extending the Schwarzschild spacetime across the event horizon.
A useful alternative are the EF coordinates $`(v,r,\theta ,\varphi )`$, in which the metric takes the form
$$ds^2=fdv^2+2dvdr+r^2d\mathrm{\Omega }^2.$$
(4)
The new time coordinate $`v`$ is constant on ingoing, radial, null geodesics ($`r`$ decreases, $`\theta `$ and $`\varphi `$ are constant); it is related to the Schwarzschild time $`t`$ by $`v=t+r^{}`$, where
$$r^{}=\frac{dr}{f}=r+2M\mathrm{ln}\left|\frac{r}{2M}1\right|.$$
(5)
The metric of Eq. (4) is regular across the event horizon. While its nondiagonal structure makes it slightly harder to work with than the metric of Eq. (LABEL:1.1), the fact that $`r`$ appears explicitly as one of the coordinates makes it much more convenient than the KS version of the Schwarzschild metric. We believe that in a pedagogical context, the Eddington-Finkelstein coordinates should be introduced before the KS coordinates.
Our first objective in this article is to popularize another set of coordinates for Schwarzschild spacetime, and propose this system as a useful alternative to the EF coordinates. These are the Painlevé-Gullstrand (PG ) coordinates $`(T,r,\theta ,\varphi )`$. They are constructed and discussed in Sec. II. Our second objective is to provide generalizations of this coordinate system. In Sec. III we discuss a one-parameter family of PG-like coordinates for Schwarzschild spacetime. To the best of our knowledge this family was first discovered by Kayll Lake in 1994 , but a related family of coordinates was previously discussed by Gautreau and Hoffmann . We show in Sec. III that the PG and EF coordinates are both members of Lake’s family. In Sec. IV we generalize this family of coordinate systems to other spherical (and static) spacetimes; equivalent coordinates were constructed, in a two-dimensional context, by Corley and Jacobson . In Sec. V we look back at our coordinates, and offer some additional comments regarding their construction. In the Appendix we relate these coordinate systems to the KS coordinates, and provide details regarding the spacetime diagrams of Figs. 1 and 2.
## II Painlevé-Gullstrand coordinates
It is often a good strategy, when looking for regular coordinate systems, to anchor the coordinates to a specific family of freely moving observers . We shall employ this strategy throughout this paper. The following derivation of the PG coordinates can be found in the book by Robertson and Noonan . Other derivations can be found in Refs. , in which the PG coordinates were independently rediscovered.
We consider observers which move along ingoing, radial, timelike geodesics of the Schwarzschild spacetime ($`r`$ decreases, $`\theta `$ and $`\varphi `$ are constant). It is easy to check that in the standard coordinates of Eq. (LABEL:1.1), the geodesic equations can be expressed in first-order form as
$$\dot{t}=\frac{\stackrel{~}{E}}{f},\dot{r}^2+f=\stackrel{~}{E}^2,$$
(6)
where an overdot denotes differentiation with respect to the observer’s proper time, and $`\stackrel{~}{E}=E/m`$ is the observer’s (conserved) energy per unit rest mass. (For a derivation, see Chap. 11 of Ref. , Chap. 25 of Ref. , or Chap. 6 of Ref. .) We assume that $`\dot{r}<0`$, and the energy parameter is related to the observer’s initial velocity $`v_{\mathrm{}}`$ — the velocity at $`r=\mathrm{}`$ — by
$$\stackrel{~}{E}=\frac{1}{\sqrt{1v_{\mathrm{}}^{}{}_{}{}^{2}}}.$$
(7)
In this section we specialize to the particular family of observers characterized by $`\stackrel{~}{E}=1`$; our observers are all starting at infinity with a zero initial velocity: $`v_{\mathrm{}}=0`$. For these observers, the geodesic equations reduce to $`\dot{t}=1/f`$ and $`\dot{r}=\sqrt{1f}`$. We notice that $`u_\alpha `$, the covariant components of the observer’s four-velocity, whose contravariant components are $`u^\alpha =(\dot{t},\dot{r},0,0)`$, is given by $`u_\alpha =(1,\sqrt{1f}/f,0,0)`$. This means that $`u_\alpha `$ is equal to the gradient of some time function $`T`$:
$$u_\alpha =_\alpha T,$$
(8)
where
$$T=t+\frac{\sqrt{1f}}{f}𝑑r.$$
(9)
Integration of the second term is elementary, and we obtain
$$T=t+4M\left(\sqrt{r/2M}+\frac{1}{2}\mathrm{ln}\left|\frac{\sqrt{r/2M}1}{\sqrt{r/2M}+1}\right|\right).$$
(10)
This shall be our new time coordinate, and $`(T,r,\theta ,\varphi )`$ are nothing but the PG coordinates. It should be clear that the key to the construction of the PG coordinates is the fact that the four-velocity can be expressed as in Eq. (8). In Sec. V we will explain how this equation comes about.
Going back to Eq. (9), we see that $`dt=dTf^1\sqrt{2M/r}dr`$. Substituting this into Eq. (LABEL:1.1) gives
$$ds^2=fdT^2+2\sqrt{2M/r}dTdr+dr^2+r^2d\mathrm{\Omega }^2.$$
(11)
This is the Schwarzschild metric in the PG coordinates. An equivalent way of expressing this is
$$ds^2=dT^2+\left(dr+\sqrt{2M/r}dT\right)^2+r^2d\mathrm{\Omega }^2.$$
(12)
This metric is manifestly regular at $`r=2M`$, in correspondence with the fact that our observers do not consider this surface to be in any way special. (The metric is of course still singular at $`r=0`$.) While the metric is now nondiagonal, it has a remarkably simple form. It is much simpler than the Kruskal-Szekeres metric, and we believe that it provides a useful alternative to the Eddington-Finkelstein form of the metric, Eq. 1.2.
In Fig. 1 we show several surfaces $`T=\text{constant}`$ in a Kruskal diagram. The construction is detailed in the Appendix. The diagram makes it clear that the PG coordinates do not extend inside the past horizon of the Schwarzschild spacetime — the “white-hole region” is not covered. The reason for this is that our observers fall inward from infinity and end up crossing the future, but not the past, horizon. By reversing the motion (choosing the opposite sign for $`\dot{r}`$), we would obtain alternative coordinates that extend within the past horizon but do not cover the black-hole region of the spacetime. While the PG coordinates do not cover the entire KS manifold, they do cover the two most interesting regions of the maximally extended Schwarzschild spacetime.
Equations (11) and (12) reveal the striking property that the surfaces $`T=\text{constant}`$ are intrinsically flat: Setting $`dT=0`$ returns $`ds^2=dr^2+r^2d\mathrm{\Omega }^2`$, which is the metric of flat, three-dimensional space in spherical polar coordinates. The information about the spacetime curvature is therefore entirely encoded in the “shift vector”, the off-diagonal component of the metric tensor. We consider this aspect of the PG coordinates to be their most interesting property.
We note that it is possible to construct PG-like coordinates for the nonspherical Kerr spacetime. This was carried out by C. Doran in a recent paper .
## III Generalization to other observers
It is easy to generalize the preceding construction to other families of freely moving observers. In this section we consider families such that $`\stackrel{~}{E}`$ is the same for all observers within the family, but not equal to unity (as in Sec. II). Each family is therefore characterized by its unique value of the energy parameter. We find it convenient to use instead the parameter $`p`$, related to the energy and initial-velocity parameters by
$$p=\frac{1}{\stackrel{~}{E}^2}=1v_{\mathrm{}}^{}{}_{}{}^{2}.$$
(13)
We take $`p`$ to be in the interval $`0<p1`$, with $`p=1`$ taking us back to the PG coordinates . To each value of $`p`$ in this interval corresponds a family of observers, and a distinct coordinate system. We are therefore constructing a one-parameter family of PG-like coordinates for Schwarzschild spacetime.
With the geodesic equations now given by $`\dot{t}=1/(\sqrt{p}f)`$ and $`\dot{r}=\sqrt{1pf}/\sqrt{p}`$, we find that $`u_\alpha `$ is now equal to a constant times the gradient of a time function $`T`$:
$$u_\alpha =\frac{1}{\sqrt{p}}_\alpha T,$$
(14)
with
$$T=t+\frac{\sqrt{1pf}}{f}𝑑r.$$
(15)
Integration of the second term doesn’t present any essential difficulties, and we obtain
$`T`$ $`=`$ $`t+2M({\displaystyle \frac{1pf}{1f}}+\mathrm{ln}\left|{\displaystyle \frac{1\sqrt{1pf}}{1+\sqrt{1pf}}}\right|`$ (17)
$`{\displaystyle \frac{1p/2}{\sqrt{1p}}}\mathrm{ln}\left|{\displaystyle \frac{\sqrt{1pf}\sqrt{1p}}{\sqrt{1pf}+\sqrt{1p}}}\right|).`$
This shall be our new time coordinate. In Sec. V we will return to the question of the origin of Eq. (14).
With $`dt`$ now equal to $`dTf^1\sqrt{1pf}dr`$, we find that the Schwarzschild metric takes the form
$$ds^2=fdT^2+2\sqrt{1pf}dTdr+pdr^2+r^2d\mathrm{\Omega }^2,$$
(18)
or
$$ds^2=\frac{1}{p}dT^2+p\left(dr+\frac{1}{p}\sqrt{1pf}dT\right)^2+r^2d\mathrm{\Omega }^2.$$
(19)
This metric is still regular at $`r=2M`$, although it is now slightly more complicated than the PG form.
In Fig. 2 we show several surfaces $`T=\text{constant}`$ in a Kruskal diagram, for several values of $`p`$. This construction is detailed in the Appendix.
In this generalization of the PG coordinates, the surfaces $`T=\text{constant}`$ are no longer intrinsically flat. The induced metric is now $`ds^2=pdr^2+r^2d\mathrm{\Omega }^2`$, and although the factor of $`p`$ in front of $`dr^2`$ looks innocuous, it is sufficient to produce a curvature. It may indeed be checked that the Riemann tensor associated with this metric is nonzero. The only nonvanising component is $`R_{\theta \varphi \theta }^\varphi =(1p)/p`$, and $`R^{abcd}R_{abcd}=4(1p)^2/(pr^2)^2`$.
It is instructive to go back to Eq. (17) and check that in the limit $`p1`$, $`T`$ reduces to the expression of Eq. (10). (This must be done as a limiting procedure, because $`T`$ is ambiguous for $`p=1`$.) Taking the limit gives
$$\underset{p1}{lim}T=t+2M\left(\frac{2}{\sqrt{1f}}+\mathrm{ln}\left|\frac{1\sqrt{1f}}{1+\sqrt{1f}}\right|\right),$$
(20)
which is indeed equivalent to Eq. (10). The PG coordinates are therefore a member of our one-parameter family.
Another interesting limit is $`p0`$, which corresponds to $`\stackrel{~}{E}\mathrm{}`$, or $`v_{\mathrm{}}1`$. In this limit, our observers start at infinity with a velocity nearly equal to the speed of light. Starting from Eq. (17) we have
$$\underset{p0}{lim}T=t+2M\left(\frac{1}{1f}+\mathrm{ln}\left|\frac{f}{1f}\right|\right)=t+r^{},$$
(21)
where we have compared with Eq. (5). Thus, $`T=v`$ in the limit $`p0`$, and our generalized PG coordinates reduce to the Eddington-Finkelstein coordinates of Eq. (4). This is not entirely surprising, in view of the fact that our observers become light-like in this limit. The EF coordinates, therefore, are also a (limiting) member of our one-parameter family.
We have constructed an interpolating family of coordinate systems for Schwarzschild spacetime; as the parameter $`p`$ varies from 1 to 0, the coordinates go smoothly from the Painlevé-Gullstrand coordinates to the Eddington-Finkelstein coordinates. This one-parameter family of coordinate systems was first discovered by Kayll Lake , but a related family of coordinates, corresponding to $`p>1`$, were previously introduced by Gautreau and Hoffmann . Lake obtained the new coordinates by solving the Einstein field equations for a vacuum, spherical spacetime in a coordinate system involving $`r`$ and an arbitrary time $`T`$. The intimate relation between his coordinates and our families of freely moving observers remained unnoticed by him.
## IV Generalization to other spacetimes
The coordinates constructed in the previous two sections can be generalized to other static and spherically symmetric spacetimes. In the usual Schwarzschild coordinates, we write the metric as
$$ds^2=e^{2\psi }fdt^2+f^1dr^2+r^2d\mathrm{\Omega }^2,$$
(22)
where $`f`$ and $`\psi `$ are two arbitrary functions of $`r`$. Under the stated symmetries, Eq. (22) gives the most general form for the metric. We assume that the spacetime is asymptotically flat, so that $`f1`$ and $`\psi 0`$ as $`r\mathrm{}`$. If the spacetime possesses a regular event horizon at $`r=r_0`$, then $`f(r_0)=0`$ and $`\psi `$ must be nonsingular for all values of $`r0`$.
The geodesic equations are now
$$\dot{t}=\frac{\stackrel{~}{E}}{e^{2\psi }f},\dot{r}^2+f=e^{2\psi }\stackrel{~}{E}^2,$$
(23)
where $`\stackrel{~}{E}`$ is still the conserved energy per unit rest mass. Re-introducting $`p=1/\stackrel{~}{E}^2`$, we find that the covariant components of the four-velocity can be again expressed as in Eq. (14), with a time function $`T`$ now given by
$$T=t+\frac{\sqrt{e^{2\psi }pf}}{f}𝑑r.$$
(24)
The second term can be integrated if $`f`$ and $`\psi `$ are known. Rewriting the metric of Eq. (22) in terms of $`dT`$ yields
$`ds^2`$ $`=`$ $`fe^{2\psi }dT^2+2e^{2\psi }\sqrt{e^{2\psi }pf}dTdr`$ (26)
$`+pe^{2\psi }dr^2+r^2d\mathrm{\Omega }^2,`$
or
$`ds^2`$ $`=`$ $`{\displaystyle \frac{1}{p}}dT^2+pe^{2\psi }\left(dr+{\displaystyle \frac{1}{p}}\sqrt{e^{2\psi }pf}dT\right)^2`$ (28)
$`+r^2d\mathrm{\Omega }^2.`$
This metric is manifestly regular at an eventual event horizon, at which $`f`$ vanishes.
The surfaces $`T=\text{constant}`$ have an induced metric given by $`ds^2=pe^{2\psi }dr^2+r^2d\mathrm{\Omega }^2`$. Unless $`\psi =0`$ and $`p=1`$, these surfaces are not intrinsically flat .
## V Final comments
In all the cases considered in Secs. II, III, and IV, the construction of our coordinate systems relied on the key fact that the four-velocity could be expressed as $`u_\alpha =_\alpha T/\sqrt{p}`$, with $`p`$ a constant. \[This is Eq. (14), and in Sec. II, $`p`$ was set equal to unity.\] This property is remarkable, and it seems to follow quite accidentally from the equations of motion. There is of course no accident, but the point remains that not every four-velocity vector can be expressed in this form.
A standard theorem of differential geometry (for example, see Appendix B of Ref. ) states that for $`u_\alpha `$ to admit the form of Eq. (14), it must satisfy the equations $`u_{;\beta }^\alpha u^\beta =0`$ and $`u_{[\alpha ;\beta }u_{\gamma ]}=0`$, in which a semicolon denotes covariant differentiation and the square brackets indicate complete antisymmetrization of the indices. The second equation states that the world lines are everywhere orthogonal to a family of spacelike hypersurfaces, the surfaces of constant $`T`$. This ensures that the four-velocity can be expressed as $`u_\alpha =\mu _\alpha T`$, for some function $`\mu (x^\alpha )`$. In general, this function is not a constant, and we do yet have Eq. (14). For this we need to impose also the first equation, which states that the motion is geodesic. When both equations hold we find that $`\mu =\text{constant}`$, and this gives us Eq. (14).
In our constructions, we have enforced the geodesic equation by selecting freely moving observers. By selecting radial observers, we have also enforced the condition that the geodesics be hypersurface orthogonal. Our strategy for constructing coordinate systems is therefore limited to radial, freely moving observers in static, spherically-symmetric spacetimes; it may not work for more general motions and/or more general spacetimes.
## Acknowledgments
This work was supported by the Natural Sciences and Engineering Research Council of Canada. We are grateful to Kayll Lake, Ted Jacobson, and an anonymous referee for discussions and comments on the manuscript.
## A Kruskal diagrams
The Kruskal diagrams of Figs. 1 and 2 are constructed as follows.
From the Schwarzschild coordinates $`t`$ and $`r`$ we define two null coordinates, $`u=tr^{}`$ and $`v=t+r^{}`$, where $`r^{}`$ is given by Eq. (5). From these we form the null KS coordinates, $`V=e^{v/4M}`$ and $`U=e^{u/4M}`$, where the upper sign refers to the region $`r>2M`$ of the Schwarzschild spacetime, while the lower sign refers to $`r<2M`$. From this we derive the relations
$$UV=e^{r/2M}\left(\frac{r}{2M}1\right)$$
(A1)
and
$$\frac{V}{U}=e^{t/2M}.$$
(A2)
Timelike and spacelike KS coordinates are then defined by $`V=\tau +\rho `$ and $`U=\tau \rho `$. In our spacetime diagrams, the $`\tau `$ axis runs vertically, while the $`\rho `$ axis runs horizontally. The future horizon is located at $`U=0`$, and the past horizon is at $`V=0`$. The curvature singularity is located at $`UV=1`$.
We express the time function of Eq. (17) as
$$T=t+r^{}+S(r),$$
(A3)
where $`S(r)`$ is the function of $`r`$ that results when the second term of Eq. (17) is shifted by $`r^{}`$, as given in Eq. (5); this function is regular at $`r=2M`$. With this definition we have $`v=TS`$, $`u=TS2r^{}`$, as well as
$$V=e^{T/4M}e^{S/4M}$$
(A4)
and
$$U=e^{r/2M}(r/2M1)e^{T/4M}e^{S/4M}.$$
(A5)
The surfaces $`T=\text{constant}`$ give rise to parametric equations of the form $`V(r)`$ and $`U(r)`$, which are obtained from Eqs. (A4) and (A5) by explicitly evaluating the function $`S(r)`$. In these equations, $`r`$ can be varied from zero to an arbitrarily large value without difficulty. The diagrams of Figs. 1 and 2 are then produced by switching to the coordinates $`\overline{t}`$ and $`\overline{r}`$ and plotting the parametric curves. |
warning/0001/math0001007.html | ar5iv | text | # Frobenius∞ invariants of homotopy Gerstenhaber algebras, I
## 1 Introduction
Frobenius manifolds play a central role in the usual formulation of Mirror Symmetry, as may be seen in the following diagram,
where morphisms in all categories are just diffeomorphisms preserving relevant structures, and $`𝖦W`$ and $`𝖡K`$ stand, respectively, for the Gromov-Witten (see, e.g., \[Ma2\]) and Barannikov-Kontsevich \[BK, Ba\] functors. A pair $`(\stackrel{~}{M},M)`$ consisting of a symplectic manifold $`\stackrel{~}{M}`$ and a Calabi-Yau manifold $`M`$ is said to be Mirror if $`𝖦W(\stackrel{~}{M})=𝖡K(M)`$. According to Kontsevich \[Ko1\], this equivalence is a shadow of a more fundamental equivalence of natural $`A_{\mathrm{}}`$-categories attached to $`M`$ and $`\stackrel{~}{M}`$.
This paper is much motivated by the Barannikov-Kontsevich construction \[BK, Ba\] of the functor from the right in the above diagram, and by Manin’s comments \[Ma1\] on their construction. The roots of the $`𝖡K`$ functor lie in the extended deformation theory of complex structures on $`M`$, more precisely in very special properties of the (differential) Gerstenhaber algebra $`𝔤`$ “controlling” such deformations. One of the miracle features of Calabi-Yau manifolds, the one which played a key role in the $`𝖡K`$ construction, is that deformations of their complex structures are non-obstructed, always producing a smooth versal moduli space<sup>1</sup><sup>1</sup>1A similar phenomenon occurs in the extended deformation theory of Lefschetz symplectic structures which also produces, via the same $`𝖡K`$ functor, Frobenius manifolds \[Me1\]. These should not be confused with GW.. In the language of Gerstenhaber algebras, the exceptional algebraic properties necessary to produce a Frobenius manifold out of $`𝔤`$, have been axiomatized in \[Ma1, Ma2\]. As a result, the functor
is now well understood.
One of our purposes in this paper is to extend the BK functor from the category of Calabi-Yau manifolds to the category of arbitrary compact complex manifolds. Which means the study of a diagram
Generically, the extended deformation theory of complex structures is obstructed, and it would be naive to expect that the question mark above stands for the category of Frobenius manifolds. In fact, it is not, and the answer is captured in the following notion.
### 1.1. Definition.
An $`F_{\mathrm{}}`$-manifold is the data $`(,E,,[\mu _{}],e)`$, where
* $``$ is a formal pointed $``$-graded manifold,
* $`E`$ is the Euler vector field on $``$, $`Ef:=\frac{1}{2}|f|f`$, for all homogeneous functions on $``$ of degree $`|f|`$,
* $``$ is an odd homological (i.e. $`^2=0`$) vector field on $``$ such that $`[E,]=\frac{1}{2}`$ and $`II^2`$, $`I`$ being the ideal of the distinguished point in $``$,
* $`[\mu _n:^n𝒯_{}𝒯_{}]`$, $`n`$, is a homotopy class of smooth unital strong homotopy commutative ($`C_{\mathrm{}}`$) algebras defined on the tangent sheaf, $`𝒯_{}`$, to $``$, such that $`Lie_E\mu _n=\frac{1}{2}n\mu _n`$, for all $`n`$, and $`\mu _1`$ is given by
$$\begin{array}{cccc}\hfill \mu _1:& 𝒯_{}& & 𝒯_{}\\ & X& & \mu _1(X):=[,X].\end{array}$$
* $`e`$ is the unit, i.e. an even vector field on $``$ such that $`[,e]=0`$, $`\mu _2(e,X)=X`$, $`X𝒯_{}`$, and $`\mu _n(\mathrm{},e,\mathrm{})=0`$ for all $`n3`$.
Clearly, the category of $`F_{\mathrm{}}`$-manifolds containes Frobenius manifolds as a subcategory.
On any $`F_{\mathrm{}}`$-manifold the vector field $``$ defines an integrable distribition, $`𝖨m\mu _1`$, which is tangent to its subspace of zeros, $`\mathrm{`}\mathrm{`}𝗓eros()\text{}`$. The structures (i)-(v) make the tangent sheaf to the smooth part of the associated quotient, $`\mathrm{`}\mathrm{`}𝗓eros()\text{}/𝖨m\mu _1`$, into a sheaf of graded unital associative algebras.
### Theorem A.
For any differential unital (graded commutative) Gerstenhaber algebra $`𝔤`$, its cohomology, $`𝐇(𝔤)`$, if finite-dimensinal, is canonically an $`F_{\mathrm{}}`$-manifold.
The resulting diagram
implies, in turn, a diagram,
through the Gerstenhaber algebras controlling extended deformations of symplectic, complex and holomorphic vector bundle structures. Moreover, the $`F_{\mathrm{}}`$-functor enjoys the correct “classical limit”: when restricted to exceptional Gerstenhaber algebras (i.e. the ones satisfying Manin’s axioms \[Ma1\]), the $`F_{\mathrm{}}`$-functor coincides precisely with the Barannikov-Kontsevich construction \[BK\], and hence takes values in the subcategory of Frobenius manifolds.
Let us emphasize once again that, in the above diagram, $`F_{\mathrm{}}(𝗌ymplecticmanifolds)`$ has nothing to do with Gromov-Witten invariants<sup>2</sup><sup>2</sup>2At best, this is a very weak shadow of the Mirror Symmetry, see below.. Nevertheless, to rather different mathematical objects we can canonically attach invariants lying in one and the same geometric category. Hence we can use these $`F_{\mathrm{}}`$-invariants for a classification, and even speak about dull mirror symmetry when
$$F_{\mathrm{}}(𝖮bject)=F_{\mathrm{}}(\stackrel{~}{𝖮bject}).$$
Such a relation may be a shadow of something conceptually more interesting (cf. \[Ko1\]).
Theorem A is explained and generalized by the following
### $`𝐓heoremB`$.
There is a canonical functor, $`F_{\mathrm{}}`$, from the derived category of unital homotopy Gerstenhaber ($`G_{\mathrm{}}`$) algebras with finite-dimensional cohomology to the category of $`F_{\mathrm{}}`$-manifolds.
This result implies that the cohomology space of any homotopy Gerstenhaber algebra is, if finite-dimensional, canonically an $`F_{\mathrm{}}`$-manifold,
The recent proof of Deligne’s conjecture \[Ta1, Ko3, V, MS\] gives the following diagrammatic corollary of Theorem B,
which, probably, has a direct relevance to the Mirror symmetry through the following specializations, C and D, of statement B.
Any $`G_{\mathrm{}}`$-algebra $`𝔤`$ is, in particular, a $`L_{\mathrm{}}`$-algebra so that its cohomology, $`𝐇(𝔤)`$, has the induced structure, $`[,]_{\mathrm{i}nd}`$, of Lie algebra. If there exists a quasi-isomorphism of $`L_{\mathrm{}}`$-algebras,
$$(𝔤,L_{\mathrm{}}\mathrm{c}omponentoftheG_{\mathrm{}}\mathrm{s}tructure)\stackrel{F}{}(𝐇(𝔤),[,]_{\mathrm{i}nd}),$$
then $`𝔤`$ is said to be $`L_{\mathrm{}}`$-formal, and $`F`$ is called a formality map. In terms of the associated $`F_{\mathrm{}}`$-invariant, the $`L_{\mathrm{}}`$-formality of a $`G_{\mathrm{}}`$-algebra $`𝔤`$ gets translated into a canonical flat structure, the Gauss-Manin connection $``$, such that $`_Xe=0`$ and $`_X_Y_Z=0`$ for any flat vector fields $`X,Y`$ and $`Z`$. This specialization of $`F_{\mathrm{}}`$-structure is called pre-$`Frobenius_{\mathrm{}}`$ structure.
### Theorem C.
There is a canonical functor from the category of pairs $`(𝔤,F)`$, where $`𝔤`$ is a $`L_{\mathrm{}}`$-formal homotopy Gerstenhaber algebra and $`F`$ a formality map, to the category of pre-$`Frobenius_{\mathrm{}}`$ manifolds.
In fact, a pre-$`\mathrm{F}robenius_{\mathrm{}}`$ manifold (a Frobenius manifold, in particular) is itself a homotopy Gerstenhaber algebra.
According to Kontsevich’s celebrated Formality Theorem \[Ko2\], for any compact complex manifold $`M`$, the Hochschild differential Lie algebra, $`C^{}(𝒜,𝒜)`$, associated to the algebra of Dolbeault forms, $`𝒜=(\mathrm{\Gamma }(M,\mathrm{\Omega }_M^{0,}),\overline{})`$, is $`L_{\mathrm{}}`$-formal so that Theorem C has a wide area of applications.
### Theorem D.
If a homotopy Gerstenhaber algebra $`𝔤`$ is quasi-isomorphic, as a $`L_{\mathrm{}}`$-algebra, to an Abelian differential Lie algebra, then the tangent sheaf, $`𝒯_{}`$, to its cohomology $`𝐇(𝔤)`$ viewed as a linear supermanifold is canonically a sheaf of unital graded commutative associative algebras.
The point is that the Hochschild complex built out of the Dolbeault algebra, $`𝒜=(\mathrm{\Gamma }(,\mathrm{\Omega }_M^{0,}),\overline{})`$, of a Calabi-Yau manifold $``$, satisfies the conditions of Theorem D. In fact, the canonically induced associative product on the tangent sheaf to the associated linear supermanifold, $`𝐇^{}(,^{}T_{})`$, is, for an appropriate formality map, potential and satisfies the WDVV equations. The resulting composition,
together with its analogue for the de Rham algebra of a compact Lefschetz symplectic manifold, will be discussed in the second part of this paper.
The paper is organized as follows:
* the origin of the data (i)-(iii) in Definition 1.1 of an $`F_{\mathrm{}}`$-manifold is explained via the deformation theory. Here we use only the $`L_{\mathrm{}}`$-component of a $`G_{\mathrm{}}`$-structure. The main technical tool is a modified version of the classical deformation functor which is proved to be non-obstructed. This part of the story is, probably, of independent interest.
* we use a homotopy techique to explain the origin of the data (iv)-(v) in Definition 1.1, and to prove Theorems A-D.
* we give second proofs of the main claims of this paper using perturbative solutions of algebro-differential equations.
## 2 Deformation functors
### 2.1. Odd Lie superalgebras.
Let $`k`$ be a field with characteristic $`2`$. An odd Lie superalgebra over $`k`$ is a vector superspace $`𝔤=𝔤_{\stackrel{~}{0}}𝔤_{\stackrel{~}{1}}`$ equipped with an odd $`k`$-linear map \[Ma2\]
$$\begin{array}{cccc}\hfill []:& 𝔤𝔤& & 𝔤\\ & ab& & [ab],\end{array}$$
which satisfies the following conditions
* odd skew-symmetry: $`[ab]=(1)^{(\stackrel{~}{a}+1)(\stackrel{~}{b}+1)}[ba]`$,
* odd Jacobi identity:
$$[a[bc]]=[[ab]c]+(1)^{(\stackrel{~}{a}+1)(\stackrel{~}{b}+1)}[b[ac]],$$
for all $`a,b,c𝔤_{\stackrel{~}{0}}𝔤_{\stackrel{~}{1}}`$.
The parity change functor transforms this structure into the usual Lie superalgebra brackets, $`[,]`$, on $`\mathrm{\Pi }𝔤`$. Thus odd Lie superlgebras are nothing but Lie superalgebras in the “awkwardly” chosen $`_2`$-grading. For this reason we sometimes omit the prefix odd, and treat $`(𝔤,d,[])`$ and $`(\mathrm{\Pi }𝔤,d,[,])`$ as different representations of one and the same object. One advantage of of working with $`[]`$ rather than with the usual Lie brackets $`[,]`$ is that this awkward $`_2`$ grading induces the correct, for our purposes, $`_2`$-grading on the associated cohomology supermanifold (see below). Another advantage will become clear below, when we introduce on $`𝔤`$ one more algebraic structure (an even, in this awkward $`_2`$-grading, associative product) making $`𝔤`$ into a Gerstenhaber algebra (cf. \[Ma1, Ma2\]).
### 2.2. Cohomology as a formal supermanifold.
A data $`(𝔤,[],d)`$ with $`(𝔤,[])`$ being a Lie superalgebra and
$$d:𝔤𝔤$$
an odd $`k`$-linear map satisfying
$`d[ab]`$ $`=`$ $`[dab](1)^{\stackrel{~}{a}}[adb]`$
is called a differential Lie superalgebra, or shortly dLie-algebra. This triple $`(𝔤,[],d)`$ is often abbreviated to $`𝔤`$.
The cohomology of $`𝔤`$,
$$𝐇(𝔤):=𝖪erd/𝖨md,$$
inherits the structure of Lie superalgebra. We always assume in this paper that $`𝐇(𝔤)`$, which we often abbreviate to $`𝐇`$, is a finite dimensional superspace, say $`dim𝐇=p|q`$. Let $`\{[e_i],i=1,\mathrm{},p+q\}`$ be a basis consisting of homogeneous elements with parity denoted by $`\stackrel{~}{i}`$, and $`\{t^i,i=1,\mathrm{},p+q\}`$ the associated dual basis in $`𝐇^{}`$. The supercommutative ring, $`k[[t^1,\mathrm{},t^{p+q}]]`$, of formal power series will be abbreviated to $`k[[t]]`$. The (purely) notational advantage of working with $`k[[t]]`$ rather than with the invariantly defined object $`^{}𝐇^{}`$ is that we shall want viewing
* $`𝐇`$ as a smooth formal pointed $`(p|q)`$-dimensional supermanifold denoted (to emphasize this change of thought) by $``$ or $`_𝔤`$,
* $`\{t^i\}`$ as linear coordinates on $``$,
* $`k[[t]]`$ as the space of global sections of the structure sheaf, $`𝒪_{}`$, on $``$.
The ideal sheaf of the origin, $`0`$, will be denote by $`I`$.
There is a canonical map
$$\begin{array}{cccc}\hfill s:& k[[t]]𝐇& & H^0(,𝒯_{})\\ & a^i(t)[e_i]& & a^i(t)\frac{}{t^i},\end{array}$$
where $`H^0(,𝒯_{})`$ stands for the space of global sections of the sheaf, $`𝒯_{}`$, of formal vector fields on $``$. There is a well defined action of $`H^0(,𝒯_{})`$ on both $`k[[t]]𝐇`$ and $`k[[t]]𝔤`$ through the first factor. If X is a formal vector field on $``$ and $`\mathrm{\Gamma }`$ is an element of $`k[[t]]𝔤`$ (or of $`k[[t]]𝐇`$), then the result of this action is denoted by $`\stackrel{}{\text{X}}\mathrm{\Gamma }`$.
Any element $`\mathrm{\Gamma }`$ in $`k[[t]]𝔤`$ (or in $`k[[t]]𝐇`$) can be uniquely decomposed,
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{[0]}+\mathrm{\Gamma }_{[1]}+\mathrm{}+\mathrm{\Gamma }_{[n]}+\mathrm{}$$
into homogeneous polynomials, $`\mathrm{\Gamma }_{[n]}`$, of degree $`n`$ in the variables $`t^i`$. The sum of the first $`n`$ terms in the above decomposition is denoted by $`\mathrm{\Gamma }_{(n)}`$, i.e. $`\mathrm{\Gamma }_{(n)}=\mathrm{\Gamma }modI^{n+1}`$.
We shall call an element $`\mathrm{\Gamma }k[[t]]𝔤`$ versal if it is even, $`\mathrm{\Gamma }modI=0`$, $`\mathrm{\Gamma }_{[1]}=\mathrm{\Gamma }modI^2I𝖪erd`$ and $`\mathrm{\Gamma }_{[1]}mod𝖨md=_{i=1}^{p+q}t^i[e_i]`$.
The sheaf $`𝒯_{}`$ comes canonically equipped with a flat torsion-free affine connection $``$ whose horizontal sections are, by definition, the linear span of $`s([e_i])`$, $`i=1,\mathrm{},p+q`$, i.e.
$$\text{X}=0$$
if and only of $`s^1(\text{X})`$ is a “constant” (independent of $`t^i`$) element in $`k[[t]]𝐇`$. This connection memorizes the origin of $``$ as a vector superspace.
It will be important, in this paper, to ignore sometimes the flat structure and view $``$ only as a smooth formal supermanifold with a distinguished point $`0`$ but no preferred coordinate system. To avoid possible confusion, we adopt from now on this latter viewpoint unless the flat connection $``$ is explicitly mentioned.
### 2.2.1. $``$-grading.
Differential Lie superalgebras $`(𝔤,d,[])`$ which we often encounter in geometry have their $`_2`$-grading induced from a finer structure, $``$-grading, which is, by definition, a decomposition of $`𝔤`$ into a direct sum,
$$𝔤=\underset{i}{}𝔤^i,$$
with the following consistency conditions
* $`d𝔤^i𝔤^{i+1}`$, and
* $`[𝔤^i𝔤^j]𝔤^{i+j1}`$.
The $`_2`$-grading associated to this structure is then simply $`𝔤_{\stackrel{~}{0}}:=_{i2}𝔤^i`$ and $`𝔤_{\stackrel{~}{1}}:=_{i2+1}𝔤^i`$.
Clearly, there is an induced $``$-grading on the cohomology Lie superalgebra, $`𝐇=_i𝐇^i(𝔤)`$, as well as on the structure sheaf, $`𝒪_{}`$, of the associated cohomology supermanifold.
### 2.3. Classical deformation functor.
One of the approaches to constructing a (versal) deformation space of a given mathematical structure $`𝒜`$ consists of the following steps (see, e.g. \[GM, Ko2, Ba\], and references therein):
* Associate to $`𝒜`$ a “controlling” differential $``$-graded Lie algebra $`(𝔤=_k𝔤^k,d,[])`$ over a field $`k`$ (which is usually $``$ or $``$).
* Define the deformation functor
$$\begin{array}{cccc}\hfill \mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}:& \left\{\begin{array}{c}\text{the category of Artin}\hfill \\ k\text{-local algebras}\hfill \end{array}\right\}& & \left\{\text{the category of sets}\right\}\end{array}$$
as follows
$$\mathrm{𝖣𝖾𝖿}_𝔤^0()=\{\mathrm{\Gamma }(𝔤m_{})^2d\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}/\mathrm{exp}(𝔤m_{})^1,$$
where $`m_{}`$ is the maximal ideal of the Artin algebra $``$, the latter is viewed as a $``$-graded algebra concentrated in degree zero (so that $`(𝔤m_{})^i=𝔤^im_{}`$), and the quotient is taken with respect to the following representation of the gauge group $`\mathrm{exp}(𝔤m_{})^1`$,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{ad_g}1}{\mathrm{a}d_g}dg,g(𝔤m_{})^1,$$
where $`\mathrm{a}d`$ is just the usual internal automorphism of $`𝔤`$, $`\mathrm{a}d_g\mathrm{\Gamma }:=[g\mathrm{\Gamma }]`$.
* Try to represent the deformation functor by a topological (pro-Artin) algebra $`𝒪_{}`$ so that
$$\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}()=\text{Hom}_{\mathrm{c}ont}(𝒪_{},).$$
This associates to the mathematical structure $`𝒜`$ the formal moduli space $``$ whose “ring of functions” is $`𝒪_{}`$.
In geometry, one often continues with a fourth step by constructing a cohomological splitting of $`𝔤`$ and applying the Kuranishi method \[Ku, GM\] to represent versally the deformation functor by the ring of analytic (rather than formal) functions on the Kuranishi space.
The tangent space, $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}(k[\epsilon ]/\epsilon ^2)`$, to the functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$ is isomorphic to the cohomology group $`𝐇^2(𝔤)`$ of the complex $`(𝔤,d)`$. If one extends in the obvious way the above deformation functor to the category of arbitrary $``$-graded $`k`$-local Artin algebras (which may not be concentrated in degree 0), one gets the functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ with the tangent space isomorphic to the full cohomology group $`_i𝐇^i(𝔤)`$.
When working with the extended deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ it is often no loss of essential information to forget the $``$-grading on $`𝔤`$ and keep only the associated $`_2`$-grading. One gets then the following equivalent definition of $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$:
$$\begin{array}{cccc}\hfill \mathrm{𝖣𝖾𝖿}_𝔤^{}:& \left\{\begin{array}{c}\text{the category of Artin}\hfill \\ k\text{-local superalgebras}\hfill \end{array}\right\}& & \left\{\text{the category of sets}\right\}\end{array}$$
$$\mathrm{𝖣𝖾𝖿}_𝔤^{}():=\{\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}d\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}/\mathrm{exp}(𝔤m_{})_{\stackrel{~}{1}}.$$
This functor is representable by a smooth formal moduli space $``$ if there exists a versal (in the sense of Sect. 2.2) solution,
$$\mathrm{\Gamma }=\underset{a}{}e_at^a+\underset{a_1,a_2}{}\mathrm{\Gamma }_{a_1a_2}t^{a_1}t^{a_2}+\mathrm{}(𝔤k[[t]])_{\stackrel{~}{0}},$$
(1)
to the so-called Maurer-Cartan equation,
$$d\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0.$$
Due to versality of $`\mathrm{\Gamma }`$, any other solution over an arbitrary Artin algebra $``$ is equivalent to this one by a base change $`k[[t]]`$.
### 2.3.1. $`L_{\mathrm{}}`$-morphisms, part I.
Let $`𝔤_1`$ and $`𝔤_2`$ be two dLie-algebras. To formulate the basic theorem of the classical deformation theory, we shall need the following notion: a sequence of linear maps
$$F_n:^n𝔤_1𝔤_2,n=1,2,\mathrm{},\stackrel{~}{F}_n=\stackrel{~}{n}+1,$$
defines a $`L_{\mathrm{}}`$-morphism from $`𝔤_1`$ to $`𝔤_2`$ if
$$dF_n(\gamma _1,\gamma _2,\mathrm{},\gamma _n)+\underset{i=1}{\overset{n}{}}\pm F_n(\gamma _1,\mathrm{},d\gamma _i,\mathrm{},\gamma _n)=$$
$$=\frac{1}{2}\underset{\genfrac{}{}{0pt}{}{k+l=n}{k,l1}}{}\frac{1}{k!l!}\underset{\sigma \mathrm{\Sigma }_n}{}\pm \left[F_k(\gamma _{\sigma _1},\mathrm{},\gamma _{\sigma _k})F_l(\gamma _{\sigma _{k+1}},\mathrm{},\gamma _{\sigma _n})\right]+$$
$$+\underset{i<j}{}\pm F_{n1}([\gamma _i\gamma _j],\gamma _1,\mathrm{},\gamma _n),$$
for arbitrary $`\gamma _1,\mathrm{},\gamma _n𝔤_1`$.
In particular, the first map $`F_1`$ is a morphism of complexes which respects the Lie brackets up to homotopy defined by the second map $`F_2`$.
A $`L_{\mathrm{}}`$-morphism $`F=\{F_n\}:𝔤_1𝔤_2`$ is called a quasi-isomorphism if its linear part $`F_1`$ induces an isomorphism, $`𝐇(𝔤_1)𝐇(𝔤_2)`$, of associated cohomology groups.
### 2.3.2. Basic Theorem of $`\mathrm{𝖣𝖾𝖿}`$ormation Theory \[Ko2\] .
An $`L_{\mathrm{}}`$-morphism $`F=\{F_n\}:𝔤_1𝔤_2`$ defines a natural transformation of the functors
$$\begin{array}{cccc}\hfill F_{}:& \mathrm{𝖣𝖾𝖿}_{𝔤_1}^{}& & \mathrm{𝖣𝖾𝖿}_{𝔤_2}^{}\\ & \mathrm{\Gamma }& & F_{}(\mathrm{\Gamma }):=_{n=1}^{\mathrm{}}\frac{1}{n!}F_n(\mathrm{\Gamma },\mathrm{},\mathrm{\Gamma }),\end{array}$$
i.e. if $`\mathrm{\Gamma }`$ is a solution to Maurer-Cartan equations in $`(𝔤_1m_{})_{\stackrel{~}{0}}`$, then $`F_{}(\mathrm{\Gamma })`$ is a solution to Maurer-Cartan equations in $`(𝔤_2m_{})_{\stackrel{~}{0}}`$.
Moreover, if $`F`$ is a quasi-isomorphism, then $`F_{}`$ is an isomorphism.
### 2.3.3. Corollary.
If a dLie-algebra $`𝔤`$ is quasi-isomorphic to an Abelian dLie-algebra, then $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ is versally representable by a smooth formal pointed supermanifold $`_𝔤`$ (the cohomology supermanifold of $`𝔤`$).
Proof. If $`𝔥`$ is an Abelian dLie-algebra, then, in the notations of Sect. 2.2,
$$\mathrm{\Gamma }=\underset{i=1}{\overset{dim𝐇(𝔥)}{}}t^i[e_i]$$
is a versal solution of Maurer-Cartan equations. Hence $`\mathrm{𝖣𝖾𝖿}_𝔥^{}`$ is representable by $`_𝔥`$.
If $`𝔤`$ is quasi-isomorphic to $`𝔥`$, the required statement follows from Theorem 2.3.2 and the isomorphism $`_𝔤=_𝔥`$. $`\mathrm{}`$
There are two remarkable examples, one dealing with extended deformations of complex structures on a Calabi-Yau manifold \[BK\] and another with extended deformations of the symplectic structure on a Lefschetz manifold \[Me1\], when the rather strong condition of Corollary 2.3.3 holds true<sup>3</sup><sup>3</sup>3One more example (of a different technical origin though of the same mirror symmetry flavor) when a naturally extended deformation problem gives rise to a smooth extended moduli space is discussed in \[Me2\].. In general, however, there will be obstructions to constructing a versal solution to the Maurer-Cartan equations, and we will have to resort to other technical means such as modifying the deformation functor as explained below in Sect. 2.4 below or further extending the deformation problem to the category of $`L_{\mathrm{}}`$-algebras.
### 2.3.4. Example (deformations of complex manifolds).
It is well known that the total space of the cotangent bundle, $`\mathrm{\Omega }_{}^1`$, to a real $`n`$-dimensional manifold $`M`$ carries a natural Poisson structure $`\{,\}`$ making the structure sheaf $`𝒪_{\mathrm{\Omega }_{}^1}`$ into a sheaf of Lie algebras. In a natural local coordinate system $`(x^a,p_a:=/x^a)`$,
$$\{f,g\}=\frac{f}{p_a}\frac{g}{x^a}\frac{f}{x^a}\frac{g}{p_a}.$$
If we now change the parity of the fibers of the natural projection $`\mathrm{\Omega }_{}^1M`$ (which is allowed since they are vector spaces), we will get an $`(n|n)`$-dimensional supermanifold $`\mathrm{\Pi }\mathrm{\Omega }_{}^1`$ equipped with a natural odd Poisson structure $`\{\}`$ making the structure sheaf $`𝒪_{\mathrm{\Pi }\mathrm{\Omega }_{}^1}`$ into a sheaf of odd Lie superalgebras. In a natural local coordinate system $`(x^a,\psi _a:=\mathrm{\Pi }/x^a)`$ on $`\mathrm{\Pi }\mathrm{\Omega }_{}^1`$,
$$\{fg\}=\frac{f}{\psi _a}\frac{g}{x^a}+(1)^{\stackrel{~}{g}(\stackrel{~}{f}+1)}\frac{g}{\psi _a}\frac{f}{x^a}.$$
The smooth functions on $`\mathrm{\Pi }\mathrm{\Omega }_{}^1`$ have a simple geometric interpretation in term of the underlying manifold $`M`$ — they are just smooth polyvector fields. Indeed, a standard power series decomposition in odd variables gives
$$f=\underset{k=0}{\overset{n}{}}\underset{a_1,\mathrm{},a_k}{}f^{a_1\mathrm{}a_k}(x)\psi _{a_1}\mathrm{}\psi _{a_k}$$
implying the isomorphism of sheaves $`𝒪_{\mathrm{\Pi }\mathrm{\Omega }_{}^1}=\mathrm{\Lambda }^{}T_{}`$, where $`T_{}`$ is the real tangent bundle to $`M`$.
Therefore, the sheaf $`\mathrm{\Lambda }^{}T_{}`$ with the $`_2`$-grading,
$$(\mathrm{\Lambda }^{}T_{})_{\stackrel{~}{0}}:=\mathrm{\Lambda }^{\mathrm{e}ven}T_{},(\mathrm{\Lambda }^{}T_{})_{\stackrel{~}{1}}:=\mathrm{\Lambda }^{\mathrm{o}dd}T_{},$$
induced from that on $`𝒪_{\mathrm{\Pi }\mathrm{\Omega }_{}^1}`$, is naturally a sheaf of odd Lie superalgebras. The odd Poisson bracket $`\{\}`$ is called, in this incarnation, the Schouten bracket and is often denoted by $`[]_{\mathrm{S}ch}`$.
If $`M`$ is a complex manifold, then the canonical odd Poison structure on the parity changed holomorphic cotangent bundle, $`\mathrm{\Pi }\mathrm{\Omega }_M^1`$, is itself holomorphic giving rise thereby to the structure of odd Lie superalgebra on the sheaf, $`\mathrm{\Lambda }^{}T_M`$, of holomorphic polyvector fields. This can be used to make the vector space
$$𝔤=\underset{k}{}𝔤^k,𝔤^k:=\underset{i+j=k}{}\mathrm{\Gamma }(M,\mathrm{\Lambda }^iT_M\mathrm{\Lambda }^j\overline{T}_M^{}),$$
into a $``$-graded differential algebra by taking $`\overline{}`$, the $`(0,1)`$ part of the de Rham operator, as a differential, and the map,
$$\begin{array}{cccc}\hfill []:& \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_1}T_M\mathrm{\Lambda }^{j_1}\overline{T}_M^{})\times \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_2}T_M\mathrm{\Lambda }^{j_2}\overline{T}_M^{})& & \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_1+i_21}T_M\mathrm{\Lambda }^{j_1+j_2}\overline{T}_M^{})\\ & \text{X}_1\overline{w}_1\times \text{X}_2\overline{w}_2& & [\text{X}_1\overline{w}_1\text{X}_2\overline{w}_2],\end{array}$$
given by
$$[\text{X}_1\overline{w}_1\text{X}_2\overline{w}_2]:=(1)^{\stackrel{~}{j}_1\stackrel{~}{i}_2}[\text{X}_1\text{X}_2]_{\mathrm{S}ch}(\overline{w}_1\overline{w}_2).$$
as the (odd) Lie brackets.
The importance of this dLie-algebra stems from the fact that the associated deformation functors $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$ and $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ describe, respectively, ordinary and extended deformations of the given complex structure on a smooth manifold $`M`$. Indeed, a complex structure on a real $`2n`$-dimensional manifold $`M`$ is a decomposition, $`T_{}=T_M\overline{T}_M`$, of the complexified real tangent bundle into a direct sum of complex integrable distributions, $`T_M`$ and its complex conjugate $`\overline{T}_M`$. Another decomposition like that, $`T_{}=T_M^{}\overline{T}_M^{}`$, can be described in terms of the original complex structure by the graph, $`\overline{T}_M^{}`$, of a linear map $`\mathrm{\Gamma }:\overline{T}_MT_M`$, i.e. by an element $`\mathrm{\Gamma }𝔤^2`$. The integrability of $`T_M^{}`$ amounts then to the Maurer-Cartan equation in $`𝔤^2`$,
$$\overline{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]_{\mathrm{S}ch}=0.$$
By solving (if possible) the above equation in the full Lie algebra $`𝔤`$ rather than in its subalgebra $`𝔤^2`$ (and taking the quotient by the gauge group describing equivalent deformations), one gets a so-called extended complex structure on $`M`$ whose geometric meaning is not yet fully understood. It is understood \[Ko1\], however, that this structure does have an important Mirror Symmetry aspect (at least for Calabi-Yau manifolds): Kontsevich noticed that his Formality Theorem \[Ko2\] identifies the moduli space of extended complex structures on a given complex manifold $`M`$ with the moduli space of $`A_{\mathrm{}}`$-deformations of the derived category of coherent sheaves on $`M`$, mirror counterpart of the conjectured Fukaya category built out of a dual complex manifold $`\stackrel{~}{M}`$.
If $`M`$ is a Calabi-Yau manifold, then, as was shown by Barannikov and Kontsevich \[BK\], the Maurer-Cartan equations in $`𝔤`$ admit a versal solution of the form (1) implying that the moduli space of extended deformations of complex structures on $`M`$ is smooth, and is isomorphic<sup>4</sup><sup>4</sup>4Strictly speaking, this isomorphism holds true in the category of formal manifolds, in which we work in this paper. It is no problem to choose the power series (1) convergent thereby inducing on the extended moduli space a smooth analytic structure. The latter is then analytically isomorphic to an open neighbourhood of zero in $`𝐇`$ which we denote by the same symbol $``$ (and continue doing this every time the analyticity aspect emerges). to $``$. In fact they have shown much more \[BK, Ba\]: in this case $``$ has an induced structure of Frobenius manifold which conjecturally coincides with the Frobenius manifold structure on $`H_{}(\stackrel{~}{M},)`$ constructed via the Gromov-Witten invariants of the dual Calabi-Yau manifold $`\stackrel{~}{M}`$. Barannikov \[Ba\] checked this conjecture for projective complete intersections Calabi-Yau manifolds.
For a general complex manifold $`M`$, the (extended) deformations are obstructed and the Maurer-Cartan equations associated to $`𝔤`$ have no versal solutions of the form (1). It is one of the main tasks of this paper to understand what happens to the Barannikov-Kontsevich’s Frobenius structure on $`𝒯_{}`$ in the presence of obstructions.
### 2.3.5. Example (deformations of Poisson and symplectic manifolds).
It is well known that a 2-vector field, $`\nu _0\mathrm{\Gamma }(M,\mathrm{\Lambda }^2T_{})`$, defines a Poisson structure,
$$\{f,g\}=\nu _0(dfdg),f,g𝒪_M,$$
on a smooth real $`n`$-dimensional manifold $`M`$ if and only if
$$[\nu _0\nu _0]_{\mathrm{S}ch}=0.$$
Then a deformed 2-vector field, $`\nu _0+\nu \mathrm{\Gamma }(M,\mathrm{\Lambda }^2T_{})`$, is again a Poisson structure if and only if $`\nu `$ satisfies the Maurer-Cartan equation,
$$d\nu +\frac{1}{2}[\nu \nu ]_{\mathrm{S}ch}=0,$$
in the differential $``$-graded algebra
$$\left(𝔤=\underset{i=0}{\overset{n}{}}\mathrm{\Gamma }(M,\mathrm{\Lambda }^iT_{}),[]_{\mathrm{S}ch},d:=[\nu _0\mathrm{}]_{\mathrm{S}ch}\right).$$
Hence the associated deformation functors $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$/$`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ describe (extended) deformations of the given Poisson structure $`\nu _0`$ on $`M`$.
For generic $`\nu _0`$, the associated cohomology group $`_i𝐇^i(𝔤)`$ may not be finite-dimensional even for compact manifolds. If, however, $`\nu _0`$ comes from a symplectic $`2`$-form $`\omega `$ on $`M`$, the situation is very different. In this case one may use the “lowering indices map” $`\omega :T_{}\mathrm{\Omega }_{}^1`$ to identify $`(𝔤,d)`$ with the de Rham complex on $`M`$ — it is not hard to check that this map sends the differential $`[\nu _0\mathrm{}]_{\mathrm{S}ch}`$ on $`\mathrm{\Lambda }^{}T_{}`$ into the usual de Rham differential on $`\mathrm{\Omega }_{}^{}`$. The (odd) Lie brackets induced on $`\mathrm{\Omega }_{}^{}`$ from the Schouten brackets on $`\mathrm{\Lambda }^{}T_{}`$ we denote by $`[]_\omega `$ to emphasize its dependence on the symplectic structure. Hence the deformation functors $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$/$`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ associated with the dLie-algebra
$$\left(𝔤=\underset{i=0}{\overset{n}{}}\mathrm{\Gamma }(M,\mathrm{\Omega }_{}^i),[]_\omega ,d=\text{de Rham differential}\right)$$
describe (extended) deformations of a symplectic structure $`\omega `$ on $`M`$. Its cohomology $`𝐇`$ is nothing but the de Rham cohomology of $`M`$.
A compact symplectic manifold $`(M,\omega )`$ is called Lefschetz if the the natural cup product on the de Rham cohomology,
$$[\omega ^k]:H^{mk}(M,)H^{m+k}(M,)$$
is an isomorphism for any $`km:=\frac{1}{2}dimM`$. This class of manifolds (which includes the class of Kähler manifolds by the Hard Lefschetz Theorem) is of interest to us for the extended deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ associated with an arbitrary Lefschetz symplectic manifold is non-obstructed and is representable by a smooth moduli space isomorphic to $``$; moreover, this moduli space of “extended symplectic structures” is always a Frobenius manifold \[Me1\]. This result is more than parallel to the Barannikov-Kontsevich’s construction of extended moduli spaces/Frobenius structures for Calabi-Yau manifolds — it is just another example when their beautiful machinery works (see a nice exposition of Manin \[Ma2\]).
The extended deformation functor associated with a generic compact symplectic manifold seems to be obstructed, and one should employ different techniques (see below) to study geometric structures induced on moduli spaces of extended symplectic structures.
### 2.3.6. Example (deformations of holomorphic vector bundles).
Let $`EM`$ be a holomorphic vector bundle on a complex $`n`$-dimensional manifold $`M`$. There is an associated differential $``$-graded Lie algebra
$$\left(𝔤=\mathrm{\Gamma }(M,\mathrm{E}ndE\overline{\mathrm{\Omega }}_M^{}[1]),[],\overline{}\right)$$
with the Lie brackets,
$$\begin{array}{cccc}\hfill []:& \mathrm{\Gamma }(M,\mathrm{E}ndE\overline{\mathrm{\Omega }}_M^{i_11})\times \mathrm{\Gamma }(M,\mathrm{E}ndE\overline{\mathrm{\Omega }}_M^{i_21})& & \mathrm{\Gamma }(M,\mathrm{E}ndE\overline{\mathrm{\Omega }}_M^{(i_1+i_21)1})\\ & A_1\overline{w}_1\times A_2\overline{w}_2& & [A_1\overline{w}_1A_2\overline{w}_2],\end{array}$$
given by
$$[A_1\overline{w}_1A_2\overline{w}_2]:=(A_1A_2A_2A_1)\left(\overline{w}_1\overline{w}_2\right).$$
The deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{\mathrm{\hspace{0.17em}0}}`$ associated with this algebra describes deformations of the holomorphic structure in the vector bundle $`E`$. It is tempting to view its extension $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ as a tool for studying extended deformations, but we reserve this role for the functor $`\mathrm{𝖣𝖾𝖿}_{\mathrm{}}^{}`$ associated with a larger differential algebra constructed in 3.1.5 below.
In general, all these functors are obstructed.
One may combine this differential Lie algebra (or its extension 3.1.5) together with the one of Example 2.3.4 into their natural semi-direct product to study joint (extended) deformations of the pair $`EM`$.
### 2.4. $`L_{\mathrm{}}`$-algebras.
These algebras will play only an auxiliary, purely technical, role in this paper.
By definition, a strong homotopy Lie algebra, or shortly $`L_{\mathrm{}}`$-algebra, is a vector superspace $`𝔥`$ equipped with linear maps,
$$\begin{array}{ccccc}\hfill \mu _k:& \mathrm{\Lambda }^k𝔥& & 𝔥& \\ & v_1\mathrm{}v_k& & \mu _k(v_1,\mathrm{},v_k),& k1,\stackrel{~}{\mu }_k=\stackrel{~}{k},\end{array}$$
satisfying, for any $`n1`$ and arbitrary $`v_1,\mathrm{},v_n𝔥_{\stackrel{~}{0}}𝔥_{\stackrel{~}{1}}`$, the following higher order Jacobi identities,
$$\underset{k+l=n+1}{}\underset{\sigma Sh(k,n)}{}(1)^{\stackrel{~}{\sigma }+k(l1)}e(\sigma ;v_1,\mathrm{},v_n)\mu _l(\mu _k(v_{\sigma (1)},\mathrm{},v_{\sigma (k)}),v_{\sigma (k+l)},\mathrm{},v_{\sigma (n)})=0,$$
where $`Sh(k,n)`$ is the set of all permutations $`\sigma :\{1,\mathrm{},n\}\{1,\mathrm{},n\}`$ which satisfy $`\sigma (1)<\mathrm{}<\sigma (k)`$ and $`\sigma (k+1)<\mathrm{}<\sigma (n)`$. The symbol $`e(\sigma ;v_1,\mathrm{},v_n)`$ (which we abbreviate from now on to $`e(\sigma )`$) stands for the Koszul sign defined by the equality
$$v_{\sigma (1)}\mathrm{}v_{\sigma (n)}=(1)^{\stackrel{~}{\sigma }}e(\sigma )v_1\mathrm{}v_n,$$
$`\stackrel{~}{\sigma }`$ being the parity of the permutation $`\sigma `$.
The $``$-graded version of this definition would require $`\mu _k`$ to be homogeneous of degree $`2k`$.
This notion as well as the associated notion of $`A_{\mathrm{}}`$-algebra (reminded below) are due to Stasheff \[S\].
The first three higher order Jacobi identities have the form
$`d^2=0`$,
$`d[v_1,v_2]=[dv_1,v_2]+(1)^{\stackrel{~}{v}_1}[v_1,dv_2]`$,
$`[[v_1,v_2],v_3]+(1)^{(\stackrel{~}{v}_1+\stackrel{~}{v}_2)\stackrel{~}{v}_3}[[v_3,v_1],v_2]+(1)^{\stackrel{~}{v}_1(\stackrel{~}{v}_2+\stackrel{~}{v}_3)}[[v_2,v_3],v_1]`$ = $`d\mu _3(v_1,v_2,v_3)\mu _3(dv_1,v_2,v_3)(1)^{\stackrel{~}{v}_1}\mu _3(v_1,dv_2,v_3)(1)^{\stackrel{~}{v}_1+\stackrel{~}{v}_2}\mu _3(v_1,v_2,dv_3)`$,
where we denoted $`dv_1:=\mu _1(v_1)`$ and $`[v_1,v_2]:=\mu _2(v_1,v_2)`$.
Therefore, $`L_{\mathrm{}}`$-algebras with $`\mu _k=0`$ for $`k3`$ are nothing but the usual differential Lie superalgebras with the differential $`\mu _1`$ and the Lie bracket $`\mu _2`$. If, furthermore, $`\mu _1=0`$, one gets the class of usual Lie superalgebras.
### 2.4.1. Odd $`L_{\mathrm{}}`$-algebras.
To make the above picture consistent with the choices made in Sect. 2.1, we should change the parity of $`𝔥`$. Hence we shall work from now on with $`𝔤:=\mathrm{\Pi }𝔥`$ and denote
$$\mu _n(v_1v_2\mathrm{}v_n):=\mathrm{\Pi }\mu _n(\mathrm{\Pi }v_1,\mathrm{\Pi }v_2,\mathrm{},\mathrm{\Pi }v_n),v_1,\mathrm{},v_n𝔤,$$
for all $`n1`$. This change of grading also unveils, through the following three observations, a rather compact image of the $`L_{\mathrm{}}`$-structure itself:
* The vector superspace $`^{}𝔤=_{n=1}^{\mathrm{}}^n𝔤`$ has a natural structure of cosymmetric coalgebra,
$$\mathrm{\Delta }(w_1\mathrm{}w_n)=\underset{i=1}{\overset{n}{}}\underset{\sigma Sh(i,n)}{}e(\sigma )\left(w_{\sigma (1)}\mathrm{}w_{\sigma (i)}\right)\left(w_{\sigma (i+1)}\mathrm{}w_{\sigma (n)}\right).$$
* Every coderivation of this coalgebra, i.e. an odd map $`Q:^{}𝔤^{}𝔤`$ satisfying $`\mathrm{\Delta }Q=\left(Q\text{Id}+\text{Id}Q\right)\mathrm{\Delta }`$, is equivalent to an arbitrary series of odd linear maps, $`\mu _n:^n𝔤𝔤`$.
* A codifferential $`Q=\{\mu _{}\}`$ is a differential, i.e. $`Q^2=0`$, if and only if $`\mu _n`$ satisfy the higher order Jacobi identities.
In conclusion, an (odd) $`L_{\mathrm{}}`$-structure on $`𝔤`$ is equivalent to a codifferential on $`(^{}𝔤,\mathrm{\Delta })`$.
### 2.4.2. $`L_{\mathrm{}}`$-morphisms, part II.
Given two $`L_{\mathrm{}}`$-algebras, $`(𝔤,\mu _{})`$ and $`(𝔤^{},\mu _{}^{})`$. A $`L_{\mathrm{}}`$-morphism $`F`$ from the first one to the second is, by definition, a differential coalgebra homomorphism
$$F:(^{}𝔤,\mathrm{\Delta },Q)(^{}𝔤^{},\mathrm{\Delta },Q^{}),$$
i.e. a linear map $`F:^{}𝔤^{}𝔤^{}`$ satisfying $`(FF)\mathrm{\Delta }=\delta ^{}F`$ and $`FQ=Q^{}F`$. The first of these equations says that $`F`$ is completely determined by a set of linear maps $`F_n:^n𝔤𝔤^{}`$ of parity $`\stackrel{~}{n}+\stackrel{~}{1}`$, while the second one imposes on these $`F_n^{}`$ a sequence of linear equations. If both input and output of $`F`$ are usual differential Lie superalgebras, these equations are presicely the ones written down in Sect. 2.3.1.
A $`L_{\mathrm{}}`$-morphism $`F:(𝔤,\mu _{})(𝔤^{},\mu _{}^{})`$ is called a quasi-isomorphism if its first component $`F_1:𝔤𝔤^{}`$ induces an isomorphism between cohomology groups of complexes $`(𝔤,\mu _1)`$ and $`(𝔤^{},\mu _1^{})`$. It is called a homotopy of the $`L_{\mathrm{}}`$-algebras, if $`F_1:𝔤𝔤^{}`$ is an isomorphism of underlying vector graded superspaces.
If the $`L_{\mathrm{}}`$-morphism $`F:(𝔤,\mu _{})(𝔤^{},\mu _{}^{})`$ is a quasi-isomorphism, then, as was proved in \[Ko2\], there exists a $`L_{\mathrm{}}`$-morphism $`F^{}:(𝔤^{},\mu _{}^{})(𝔤,\mu _{})`$ which induces the inverse isomorphism between cohomology groups of complexes $`(𝔤,\mu _1)`$ and $`(𝔤^{},\mu _1^{})`$.
Two $`L_{\mathrm{}}`$-morphisms, $`F,G:(𝔤,\mu _{})(𝔤^{},\mu _{}^{})`$, are said to be homotopy equivalent if there is an odd linear map $`h:^{}𝔤𝔤^{}`$ such that
$$\mathrm{\Delta }^{}h=Fh+hG,\mathrm{a}ndF=G+Q^{}h+hQ.$$
This map is completely determined by a set of linear maps, $`\{h_n:^n𝔤𝔤^{},\stackrel{~}{h_n}=\stackrel{~}{n},n=1,2,\mathrm{}\}`$, and is called a homotopy of morphisms. The resulting homotopy relation on the set of all $`L_{\mathrm{}}`$-morphisms from $`(𝔤,\mu _{})`$ to $`(𝔤^{},\mu _{}^{})`$ is an equivalence relation.
### 2.4.3. A geometric interpretation of a $`L_{\mathrm{}}`$-algebra $`𝔤`$.
The dual of the free cocommutative coalgebra $`^{}𝔤`$ can be identified with the algebra of formal power series on the vector superspace $`𝔤`$ viewed as a formal pointed supermanifold (to emphasize this change of thought we denote the supermanifold structure on $`𝔤`$ by $`M_𝔤`$). With this identification, the $`L_{\mathrm{}}`$-structure $`\mu _{}`$ on $`𝔤`$, that is the codifferential $`Q`$ on $`^{}𝔤`$, goes into an odd vector field $`Q`$ on $`M_𝔤`$ satisfying \[Ko2\]
* $`Q^2=0`$,
* $`QII`$,
where $`I`$ is the ideal of the distinguished point, $`0M_𝔤`$. (An odd vector field on a formal pointed superspace satisfying the above two conditions is usually called homological.)
A $`L_{\mathrm{}}`$-morphism $`F`$ between two $`L_{\mathrm{}}`$-algebras $`(𝔤,\mu _{})`$ and $`(𝔤^{},\mu _{}^{})`$ is nothing but a $`Q`$-equivariant map between the associated formal pointed homological supermanifolds, $`(M_𝔤,Q)`$ and $`(M_𝔤^{},Q^{})`$.
### 2.5. A modified deformation functor.
For a general dLie-algebra $`𝔤`$, the classical deformation functor $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ is not representable by a smooth versal moduli space. At best one can use Kuranishi technique to represent $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ by a singular analytic space. There is, however, a simple geometric way to keep track of versality and smoothness. The idea is as follows.
First, we extend the input category used in the construction of $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$ in Sect. 2.3 to a category of differential Artin superalgebras whose $`\mathrm{𝖮𝖻}`$s are pairs, $`(,)`$, consisting of an Artin superalgebra $``$ together with a differential $`:`$ satisfying $`m_{}m_{}^2`$, and whose $`\mathrm{𝖬𝗈𝗋}`$s are morphisms of Artin superalgebras commuting with the differentials. A $``$-graded version of this definition would require $``$ to have degree $`+1`$.
Second, to the “controlling” differential Lie algebra $`𝔤`$ we associate a new deformation functor,
$$\begin{array}{cccc}\hfill \text{Def}_𝔤^{}:& \left\{\begin{array}{c}\text{the category of differential}\\ \text{Artin superalgebras}\end{array}\right\}& & \left\{\text{the category of sets}\right\}\\ & (,)& & \text{Def}_𝔤^{}(,)\end{array}$$
by setting
$$\text{Def}_𝔤^{}(,)=\{\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}/\mathrm{exp}(𝔤m_{})_{\stackrel{~}{1}}.$$
Here the quotient is taken with respect to the following representation of the gauge group,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{ad_g}1}{\mathrm{a}d_g}(d+\stackrel{}{})g,g(𝔤m_{})_{\stackrel{~}{1}}.$$
On the subcategory $`(,0)`$ the deformation functor $`\text{Def}_𝔤^{}`$ coincides precisely with the classical one $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$.
### 2.5.1. Remark.
If, for a derivation $`:`$, an element $`\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}`$ satisfies the equation (which we sometimes call the Master equation),
$$d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0,$$
then
$`0`$ $`=`$ $`d\left(d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+{\displaystyle \frac{1}{2}}[\mathrm{\Gamma }\mathrm{\Gamma }]\right)`$
$`=`$ $`\stackrel{}{}d\mathrm{\Gamma }+[d\mathrm{\Gamma }\mathrm{\Gamma }]`$
$`=`$ $`\stackrel{}{}d\mathrm{\Gamma }[\stackrel{}{}\mathrm{\Gamma }\mathrm{\Gamma }]{\displaystyle \frac{1}{2}}\left[[\mathrm{\Gamma }\mathrm{\Gamma }]\mathrm{\Gamma }\right]`$
$`=`$ $`\stackrel{}{}\left(d\mathrm{\Gamma }+{\displaystyle \frac{1}{2}}\left[\mathrm{\Gamma }\mathrm{\Gamma }\right]\right)`$
$`=`$ $`\stackrel{}{}^{\mathrm{\hspace{0.17em}2}}\mathrm{\Gamma },`$
motivating our assumption above that $``$ is a differential in $``$ rather than merely a derivation.
### 2.5.2. $`L_{\mathrm{}}`$-extension of $`\text{Def}^{}`$.
This extension will be used later only as a technical tool in the study of $`\text{Def}_𝔤^{}`$ for usual differential Lie superalgebras $`𝔤`$.
Given a $`L_{\mathrm{}}`$-algebra $`(𝔤,\mu _{})`$, we define,
$$\begin{array}{cccc}\hfill \text{Def}_𝔤^{}:& \left\{\begin{array}{c}\text{the category of differential}\\ \text{Artin superalgebras}\end{array}\right\}& & \left\{\text{the category of sets}\right\}\\ & (,)& & \text{Def}_𝔤^{}(,)\end{array}$$
by setting
$$\text{Def}_𝔤^{}(,)=\{\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}\stackrel{}{}\mathrm{\Gamma }=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n(n+1)/2}}{n!}\mu _n(\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma })\}/$$
Here the quotient is taken with respect to the gauge equivalence, $``$, which is best described using the following geometric model of the $`\text{Def}\mathrm{o}rmation`$ functor.
In $`𝖢ategory^{op}`$, both the differential Artin superalgebra, $`(,)`$, and the $`L_{\mathrm{}}`$-algebra, $`(𝔤,\mu _{})`$, are represented by formal pointed analytic homological superspaces, $`(M_{},0,)`$ and, respectively, $`(M_𝔤,0,Q)`$. Then the set
$$S=\{\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}\stackrel{}{}\mathrm{\Gamma }=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n(n+1)/2}}{n!}\mu _n(\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma })\}$$
is just the set of all formal maps of pointed supermanifolds,
$$\mathrm{\Gamma }:(M_{},0)(M_𝔤,0),$$
satisfying the equivariency condition
$$d\mathrm{\Gamma }()=\mathrm{\Gamma }^{}(Q).$$
The latter is precisely the ($`L_{\mathrm{}}`$-generalization of) the Master equation.
Both superpaces, $`(M_{},0,)`$ and $`(M_𝔤,0,Q)`$, are foliated by integrable distributions,
$`𝒟_{}`$ $`:=`$ $`\{XTM_{}X=[,Y]\mathrm{f}orsomeYTM_{}\},`$
$`𝒟_Q`$ $`:=`$ $`\{X^{}TM_𝔤X^{}=[Q,Y^{}]\mathrm{f}orsomeY^{}TM_𝔤\},`$
and $`d\mathrm{\Gamma }(𝒟_{})\mathrm{\Gamma }^{}(𝒟_Q)`$ for any $`\mathrm{\Gamma }S`$. Hence any such $`\mathrm{\Gamma }`$ defines a map, $`\widehat{\mathrm{\Gamma }}`$, through the following commutative diagram,
We say that two elements in $`S`$ are gauge equivalent, $`\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$, if $`\widehat{\mathrm{\Gamma }}_1=\widehat{\mathrm{\Gamma }}_2`$. Infinitesimally, the gauge equivalence is given by
$$\mathrm{\Gamma }\mathrm{\Gamma }+(d+\stackrel{}{})g\underset{n=2}{\overset{\mathrm{}}{}}\frac{(1)^{n(n+1)/2}}{(n1)!}\mu _n(g\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }),g(𝔤m_{})_{\stackrel{~}{1}}.$$
### 2.5.4. Remark.
If, for a derivation $`:`$, an element $`\mathrm{\Gamma }(𝔤m_{})_{\stackrel{~}{0}}`$ satisfies the Master equation,
$$\stackrel{}{}\mathrm{\Gamma }=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n(n+1)/2}}{n!}\mu _n(\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }),$$
then, using the higher Jacobi identities (as in 2.5.1), one gets an implication,
$$\stackrel{}{}^{\mathrm{\hspace{0.17em}2}}\mathrm{\Gamma }=0.$$
### 2.5.5. Basic Theorem of Deformation Theory.
Let $`(𝔤_1,Q_1)`$ and $`(𝔤_2,Q_2)`$ be two $`L_{\mathrm{}}`$-algebras (in partucilar, dLie-algebras). An $`L_{\mathrm{}}`$-morphism $`F=\{F_n\}:𝔤_1𝔤_2`$ defines a natural transformation of the functors,
$$\begin{array}{cccc}\hfill F_{}:& \text{Def}_{𝔤_1}^{}& & \text{Def}_{𝔤_2}^{}\\ & \mathrm{\Gamma }& & F_{}(\mathrm{\Gamma }):=_{n=1}^{\mathrm{}}\frac{1}{n!}F_n(\mathrm{\Gamma },\mathrm{},\mathrm{\Gamma }).\end{array}$$
Moreover, if $`F`$ is a quasi-isomorphism, then $`F_{}`$ is an isomorphism.
Proof. Assume $`\mathrm{\Gamma }(𝔤_1m_{})_{\stackrel{~}{0}}`$ satisfies the Master equation,
$$d\mathrm{\Gamma }()=\mathrm{\Gamma }^{}(Q_1).$$
The $`L_{\mathrm{}}`$-morphism $`F`$, when viewed as map, $`M_{𝔤_1}M_{𝔤_2}`$, of pointed formal manifold, satisfies,
$$dF(Q_1)=F^{}(Q_2).$$
Hence $`F_{}(\mathrm{\Gamma })`$, which the same as $`F\mathrm{\Gamma }`$, obviously satisfies the Master equation in $`(𝔤_2,Q_2)`$. $`\mathrm{}`$
### 2.5.6. Smoothness Theorem.
The deformation functor $`\text{Def}^{}`$ is unobstructed, i.e. for any dLie-algebra $`𝔤`$ with finite-dimensional cohomology $`𝐇(𝔤)`$, the functor $`\text{Def}_𝔤^{}`$ is versally representable by a smooth pointed formal $`dim𝐇(𝔤)`$-dimensional homological supermanifold $`(_𝔤,)`$.
Moreover, the diffeomorphism class of $`(_𝔤,)`$ is an invariant of $`𝔤`$.
Proof. It is enough to show that
* there exists a versal element $`\mathrm{\Gamma }k[[t]]𝔤`$ and a differential $`:k[[t]]k[[t]]`$ satisfying the Master equation
$$d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0,$$
(2)
* the differential $``$, when viewed as a vector field on the cohomological supermanifold $`_𝔤`$, is an invariant of $`𝔤`$.
Unless $`𝔤`$ is formal, there is no quasi-isomorphism from $`(𝔤,[],d)`$ to its cohomology $`(𝐇(𝔤),[]_{\mathrm{i}nd},0)`$. However, there always exists a $`L_{\mathrm{}}`$-structure, $`\{\mu _{},\mathrm{w}ith\mu _1=0\}`$, on $`𝐇(𝔤)`$ which is quasi-isomorphic, via some $`L_{\mathrm{}}`$-morphism $`F`$, to $`(𝔤,[],d)`$. Moreover, this structure is defined uniquely up to a homotopy.
Setting $`\mathrm{\Gamma }_{[1]}=_it^i[e_i]`$, in the notations of Sect. 2.2, we define a derivation, $`:k[[t]]k[[t]]`$, by the formula
$$\stackrel{}{}\mathrm{\Gamma }_{[1]}=\underset{n=2}{\overset{\mathrm{}}{}}\frac{(1)^{n(n+1)/2}}{n!}\mu _n(\mathrm{\Gamma }_{[1]}\mathrm{}\mathrm{\Gamma }_{[1]}).$$
By Remark 2.5.4, this derivation satisfies $`^2=0`$. Hence $`(\mathrm{\Gamma }_{[1]},)`$ is a versal solution of the Master equation in $`(𝐇(𝔤),\mu _{})`$, while $`(F_{}(\mathrm{\Gamma }_{[1]}),)`$ is, by Theorem 2.5.5, a versal solution of the Master equation in $`(𝔤,[],d)`$. This proves claim (i).
The $`L_{\mathrm{}}`$-structure $`\{\mu _{}\}`$ on $`𝐇^{p|q}`$ is well-defined only up to a homotopy, $`\{\eta _{(n)}:^n𝐇𝐇,n2,\stackrel{~}{\eta _n}=\stackrel{~}{0}\}`$. Is is easy to check that a homotopy change of the induced $`L_{\mathrm{}}`$-structure
$$\mu _{}\stackrel{\eta _{()}}{}\mu _{}^{},$$
does change the differential,
$$^{},$$
but in a remarkably geometric way,
$$^{}=d\eta (),$$
where $`\eta :^{p|q}^{p|q}`$ is just a formal change of coordinates
$$t^it^i=t^i+\underset{j,k}{}\pm \eta _{(2)jk}^it^jt^k+\underset{j,k,l}{}\pm \eta _{(3)jkl}^it^jt^kt^l+\mathrm{}.$$
Put another way, a homotopy change of $`\mu _{}`$ affects only the coordinate representation of the vector field $``$ on $`_𝔤`$. As a geometric entity, this is an invariant of $`𝔤`$. $`\mathrm{}`$
### 2.5.7. Corollary.
The derived category of $`L_{\mathrm{}}`$-algebras with finite-dimensional cohomology is canonically equivalent to the (purely geometric) category of pointed formal homological supermanifolds, $`(,,0)`$, with $``$ satisfying $`II^2`$, $`I`$ being the ideal of the distinguished point $`0`$.
Proof. It is well-known that each quasi-isomorphism of $`L_{\mathrm{}}`$-algebras is a homotopy equivalence. Thus the derived category of $`L_{\mathrm{}}`$-algebras is equivalent to the their homotopy category. Then the required statement follows immediately from an observation made in the proof of Theorem 2.5.6 that, for any homotopy class of $`L_{\mathrm{}}`$-algebras $`[𝔤]`$, the associated homotopy class of $`L_{\mathrm{}}`$-structures induced on the cohomology $`𝐇(𝔤)`$ is isomorphically mapped into one and the same homological manifold $`(_𝔤,,0)`$. $`\mathrm{}`$
### 2.5.8. Remarks.
(i) The origin of the vector field $``$ in Theorem 2.5.6 can be traced back to Chen’s power series connection \[C\]. This will be made apparent in Section 4 where we give another, perturbative, proof of the above Theorem. From now on we call $``$ the Chen’s differential or Chen’s vector field.
(ii) The higher order tensors $`\mu _{}`$ induced on the cohomology $`𝐇(𝔤)`$ by a $`L_{\mathrm{}}`$-quasi-isomorphism from a dLie-algebra $`𝔤`$ coincide precisely with the Massey products when they are well-defined and univalued. Thus the Chen’s differential gives a compact (and invariant) representation of the homotopy class of Massey products.
### 2.5.9. Extended Kuranishi moduli space.
Since the Chen’s vector field $``$ on $``$ is homological, the distribution
$$𝒟_{}=\{\text{X}𝒯_{}\text{X}=[,Y]\text{for some}Y𝒯_{}\}$$
is integrable (cf. Sect. 2.5.3). Indeed, the Jacobi identities imply
$$[[,\text{X}],[,\text{Y}]]=[,[\text{X},[,\text{Y}]]].$$
Consider an affine subscheme,
$$\mathrm{`}\mathrm{`}𝖹eros()\text{}:=𝖲peck[[t]]/<t^1,\mathrm{},t^{p+q}>,$$
of zeros of the vector field $``$. The distribution $`𝒟_{}`$ is tangent to $`\mathrm{`}\mathrm{`}𝖹eros()\text{}`$ since $`[,[,\text{X}]]=0`$. We define the extended Kuranishi space, $`𝒦_𝔤`$, as the so called non-linear homology \[BK, Ma2, Ba\] of the Chen differential, i.e. as the quotient $`\mathrm{`}\mathrm{`}𝖹eros()\text{}/𝒟_{}`$. (For our purposes it is enough to understand the latter as $`𝖲peck[[t]]𝖪er/<t^1,\mathrm{},t^{p+q}>`$.) This passage from $`(,)`$ to $`𝒦_𝔤`$ establishes a clear link between the unobstructed deformation functor $`\text{Def}_𝔤^{}`$ and the classical one $`\mathrm{𝖣𝖾𝖿}_𝔤^{}`$.
Kuranishi spaces, $`𝖪_𝔤`$, originally emerged \[Ku, GM\] in the context of the deformation functor $`\mathrm{D}ef_𝔤^0`$ associated to a cohomologically split $``$-graded dLie algebra $`𝔤`$. It is not hard to see that $`𝖪_𝔤`$ is a proper subspace of the extended Kuranishi space $`𝒦_𝔤`$. We will see below that for a rich class of dLie algebras $`𝔤`$ — the so-called (homotopy) Gerstenhaber algebras — the tangent sheaves to the smooth parts, $`𝒦_{\mathrm{s}mooth}`$, of the associated extended Kuranishi spaces are canonically sheaves of associative algebras. This structure is not visible if working in the category of original (non-extended) Kuranishi spaces $`𝖪`$ only.
### 2.6. Cohomological splitting.
It is very easy to compute Chen’s differential once a cohomological decomposition of a dLie-algebra $`𝔤`$ under investigation is chosen. The latter means the data $`(i,p,Q)`$, where $`i:𝐇(𝔤)𝔤`$ is a linear injection, $`p:𝔤𝐇(𝔤)`$ a linear surjection, and $`Q:𝔤𝔤`$ an odd linear operator, all satisfying the conditions,
$$pi=\text{Id}=ipdQQd,$$
in $`\text{End}_k(𝔤)`$.
Such a decomposition of $`𝔤`$ often occurs in (complex) differential geometry \[K, Ku\], where typical dLie-algebras come equipped with a norm $`||||`$ and their cohomologies $`𝐇`$ get identified with harmonic subspaces, $`\text{Harm}:=𝖪erd𝖪erd^{}𝔤`$, $`d^{}`$ being the adjoint of $`d`$ with respect to $`||||`$. The operator $`Q`$ is then $`Gd^{}`$, where $`G`$ is the Green function of the Laplacian $`\mathrm{}=dd^{}+d^{}d`$. In this situation the formal power series solution, $`\mathrm{\Gamma }`$, of the Master equation as well as the associated Chen’s vector field can be chosen to be convergent inducing, thereby, the structure of analytic (rather than formal) homological supermanifold on $``$.
It is not hard to check that, given a cohomological splitting of $`𝔤`$, the pair, $`(\mathrm{\Gamma },)`$, given recursively by
$`\mathrm{\Gamma }_{[1]}`$ $`=`$ $`{\displaystyle \underset{i}{}}t^ie_i𝖪erQ𝖪erd`$
$`\mathrm{\Gamma }_{[2]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Q[\mathrm{\Gamma }_{[1]}(t)\mathrm{\Gamma }_{[1]}(t)],`$
$`\mathrm{\Gamma }_{[3]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Q\left([\mathrm{\Gamma }_{[1]}(t)\mathrm{\Gamma }_{[2]}(t)]+[\mathrm{\Gamma }_{[2]}(t)\mathrm{\Gamma }_{[1]}(t)]\right),`$
$`\mathrm{}`$
$`\mathrm{\Gamma }_{[n]}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Q\left({\displaystyle \underset{k=1}{\overset{n1}{}}}[\mathrm{\Gamma }_{[k]}(t)\mathrm{\Gamma }_{[nk]}(t)]\right)`$ (3)
$`\mathrm{}`$
and
$$\stackrel{}{}p(\mathrm{\Gamma }_{[1]}):=\frac{1}{2}p\left([\mathrm{\Gamma }\mathrm{\Gamma }]\right),$$
give an explicit versal solution of the Master equation in $`𝔤`$.
The above power series for $`\mathrm{\Gamma }`$ is well known in the classical Deformation Theory \[K, Ku\] where it plays a key role in constructing Kuranishi analytic moduli spaces. This series is essentially an inversion of the Kuranishi map \[Ku\] in the category of $`L_{\mathrm{}}`$-algebras.
### 2.7. Formality and flat structures.
If an algebra $`(𝔤,d,[])`$ is formal, then, as follows from the proof of Theorem 2.5.6, the associated homological supermanifold $`(_𝔤,,0)`$ has a canonical flat structure $``$. In the associated flat coordinates $`t^i`$, the Chen’s vector field $``$ has coefficients which are polynomials in $`t^i`$ of order $`2`$. (This observation can, in fact, be made into a geometric criterion of formality.) More precisely, the following is true.
2.7.1. Theorem. For any formal dLie-algebra $`𝔤`$ there is a canonical isomorphism of sets,
$$\left\{\begin{array}{c}𝖥latstructureson_𝔤𝗌uchthat_X_Y_Z=0\\ 𝖿oranyhorizontalvectorfieldsX,Y𝖺ndZ\end{array}\right\}\left\{\begin{array}{c}𝖧omotopyclasses\\ 𝗈fformalitymaps\end{array}\right\}$$
Let us choose a basis, $`\{s_i,i=1,\mathrm{},p+q\}`$, in the $`(p|q)`$-dimensional vector superspace $`𝐇(𝔤)`$, and let $`\{t^i\}`$ be the associated linear coordinates. We shall need, for a short time, a category $`𝖠rtin_{k[[t]]}`$ consisting of Artin superalgebras of the form
$$𝒜_N:=k[[t^1,\mathrm{},t^{p+q}]]/<(t^1)^{N_1}\mathrm{}(t^{p+q})^{N_{p+q}}>.$$
Denoting the maximal ideal of such a superalgebra by $`m_N`$, we set $`(𝔤m_N)_{𝗏ersal}`$ to be a linear subspace in $`𝔤m_N`$ consisting of even elements, $`\mathrm{\Gamma }`$, satisfying $`\mathrm{\Gamma }modm_N=0`$, $`\mathrm{\Gamma }modm_N^2𝖪erd`$, and
$$\left(\mathrm{\Gamma }modm_N^2\right)mod𝖨md=\underset{i=1}{\overset{p+q}{}}t^is_i.$$
This set is invariant under the action of the gauge group $`\mathrm{exp}(𝔤m_N)_{\stackrel{~}{1}}`$ (see Sect. 2.5).
2.7.2. Lemma. For any formal dLie-algebra $`𝔤`$ there is a canonical isomorphism of sets,
$$\underset{}{lim}\frac{\{\mathrm{\Gamma }(𝔤m_N)_{𝗏ersal}d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}}{𝗀augegroup}=\frac{𝖥ormalitymaps}{𝗁omotopyequivalence},$$
where the projective limit is taken over the category $`𝖠rtin_{k[[t]]}`$.
Proof. For any $`𝒜_N𝖠rtin_{k[[t]]}`$, the Master equation in the Lie algebra $`(𝐇(𝔤)𝒜_N,[]_{\mathrm{i}nd})`$ has a canonical versal solution,
$$\mathrm{\Gamma }_0=\underset{i=1}{\overset{p+q}{}}t^is_i,=\underset{i,j,k=1}{\overset{p+q}{}}(1)^{\stackrel{~}{j}(\stackrel{~}{i}+1)}t^it^jC_{ij}^k\frac{}{t^k},$$
where $`C_{ij}^k`$ are the structure constants of $`[]_{\mathrm{i}nd}`$.
If $`𝔤`$ is formal, and $`F=\{F_n:^n𝐇(𝔤)𝔤,n=1,2,\mathrm{}\}`$ is a formality map, then, by Theorem 2.5.5,
$$\mathrm{\Gamma }:=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}F_n(\mathrm{\Gamma }_0,\mathrm{},\mathrm{\Gamma }_0),𝗍hesame,$$
is a versal solution of the Master equation in $`𝔤𝒜_N`$.
It is easy to check that an arbitrary homotopy change,
$$FF^h,$$
of the formality map, say the one induced by a set of linear maps $`h=\{h_n:^n𝐇(𝔤)𝔤,\stackrel{~}{h}_n=\stackrel{~}{n},n=1,2,\mathrm{}\}`$, change the versal solution $`\mathrm{\Gamma }`$ into a gauge equivalent one,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g,$$
where
$$g=\underset{i}{}h_1(e_i)t^i+\underset{i,j}{}\pm h_2(e_i,e_j)t^it^j+\underset{i,j,k}{}\pm h_3(e_i,e_j,e_k)t^it^jt^k+\mathrm{}.$$
Hence there is a canonical map,
$$\begin{array}{ccc}\frac{𝖥ormalitymaps}{𝗁omotopyequivalence},& & \frac{\{\mathrm{\Gamma }(𝔤m_N)_{𝗏ersal}d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}}{𝗀augegroup}\\ F=\{F_n:^n𝐇(𝔤)𝔤,n=1,2,\mathrm{}\}/& & _{n=1}^{\mathrm{}}\frac{1}{n!}F_n(\mathrm{\Gamma }_0,\mathrm{},\mathrm{\Gamma }_0)/\end{array}$$
which implies (almost immediately) the desired result. $`\mathrm{}`$
2.7.3. Proof of Theorem 2.7.1. Let $`𝖣iff_0`$ be the group of all formal diffeomorphisms of $`_𝔤=𝐇(𝔤)`$ into itself preserving the origin, and set
$$𝖣iff_{0,}:=\{\varphi 𝖣iff_0\varphi _{}()𝗂squadraticint^i\}.$$
Note that $`𝖣iff_{0,}=𝖣iff_0`$ if the Chen’s vector field $``$ vanishes. In general,
$$GL(p+q)𝖣iff_{0,}𝖣iff_0.$$
There is an obvious isomorphism,
$$\left\{\begin{array}{c}𝖥latsructureson_𝔤𝗌uchthat_X_Y_Z=0\\ 𝖿oranyhorizontalvectorfieldsX,Y𝖺ndZ\end{array}\right\}=\frac{𝖣iff_{0,}}{GL(p+q)}.$$
On the other hand, by Theorem 4.2.2 (see below),
$$\frac{𝖣iff_{0,}}{GL(p+q)}=\underset{}{lim}\frac{\{\mathrm{\Gamma }(𝔤m_N)_{𝗏ersal}d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0\}}{𝗀augegroup}.$$
The final link in the chain of canonical isomorphisms is provided by Lemma 2.7.2. $`\mathrm{}`$
2.7.4. Corollary. For any compact Calabi-Yau manifold, there is a canonical isomorphism of sets,
$$\left\{\begin{array}{c}𝖥latconnectionsonBarannikovKontsevich^{}s\\ 𝗆odulispaceofextendedcomplexstructures\end{array}\right\}\left\{\begin{array}{c}𝖧omotopyclasses\\ 𝗈fformalitymaps\end{array}\right\}$$
## 3 Homotopy Gerstenhaber algebras
### 3.1. Differential Gerstenhaber algebras.
A differential Gerstenhaber algebra, or shortly, a dG-algebra, is the data $`(𝔤,d,[],)`$ where
* $`(𝔤,d,[])`$ is a $``$-graded dLie-algebra as defined in Sect. 2.2.1;
* $`(𝔤,d,)`$ is a differential $``$-graded associative algebra with the product
$$\begin{array}{cccc}\hfill :& 𝔤𝔤& & 𝔤\\ & ab& & ab,\end{array}$$
having degree $`0`$;
* the binary operations $`[]`$ and $``$ satisfy the odd Poisson identity,
$$[a(bc)]=[ab]c+(1)^{(\stackrel{~}{a}+1)\stackrel{~}{b}}b[ac],$$
for all homogeneous $`a,b,c𝔤`$.
A dG-algebra is called graded commutative if such is the dot product.
The identity in $`𝔤`$ is an even element $`e_0`$ such that $`de_0=0`$, $`e_0a=ae_0=a`$ and $`[e_0a]=0`$ for any $`a𝔤`$. It defines a cohomology class $`[e_0]`$ in $`𝐇`$, and a constant vector field on $``$ which we denote by $`e`$.
### 3.1.1. Remark.
If $`𝔤`$ is a unital dG-algebra, then a versal solution, $`\mathrm{\Gamma }`$, of the Master equation (2) in $`𝔤`$ can (and will) be always normalized in such a way that
$$\stackrel{}{e}\mathrm{\Gamma }=e_0.$$
### 3.1.2. Differential Gerstenhaber-Batalin-Vilkovisky algebras.
Let $`(𝔤,)`$ be a $``$-graded commutative associative algebra over a field $`k`$. Let us say that the zero operator, $`0:𝔤𝔤`$, is of order $`1`$, and let us denote the linear operator, $`xax`$, of left multiplication by an element $`a𝔤`$ by $`l_a`$. A homogeneous linear operator, $`D:𝔤𝔤`$, is said to be an operator of order $`k`$ if the operator $`[D,l_a]`$ is of order $`k1`$ for any homogeneous $`a`$ in $`𝔤`$.
Assume now that $`(𝔤,)`$ comes equipped with
* a degree $`+1`$ linear operator, $`d:𝔤𝔤`$, of order 1, and
* a degree $`1`$ linear operator, $`\mathrm{\Delta }:𝔤𝔤`$, of order $`2`$,
satisfying the conditions,
$$d^2=0,\mathrm{\Delta }^2=0,d\mathrm{\Delta }+\mathrm{\Delta }d=0.$$
In this case the data
$$(𝔤,d,,[])$$
with
$$[ab]:=(1)^{\stackrel{~}{a}}\mathrm{\Delta }(ab)(1)^{\stackrel{~}{a}}(\mathrm{\Delta }a)ba(\mathrm{\Delta }b),a,b𝔤,$$
defines a dG-algebra \[Ma2\]. The dG-algebras arising in this way are often called exact or dGBV-algebras.
### 3.1.3. Example (complex manifolds).
For any $`n`$-dimensional comlex manifold $`M`$ the differential Lie algebra of Example 2.3.1,
$$𝔤=(\mathrm{\Gamma }(M,\mathrm{\Lambda }^{}T_M\mathrm{\Lambda }^{}\overline{T}_M^{}),\overline{},[]),$$
equipped with a supercommutative product,
$$\begin{array}{cccc}\hfill :& \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_1}T_M\mathrm{\Lambda }^{j_1}\overline{T}_M^{})\times \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_2}T_M\mathrm{\Lambda }^{j_2}\overline{T}_M^{})& & \mathrm{\Gamma }(M,\mathrm{\Lambda }^{i_1+i_2}T_M\mathrm{\Lambda }^{j_1+j_2}\overline{T}_M^{})\\ & \text{X}_1\overline{w}_1\times \text{X}_2\overline{w}_2& & (1)^{\stackrel{~}{j}_1\stackrel{~}{i}_2}(\text{X}_1\text{X}_2)\overline{w}_1\overline{w}_2,\end{array}$$
is a unital $``$-graded commutative dG-algebra.
If $`M`$ admits a nowhere vanishing global holomorphic volume form, $`\mathrm{\Omega }\mathrm{\Gamma }(M,\mathrm{\Omega }_M^n)`$, then the above dG-algebra is actually exact \[Ti, To, BK\] with $`\mathrm{\Delta }`$ being the composition,
$$\mathrm{\Delta }:\mathrm{\Lambda }^iT_M\stackrel{i_\mathrm{\Omega }}{}\mathrm{\Omega }_M^{ni}\stackrel{}{}\mathrm{\Omega }_M^{ni+1}\stackrel{i_\mathrm{\Omega }^1}{}\mathrm{\Lambda }^{i1}T_M.$$
Here $`i_\mathrm{\Omega }:\mathrm{\Lambda }^{}T_M\mathrm{\Omega }_M^{}`$ is the natural isomorphism given by contraction with the holomorphic volume form, and $``$ is the $`(1,0)`$-part of the de Rham operator.
### 3.1.4. Example (symplectic manifolds).
For any symplectic manifold $`(M,\omega )`$ the dLie algebra of Example 2.3.5,
$$𝔤=(\mathrm{\Gamma }(M,\mathrm{\Omega }_{}^{}),d,[]_\omega ),$$
together with a graded commutative product, $`ab:=ab`$, is a unital $``$-graded dG-algebra. Moreover, it is a dGBV-algebra with the 2-nd order differential given by
$$\mathrm{\Delta }|_{\mathrm{\Omega }_{}^k}=(1)^{k+1}d.$$
Here $`:\mathrm{\Omega }_{}^k\mathrm{\Omega }_{}^{2mk}`$ is the symplectic analogue of the Hodge duality operator defined by the condition, $`\beta (\alpha )=\beta ,\alpha \omega ^m/m!`$, with $`,`$ being the pairing between $`k`$-forms induced by the symplectic form.
### 3.1.5. Example (vector bundles).
Let $`M`$ be a complex manifold, and $`\pi :EM`$ a holomorphic vector bundle. There is a complex of $``$-graded sheaves $`(^{}E\mathrm{\Lambda }^{}E^{},\mathrm{\Delta })`$,
$$\mathrm{}\stackrel{\mathrm{\Delta }}{}^{k+1}E\mathrm{\Lambda }^{l+1}E^{}\stackrel{\mathrm{\Delta }}{}^kE\mathrm{\Lambda }^lE^{}\stackrel{\mathrm{\Delta }}{}^{k1}E\mathrm{\Lambda }^{l1}E^{}\stackrel{\mathrm{\Delta }}{}\mathrm{},$$
where the differential $`\mathrm{\Delta }`$ is just the contraction, the $``$-grading is induced from the one on $`\mathrm{\Lambda }^{}E`$, and we set $`^kE=\mathrm{\Lambda }^kE^{}=0`$ for $`k<0`$. It is easy to see that $`\mathrm{\Delta }`$ is a linear operator of oder 2 with respect to the natural supercommutative product,
$$\begin{array}{ccc}^{}E\mathrm{\Lambda }^{}E^{}\times ^{}E\mathrm{\Lambda }^{}E^{}& \stackrel{}{}& ^{}E\mathrm{\Lambda }^{}E^{}\\ (a_1b_1^{})\times (a_2b_2^{})& & (a_1a_2)(b_1^{}b_2^{}).\end{array}$$
Hence the data
$$𝔤=(\underset{k}{}𝔤^k,𝔤^k:=\underset{i+j=k}{}\mathrm{\Gamma }(M,^{}E\mathrm{\Lambda }^iE^{}\overline{\mathrm{\Omega }}_M^j),\mathrm{\Delta },\overline{})$$
is a unital dGBV-algebra. It extends the dLie algebra of Example 2.3.6, as the following calculation shows.
### Proposition.
The bracket $`[]`$ on $`^{}E\mathrm{\Lambda }^{}E^{}`$, when restricted to $`EE^{}`$, coincides, up to a sign, with the usual commutator of morphisms.
Proof. Let us consider a pair of germs, $`C_1=a_1b_1^{}`$ and $`C_2=a_2b_2^{}`$, in the same stalk of $`EE^{}`$. Then
$`[C_1C_2]`$ $`=`$ $`(1)^{\stackrel{~}{C}_1}\mathrm{\Delta }(C_1C_2)(1)^{\stackrel{~}{C}_1}\mathrm{\Delta }(C_1)C_2C_1\mathrm{\Delta }(C_2)`$
$`=`$ $`\mathrm{\Delta }\left((a_1a_2)(b_1^{}b_2^{})\right)+\mathrm{\Delta }(a_1b_1^{})(a_2b_2)(a_1b_1)\mathrm{\Delta }(a_2b_2^{})`$
$`=`$ $`a_1,b_1^{}a_2b_2^{}+a_1,b_2^{}a_2b_1^{}a_2,b_1^{}a_1b_2^{}`$
$`+a_2,b_2^{}a_1b_1^{}+a_1,b_1^{}a_2b_2^{}a_2,b_2^{}a_1b_1^{}`$
$`=`$ $`a_1,b_2^{}a_2b_1^{}a_2,b_1^{}a_1b_2^{}`$
$`=`$ $`\left(C_1C_2C_2C_1\right),`$
where angular brackets stand for the usual pairing between a vector and a 1-form. $`\mathrm{}`$
There is a problem with the constructed dGBV-algebra — its cohomology may not be finite dimensional even for compact manifolds. It is can be resolved by passing to its dGBV-subalgebra,
$$𝔤_E=(\underset{k}{}𝔤^k,𝔤^k:=\underset{i+j=k}{}\mathrm{\Gamma }(M,^iE\mathrm{\Lambda }^iE^{}\overline{\mathrm{\Omega }}_M^j),\mathrm{\Delta },\overline{}).$$
If necessary, the asymmetry of $`E`$ and $`E^{}`$ can be eliminated by taking the tensor product $`𝔤_E𝔤_E^{}`$.
### 3.1.6. Example (Hochschild cohomology).
Let $`A`$ be an associative algebra over a field $`k`$. The $``$-graded vector space of Hochschild cochains,
$$C^{}(A,A):=\underset{n=0}{\overset{\mathrm{}}{}}\text{Hom}_k(A^n,A),$$
can be made into a Hochschild complex with the differential, $`d:C^n(A,A)C^{n+1}(A,A)`$, given by
$`(df)(a_1\mathrm{}a_{n+1})`$ $`:=`$ $`a_1f(a_2\mathrm{}a_{n+1})+{\displaystyle \underset{i=1}{\overset{n}{}}}(1)^If(a_1\mathrm{}a_ia_{i+1}\mathrm{}a_{n+1})`$
$`+(1)^{n+1}f(a_1\mathrm{}a_n)a_{n+1}`$
for any $`fC^n(A,A)`$.
One can define two binary operations, $`C^{}(A,A)C^{}(A,A)C^{}(A,A)`$, the degree 0 dot product,
$$(fg)(a_1\mathrm{}a_{k+l}):=(1)^{kl}f(a_1\mathrm{}a_k)g(a_1\mathrm{}a_l),fC^k(A,A),gC^l(A,A),$$
and the degree $`1`$ bracket,
$$[fg]:=fg(1)^{(k+1)(l+1)}gf,$$
where
$$(fg)(a_1\mathrm{}a_{k+l1}):=\underset{i=1}{\overset{k1}{}}(1)^{(i+1)(l+1)}f(a_1\mathrm{}a_ig(a_{i+1}\mathrm{}a_{i+l})\mathrm{}a_{k+l1}).$$
These two make the Hochschild complex into a $``$-graded differential associative algebra and a differential (odd) Lie algebra respectively. Though $`(C^{}(A,A),d,[],)`$ is not a dG-algebra, it is a remarkabale fact that the associated Hochschild cohomology,
$$\mathrm{H}och^{}(A,A)=\frac{𝖪erd}{𝖨md},$$
carries the structure of graded commutative dG-algebra with respect to the naturally indiced dot product, Lie bracket and the zero differential.
### 3.2. $`A_{\mathrm{}}`$-algebras.
A strong homotopy algebra, or shortly $`A_{\mathrm{}}`$-algebra, is by definition a vector superspace $`V`$ equipped with linear maps,
$$\begin{array}{ccccc}\hfill \mu _k:& ^kV& & V& \\ & v_1\mathrm{}v_k& & \mu _k(v_1,\mathrm{},v_k),& k1,\end{array}$$
of parity $`\stackrel{~}{k}`$ satisfying, for any $`n1`$ and any $`v_1,\mathrm{},v_nV`$, the following higher order associativity conditions,
$$\underset{k+l=n+1}{}\underset{j=0}{\overset{k1}{}}(1)^r\mu _k(v_1,\mathrm{},v_j,\mu _l(v_{j+1},\mathrm{},v_{j+l}),v_{j+l+1},\mathrm{},v_n)=0,$$
(4)
where $`r=\stackrel{~}{l}(\stackrel{~}{v}_1+\mathrm{}+\stackrel{~}{v}_j)+\stackrel{~}{j}(\stackrel{~}{l}1)+(\stackrel{~}{k}1)\stackrel{~}{l}`$ and $`\stackrel{~}{v}`$ denotes the parity of $`vV`$.
Denoting $`dv_1:=\mu _1(v_1)`$ and $`v_1v_2:=\mu _2(v_1,v_2)`$, we can spell the first three conditions from the above infinite series as follows,
$`d^2=0`$,
$`d(v_1v_2)=(dv_1)v_2+(1)^{\stackrel{~}{v}_1}v_1(dv_2)`$,
$`v_1(v_2v_3)(v_1v_2)v_3=d\mu _3(v_1,v_2,v_3)+\mu _3(dv_1,v_2,v_3)+(1)^{\stackrel{~}{v}_1}\mu _3(v_1,dv_2,v_3)+\text{ }(1)^{\stackrel{~}{v}_1+\stackrel{~}{v}_2}\mu _3(v_1,v_2,dv_3)`$,
Therefore $`A_{\mathrm{}}`$-algebras with $`\mu _k=0`$ for $`k3`$ are nothing but the differential associative superalgebras with the differential $`\mu _1`$ and the associative multiplication $`\mu _2`$. If, furthermore, $`\mu _1=0`$, one recovers the usual associative superalgebras.
There is a (finer) $``$-graded version of the above definition in which the maps $`\mu _n`$ are required to be homogeneous (usually of degree $`n2`$) with respect to the given $``$-grading on $`V`$.
### 3.2.1. Identity.
An element $`e`$ in the $`A_{\mathrm{}}`$-algebra is called the identity if $`\mu _1(e)=0`$, $`\mu _2(e,v)=\mu _2(v,e)=v`$ and $`\mu _n(v_1,\mathrm{},e,\mathrm{},v_{n1})=0`$ for all $`n3`$ and arbitrary $`v,v_1,\mathrm{},v_{n1}V`$.
### 3.2.2. Homotopy classes of $`A_{\mathrm{}}`$-algebras.
For a pair of $`A_{\mathrm{}}`$-algebras, $`(V,\mu _{})`$ and $`(\stackrel{~}{V},\stackrel{~}{\mu }_{})`$, there is a natural notion of a $`A_{\mathrm{}}`$-morphism from $`V`$ to $`\stackrel{~}{V}`$ which is, by definition, a set of linear maps
$$F=\{f_n:V^n\stackrel{~}{V},n1\},$$
of parity $`\stackrel{~}{n}+1`$ (or of degree $`1n`$ in the $``$-graded case) which satisfy
$$\underset{1k_1<k_2<\mathrm{}<k_i=n}{}(1)^{i+r}\stackrel{~}{\mu }_i(f_{k_1}(v_1,\mathrm{},v_{k_1}),f_{k_2k_1}(v_{k_1+1},\mathrm{},v_{k_2}),\mathrm{},f_{nk_{i1}}(v_{k_{i1}+1},\mathrm{},v_n))$$
$$=\underset{k+l=n+1}{}\underset{j=0}{\overset{k1}{}}(1)^{l(\stackrel{~}{v}_1+\mathrm{}+\stackrel{~}{v}_j+n)+j(l1)}f_k(v_1,\mathrm{},v_j,\mu _l(v_{j+1},\mathrm{},v_{j+l}),v_{j+l+1},\mathrm{},v_n).$$
The first three floors in the above infinite tower are
$`\stackrel{~}{\mu }_1=\mu _2=:d`$,
$`\stackrel{~}{\mu }_2(v_1,v_2)=\mu _2(v_1,v_2)+(df_2)(v_1,v_2)`$,
$`\stackrel{~}{\mu }_3(v_1,v_2,v_3)+\stackrel{~}{\mu }_2(f_2(v_1,v_2),v_3)(1)^{\stackrel{~}{v}_1}\stackrel{~}{\mu }_2(a_1,f_2(v_3,v_4))=`$
$`\mu _3(v_1,v_2,v_3)f_2(\mu _2(v_1,v_2),v_3)+f_2(v_1,\mu _2(v_3,v_4))+(df_3)(v_1,v_2,v_3),`$
where we naturally extended the differential $`d:VV`$ to $`d:^kV^{}V^kV^{}V`$ (so that, for example, $`(df_2)(v_1,v_2)=df_2(v_1,v_2)+f_2(dv_1,v_2)+(1)^{\stackrel{~}{v}_1}f_2(v_1,dv_2)`$)
A morphism $`F=\{f_n\}`$ of the $`A_{\mathrm{}}`$-algebra $`(V,\mu _{})`$ to itself is called a homotopy if $`f_1`$ is an isomorphism. If $`(V,\mu _{})`$ has the identity $`e`$, then by a homotopy of $`(V,\mu _{},e)`$ we understand a homotopy of $`(V,\mu _{})`$ satisfying the additional conditions, $`f_n(v_1,\mathrm{},e,\mathrm{},v_{n1})=0`$ for all $`n2`$ and arbitrary $`v_1,\mathrm{},v_{n1}V`$.
It is not hard to see that homotopy defines an equivalence relation in the set of all possible (unital) $`A_{\mathrm{}}`$-structures on a given vector superspace $`V`$.
### 3.2.3. Remark.
For future reference we rewrite the $`n`$-th order associativity condition (4) as
$$\mathrm{\Lambda }_n(v_1,\mathrm{},v_n)=(d\mu _n)(v_1,\mathrm{},v_n)$$
where
$$\mathrm{\Lambda }_n(v_1,\mathrm{},v_n):=\underset{\genfrac{}{}{0pt}{}{k+l=n1}{k,l1}}{}(1)^r^{}\mu _{k+1}(v_1,\mathrm{},v_j,\mu _{l+1}(v_{j+1},\mathrm{},v_{j+l+1}),v_{j+l+2},\mathrm{},v_n)$$
and $`r^{}=(l+1)(\stackrel{~}{v}_1+\mathrm{}+\stackrel{~}{v}_j)+jl+k(l+1)+1`$.
### 3.2.4. Remark.
It follows from (4) for $`n=3`$ that the cohomology,
$$H(V):=\frac{𝖪er\mu _1}{𝖨m\mu _1}$$
of a (unital) $`A_{\mathrm{}}`$-algebra $`(V,\mu _{})`$ is canonically a (unital) associative algebra. Moreover, a homotopy class of (unital) $`A_{\mathrm{}}`$-structures on $`V`$ induces one and the same structure of (unital) associative algebra on $`H(V)`$.
### 3.2.5. The bar construction.
There is a conceptually better interpretation \[S\] of an $`A_{\mathrm{}}`$-structure on the vector superspace $`V`$ as a co-differential on the bar-construction of $`V`$. Here are the details:
* The vector space
$$𝖡(V):=\underset{n=1}{\overset{\mathrm{}}{}}\left(V[1]\right)^n$$
is naturally a co-algebra with the co-product given by
$$\mathrm{\Delta }(w_1\mathrm{}w_n)=\underset{i=1}{\overset{n}{}}\left(w_1\mathrm{}w_i\right)\left(w_i\mathrm{}w_n\right).$$
* A linear map $`Q:𝖡(V)𝖡(V)`$ is said to be a co-derivation if $`\mathrm{\Delta }Q=Q\text{Id}+\text{Id}Q`$. There is a one-to-one correspondence between such co-derivations and Hochschild cochains understood as elements of $`\text{Hom}(𝖡(V),V)`$.
* A homogeneous (of degree $`2`$) Hochschild cochain $`\mu _{}:𝖡(V)V`$ defines an $`A_{\mathrm{}}`$-structure on $`V`$ if and only if the associated co-derivation $`Q`$ is a co-differential, i.e. satisfies $`Q^2=0`$.
In this setup, a morphism $`(V,\mu _{})(\stackrel{~}{V},\stackrel{~}{\mu }_{})`$ as in 3.8.1 is precisely a morphism of the associated bar-constructions respecting co-differentials.
### 3.3. $`C_{\mathrm{}}`$-algebras.
This notion is a supercommutative analogue of the notion of $`A_{\mathrm{}}`$-algebra.
Let $`V`$ be a $``$-graded vector space and $`𝖡(V)`$ its bar construction. One can make the latter into an associative and graded commutative algebra by defining the shuffle tensor product, $`:𝖡(V)𝖡(V)𝖡(V)`$, as follows
$$(w_1\mathrm{}w_k)(w_{k+1}\mathrm{}w_n):=\underset{\sigma Sh(k,n)}{}e(\sigma ;w_1,\mathrm{},w_n)w_{\sigma (1)}\mathrm{}w_{\sigma (n)}.$$
Here we used the notations explained in Sect. 2.4.
By definition \[GJ\], a strong homotopy commutative algebra, or shortly, $`C_{\mathrm{}}`$-algebra is an $`A_{\mathrm{}}`$-algebra $`(V,\mu _{})`$ such that the associated Hochschild cochain $`\mu _{}:𝖡(V)V`$ factors through the composition<sup>5</sup><sup>5</sup>5Such cochains are often called Harrison cochains.
This implies, in particular, that
$$\mu _2(v_1,v_2)=(1)^{\stackrel{~}{v}_1\stackrel{~}{v}_2}\mu _2(v_2,v_1)$$
for any $`v_1,v_2V`$.
One defines notions of unital $`C_{\mathrm{}}`$-algebras, of a morphism of $`C_{\mathrm{}}`$-algebras, of their homotopy etc. in the same way as in the $`A_{\mathrm{}}`$-case.
### 3.4. $`G_{\mathrm{}}`$-algebras.
Let $`V`$ be a $``$-graded vector space and let
$$𝖫ie(V[1]^{})=\underset{k=1}{\overset{\mathrm{}}{}}𝖫ie^k(V[1]^{})$$
the free graded Lie algebra generated by the shifted dual vector space $`V[1]^{}`$, i.e.
$$𝖫ie^1(V[1]^{}):=V[1]^{},𝖫ie^k(V[1]^{}):=[V[1]^{},𝖫ie^{k1}(V[1]^{})].$$
The Lie bracket on $`𝖫ie(V[1]^{})`$ extends in a usual way to the skew-symmetric associative algebra
$$^{}𝖫ie(V[1]^{})=\underset{k=0}{\overset{\mathrm{}}{}}^k𝖫ie(V[1]^{}),$$
making the latter into a Gerstenhaber algebra.
### 3.4.1. Definition \[Ta1, TT\].
A homotopy Gerstenhaber algebra, or shortly $`G_{\mathrm{}}`$-algebra is a graded vector space $`V`$ together with a degree one linear operator
$$Q:^{}𝖫ie(V[1]^{})^{}𝖫ie(V[1]^{})$$
such that $`Q^2=0`$ and $`Q`$ is a derivation with respect to both the product and the bracket.
A $`G_{\mathrm{}}`$-morphism, $`VV^{}`$, of $`G_{\mathrm{}}`$-algebras is by definition a morphism, $`(^{}𝖫ie(V[1]^{}),Q)(^{}𝖫ie(V^{}[1]^{}),Q^{})`$, of associated differential Gerstenhaber algebras.
The definition 3.4.1 makes sense only in the case when $`V`$ is finite-dimensional. However, an obvious dualization fixes the problem \[TT\]:
* The dual of $`𝖫ie(V[1]^{})`$ can be identified with the quotient $`𝖡(V)/𝖡(V)𝖡(V)`$, $``$ being the shuffle tensor product.
* Derivations of $`^{}𝖫ie(V[1]^{})`$ can be identified with arbitrary collections of linear maps,
$$m_{k_1,\mathrm{},k_n}^{}:V[1]^{}𝖫ie^{k_1}(V[1]^{})\mathrm{}𝖫ie^{k_n}(V[1]^{}),$$
which upon dualization go into linear homogeneous maps,
$$m_{k_1,\mathrm{},k_n}:\frac{V^{k_1}}{\mathrm{s}huffleproducts}\mathrm{}\frac{V^{k_n}}{\mathrm{s}huffleproducts}V,$$
of degree $`3nk_1\mathrm{}k_n`$.
* the condition $`Q^2=0`$ translates into a well-defined set of quadratic equations for $`m_{k_1,\mathrm{},k_n}`$ which say, in particular, that $`m_1`$ is a differential on $`V`$ and that the product, $`v_1v_2:=(1)^{\stackrel{~}{v}_1}m_2(v_1,v_2)`$, together with the Lie bracket, $`[v_1v_2]:=(1)^{\stackrel{~}{v}_1}m_{1,1}(v_1,v_2)`$, satisfy the Poisson identity up to a homotopy given by $`m_{2,1}`$. Hence the associated cohomology space $`𝐇`$ is a graded commutative Gerstenhaber algebra with respect to the binary operations induced by $`m_2`$ and $`m_{1,1}`$.
The identity in a $`G_{\mathrm{}}`$-algebra $`V`$ is an even element $`e`$ such that all $`m_{k_1,\mathrm{},k_n}(\mathrm{},e,\mathrm{})`$ vanish except $`m_2(e,v)=v`$.
### 3.4.2. Theorem-construction.
There is a canonical functor from the derived category of unital $`G_{\mathrm{}}`$-algebras with finite-dimensional cohomology to the category of $`F_{\mathrm{}}`$-manifolds.
Proof. Since each quasi-isomorphism of $`G_{\mathrm{}}`$-algebras is an equivalence relation, the derived category of $`G_{\mathrm{}}`$-algebras coincides with their homotopy category.
We construct the desired functor,
in two steps.
Step 1. Suppose we are given a homotopy class, $`[]`$, of $`G_{\mathrm{}}`$-structures on a graded vector space $`V`$. By Kontsevich’s Lemma 1 in \[Ko3\], a cohomological splitting of the complex $`(V,m_1)`$ transfers $`[]`$ into a homotopy class, $`[^{}𝖫ie(𝐇[1]^{}),Q]`$, of minimal $`G_{\mathrm{}}`$-algebras on the finite-dimensional cohomology space of the above complex. Moreover, this class does not depend on the choice of a particular cohomological splitting, and it is homotopy equivalent to the original one.
Step 2. Let $``$ be the multiplicative ideal in $`^{}𝖫ie(𝐇[1]^{})`$ generated by the commutant of $`𝖫ie(𝐇[1]^{})`$. Any differential $`Q`$ from the induced homotopy class preserves this ideal and induces, through the quotient $`^{}𝖫ie(𝐇[1]^{})/`$, a homotopy class of $`L_{\mathrm{}}`$-structures on $`V`$ which, by Corollary 2.5.7, can be identified with an odd vector field $``$ on the associated cohomological supermanifold $``$ satisfying $`II^2`$ and $`[E,]=`$, $`I`$ being the ideal of the distinguished point $`0`$ and $`E`$ the Euler vector field. We claim that the rest of the data listed in Definition 1.1 gets induced on $``$ through the quotient $`^{}𝖫ie(𝐇[1]^{})/^2`$. Indeed, what is left of a differential $`Q`$ on this quotient can be described as a collection of tensors, $`m_{k,1,\mathrm{},1}`$, which, in a basis $`\{e_a\}`$ of $`𝐇`$, are represented by their components, $`\mu _{b_1\mathrm{}b_k,c_1,\mathrm{},c_l}^a`$, $`k1,l0`$. The Chen’s vector field $``$ and the tensors $`\mu _k`$ defining the structure of a $`C_{\mathrm{}}`$-algebra on the tangent sheaf $`𝒯_{}`$ are then given by formal power series,
$$=\underset{l0}{}\pm \mu _{b_1,c_1,\mathrm{},c_l}^at^{b_1}t^{c_1}\mathrm{}t^{c_l}\frac{}{t^a}$$
and
$$\mu _{b_1\mathrm{}b_k}^a=\underset{l0}{}\pm \mu _{b_1\mathrm{}b_k,c_1,\mathrm{},c_l}^at^{c_1}\mathrm{}t^{c_l}.$$
where $`t^a`$ are the associated linear coordinates on $``$ to which we assign degree $`2|e_a|`$. It is easy to see that the $`G_{\mathrm{}}`$-identities for $`m_{k,1,\mathrm{},1}`$ get transformed into the right identities for the tensor fields $``$ and $`\mu _k`$ on $``$. This completes the construction. $`\mathrm{}`$
### 3.4.3. Corollary.
For any unital $`G_{\mathrm{}}`$-algebra with finite dimensional cohomology, the tangent sheaf to the smooth part of the extended Kuranishi space, $`=\mathrm{`}\mathrm{`}𝗓eros()\text{}/𝖨m\mu _1`$, is canonically a sheaf of induced (unital) associative algebras.
It will be interesting to find out when $`_{\mathrm{s}mooth}`$ with its canonically induced structure 3.4.3 is an $`F`$-manifold in the sense of Hertling and Manin \[HM\].
### 3.5. Remark.
Different “resolutions” of the chain operad in the little disk operad give different notions of homotopy Gerstenhaber algebra \[V\]. The definition 3.4.1 is the most canonical one. However, the functor $`F_{\mathrm{}}`$ is not an equivalence in this case.
The proof of Theorem 3.4.2 suggests one more version: a reduced homotopy Gerstenhaber algebra is a graded vector space $`V`$ together with the structure of $`G_{\mathrm{}}`$-algebra such that all composition maps $`m_{k_1,\mathrm{},k_n}`$ vanish except $`m_{k_1,1,\mathrm{},1}`$. The derived category of such algebras is equivalent to the category of $`F_{\mathrm{}}`$-manifolds (cf. Theorem 2.5.7).
### 3.6. Formality and Gauss-Manin connections.
A pre-$`Frobenius_{\mathrm{}}`$ manifold is the data $`(,E,,,[\mu _{}],e)`$, where
* $``$ is a formal pointed $``$-graded manifold,
* $`E`$ is the Euler vector field on $``$, $`Ef:=\frac{1}{2}|f|f`$, for all homogeneous functions on $``$ of degree $`|f|`$,
* $``$ is a flat torsion-free affine connection, called the Gauss-Manin connection, on $``$,
* $``$ is an odd homological (i.e. $`^2=0`$) vector field on $``$ such that $`[E,]=`$, $`_X_Y_Z=0`$ for any horizontal vector fields, $`X,Y`$ and $`Z`$ on $``$, and $`II^2`$, $`I`$ being the ideal of the distinguished point in $``$,
* $`[\mu _n:^n𝒯_{}𝒯_{}]`$, $`n`$, is a homotopy class of smooth unital strong homotopy commutative ($`C_{\mathrm{}}`$) algebras defined on the tangent sheaf, $`𝒯_{}`$, to $``$, such that $`Lie_E\mu _n=\frac{1}{2}n\mu _n`$, for all $`n`$, and $`\mu _1`$ is given by
$$\begin{array}{cccc}\hfill \mu _1:& 𝒯_{}& & 𝒯_{}\\ & X& & \mu _1(X):=[,X].\end{array}$$
* $`e`$ is the flat unit, i.e. an even vector field on $``$ such that $`[,e]=0`$, $`e=0`$, $`\mu _2(e,X)=X`$, $`X𝒯_{}`$, and $`\mu _n(\mathrm{},e,\mathrm{})=0`$ for all $`n3`$.
### 3.6.1. Theorem.
There is a canonical functor from the category of pairs $`(𝔤,F)`$, where $`𝔤`$ is $`L_{\mathrm{}}`$-formal unital homotopy Gerstenhaber algebra and $`F`$ a formality map, to the category of pre-$`Frobenius_{\mathrm{}}`$ manifolds.
Proof. The desired statement follows immediately from Theorem 2.7.1 and a version of Theorem-Construction 3.4.2 where the formality map $`F`$ is used to transfer the $`G_{\mathrm{}}`$-structure from the algebra to its cohomology. $`\mathrm{}`$
### 3.6.2. Theorem.
If a homotopy Gerstenhaber algebra $`𝔤`$ is quasi-isomorphic, as a $`L_{\mathrm{}}`$-algebra, to an Abelian dLie algebra, then the tangent sheaf, $`𝒯_{}`$, to its cohomology viewed as a linear supermanifold is canonically a sheaf of unital graded commutative associative algebras.
Proof. In this case $`=0`$ and $`\mu _2`$, which is now defined uniquely, makes $`𝒯_{}`$ into a sheaf of unital graded commutative associative algebras. $`\mathrm{}`$
## 4 Perturbative construction of $`F_{\mathrm{}}`$-invariants
The purpose of this section is to give second “down-to-earth” proofs of some of the main claims of this paper. Our approach here is based on perturbative solutions of algebro-differential equations rather than on the homotopy technique used in the two previous Sections.
First comes a perturbative proof of the Smoothness Theorem 2.5.6.
### 4.1. Theorem (Chen’s construction).
For any differential Lie superalgebra $`𝔤`$, there exists a versal element, $`\mathrm{\Gamma }k[[t]]𝔤`$, and an odd derivation, $`:k[[t]]k[[t]]`$, such that $`^2=0`$ and the equation,
$$d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0$$
holds. Moreover, for any quasi-isomorphism of complexes of vector spaces, $`\varphi :(𝔤,d)(𝐇,0)`$, $`\mathrm{\Gamma }`$ may be normalized so that $`\varphi (\mathrm{\Gamma }_{[n]})=0`$ for all $`n2`$.
We shall prove this Theorem by induction using (twice) the following Lemma which is merely a trancated version of Remark 2.5.1.
### 4.1.1. Lemma.
Assume the elements $`\mathrm{\Gamma }_{(n)}=_{k=0}^n\mathrm{\Gamma }_{[k]}k[[t]]𝔤`$ and $`_{(n)}=_{k=0}^n_{[n]}𝖣erk[[t]]`$ satisfy
$$d\mathrm{\Gamma }_{(n)}+\stackrel{}{}_{(n)}\mathrm{\Gamma }_{(n)}+\frac{1}{2}\left[\mathrm{\Gamma }_{(n)}\mathrm{\Gamma }_{(n)}\right]=0modI^{n+1}.$$
Then
$$\psi _{[n+1]}:=d\mathrm{\Gamma }_{(n)}+\stackrel{}{}_{(n)}\mathrm{\Gamma }_{(n)}+\frac{1}{2}\left[\mathrm{\Gamma }_{(n)}\mathrm{\Gamma }_{(n)}\right]modI^{n+2}$$
satisfies
$$d\psi _{[n+1]}=\stackrel{}{}_{(n)}^{\mathrm{\hspace{0.17em}2}}\left(\mathrm{\Gamma }_{(n)}\right)modI^{n+2}.$$
4.1.2. Proof of the Theorem. Let
$$\varphi :(𝔤,d)(𝐇,0),$$
be a quasi-isomorphism, i.e. a morphism of complexes inducing an isomorphism on cohomology. Since $`𝔤`$ is defined over a field, such a quasi-isomorphism always exists (note that we do not ask for any sort of a relationship between $`\varphi `$ and the Lie brackets).
Let $`e_i`$ be any representatives of the cohomology classes $`[e_i]`$ in $`𝖪erd𝔤`$. We may assume without loss of generality that $`\varphi (e_i)=[e_i]`$. Then choosing $`\mathrm{\Gamma }_{[0]}=0`$, $`\mathrm{\Gamma }_{[1]}:=_{i=1}^{p+q}t^ie_i`$, and $`_{[0]}=_{[1]}=0`$ we get the data $`(\mathrm{\Gamma }_{(1)},_{(1)})`$ satisfying the Master equation modulo terms in $`I^2`$ and the nilpotency condition $`^2=0`$ modulo terms in $`I^3`$.
Assume we have constructed a versal element $`\mathrm{\Gamma }_{(n)}=_{k=1}^n\mathrm{\Gamma }_{[k]}k[[t]]𝔤`$ and an odd vector field $`_{(n)}=_{k=2}^n_{[n]}`$ on $``$ such that the equations
$$P_n:\{\begin{array}{c}d\mathrm{\Gamma }_{(n)}+\stackrel{}{}_{(n)}\mathrm{\Gamma }_{(n)}+\frac{1}{2}\left[\mathrm{\Gamma }_{(n)}\mathrm{\Gamma }_{(n)}\right]=0modI^{n+1}\hfill \\ _{(n)}^2=0modI^{n+2}.\hfill \end{array}$$
are satisfied.
Let us show that there exists $`\mathrm{\Gamma }_{[n+1]}k[[t]]𝔤`$ and $`_{[n+1]}H^0(𝒯M_𝐇)`$ such that
$$\mathrm{\Gamma }_{(n+1)}=\mathrm{\Gamma }_{(n)}+\mathrm{\Gamma }_{[n+1]},_{(n+1)}=_{(n)}+_{[n+1]},$$
satisfy the equations $`P_{n+1}`$.
Note that, in the notations of Lemma 4.4.1, one has
$$d\mathrm{\Gamma }_{(n+1)}+\stackrel{}{}_{(n+1)}\mathrm{\Gamma }_{(n+1)}+\frac{1}{2}\left[\mathrm{\Gamma }_{(n+1)}\mathrm{\Gamma }_{(n+1)}\right]modI^{n+2}=d\mathrm{\Gamma }_{[n+1]}+\psi _{[n+1]}+\stackrel{}{}_{[n+1]}\mathrm{\Gamma }_{[1]}.$$
Let us now define $`\stackrel{}{}_{[n+1]}`$ by setting
$$\stackrel{}{}_{[n+1]}\mathrm{\Gamma }_{[1]}:=\varphi (\psi _{[n+1]}).$$
As $`d\psi _{[n+1]}=0`$ by Lemma 4.4.1 and the second equation of $`P_n`$, we conclude that
$$\psi _{[n+1]}+\stackrel{}{}_{[n+1]}\mathrm{\Gamma }_{[1]}(𝖪er\varphi 𝖪erd)k[[t]]_{[n+1]}.$$
Since $`\varphi `$ is a quasi-isomorphism, $`𝖪er\varphi 𝖪erd=𝖨md`$. Hence, there exists $`\mathrm{\Gamma }_{[n+1]}k[[t]]𝔤`$ such that
$$d\mathrm{\Gamma }_{[n+1]}=\psi _{[n+1]}\stackrel{}{}_{[n+1]}\mathrm{\Gamma }_{[1]}.$$
Thus the first equation of the system $`P_{n+1}`$ holds. This implies, by Lemma 2.4.2,
$`d\psi _{[n+2]}`$ $`=`$ $`\stackrel{}{}_{(n+1)}^{\mathrm{\hspace{0.17em}2}}\mathrm{\Gamma }_{(n+1)}modI^{n+3}`$
$`=`$ $`\stackrel{}{}_{(n+1)}^{\mathrm{\hspace{0.17em}2}}\mathrm{\Gamma }_{[1]}modI^{n+3}.`$
Applying $`\varphi `$ to both sides of this equation, we get
$$\stackrel{}{}_{(n+1)}^{\mathrm{\hspace{0.17em}2}}\varphi (\mathrm{\Gamma }_{[1]})=0$$
implying the second equation of the system $`P_{n+1}`$,
$$\stackrel{}{}_{(n+1)}^{\mathrm{\hspace{0.17em}2}}=0modI^{n+3},$$
and completing thus the inductive procedure.
Finally, we note that $`𝖪erd+𝖪er\varphi =𝔤`$ for $`\varphi `$ is a quasi-isomorphism. Hence we can always adjust $`\mathrm{\Gamma }_{[n+1]}`$, $`n1`$, so that it lies in $`𝖪er\varphi `$. $`\mathrm{}`$.
### 4.1.3. Remarks.
(i) The role of $``$ in the Chen’s construction is to absorb all the obstructions so that constructing a versal solution to the Master equation poses no problem (cf. Smoothness Theorem 2.5.6).
(ii) The Chen’s differential $``$ is completely determined by $`\mathrm{\Gamma }`$. Indeed, the Master equations imply,
$$\stackrel{}{}\varphi (\mathrm{\Gamma })=\frac{1}{2}\varphi \left([\mathrm{\Gamma }\mathrm{\Gamma }]\right).$$
Decomposing,
$$\varphi (\mathrm{\Gamma })=\underset{i}{}f^i(t)[e_i],$$
we note that $`f^i(t)=t^imodI^2`$. Hence the functions $`f^i(t)`$ define a coordinate system on $``$ and the values, $`f^i(t)`$, completely determine the differential $``$.
In particular, if $`\mathrm{\Gamma }`$ is $`\varphi `$-normalized, i.e. $`\varphi (\mathrm{\Gamma }_{[n2]})=0`$, then $``$ can be computed by the formula
$$\stackrel{}{}\left(\underset{i=1}{\overset{p+q}{}}t^i[e_i]\right)=\frac{1}{2}\varphi \left([\mathrm{\Gamma }\mathrm{\Gamma }]\right).$$
(iii) We shall understand from now on a versal solution, $`\mathrm{\Gamma }`$, of the Master equations and the associated Chen differential $``$ as, respectively, global sections of the sheaves $`𝔤𝒪_{}`$ and $`𝒯_{}`$ on $``$ (in practical terms, this essentially fixes their transformation properties under arbitrary changes of coordinates on the cohomology supermanifold).
We call sometimes $`\mathrm{\Gamma }`$ a Master function.
(iv) The argument in (ii) also downplays the role of the quasi-isomorphism $`\varphi `$ used in the Chen construction. If $`\mathrm{\Gamma }`$ is normalised with respect to a quasi-isomorphism $`\varphi :(𝔤,d)(𝐇,0)`$, then, for any other quasi-isomorphism $`\varphi ^{}`$, the same$`\mathrm{\Gamma }`$ can be viewed as $`\varphi ^{}`$-normalized, but in a new coordinate system $`t_{}^{}{}_{}{}^{i}=f^i(t^j)`$ given by $`\varphi ^{}(\mathrm{\Gamma })=f^i(t^j)[e_i]`$. Thus varying quasi-isomorphism $`\varphi `$ used in the construction of $`\mathrm{\Gamma }`$ amounts to varying flat structure on the pointed supermanifold $``$.
(v) Chen has actually invented his differential $``$ in the context of differential associative algebras \[C\]. Its Lie algebra analogue, Theorem 4.1, is due to Hain \[H\].
### 4.2. Gauge equivalence.
Let us consider the following action, called a gauge transformation, of $`𝔤_{\stackrel{~}{1}}I`$ on $`𝔤k[[t]]`$:
$$\begin{array}{ccc}𝔤I\times 𝔤k[[t]]& & 𝔤k[[t]]\\ g\mathrm{\Gamma }& & \mathrm{\Gamma }^g:=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{\mathrm{a}d_g}1}{\mathrm{a}d_g}(d+\stackrel{}{})g.\end{array}$$
### 4.2.1. Lemma.
If $`\mathrm{\Gamma }𝔤k[[t]]`$ is a Master function, then, for any $`g𝔤_{\stackrel{~}{1}}I`$, the function $`\mathrm{\Gamma }^g`$ is also a Master function, and both these share the same Chen differential.
Proof. We have to show that the equation $`d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0`$ implies
$$d\mathrm{\Gamma }^g+\stackrel{}{}\mathrm{\Gamma }^g+\frac{1}{2}[\mathrm{\Gamma }^g\mathrm{\Gamma }^g]=0.$$
This follows immediately from the well-known formulae \[GM\],
$$e^{\mathrm{a}d_g}de^{\mathrm{a}d_g}=d\mathrm{a}d_{\frac{e^{\mathrm{a}d_g}1}{\mathrm{a}d_g}dg},e^{\mathrm{a}d_g}e^{\mathrm{a}d_g}=\mathrm{a}d_{\frac{e^{\mathrm{a}d_g}1}{\mathrm{a}d_g}\stackrel{}{}g},e^{\mathrm{a}d_g}\mathrm{a}d_\mathrm{\Gamma }e^{\mathrm{a}d_g}=\mathrm{a}d_{e^{\mathrm{a}d_g}\mathrm{\Gamma }},$$
and
$$e^{\mathrm{a}d_g}\left[(\mathrm{})(\mathrm{})\right]=\left[e^{\mathrm{a}d_d}(\mathrm{})e^{\mathrm{a}d_g}(\mathrm{})\right].$$
$`\mathrm{}`$
### 4.2.2. Theorem.
Let $`𝔤`$ be a differential Lie algebra. For any two Master functions on $``$, $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{}`$, there is a gauge function $`g\mathrm{\Gamma }(,𝔤I)`$ and a diffeomorphism $`f:(,0)(,0)`$ such that $`\mathrm{\Gamma }^{}=f^{}(\mathrm{\Gamma }^g)`$ and $`=f_{}(^{})`$.
A sketch of the proof. Let us fix a quasi-isomorphism $`\varphi :(𝔤,d)(𝐇,0)`$ of complexes of Abelian groups, and a coordinate system on $``$ in which $`\mathrm{\Gamma }^{}`$ is $`\varphi `$-normalized.
We have $`\mathrm{\Gamma }_{[1]}^{}=\mathrm{\Gamma }_{[1]}dg_{[1]}`$, for some $`g_{[1]}\mathrm{\Gamma }(,𝔤I)`$, and $`_{[1]}^{}=_{[1]}=0`$. Hence $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }^{g_{[1]}}modI^2`$ and there is a unique diffeomorphism $`f_1:`$ such that the Master function $`\mathrm{\Gamma }^{\prime \prime }:=f_1^{}(\mathrm{\Gamma }^{g_{[1]}})`$ is $`\varphi `$-normalized and $`\mathrm{\Gamma }_{[1]}^{\prime \prime }=\mathrm{\Gamma }_{[1]}^{}`$. Hence,
$$d(\mathrm{\Gamma }_{[2]}^{}\mathrm{\Gamma }_{[2]}^{\prime \prime })+\stackrel{}{(_{[2]}^{}_{[2]}^{\prime \prime })}\mathrm{\Gamma }_{[1]}=0$$
implying $`_{[2]}^{}=_{[2]}^{\prime \prime }`$ and $`d(\mathrm{\Gamma }_{[2]}^{}\mathrm{\Gamma }_{[2]}^{\prime \prime })=0`$. Since $`\varphi (\mathrm{\Gamma }_{[2]}^{}\mathrm{\Gamma }_{[2]}^{\prime \prime })=0`$ and $`\varphi `$ is a quasi-isomorphism, there exists $`g_{[2]}𝔤I^2`$ such that $`\mathrm{\Gamma }_{[2]}^{}\mathrm{\Gamma }_{[2]}=dg_{[2]}`$. Hence
$$\mathrm{\Gamma }^{}=(\mathrm{\Gamma }^{\prime \prime })^{g_{[2]}}=f_1^{}(\mathrm{\Gamma }^{g_{(2)}})modI^3.$$
Continuing by induction and using Lemma 4.2.1 one easily obtains the desired result. $`\mathrm{}`$
### 4.2.3. Corollary.
The Chen’s vector field $``$ on $``$ is an invariant of $`𝔤`$.
### 4.3. Differential on $`𝒯_{}`$.
We fix from now on a dG-algebra $`𝔤`$ and a Master function $`\mathrm{\Gamma }`$ on $``$. The latter is not defined canonically, though the associated Chen differential $``$ is.
We also fix a quasi-isomorphism, $`\varphi :(𝔤,d)(𝐇,0)`$, of complexes of Abelian groups. This puts no restriction whatsoever on the dG-algebra under consideration. Moreover, our main results will not depend on the particular choices of $`\mathrm{\Gamma }`$ and $`\varphi `$ we have made — these two are no more than working tools.
The global vector field $``$ on $``$ makes $`𝒯_{}`$ into a sheaf of complexes with the differential
$$\begin{array}{cccc}\hfill \delta :& 𝒯_{}& & 𝒯_{}\\ & \text{X}& & \delta \text{X}:=[,\text{X}],\end{array}$$
where $`[,]`$ stands for the usual commutator of (germs) of vector fields. Indeed,
$$(\delta )^2\text{X}=[,[,\text{X}]]=\frac{1}{2}[[,],\text{X}]=[^2,\text{X}]=0,$$
where we have used the Jacobi identity and the fact that $`^2=0`$.
This, of course, induces a differential on the sheaf of tensor products, $`𝒯_{}^m(𝒯_{}^{})^n`$ (and on the associated vector space of global sections), which we denote by the same symbol $`\delta `$.
### 4.4. Deformed dG-algebra.
It is easy to check that the map
$$\begin{array}{cccc}\hfill d^\mathrm{\Gamma }:& k[[t]]𝔤& & k[[t]]𝔤\\ & a& & d^\mathrm{\Gamma }a:=da+\stackrel{}{}a+[\mathrm{\Gamma }a].\end{array}$$
satisfies
* $`(d^\mathrm{\Gamma })^2=0`$,
* $`d^\mathrm{\Gamma }(ab)=(d^\mathrm{\Gamma }a)b+(1)^{\stackrel{~}{a}}ad^\mathrm{\Gamma }b`$
* $`d^\mathrm{\Gamma }\left[ab\right]=\left[d^\mathrm{\Gamma }ab\right](1)^{\stackrel{~}{a}}\left[ad^\mathrm{\Gamma }b\right]`$
implying that the data $`(k[[t]]𝔤,[],,d^\mathrm{\Gamma })`$ is a dG-algebra.
The differentials $`d^\mathrm{\Gamma }`$ and $`\delta `$ make the sheaf $`𝔤(𝒯_{}^{})^k`$ on $``$ into a sheaf of complexes with the differential which we denote by $`D^\mathrm{\Gamma }`$. For example, for any germ $`\mathrm{\Phi }𝔤𝒯_{}^{}`$ and any germ $`\text{X}𝒯_{}`$ over the same point in $``$,
$$(D^\mathrm{\Gamma }\mathrm{\Phi })(\text{X}):=d^\mathrm{\Gamma }\mathrm{\Phi }(\text{X})(1)^{\stackrel{~}{\mathrm{\Phi }}}\mathrm{\Phi }(\delta \text{X}).$$
The vector space $`\text{Hom}(𝒯_{}^k,𝔤𝒪_{})`$ is also a complex with the differential denoted by the same symbol $`D^\mathrm{\Gamma }`$.
### 4.5. Morphism of sheaves of complexes.
The versal solution $`\mathrm{\Gamma }`$ gives rise to a morphism of $`𝒪_{}`$-modules,
$$\begin{array}{cccc}\hfill \mathrm{{\rm Y}}:& 𝒯_𝐇& & 𝒪_{}𝔤\\ & \text{X}& & \mathrm{{\rm Y}}(\text{X}):=\stackrel{}{\text{X}}\mathrm{\Gamma }.\end{array}$$
It is not hard to check that $`\mathrm{{\rm Y}}`$ is a monomorphism.
### 4.5.1. Lemma.
The element$`\mathrm{{\rm Y}}\text{Hom}(𝒯_{},𝔤𝒪_{})`$ is cyclic, i.e.
$$D^\mathrm{\Gamma }\mathrm{{\rm Y}}=0.$$
Proof. Applying $`\text{X}𝒯_{}`$ to both sides of the equation
$$d\mathrm{\Gamma }+\stackrel{}{}\mathrm{\Gamma }+\frac{1}{2}[\mathrm{\Gamma }\mathrm{\Gamma }]=0$$
we get
$$(1)^{\stackrel{~}{X}}d^\mathrm{\Gamma }(\stackrel{}{\text{X}}\mathrm{\Gamma })+\stackrel{}{[\text{X},]}\mathrm{\Gamma }=0,$$
implying $`(D^\mathrm{\Gamma }\mathrm{{\rm Y}})(\text{X})=0`$. $`\mathrm{}`$
### 4.5.2. Corollary.
For any$`\text{X}𝒯_{}`$, $`d^\mathrm{\Gamma }(\stackrel{}{\text{X}}\mathrm{\Gamma })=\stackrel{}{\delta \text{X}}\mathrm{\Gamma }`$.
### 4.5.3 Corollary.
For any$`\chi \text{Hom}(^k𝒯_{},𝒯_{})`$ one has
$$D^\mathrm{\Gamma }(\mathrm{{\rm Y}}\chi )=\mathrm{{\rm Y}}(\delta \chi ).$$
Proof. We have, using Corollary 4.5.2,
$`D^\mathrm{\Gamma }(\mathrm{{\rm Y}}\chi )(\text{X}_1,\mathrm{},\text{X}_k)`$ $`=`$ $`d^\mathrm{\Gamma }\left(\stackrel{}{\chi (\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }\right)(1)^{\stackrel{~}{\chi }}\stackrel{}{\chi (\delta \text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }`$
$`\mathrm{}(1)^{\stackrel{~}{\chi }+\stackrel{~}{X}_1+\mathrm{}+\stackrel{~}{X}_{k1}}\stackrel{}{\chi (\text{X}_1,\mathrm{},\delta \text{X}_k)}\mathrm{\Gamma }`$
$`=`$ $`\left(\stackrel{}{\delta \chi (\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }\right)(1)^{\stackrel{~}{\chi }}\stackrel{}{\chi (\delta X_1,\mathrm{},X_k)}\mathrm{\Gamma }`$
$`\mathrm{}(1)^{\stackrel{~}{\chi }+\stackrel{~}{X}_1+\mathrm{}+\stackrel{~}{X}_{k1}}\stackrel{}{\chi (\text{X}_1,\mathrm{},\delta \text{X}_k)}\mathrm{\Gamma }`$
$`=`$ $`\left(\mathrm{{\rm Y}}\delta \chi \right)(\text{X}_1,\mathrm{},\text{X}_k)`$
for arbitrary $`\text{X}_1,\mathrm{},\text{X}_k𝒯_{}`$. $`\mathrm{}`$
Therefore, the map
$$\mathrm{{\rm Y}}:(𝒯_{},\delta )(𝔤𝒪_{},d^\mathrm{\Gamma })$$
is a morphism of sheaves of complexes. Note that the “projection” map $`s\varphi :𝔤𝒪_{}𝒯_{}`$ satisfies $`s\varphi \mathrm{{\rm Y}}=\text{Id}`$ but does not, in general, respect the differentials.
Analogously one shows that the morphism $`\mathrm{{\rm Y}}\mathrm{{\rm Y}}`$ defined by the commutative diagram
defines a cyclic element in $`(\text{Hom}(𝒯_{}^2,𝔤𝒪_{}),D^\mathrm{\Gamma })`$. In a similar way one uses muliplicative structure in $`𝔤`$ to construct cyclic elements $`\mathrm{{\rm Y}}\mathrm{{\rm Y}}\mathrm{{\rm Y}}`$ etc.<sup>6</sup><sup>6</sup>6The cyclicity of $`\mathrm{{\rm Y}}\mathrm{{\rm Y}}`$, etc., relies on the Poisson identity holding in $`(𝔤,[],,d)`$.
For future reference we define a morphism $`\mathrm{{\rm Y}}_{(n)}\text{Hom}(𝒯,𝔤𝒪_{})`$ by setting $`\mathrm{{\rm Y}}_{(n+1)}(\text{X}):=\stackrel{}{\text{X}}\mathrm{\Gamma }_{(n+1)}`$. Similarly one defines $`\mathrm{{\rm Y}}_{(n)}\mathrm{{\rm Y}}_{(n)}`$, etc.
### 4.6. Multiplicative structure in $`𝒯_{}`$.
We will show in this subsection that, for any dG-algebra $`𝔤`$, the associated tangent sheaf $`𝒯_{}`$ is always a sheaf of differential associative algebras (defined uniquely up to a homotopy).
### 4.6.1. Theorem.
There exists an even morphism of sheaves, $`\mu \text{Hom}(𝒯_{}^2,𝒯_{})`$, such that $`\delta \mu =0`$ and the diagram
is commutative at the cohomology level, i.e.
$$[\mathrm{{\rm Y}}\mathrm{{\rm Y}}]=[\mathrm{{\rm Y}}\mu ]$$
in the cohomology sheaf$`𝖪erD^\mathrm{\Gamma }/𝖨mD^\mathrm{\Gamma }`$ associated with the sheaf of complexes$`(Hom(𝒯_{}^2,𝔤𝒪_{}),D^\mathrm{\Gamma })`$.
Proof. We have to show that there exists $`\mu \text{Hom}(𝒯_{}^2,𝒯_{})`$ such that
$$\delta \mu (\text{X},\text{Y})=\mu (\delta \text{X},\text{Y})+(1)^{\stackrel{~}{X}}\mu (\text{X},\delta \text{Y})$$
(5)
and
$$\stackrel{}{\text{X}}\mathrm{\Gamma }\stackrel{}{\text{Y}}\mathrm{\Gamma }=\stackrel{}{\mu (\text{X},\text{Y})}\mathrm{\Gamma }+(D^\mathrm{\Gamma }A)(\text{X},\text{Y})$$
(6)
for some $`A\text{Hom}(𝒯_{}^2,𝔤𝒪_{})`$ and any $`\text{X},\text{Y}𝒯_{}`$. We shall proceed by induction and assume, without loss of generality, that the vector fields X and Y are constant, i.e. $`\text{X}=\text{Y}=0`$.
The above equations can obviously be satisfied $`modI`$: just set
$$\mu _{[0]}(\text{X},\text{Y}):=\varphi (\stackrel{}{\text{X}}\mathrm{\Gamma }_{[1]}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{[1]}).$$
Indeed,
$$\stackrel{}{\text{X}}\mathrm{\Gamma }_{[1]}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{[1]}\stackrel{}{\mu _{[0]}(\text{X},\text{Y})}\mathrm{\Gamma }_{[1]}𝖪er\varphi 𝖪erd$$
and hence this expression is $`d`$-exact. Denote it by $`dA_{[0]}(\text{X},\text{Y})`$. (We can always normalise $`A_{[0]}`$ so that it lies in $`\text{Hom}(𝒯_{}^2,\mathrm{ker}\varphi 𝒪_{})`$.) This solves (6) $`modI`$. The equation (5)$`modI`$ is trivial (recall that $`_{[<2]}=0`$).
Assume now that we have constructed $`\mu _{(n)}`$ and $`A_{(n)}`$ so that the equations
$$\delta _{(n)}\mu _{(n1)}(\text{X},\text{Y})=\mu _{(n1)}(\delta _{(n)}\text{X},\text{Y})+(1)^{\stackrel{~}{X}}\mu _{(n1)}(\text{X},\delta _{(n)}\text{Y})modI^{n+1}$$
(7)
$$\stackrel{}{\text{X}}\mathrm{\Gamma }_{(n+1)}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{(n+1)}=\stackrel{}{\mu _{(n)}(\text{X},\text{Y})}\mathrm{\Gamma }_{(n+1)}+(D_{}^{\mathrm{\Gamma }}{}_{(n)}{}^{}A_{(n)})(\text{X},\text{Y})modI^{n+1}$$
(8)
hold.
The Theorem will be proved if we find $`\mu _{[n+1]}`$ and $`A_{[n+1]}`$ satisfying
$$\delta _{(n+1)}\left(\mu _{(n)}(\text{X},\text{Y})+\mu _{[n+1])}(\text{X},\text{Y})\right)=\mu _{(n)}(\delta _{(n+1)}\text{X},\text{Y})+(1)^{\stackrel{~}{X}}\mu _{(n)}(\text{X},\delta _{(n+1)}\text{Y})modI^{n+2}$$
and
$$\stackrel{}{\text{X}}\mathrm{\Gamma }_{(n+2)}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{(n+2)}\stackrel{}{\mu _{(n)}(\text{X},\text{Y})}\mathrm{\Gamma }_{(n+2)}\stackrel{}{\mu _{[n+1]}(\text{X},\text{Y})}\mathrm{\Gamma }_{[1]}(D_{}^{\mathrm{\Gamma }}{}_{(n+1)}{}^{}A_{(n)})(\text{X},\text{Y})=$$
$$=dA_{[n+1]}(\text{X},\text{Y})modI^{n+2}$$
Defining
$$\mu _{[n+1]}(\text{X},\text{Y}):=\varphi \left(\stackrel{}{\text{X}}\mathrm{\Gamma }_{(n+2)}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{(n+2)}\stackrel{}{\mu _{(n)}(\text{X},\text{Y})}\mathrm{\Gamma }_{(n+2)}(D_{}^{\mathrm{\Gamma }}{}_{(n+1)}{}^{}A_{(n)})(\text{X},\text{Y})\right)$$
we ensure that the morphism
$$\lambda _{[n+1]}(\text{X},\text{Y}):=(\mathrm{{\rm Y}}_{(n+1)}\mathrm{{\rm Y}}_{(n+1)}\mathrm{{\rm Y}}_{(n+1)}\mu _{(n)}D_{}^{\mathrm{\Gamma }}{}_{(n+1)}{}^{}A_{(n)})\mathrm{{\rm Y}}_{(0)}\mu _{[n+1]})(\text{X},\text{Y})modI^{n+2}$$
take values in the sheaf in $`𝖪er\varphi 𝒪_{}`$. Since it vanishes modulo $`I^{n+1}`$, we have, modulo $`I^{n+2}`$,
$`d\lambda _{[n+1]}(\text{X},\text{Y})`$ $`=`$ $`(D_{(n+1)}^\mathrm{\Gamma }\lambda _{[n+1]})(\text{X},\text{Y})`$
$`=`$ $`\left(D_{(n+1)}^\mathrm{\Gamma }(\mathrm{{\rm Y}}_{(n+1)}\mu _{(n)})\right)(\text{X},\text{Y})`$
$`=`$ $`d_{n+1}^\mathrm{\Gamma }(\stackrel{}{\mu _{(n)}(\text{X},\text{Y})}\mathrm{\Gamma })+\stackrel{}{\mu _{(n)}(\delta _{(n+1)}\text{X},\text{Y})}\mathrm{\Gamma }+(1)^{\stackrel{~}{X}}\stackrel{}{\mu _{(n)}(\text{X},\delta _{(n+1)}\text{Y})}\mathrm{\Gamma }`$
$`=`$ $`\stackrel{}{(\delta _{(n+1)}\mu _{(n)}(\text{X},\text{Y})\mu _{(n)}(\delta _{(n+1)}\text{X},\text{Y})(1)^{\stackrel{~}{X}}\mu _{(n)}(\text{X},\delta _{(n+1)}\text{Y}))}\mathrm{\Gamma }`$
where we have used Corollary 4.5.2 and the fact that $`D^\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{{\rm Y}})=0`$. Applying $`\varphi `$ to the last equation, we get
$$\delta _{(n+1)}\mu _{(n)}(\text{X},\text{Y})=\mu _{(n)}(\delta _{(n+1)}\text{X},\text{Y})+(1)^{\stackrel{~}{X}}\mu _{(n)}(\text{X},\delta _{(n+1)}\text{Y})modI^{n+2}$$
and hence
$$d\lambda _{[n+1]}(\text{X},\text{Y})=0.$$
Since $`𝖪er\varphi 𝖪erd𝖨md`$, there exists $`A_{[n+1]}(\text{X},\text{Y})`$ (which can be chosen to lie in $`𝖪er\varphi 𝒪_{}`$) such that
$$\lambda _{[n+1]}(\text{X},\text{Y})=dA_{[n+1]}(\text{X},\text{Y}).$$
This completes the inductive procedure and hence the proof of the Theorem. $`\mathrm{}`$
### 4.6.2. Definition.
An even morphism of sheaves, $`\mu \text{Hom}(^2𝒯_{},𝒯_{})`$, satisfying the conditions of Theorem 4.6.1 is called induced. The associated data $`(𝒯_{},\delta ,\mu )`$ is called a sheaf of induced differential algebras.
Clearly, an induced product on $`𝒯_{}`$ is supercommutative if the product $``$ in $`𝔤`$ is supercommutative.
### 4.7. (Non)Uniqueness.
How unique is the product $`\mu `$ induced on the tangent sheaf $`𝒯_{}`$ by Theorem 4.6.1? When is it associative? To address these questions we shall need the following technical result.
### 4.7.1. Lemma.
If$`\tau \text{Hom}(^k𝒯_{},𝒯_{})`$ and $`B\text{Hom}(^k𝒯_{},𝔤𝒪_{})`$ satisfy the equation
$$\mathrm{{\rm Y}}\tau =D^\mathrm{\Gamma }B$$
then there exists$`\chi \text{Hom}(^k𝒯_{},𝒯_{})`$ and $`C\text{Hom}(^k𝒯_{},𝖪er\varphi 𝒪_{})`$ such that
* $`B=\mathrm{{\rm Y}}\chi +D^\mathrm{\Gamma }C`$,
* $`\tau =\delta \chi `$, i.e.
$`\tau (\text{X}_1,\mathrm{},\text{X}_k)`$ $`=`$ $`\delta \chi (\text{X}_1,\mathrm{},\text{X}_k)(1)^{\stackrel{~}{\chi }}\chi (\delta \text{X}_1,\mathrm{},\text{X}_k)`$
$`\mathrm{}(1)^{\stackrel{~}{\chi }+\stackrel{~}{X}_1+\mathrm{}+\stackrel{~}{X}_{k1}}\chi (\text{X}_1,\mathrm{},\delta \text{X}_k)`$
for any $`\text{X}_1,\mathrm{},\text{X}_k𝒯_{}`$.
Proof. Without loss of generality we may assume that (germs of) vectors fields $`\text{X}_1,\mathrm{},\text{X}_k`$ are constant. In view of Corollary 4.5.3 and injectivity of the map $`\mathrm{{\rm Y}}`$, it is enough to show that the equation
$$\stackrel{}{\tau (\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }=(D^\mathrm{\Gamma }B)(\text{X}_1,\mathrm{},\text{X}_k)$$
(9)
implies
$$B(\text{X}_1,\mathrm{},\text{X}_k)=\stackrel{}{\chi (\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }+(D^\mathrm{\Gamma }C)(\text{X}_1,\mathrm{},\text{X}_k)$$
for some $`\chi \text{Hom}(^k𝒯_{},𝒯_{})`$ and $`C\text{Hom}(^k𝒯_{},𝖪er\varphi 𝒪_{})`$. We shall proceed by induction.
The equation (9$`modI`$ is
$$\stackrel{}{\tau _{[0]}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{[1]}=dB_{[0]}(\text{X}_1,\mathrm{},\text{X}_k).$$
Hence $`\tau _{[0]}=0`$ and $`dB_{[0]}=0`$. Set
$$\chi _{[0]}(\text{X}_1,\mathrm{},\text{X}_k):=\varphi \left(B_{[0]}(\text{X}_1,\mathrm{},\text{X}_k)\right).$$
Then $`B_{[0]}(\text{X}_1,\mathrm{},\text{X}_k)\stackrel{}{\chi _{[0]}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{[1]}`$ lies in $`(𝖪er\varphi 𝖪erd)𝒪_{}`$ and hence equals $`dC_{[0]}(\text{X}_1,\mathrm{},\text{X}_k)`$ for some $`C_{[0]}\text{Hom}(^k𝒯_{},𝖪er\varphi 𝒪_{})`$.
Assume that $`\chi _{(n)}`$ and $`C_{(n)}`$ are constructed so that the equations
$$B_{(n)}(\text{X}_1,\mathrm{},\text{X}_k)=\stackrel{}{\chi _{(n)}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{(n+1)}+(D_{(n)}^\mathrm{\Gamma }C_{(n)})(\text{X}_1,\mathrm{},\text{X}_k)modI^{n+1}$$
holds.
Let us show that there exist $`\chi _{[n+1]}`$ and $`C_{[n+1]}`$ satisfying
$`B_{(n+1)}(\text{X}_1,\mathrm{},\text{X}_k)`$ $`=`$ $`\stackrel{}{\chi _{(n)}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{(n+2)}+\stackrel{}{\chi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{[1]}`$
$`+(D_{(n+1)}^\mathrm{\Gamma }C_{(n)})(\text{X}_1,\mathrm{},\text{X}_k)+dC_{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)modI^{n+2},`$
or, equivalently, satisfying
$$dC_{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)=\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)\stackrel{}{\chi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{[1]}modI^{n+2},$$
where we have set
$$\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k):=B_{(n+1)}(\text{X}_1,\mathrm{},\text{X}_k)\stackrel{}{\chi _{(n)}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{(n+2)}(D_{(n+1)}^\mathrm{\Gamma }C_{(n)})(\text{X}_1,\mathrm{},\text{X}_k).$$
Since $`\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)`$ vanishes $`modI^{n+1}`$, this is a monom of degree $`n+1`$ and $`t^i`$, and hence, modulo $`I^{n+2}`$,
$`d\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)`$ $`=`$ $`(D_{(n+1)}^\mathrm{\Gamma }\psi _{[n+1]})(\text{X}_1,\mathrm{},\text{X}_k)`$
$`=`$ $`(D_{(n+1)}^\mathrm{\Gamma }B_{(n+1)})(\text{X}_1,\mathrm{},\text{X}_k)D_{(n+1)}^\mathrm{\Gamma }(\mathrm{{\rm Y}}_{(n+1)}\chi _{(n)})(\text{X}_1,\mathrm{},\text{X}_k)`$
$`=`$ $`\stackrel{}{\left(\tau _{(n+1)}(\text{X}_1,\mathrm{},\text{X}_k)(\delta _{(n+1)}\chi _{(n)})(\text{X}_1,\mathrm{},\text{X}_k)\right)}\mathrm{\Gamma }.`$
Applying $`\varphi `$ to both sides of these equations we get
$$\tau _{(n+1)}(\text{X}_1,\mathrm{},\text{X}_k)=(\delta _{(n+1)}\chi _{(n)})(\text{X}_1,\mathrm{},\text{X}_k)modI^{n+1}$$
and hence
$$d\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)=0modI^{(n+1)},$$
We define $`\chi _{[n+1]}`$ by
$$\chi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k):=\varphi (\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)).$$
Then $`\psi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)\stackrel{}{\chi _{[n+1]}(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }_{[1]}𝖪erd𝖪er\varphi 𝖨md`$. This proves the existence of $`C_{[n+1]}`$ and hence completes the proof of the theorem. $`\mathrm{}`$
If $`\mu ^{},\mu ^{\prime \prime }\text{Hom}(^2𝒯_{},𝒯_{})`$ are two products as in Theorem 4.6.1, then
$$\mathrm{{\rm Y}}(\mu ^{}\mu ^{\prime \prime })=D^\mathrm{\Gamma }B$$
for some $`B\text{Hom}(^2𝒯_{},𝔤𝒪_{})`$, and hence, by Lemma 4.7.1,
$$\mu ^{}\mu ^{\prime \prime }=\delta \chi $$
for some odd $`\chi \text{Hom}(^2𝒯_{},𝒯_{})`$, i.e. $`\mu ^{}`$ and $`\mu ^{\prime \prime }`$ are what is called homotopy equivalent.
On the other hand, if $`\mu ^{\prime \prime }`$ is a product with the properties stated by Theorem 4.6.1, then, for any odd $`\chi \text{Hom}(^2𝒯_{},𝒯_{})`$, the product
$$\mu ^{}:=\mu ^{\prime \prime }+\delta \chi $$
also enjoyes the properties of Theorem 4.6.1. Indeed, by Corollary 4.6.3,
$`\mathrm{{\rm Y}}\mu ^{\prime \prime }`$ $`=`$ $`\mathrm{{\rm Y}}\mu ^{}+\mathrm{{\rm Y}}(\delta \chi )`$
$`=`$ $`\mathrm{{\rm Y}}\mu ^{}+D^\mathrm{\Gamma }(\mathrm{{\rm Y}}\chi )`$
and hence
$$[\mathrm{{\rm Y}}\mathrm{{\rm Y}}]=[\mathrm{{\rm Y}}\mu ^{}]=[\mathrm{{\rm Y}}\mu ^{}]$$
in the cohomology sheaf $`𝖪erD^\mathrm{\Gamma }/𝖨mD^\mathrm{\Gamma }`$.
Thus the set of products $`\mu `$ induced on $`𝒯_{}`$ by Theorem 4.6.1 is a principal homogeneous space over the Abelian group $`\delta \text{Hom}_{\stackrel{~}{1}}(^2𝒯_{},𝒯_{})`$. Hence all induced products on each stalk of $`𝒯_{}`$ combine into a single homotopy class which we call induced.
### 4.7.2. Theorem.
The sheaf $`𝒯_{}`$ is canonically a sheaf of induced homotopy classes of differential algebras.
Proof. We have to show that the homotopy class of products induced on $``$ is an invariant of the dG-algebra under consideration, i.e. that it is independent of the choice of a quasi-isomorphism $`\varphi `$ and on the choice of a Master function $`\mathrm{\Gamma }`$ used in its construction. In view of Remark 4.1.3(iv), it is enough to check the invariance of the product under the gauge transformations,
$$\mathrm{\Gamma }\mathrm{\Gamma }^g:=e^{\mathrm{a}d_g}\mathrm{\Gamma }\frac{e^{\mathrm{a}d_g}1}{\mathrm{a}d_g}(d+\stackrel{}{})g,g\mathrm{\Gamma }(,𝔤_{\stackrel{~}{1}}).$$
A straightforward analysis of the basic equation (5) shows that gauge transformation changes the tensor $`A`$,
$$A^g=e^{\mathrm{a}d_g}\left(A(X,Y)G\mathrm{{\rm Y}}\mathrm{{\rm Y}}G+GD^\mathrm{\Gamma }G+G\mu \right),$$
where $`G\text{Hom}(𝒯_{},𝔤𝒪_{})`$ is given by
$$G(X):=(1)^{\stackrel{~}{X}}\frac{e^{\mathrm{a}d_g}1}{\mathrm{a}d_g}(d+\stackrel{}{})\stackrel{}{\text{X}}g,$$
but leaves the product invariant, $`\mu ^g=\mu `$. $`\mathrm{}`$
### 4.8. Identity in $`𝔤`$ $``$ identity in $`𝒯_{}`$.
If the dG-algebra under consideration, $`𝔤`$, has the identity $`e_0`$, and the versal solution $`\mathrm{\Gamma }`$ is approprietly normalized (see Remark 3.1.1), then
$`\stackrel{}{\delta (e)}\mathrm{\Gamma }`$ $`=`$ $`\stackrel{}{[,e]}\mathrm{\Gamma }`$
$`=`$ $`\stackrel{}{}e_0+\stackrel{}{e}(d\mathrm{\Gamma }+{\displaystyle \frac{1}{2}}[\mathrm{\Gamma }\mathrm{\Gamma }])`$
$`=`$ $`0+de_0+[e_0\mathrm{\Gamma }]`$
$`=`$ $`0,`$
so that $`\delta (e)=0`$. We shall show next that the induced homotopy class of differential algebras on each stalk of $`𝒯_{}`$ containes a canonical subclass of unital differential algebras.
### 4.8.1. Theorem.
If $`𝔤`$ has the identity $`e_o`$, then $`𝒯_{}`$ is canonically a sheaf of induced homotopy classes of differential algebras with the identity $`e`$.
A sketch of the proof. It is enough to show that there exists a $`\delta `$-closed element, $`\mu \text{Hom}_{\stackrel{~}{0}}(^2𝒯_{},𝒯_{})`$, such that, for arbitrary (constant) $`\text{X},\text{Y}𝒯_{}`$, the equation
$$\stackrel{}{\text{X}}\mathrm{\Gamma }\stackrel{}{\text{Y}}\mathrm{\Gamma }=\stackrel{}{\mu (\text{X},\text{Y})}\mathrm{\Gamma }+(D^\mathrm{\Gamma }A)(\text{X},\text{Y})$$
holds for some $`A\text{Hom}(𝒯_{}^2,𝔤𝒪_{})`$ satisfying $`A(\text{X},e)=\text{X}`$ and $`A(e,\text{Y})=\text{Y}`$ (cf. (6)).
Recall (see the proof of Theorem 4.6.1) that at the lowest order we have,
$`\mu _{[0]}(\text{X},\text{Y})`$ $`=`$ $`\varphi (\stackrel{}{\text{X}}\mathrm{\Gamma }_{[1]}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{[1]})`$
$`dA_{[0]}(\text{X},\text{Y})`$ $`=`$ $`\stackrel{}{\text{X}}\mathrm{\Gamma }_{[1]}\stackrel{}{\text{Y}}\mathrm{\Gamma }_{[1]}\stackrel{}{\mu _{[0]}(\text{X},\text{Y})}\mathrm{\Gamma }_{[1]}.`$
and hence $`\mu _{[0]}(\text{X},e)=\mu _{[0]}(e,\text{X})=\text{X}`$ and $`dA_{[0]}(\text{X},e)=dA_{[0]}(e,\text{X})=0`$. We claim that $`A_{[0]}`$ can be chosen to satisfy $`A_{[0]}(\text{X},e)=A_{[0]}(e,\text{X})=0`$. This can be achieved by a replacement,
$$\begin{array}{ccccc}A_{[0]}(\text{X},\text{Y})& & A_{[0]}^{}(\text{X},\text{Y})& :=& A_{[0]}(\text{X},\text{Y})A_{[0]}(\text{X},e)\stackrel{}{\text{Y}}\mathrm{\Gamma }(1)^{\stackrel{~}{X}}\stackrel{}{\text{X}}\mathrm{\Gamma }A_{[0]}(e,\text{Y})\hfill \\ & & & & +(1)^{\stackrel{~}{X}}\stackrel{}{\text{X}}\mathrm{\Gamma }A_{[0]}(e,e)\stackrel{}{\text{Y}}\mathrm{\Gamma },\hfill \end{array}$$
which satisfies,
$$dA_{[0]}^{}(\text{X},\text{Y})=dA_{[0]}(\text{X},\text{Y}),A_{[0]}^{}(\text{X},e)=A_{[0]}^{}(e,\text{X})=0.$$
This observation allows us to include into the inductive procedure of the proof of Theorem 4.6.1 the additional assumptions
$$\mu _{(n)}(\text{X},e)=\mu _{(n)}(e,\text{X})=\text{X},A_{(n)}(\text{X},e)=A_{(n)}(e,\text{X})=0,$$
and show, by exactly the same argument as in the case $`n=0`$ above, that they hold true for $`n+1`$.
Thus there does exist a product $`\mu `$ from the induced homotopy class satisfying $`\mu (\text{X},e)=\mu (e,\text{X})=\text{X}`$. It is defined uniquely up to a transformation
$$\mu \mu +\delta \chi $$
with $`\chi `$ satisfying $`\chi (\text{X},e)=\chi (e,\text{X})=0`$ for arbitrary $`\text{X}𝒯_{}`$. Thus what is well-defined is the induced homotopy class of unital differential algebras. $`\mathrm{}`$
### 4.9. Theorem.
For any (unital) dG-algebra $`𝔤`$, the tangent sheaf $`𝒯_{}`$ to its cohomology supermanifold is canonically a sheaf of homotopy classes of (unital) $`A_{\mathrm{}}`$-algebras with
* $`\mu _1=[,\mathrm{}]`$, $``$ being the Chen’s vector field, and
* the homotopy class of $`\mu _2`$ being the induced homotopy class as in Theorem 4.6.1.
A sketch of the proof. By Theorem 4.6.1, there exists a product $`\mu _2\text{Hom}(^2𝒯_{},𝒯_{})`$ satisfying the equation
$$\stackrel{}{\mu _2(\text{X}_1,\text{X}_2)}\mathrm{\Gamma }=\stackrel{}{\text{X}_1}\mathrm{\Gamma }\stackrel{}{\text{X}_2}\mathrm{\Gamma }+(D^\mathrm{\Gamma }A_2)(\text{X}_1,\text{X}_2)$$
for some odd $`A_2\text{Hom}(^2𝒯_{},𝔤𝒪_{})`$ and arbitrary $`\text{X}_1,\text{X}_2𝒯_{}`$. We have, in the notations of subsection 3.2.3,
$`\stackrel{}{\mathrm{\Lambda }_3(\text{X}_1,\text{X}_2,\text{X}_3)}\mathrm{\Gamma }`$ $`=`$ $`\stackrel{}{\mu _2(\text{X}_1,\mu _2(\text{X}_2,\text{X}_3))\mu _2(\mu _2(\text{X}_1,\text{X}_2),\text{X}_3)}\mathrm{\Gamma }`$
$`=`$ $`\stackrel{}{\text{X}_1}\mathrm{\Gamma }\left(\stackrel{}{\text{X}_2}\mathrm{\Gamma }\stackrel{}{\text{X}_3}\mathrm{\Gamma }\right)+(D^\mathrm{\Gamma }A_2)(\text{X}_1,\mu _2(\text{X}_2,\text{X}_4))+\mathrm{{\rm Y}}(\text{X}_1)(D^\mathrm{\Gamma }A_2)(\text{X}_2,\text{X}_3)`$
$`\left(\stackrel{}{\text{X}_1}\mathrm{\Gamma }\stackrel{}{\text{X}_2}\mathrm{\Gamma }\right)\stackrel{}{\text{X}_3}\mathrm{\Gamma }(D^\mathrm{\Gamma }A_2)(\mu _2(\text{X}_1,\text{X}_2),\text{X}_3)(D^\mathrm{\Gamma }A_2)(\text{X}_1,\text{X}_2)\mathrm{{\rm Y}}(\text{X}_3)`$
$`=`$ $`(D^\mathrm{\Gamma }B_3)(\text{X}_1,\text{X}_2,\text{X}_3),`$
where
$`B_3(\text{X}_1,\text{X}_2,\text{X}_3)`$ $`:=`$ $`(1)^{\stackrel{~}{X}_1}\mathrm{{\rm Y}}(\text{X}_1)A_2(\text{X}_2,\text{X}_3)A_2(\text{X}_1,\text{X}_2)\mathrm{{\rm Y}}(\text{X}_3)`$
$`+A_2(\text{X}_1,\mu _2(\text{X}_2,\text{X}_3))A_2(\mu _2(\text{X}_1,\text{X}_2),\text{X}_3).`$
Here we used associativity of the dot product in $`𝔤`$, $`D^\mathrm{\Gamma }`$-closedness of $`\mathrm{{\rm Y}}`$ and $`\delta `$-closedness of $`\mu _2`$.
By Lemma 4.7.1, there exists $`\mu _3\text{Hom}_{\stackrel{~}{0}}(^3𝒯_{},𝒯_{})`$ such that the 3rd order associativity condition, $`\mathrm{\Lambda }_3=\delta \mu _3`$, is satisfied, and
$$\stackrel{}{\mu _3(\text{X}_1,\text{X}_2,\text{X}_3)}\mathrm{\Gamma }=B_3(\text{X}_1,\text{X}_2,\text{X}_3)+(D^\mathrm{\Gamma }A_3)(\text{X}_1,\text{X}_2,\text{X}_3)$$
for some $`A_3\text{Hom}_{\stackrel{~}{0}}(^3𝒯_{},𝔤𝒪_{})`$.
Exactly the same procedure constructs inductively all the higher order products $`\mu _n\text{Hom}_{\stackrel{~}{n}}(^n𝒯_{},𝒯_{})`$ satisfying the higher order associativity conditions:
* Assume that we have constructed $`\mu _k\text{Hom}_{\stackrel{~}{k}}(^k𝒯_{},𝒯_{})`$ and $`A_k\text{Hom}_{\stackrel{~}{k}+\stackrel{~}{1}}(^k𝒯_{},𝔤𝒪_{})`$, $`k=2,\mathrm{},n1`$, such that $`\mathrm{\Lambda }_k=\delta \mu _k`$ ($`k`$-th order associativity condition) and
$`\stackrel{}{\mu _k(\text{X}_1,\mathrm{},\text{X}_k)}\mathrm{\Gamma }`$ $`=`$ $`{\displaystyle \underset{i+j=k}{}}(1)^{(j+1)(\stackrel{~}{X}_1+\mathrm{}+\stackrel{~}{X}_i)+i+1}A_i(\text{X}_1,\mathrm{},\text{X}_i)A_i(\text{X}_{i+1},\mathrm{},\text{X}_k)`$
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{\genfrac{}{}{0pt}{}{i+j=k+1}{i2}}{j2}}{}}{\displaystyle \underset{l=0}{\overset{i1}{}}}(1)^rA_i(\text{X}_1,\mathrm{},\text{X}_l,\mu _j(\text{X}_{l+1},\mathrm{},\text{X}_{l+j}),\text{X}_{l+j+1},\mathrm{},\text{X}_k)`$
$`+(D^\mathrm{\Gamma }A_k)(\text{X}_1,\mathrm{},\text{X}_k)`$
$`=:`$ $`B_k(\text{X}_1,\mathrm{},X_k)+(D^\mathrm{\Gamma }A_k)(\text{X}_1,\mathrm{},\text{X}_k),`$
where we have set $`A_1:=\mathrm{{\rm Y}}`$ and $`r=\stackrel{~}{j}(\stackrel{~}{\text{X}_1}+\mathrm{}\stackrel{~}{\text{X}_1})+\stackrel{~}{l}(\stackrel{~}{j}1)+(\stackrel{~}{i}\stackrel{~}{1})\stackrel{~}{j}+1`$.
* Use the above expressions for $`\mathrm{{\rm Y}}\mu _k`$, $`k=2,\mathrm{},n1`$, to show that
$$\stackrel{}{\mathrm{\Lambda }_n(\text{X}_1,\mathrm{},\text{X}_n)}\mathrm{\Gamma }=(D^\mathrm{\Gamma }B_n)(\text{X}_1,\mathrm{},\text{X}_n).$$
* Apply Lemma 3.6.1 to conclude that there exists $`\mu _n`$ such that $`\mathrm{\Lambda }_n=\delta \mu _n`$ ($`n`$-th order associativity condition) and $`\mathrm{{\rm Y}}\mu _n=B_n+D^\mathrm{\Gamma }A_n`$ for some $`A_n\text{Hom}_{\stackrel{~}{n}+\stackrel{~}{1}}(^k𝒯_{},𝔤𝒪_{})`$
Finally, we note that at each stage of the above construction the $`n`$th product $`\mu _n`$ is defined only up to a $`\delta `$-exact term, $`\delta f_n`$. These arbitrary terms combine all together into a homotopy of the $`A_{\mathrm{}}`$-structure $`(𝒯_{},\mu _{})`$. $`\mathrm{}`$
### 4.9.1. Corollary.
The cohomology sheaf on$``$,
$$_𝐇:=\frac{𝖪er\delta }{𝖨m\delta },$$
is canonically a sheaf of induced (unital) associative algebras.
### 4.9.2. Corollary.
The tangent sheaf, $`𝒯_{\mathrm{s}mooth}`$, to the smooth part of the extended Kuranishi space is canonically a sheaf of induced (unital) associative algebras.
### 4.9.3. The Euler field.
If the dG-algebra $`𝔤`$ under consideration is $``$-graded, then the cohomology $`𝐇`$ and hence its dual $`𝐇^{}`$ are also $``$-graded. We make $`k[[t]]`$ into a $``$-graded ring my setting
$$k[[t]]=^{}𝐇^{}[2].$$
This also induces $``$-grading in the sheaves $`𝒪_H`$ and $`𝒯_{}`$ on the supermanifold $``$.
If $`\{[e_i]\}`$ is a basis in $`𝐇`$ and $`\{t^i\}`$ are the associated linear coordinates on $``$ as in Sect. 2.2, then
$$|t^i|=2|[e_i]|.$$
With this choice of $``$-grading on $`𝒪_{}`$ we ensure that $`|\mathrm{\Gamma }|=2`$ and hence $`||=1`$, $`|\delta |=1`$, and $`|\mu _n|=n`$ for all the induced higher order products on $`𝒯_{}`$.
The Euler field on $`M`$ is, by definition, the derivation $`E`$ of $`k[[t]]`$ given by
$$Ef=\frac{1}{2}|f|f,fk[[t]].$$
In coordinates,
$$E=\frac{1}{2}\underset{i}{}|t^i|t^i\frac{}{t^i}.$$
This vector field generates the scaling symmetry on $`(,\mu _{})`$ (cf. \[BK, Ma2\]). If we decompose
$$\mu _n(\frac{}{t^{i_1}},\mathrm{},\frac{}{t^{i_n}})=\underset{k}{}\mu _{i_1\mathrm{}i_n}^k(t)\frac{}{t^k},$$
then, as follows from the explicit construction of $`\mu _n`$ given in the proof of Theorem 4.9,
$$E\mu _{i_1\mathrm{}i_n}^k=\frac{1}{2}\left(|t^k||t^{i_1}|\mathrm{}|t^{i_n}|+n\right)\mu _{i_1\mathrm{}i_n}^k.$$
Note also that in the presence of identity, $`[e,E]=e`$.
### 4.9.4. The perturbative proof of Theorem A.
The required statement follows immediately from the graded commutative version of Theorem 4.9 and Sect. 4.9.3. $`\mathrm{}`$
### 4.9.5. A generalization to $`G_{\mathrm{}}`$-algebras.
In the perturbative construction of the $`F_{\mathrm{}}`$-functor for dG-algebras the odd Poisson identity was used in a few places. For example, in the construction of $`\mu _2`$ the only place where we relied on it was the cyclicity of $`\mathrm{{\rm Y}}\mathrm{{\rm Y}}`$,
$$D^\mathrm{\Gamma }(\mathrm{{\rm Y}}\mathrm{{\rm Y}})=0.$$
However, a glance at the basic equation (6) (and its higher order analogues in Sect. 4.9) shows that the perturbative argument stands if the cyclicity (and its analogues) holds only up to a homotopy. Therefore, the generalization from dG-algebras to $`G_{\mathrm{}}`$-algebras is straightforward, affecting only auxiliary tensors $`A_n`$.
Acknowledgement. This work was done during author’s visit to the Max Planck Institute for Mathematics in Bonn. Excellent working conditions in the MPIM are gratefully acknowledged. I would like to thank Yu.I. Manin for many stimulating discussions, and A.A. Voronov for valuable communications. I am also grateful to J. Stasheff for useful comments.
| Max Planck Institute for Mathematics in Bonn, and |
| --- |
| Department of Mathematics, University of Glasgow |
| sm@maths.gla.ac.uk | |
warning/0001/physics0001011.html | ar5iv | text | # The Hydration Number of Li+ in Liquid Water*footnote **footnote *LA-UR-99-3360.
## Abstract
A theoretical treatment based upon the quasi-chemical theory of solutions predicts the most probable number of water neighbors in the inner shell of a Li<sup>+</sup> ion in liquid water to be four. The instability of a six water molecule inner sphere complex relative to four-coordinated structures is confirmed by an ‘ab initio’ molecular dynamics calculation. A classical Monte Carlo simulation equilibrated 26 water molecules with a rigid six-coordinated Li(H<sub>2</sub>O)$`{}_{6}{}^{}^+`$ complex with periodic boundary conditions in aqueous solution. With that initial configuration for the molecular dynamics, the six-coordinated structure relaxed into four-coordinated arrangements within 112 fs and stabilized. This conclusion differs from prior interpretations of neutron and X-ray scattering results on aqueous solutions.
The hydration of ions in water is not only fundamental to physical chemistry but also relevant to the current issue of selectivity of biological ion channels. In the context of potassium channels, for example, the free energies for replacement of inner shell water ligands with peptide carbonyls donated by proteins of the channel structure seem decisive to the selectivity of the channel, specifically for preference of K<sup>+</sup> over Na<sup>+</sup>. Studies to elucidate the thermodynamic features of such inner shell exchange reactions require prior knowledge of the ion hydration structures and energetics.
Unfortunately, our understanding of the inner hydration shell structure of ions in water is not as clear as it might be. The simplest and most favorable case to pursue is the Li<sup>+</sup> solute. Neutron scattering measurements on LiCl solutions in liquid water have led to a firm conclusion that the Li<sup>+</sup> ion has six near-neighbor water molecule partners. That result, however, has not been entirely uniform across studies of similar aqueous solutions containing Li<sup>+</sup> ions. X-ray scattering results have been interpreted similarly to indicate a hydration number of six, again with some nonuniformity. In contrast, some spectroscopic studies have suggested tetrahedral coordination of the Li<sup>+</sup> ion in water and an array of physical chemical inferences lend some support to that conclusion. On the theoretical side, electronic structure calculations on the Li<sup>+</sup> ion with six water molecules predict a slightly, but distinctly, lower energy for a structure with four inner shell and two outer shell water molecules than for structures with six water molecules in the innermost shell; results such as those seem to be universally supported by other electronic structure efforts. Simulations have produced a range of results including both four and six inner shell water neighbors with considerable statistical dispersion. It is well recognized, of course, that simulations are typically not designed to provide a sole determination of such properties, though they do shed light on the issues determining the hydration number of ions in water.
The theoretical scheme used here to address these problems for the Li<sup>+</sup>(aq) ion is based upon the quasi-chemical organization of solution theory, which is naturally suited to these problems. The first step is the study of the reactions
$$Li^++nH_2OLi(H_2O)_n^+$$
(1)
that combine $`n`$ water molecule ligands with the Li<sup>+</sup> ion in a geometrically defined inner sphere under ideal gas conditions. At a subsequent step an approximate, physical description of the aqueous environment surrounding these complexes is included. The geometric definition of an inner sphere region enforces a physical balance in this method. The goal of this approach is to treat inner sphere ligands explicitly, in molecular detail, but at the same time to achieve a description of outer sphere hydration thermodynamics that is consistent from one complex to another. If minimum energy complex geometries were to shift different numbers of ligands to outer sphere regions, that would unbalance the thermodynamic description of the hydration of the inner sphere materials. For example, in the quantitative implementation of the quasi-chemical approach we specifically do not use the Li\[(H<sub>2</sub>O)<sub>4</sub>\]\[(H<sub>2</sub>O)<sub>2</sub>\]<sup>+</sup> complex cited above, with two water molecules outside the inner sphere, even though this structure helpfully clarifies the physical issue.
Gas-phase thermochemical data required for the equilibria in Eq. (1) were obtained by electronic structure calculations using the Gaussian98 programs with the B3LYP hybrid density functional theory approximation. All structures were fully optimized with a basis including polarization functions on Li<sup>+</sup> (6-31G\*) and both polarization and diffuse functions (6-31++G\**) on the oxygen and hydrogen centers. At the optimum geometry and with the same basis set, harmonic vibrational frequencies of the clusters were calculated and atomic charges determined using the ChelpG capability in Gaussian98. Partition functions were then calculated, thus providing a determination of the free energy changes of the equilibria in Eq. (1) due to atomic motions internal to the clusters within the harmonic approximation.
Interactions of these complexes with the external aqueous environment were treated with a dielectric model following the previous study of the hydrolysis of the ferric ion. Classic electrostatic interactions based upon the ChelpG partial atomic charges were the only solution-complex interactions treated; in particular, repulsive force (overlap) interactions were neglected based on the expectation that they make a secondary contribution to the thermodynamic properties considered here. The external boundary of the volume enclosed by spheres centered on all atoms defined the solute molecular surface. The sphere radii were those determined empirically by Stefanovich and Truong, except R$`_{Li^+}`$=2.0 Å for the lithium ion. Because the lithium ion is well buried by the inner shell waters, slight variations of the lithium radius were found to be unimportant. The value R$`_{Li^+}`$=2.0 Å was identified as slightly larger than the nearest Li-O distances and significantly smaller than the Li-O distances (3.5 – 4.0 Å) for second shell pairs.
Results of the calculations are summarized in Fig. 1. Geometry optimization of each of the $`n`$-coordinated clusters confirms that the inner shell structures used in these calculations are not necessarily the lowest energy structures for a given number of water neighbors. Although a tetrahedral cluster of inner shell water molecules is the lowest energy structure for Li(H<sub>2</sub>O)$`{}_{4}{}^{}^+`$, a cluster with five inner shell water molecules is slightly higher in energy than a cluster with one outer shell and four inner shell water molecules. Similarly, the lowest energy cluster with six water molecules contains four inner shell water molecules arranged tetrahedrally and two outer shell water molecules.
Fig. 1 shows that the $`n`$=4 inner sphere cluster has the lowest free energy for a dilute (p=1 atm) ideal gas phase. Adjustment of the concentration of water molecules to the value $`\rho _W`$ = 1 g/cm<sup>3</sup>, to match the normal density of liquid water, changes the most favored cluster to the one with $`n`$=6 inner shell water molecules. Outer sphere interactions described by the dielectric model progressively destabilize the larger clusters, as they should since larger numbers of water molecules are being treated explicitly as members of the inner shell. As a consequence of including the outer sphere contributions, the final position of minimum free energy is returned to the $`n`$=4 structure, with the $`n`$=3 complex predicted to be next most populous in liquid water at T=298.15 K and p=1 atm. The mean hydration number predicted by this calculation is $`\overline{n}`$=4.0.
The current quasi-chemical prediction for the absolute hydration free energy of the Li<sup>+</sup> ion under these conditions is -128 kcal/mol, not including any repulsive force (packing) contributions. An extreme increase of R$`_{Li^+}`$ to 2.65 Å raises this value to about -126 kcal/mol, showing that the theoretical results are insensitive to the ion radius, as remarked above. Experimental values are -113 kcal/mol, -118 kcal/mol, and -125 kcal/mol, converted to this standard state. This dispersion of experimental values for the absolute hydration free energy of the Li<sup>+</sup> (aq) ion is accurately mirrored in the dispersion of reference values adopted for the absolute hydration free energy of the H<sup>+</sup> (aq) ion. Inclusion of repulsive force contributions would reduce the present calculated value slightly. Furthermore, Li<sup>+</sup>(aq) is believed to have a strongly structured second hydration shell, which is treated only approximately in this calculation. Nevertheless, this level of agreement between calculation and experiment is satisfactory.
We additionally emphasize that the Li(H<sub>2</sub>O)$`{}_{n}{}^{}^+`$ complexes are treated in the harmonic approximation, although fully quantum mechanically. The low-$`n`$ clusters might have more entropy than is being accounted for by the harmonic approximation. If this were the case, then low-$`n`$ clusters would be more populous than currently represented. This would likely raise the theoretical value also.
To further test the $`n`$=4 prediction, ‘ab initio’ molecular dynamics calculations were carried out utilizing the VASP program. Two checks established the consistency for these problems between the electronic structure calculations described above and the energetics involved in the molecular dynamics calculations. First, the electron density functional alternative implemented in VASP was checked by comparing the electronic structure results obtained with the B3LYP hybrid electron density functional and the PW91 generalized gradient approximation exchange-correlation functional, using the Gaussian98 program and the same basis sets. As expected, satisfactory agreement was observed in the binding energies for sequential addition of a water molecule to the Li(H<sub>2</sub>O)$`{}_{n}{}^{}^+`$ clusters. Then the issues of pseudo-potentials and basis set were checked by optimizing cluster geometries with the VASP program and comparing to the results obtained for the same problems with Gaussian98. Again agreement was observed. For example, both procedures predicted the same lowest energy six-coordinated structure, the characteristic Li\[(H<sub>2</sub>O)<sub>4</sub>\]\[(H<sub>2</sub>O)<sub>2</sub>\]<sup>+</sup> cluster, with nearly identical geometries.
To initiate the ‘ab initio’ molecular dynamics calculation, the optimum $`n`$=6 inner sphere structure, rigidly constrained, was first equilibrated with 26 water molecules under conventional Monte Carlo liquid simulation conditions for liquid water, including periodic boundary conditions. This system of one Li<sup>+</sup> ion and 32 water molecules was then used as an initial configuration for the molecular dynamics calculation. As shown in Fig. 2, the initial $`n`$=6 structure relaxed to stable $`n`$=4 alternatives within 112 fs. The results of longer molecular dynamics calculations will be reported later.
The ‘ab initio’ molecular dynamics and the quasi-chemical theory of liquids exploit different approximations and produce the same conclusion here. This agreement supports the prediction that Li<sup>+</sup>(aq) has four inner shell water ligands at infinite dilution in liquid water under normal conditions. This prediction differs from interpretations of neutron and X-ray scattering data on aqueous solutions.
The conditions studied by these calculations and those targeted in the neutron scattering work do not match perfectly, particularly with regard to Li<sup>+</sup> concentration. Nevertheless, the theoretical methods are straightforward and physical, and, moreover, the distinct methods used here conform in their prediction of hydration number. Therefore, it will be of great importance for future work to fully resolve the differences between calculations and scattering experiments for these problems.
This work was supported by the US Department of Energy under contract W-7405-ENG-36 and the LDRD program at Los Alamos. |
warning/0001/math0001084.html | ar5iv | text | # The Kronecker Product of Schur Functions indexed by Two-Row Shapes or Hook Shapes.
## 1. Introduction
The aim of this paper is to derive an explicit formula for the Kronecker coefficients corresponding to partitions of certain shapes. The Kronecker coefficients, $`\gamma _{\mu \nu }^\lambda `$, arise when expressing a Kronecker product (also called inner or internal product), $`s_\mu s_\nu `$, of Schur functions in the Schur basis,
(1)
$$s_\mu s_\nu =\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda s_\lambda .$$
These coefficients can also be defined as the multiplicities of the irreducible representations in the tensor product of two irreducible representations of the symmetric group. A third way to define them is by the comultiplication expansion. Given two alphabets $`X=\{x_1,x_2,\mathrm{}\}`$ and $`Y=\{y_1,y_2,\mathrm{}\}`$
(2)
$$s_\lambda [XY]=\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda s_\mu [X]s_\nu [Y],$$
where $`s_\lambda [XY]`$ means $`s_\lambda (x_1y_1,x_1y_2,\mathrm{},x_iy_j,\mathrm{})`$. Remmel and Remmel and Whitehead have studied the Kronecker product of Schur functions corresponding to two two-row shapes, two hook shapes, and a hook shape and a two-row shape. We will use the comultiplication expansion (2) for the Kronecker coefficients, and a formula for expanding a Schur function of a difference of two alphabets due to Sergeev to obtain similar results in a simpler way. We believe that the formulas obtained using this approach are elegant and reflect the symmetry of the Kronecker product. In the three cases we found a way to express the Kronecker coefficients in terms of regions and paths in $`𝐍^2`$.
## 2. Basic definitions
A partition $`\lambda `$ of a positive integer $`n`$, written as $`\lambda n`$, is an unordered sequence of natural numbers adding to $`n`$. We write $`\lambda `$ as $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _n)`$, where $`\lambda _1\lambda _2\mathrm{}`$, and consider two such strings equal if they differ by a string of zeroes. The nonzero numbers $`\lambda _i`$ are called the parts of $`\lambda `$, and the number of parts is called the length of $`\lambda `$, denoted by $`l(\lambda )`$. In some cases, it is convenient to write $`\lambda =(1^{d_1}2^{d_2}\mathrm{}n^{d_n})`$ for the partition of $`n`$ that has $`d_i`$ copies of $`i`$. Using this notation, we define the integer $`z_\lambda `$ to be $`1^{d_1}d_1!\mathrm{\hspace{0.17em}2}^{d_2}d_2!\mathrm{}n^{d_n}d_n!`$.
We identify $`\lambda `$ with the set of points $`(i,j)`$ in $`𝐍^2`$ defined by $`1j\lambda _i`$, and refer to them as the Young diagram of $`\lambda `$. The Young diagram of a partition $`\lambda `$ is thought of as a collection of boxes arranged using matrix coordinates. For instance, the Young diagram corrresponding to $`\lambda =(4,3,1)`$ is
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To any partition $`\lambda `$ we associate the partition $`\lambda ^{}`$, its conjugate partition, defined by $`\lambda _i^{}=\left|\{j:\lambda _ji\}\right|.`$ Geometrically, $`\lambda ^{}`$ can be obtained from $`\lambda `$ by flipping the Young diagram of $`\lambda `$ around its main diagonal. For instance, the conjugate partition of $`\lambda `$ is $`\lambda ^{}=(3,2,2,1)`$, and the corresponding Young diagram is
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We recall some facts about the theory of representations of the symmetric group, and about symmetric functions. See or for proofs and details.
Let $`R(S_n)`$ be the space of class function in $`S_n`$, the symmetric group on $`n`$ letters, and let $`\mathrm{\Lambda }^n`$ be the space of homogeneous symmetric functions of degree $`n`$. A basis for $`R(S_n)`$ is given by the characters of the irreducible representations of $`S_n`$. Let $`\chi ^\mu `$ be the irreducible character of $`S_n`$ corresponding to the partition $`\mu `$. There is a scalar product $`,_{S_n}`$ on $`R(S_n)`$ defined by
$$\chi ^\mu ,\chi ^\nu _{S_n}=\frac{1}{n!}\underset{\sigma S_n}{}\chi ^\mu (\sigma )\chi ^\nu (\sigma ),$$
and extended by linearity.
A basis for the space of symmetric functions is given by the Schur functions. There exists a scalar product $`,_{\mathrm{\Lambda }^n}`$ on $`\mathrm{\Lambda }^n`$ defined by
$$s_\lambda ,s_\mu _{\mathrm{\Lambda }^n}=\delta _{\lambda \mu },$$
where $`\delta _{\lambda \mu }`$ is the Kronecker delta, and extended by linearity.
Let $`p_\mu `$ be the power sum symmetric function corresponding to $`\mu `$, where $`\mu `$ is a partition of $`n`$. There is an isometry $`ch^n:R(S_n)\mathrm{\Lambda }^n`$, given by the characteristic map,
$$\stackrel{n}{ch}(\chi )=\underset{\mu n}{}z_\mu ^1\chi (\mu )p_\mu .$$
This map has the remarkable property that if $`\chi ^\lambda `$ is the irreducible character of $`S_n`$ indexed by $`\lambda `$, then $`ch^n(\chi ^\lambda )=s_\lambda `$, the Schur function corresponding to $`\lambda `$. In particular, we obtain that $`s_\lambda =_{\mu n}z_\mu ^1\chi ^\lambda (\mu )p_\mu .`$ Hence,
(3)
$$\chi ^\lambda (\mu )=s_\lambda ,p_\mu .$$
Let $`\lambda `$, $`\mu `$, and $`\nu `$ be partitions of $`n`$. The Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ are defined by
(4) $`\gamma _{\mu \nu }^\lambda =\chi ^\lambda ,\chi ^\mu \chi ^\nu _{S_n}={\displaystyle \frac{1}{n!}}{\displaystyle \underset{\sigma S_n}{}}\chi ^\lambda (\sigma )\chi ^\mu (\sigma )\chi ^\nu (\sigma ).`$
Equation (4) shows that the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ are symmetric in $`\lambda `$, $`\mu `$, and $`\nu `$. The relevance of the Kronecker coefficients comes from the following fact: Let $`X^\mu `$ be the representation of the symmetric group corresponding to the character $`\chi ^\mu `$. Then $`\chi ^\mu \chi ^\nu `$ is the character of $`X^\mu X^\nu `$, the representation obtained by taking the tensor product of $`X^\mu `$ and $`X^\nu `$. Moreover, $`\gamma _{\mu \nu }^\lambda `$ is the multiplicity of $`X^\lambda `$ in $`X^\mu X^\nu `$.
Let $`f`$ and $`g`$ be homogeneous symmetric functions of degree $`n`$. The Kronecker product, $`fg`$, is defined by
(5)
$$fg=\stackrel{n}{ch}(uv),$$
where $`ch^nu=f`$, $`ch^nv=g`$, and $`uv(\sigma )=u(\sigma )v(\sigma )`$. To obtain (1) from this definition, we set $`f=s_\mu `$, $`g=s_\nu `$, $`u=\chi ^\mu `$, and $`v=\chi ^\nu `$ in (5).
The Kronecker product has the following symmetries:
$`s_\mu s_\nu `$ $`=s_\nu s_\mu .`$
$`s_\mu s_\nu `$ $`=s_\mu ^{}s_\nu ^{}.`$
Moreover, if $`\lambda `$ is a one-row shape
$$\gamma _{\mu \nu }^\lambda =\delta _{\mu ,\nu }.$$
We introduce the operation of substitution or plethysm into a symmetric function. Let $`f`$ be a symmetric function, and let $`X=\{x_1,x_2,\mathrm{}\}`$ be an alphabet. We write $`X=x_1+x_2+\mathrm{}`$, and define $`f[X]`$ by,
$$f[X]=f(x_1,x_2,\mathrm{}).$$
In general, if $`u`$ is any element of $`𝐐[[x_1,x_2,\mathrm{}]]`$, we write $`u`$ as $`_\alpha c_\alpha u_\alpha `$ where $`u_\alpha `$ is a monomial with coefficient $`1`$. Then $`p_\lambda [u]`$ is defined by setting
$`p_n[u]`$ $`={\displaystyle \underset{\alpha }{}}c_\alpha u_\alpha ^n`$
$`p_\lambda [u]`$ $`=p_{\lambda _1}[u]\mathrm{}p_{\lambda _n}[u]`$
for $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$. We define $`f[u]`$ for all symmetric functions $`f`$ by saying that $`f[u]`$ is linear in $`f`$.
The operation of substitution into a symmetric function has the following properties. For $`\alpha `$ and $`\beta `$ rational numbers, $`(\alpha f+\beta g)[u]=\alpha f[u]+\beta g[u].`$ Moreover, if $`c_\alpha =1`$ for all $`\alpha `$, then $`f[u]=f(\mathrm{},u_\alpha ,\mathrm{}).`$
Let $`X=x_1+x_2+\mathrm{}`$ and $`Y=y_1+y_2+\mathrm{}`$ be two alphabets. Define the sum of two alphabets by $`X+Y=x_1+x_2+\mathrm{}+y_1+y_2+\mathrm{},`$ and the product of two alphabets by $`XY=x_1y_1+\mathrm{}+x_iy_j+\mathrm{}.`$ Then
$`p_n[X+Y]`$ $`=p_n[X]+p_n[Y],`$
(6) $`p_n[XY]`$ $`=p_n[X]p_n[Y].`$
The inner product of function in the space of symmetric functions in two infinite alphabets is defined by
$$,_{_{XY}}=,__X,__Y,$$
where for any given alphabet $`Z`$, $`,__Z`$ denotes the inner product of the space of symmetric functions in $`Z`$.
For all partitions $`\rho `$, we have that $`p_\rho [XY]=p_\rho [X]p_\rho [Y].`$ If we rewrite (3) as $`p_\rho =_\lambda \chi ^\lambda (\rho )s_\lambda `$, then
(7)
$$\underset{\lambda }{}\chi ^\lambda s_\lambda [XY]=\underset{\mu ,\nu }{}\chi ^\mu \chi ^\nu s_\mu [X]s_\nu [Y].$$
Taking the coefficient of $`\chi ^\lambda `$ on both sides of the previous equation we obtain
$$s_\lambda [XY]=\chi ^\lambda ,\chi ^\mu \chi ^\nu s_\mu [X]s_\nu [Y].$$
Finally, using the definition of Kronecker coefficients (4) we obtain the comultiplication expansion (2).
###### Notation.
Let $`p`$ be a point in $`𝐍^2.`$ We say that $`(i,j)`$ can be reached from $`p`$, written $`p(i,j)`$, if $`(i,j)`$ can be reached from $`p`$ by moving any number of steps south-west or north-west. We define the weight function $`\omega `$ by
$$\omega _p(i,j)=\{\begin{array}{cc}x^iy^j,\hfill & \text{ if }p(i,j),\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
In particular, $`\sigma _{k,l}(h)=0`$ if $`h<0`$.
###### Notation.
We denote by $`x`$ the largest integer less than or equal to $`x`$ and by $`x`$ the smallest integer greater than or equal to $`x`$.
If $`f`$ is a formal power series, then $`[x^\alpha ]f`$ denotes the coefficient of $`x^\alpha `$ in $`f.`$
Following Donald Knuth we denote the characteristic function applied to a proposition $`P`$ by enclosing $`P`$ with brackets,
$$(P)=\{\begin{array}{cc}1,\hfill & \text{if proposition }P\text{ is true,}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
## 3. The case of two two-row shapes
The object of this section is to find a closed formula for the Kronecker coefficients when $`\mu =(\mu _1,\mu _2)`$ and $`\nu =(\nu _1,\nu _2)`$ are two-row shapes, and when we do not have any restriction on the partition $`\lambda `$. We describe the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ in terms of paths in $`𝐍^2`$. More precisely, we define two rectangular regions in $`𝐍^2`$ using the parts of $`\lambda `$. Then we count the number of points in $`𝐍^2`$ inside each of these rectangles that can be reached from $`(\nu _2,\mu _2+1)`$, if we are allowed to move any number of steps south-west or north-west. Finally, we subtract these two numbers.
We begin by introducing two lemmas that allow us to state Theorem 4 in a concise form.
###### Notation.
We use the coordinate axes as if we were working with matrices with first entry $`(0,0)`$. That is, the point $`(i,j)`$ belongs to the $`i`$th row and the $`j`$th column.
###### Lemma 1.
Let $`k`$ and $`l`$ be positive numbers. Let $`R`$ be the rectangle with width $`k`$, height $`l`$, and upper–left square $`(0,0)`$. Define
$$\sigma _{k,l}(h)=\left|\{(u,v)R𝐍^2:(0,h)(u,v)\}\right|$$
Then
$$\sigma _{k,l}(h)=\{\begin{array}{cc}0,\hfill & \text{if }h<0\hfill \\ (\frac{h}{2}+1)^2,\hfill & \text{if }0h<\mathrm{min}(k,l)\hfill \\ \sigma _{k,l}(s)+(\frac{hs}{2})\mathrm{min}(k,l),\hfill & \text{if }\mathrm{min}(k,l)h<\mathrm{max}(k,l)\hfill \\ \frac{kl}{2}\sigma _{k,l}(k+lh4),\hfill & \text{if }h\text{ is even and }\mathrm{max}(k,l)h\hfill \\ \frac{kl}{2}\sigma _{k,l}(k+lh4),\hfill & \text{if }h\text{ is odd and }\mathrm{max}(k,l)h\hfill \end{array}$$
where $`s`$ is defined as follows: If $`h\mathrm{min}(k,l)`$ is even, then $`s=\mathrm{min}(k,l)2`$; otherwise $`s=\mathrm{min}(k,l)1`$.
###### Proof.
If $`h`$ is to the left of the $`0`$th column, then we cannot reach any of the points in $`𝐍^2`$ inside $`R`$. Hence, $`\sigma _{k,l}(h)`$ should be equal to zero.
If $`0h\mathrm{min}(k,l)`$, then we are counting the number of points in $`𝐍^2`$ that can be reached from $`(0,h)`$ inside the square $`S`$ of side $`\mathrm{min}(k,l)`$. We have to consider two cases. If $`h`$ is odd, then we are summing $`2+4+\mathrm{}+(h+1)=(\frac{h}{2}+1)^2`$. On the other hand, if $`h`$ is even, then we are summing $`1+3+\mathrm{}+(h+1)=(\frac{h}{2}+1)^2`$.
If $`\mathrm{min}(k,l)h<\mathrm{max}(k,l)`$, then we subdivide our problem into two parts. First, we count the number of points in $`𝐍^2`$ that can be reached from $`(0,h)`$ inside the square $`S`$ by $`\sigma _{k,l}(s)`$. Then we count those points in $`𝐍^2`$ that are in $`R`$ but not in $`S`$. Since $`h<\mathrm{max}(k,l)`$ all diagonals have length $`\mathrm{min}(k,l)`$ and there are $`\frac{hs}{2}`$ of them. See Table 1.
If $`\mathrm{max}(k,l)h`$, then it is easier to count the total number of points in $`𝐍^2`$ that can be reached from $`(0,h)`$ inside $`R`$ by choosing another parameter $`\widehat{h}`$ big enough and with the same parity as $`h`$. Then we subtract those points in $`𝐍^2`$ in $`R`$ that are not reachable from $`(0,h)`$ because $`h`$ is too close. If $`h`$ is even this number is $`kl/2`$. If $`h`$ is odd this number is $`kl/2`$.
Then we subtract those points that we should not have counted. we express this number in terms of the function $`\sigma `$. The line $`y=x+h+2`$ intersects the line $`y=l1`$ at $`x=hl+3`$. This is the $`x`$ coordinate of the first point on the last row that is not reachable from $`(0,h)`$. Then to obtain the number of points that can be reached from this point by moving south-west or north-west, we subtract $`hl+3`$ to $`k1`$. We have obtained that are $`\sigma _{k,l}(k+lh4)`$ points that we should not have counted. ∎
###### Example 2.
By definition $`\sigma _{9,5}(4)`$ counts the points in $`𝐍^2`$ in Table 1 marked with $``$. Then $`\sigma _{9,5}(4)=9`$. Similarly, $`\sigma _{9,5}(8)`$ counts the points in $`𝐍^2`$ in Table 1 marked either with the symbol $``$ or with the symbol $``$. Then $`\sigma _{9,5}(8)=19.`$
$``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$
Table 1.
###### Lemma 3.
Let $`a,b,c,`$ and $`d`$ be in $`𝐍`$. Let $`R`$ be the rectangle with vertices $`(a,c)`$, $`(a+b,c)`$, $`(a,c+d)`$, and $`(a+b,c+d)`$. We define
$$\mathrm{\Gamma }(a,b,c,d)(x,y)=\left|\{(u,v)R:(x,y)(u,v)\}\right|.$$
Suppose that $`(x,y)`$ is such that $`xy`$. Then
$$\mathrm{\Gamma }(a,b,c,d)(x,y)=\{\begin{array}{cc}\sigma _{b+1,d+1}(x+yac),\hfill & 0yc\hfill \\ \sigma _{b+1,yc+1}(xa)+\sigma _{b+1,c+dy+1}(xa)\delta ,\hfill & c<y<c+d\hfill \\ \sigma _{b+1,d+1}(xy+c+da),\hfill & c+dy\hfill \end{array}$$
where $`\delta `$ is defined as follows If $`x<a`$, then $`\delta =0`$. If $`axa+b`$, then $`\delta =\frac{xa+1}{2}`$. Finally, if $`x>a+b`$ then we consider two cases: If $`xab`$ is even then $`\delta =\frac{b+1}{2}`$; otherwise, $`\delta =\frac{b+1}{2}`$.
###### Proof.
We consider three cases. If $`0yc`$ then the first position inside $`R`$ that we reach is $`(x+yac,c)`$. Therefore, we assume that we are starting at this point. Similarly, if $`yc+d`$, then the first position inside $`R`$ that we reach is $`(xy+c+da,c)`$. Again, we can assume that we are starting at this point.
On the other hand, if $`c<y<c+d`$, then we subdivide the problem in two parts. The number of position to the north of us is counted by $`\sigma _{b+1,yc+1}(xa)`$. The number of position to the south of us is counted by $`\sigma _{b+1,c+dy+1}(xa)`$. We define $`\delta `$ to be the number of points in $`𝐍^2`$ that we counted twice during this process. Then it is easy to see that $`\delta `$ is given by the previous definition. ∎
To compute the coefficient $`u_\nu `$ in the expansion $`f[X]=_\eta u_\eta s_\eta [X]`$ for $`f\mathrm{\Lambda }`$, it is enough to expand $`f[x_1+\mathrm{}+x_n]=_\eta u_\eta s_\eta [x_1+\mathrm{}+x_n]`$ for any $`nl(\nu )`$. (See \[7, section I.3\], for proofs and details.) Therefore, in this section we work with symmetric functions in a finite number of variables.
Let $`\mu `$ and $`\nu `$ be two-row partitions. Set $`X=1+x`$ and $`Y=1+y`$ in the comultiplication expansion (2) to obtain
(8)
$$s_\lambda [(1+y)(1+x)]=\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda s_\mu [1+y]s_\nu [1+x].$$
Note that the Kronecker coefficients are zero when $`l(\lambda )>4.`$
Jacobi’s definition of a Schur function on a finite alphabet $`s_\lambda [X]`$ as a quotient of alternants says that
(9) $`s_\lambda [X]=s_\lambda (x_1,\mathrm{},x_n)={\displaystyle \frac{det(x_i^{\lambda _j+nj})_{1i,jn}}{_{i<j}(x_ix_j)}}.`$
By the symmetry properties of the Kronecker product it is enough to compute the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ when $`\nu _2\mu _2`$.
###### Theorem 4.
Let $`\mu `$, $`\nu `$, and $`\lambda `$ be partitions of $`n`$, where $`\mu =(\mu _1,\mu _2)`$ and $`\nu =(\nu _1,\nu _2)`$ are two two-row partitions and let $`\lambda =(\lambda _1,\lambda _2,\lambda _4,\lambda _4)`$ be a partition of length less than or equal to $`4`$. Assume that $`\nu _2\mu _2`$. Then
$$\gamma _{\mu \nu }^\lambda =\left(\mathrm{\Gamma }(a,b,a+b+1,c)\mathrm{\Gamma }(a,b,a+b+c+d+2,c)\right)(\nu _2,\mu _2+1).$$
where $`a=\lambda _3+\lambda _4`$, $`b=\lambda _2\lambda _3`$, $`c=\mathrm{min}(\lambda _1\lambda _2,\lambda _3\lambda _4)`$ and $`d=\left|\lambda _1+\lambda _4\lambda _2\lambda _3\right|`$.
###### Proof.
We expand the polynomial $`s_\lambda [(1+y)(1+x)]=s_\lambda (1,y,x,xy)`$ in two different ways and obtain the Kronecker coefficients by equating both results. Let $`\phi `$ be the polynomial defined by $`\phi =(1x)(1y)s_\lambda (1,y,x,xy).`$ Using Jacobi’s definition of a Schur function we obtain
(10)
$$\phi =\frac{\left|\begin{array}{cccc}1& 1& 1& 1\\ y^{\lambda _1+3}& y^{\lambda _2+2}& y^{\lambda _3+1}& y^{\lambda _4}\\ x^{\lambda _1+3}& x^{\lambda _2+2}& x^{\lambda _3+1}& x^{\lambda _4}\\ (xy)^{\lambda _1+3}& (xy)^{\lambda _2+2}& (xy)^{\lambda _3+1}& (xy)^{\lambda _4}\end{array}\right|}{xy(1xy)(yx)(1x)(1y)}.$$
On the other hand, we may use Jacobi’s definition to expand $`s_\mu [1+y]`$ and $`s_\nu [1+x]`$ as quotients of alternants. Substitute this into (8):
$`s_\lambda [(1+y)(1+x)]`$ $`={\displaystyle \underset{\begin{array}{c}\mu =(\mu _1,\mu _2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda \left({\displaystyle \frac{y^{\mu _2}y^{\mu _1+1}}{1y}}\right)\left({\displaystyle \frac{x^{\nu _2}x^{\nu _1+1}}{1x}}\right)`$
(11) $`={\displaystyle \underset{\begin{array}{c}\mu =(\mu _1,\mu _2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda {\displaystyle \frac{x^{\nu _2}y^{\mu _2}x^{\nu _2}y^{\mu _1+1}x^{\nu _1+1}y^{\mu _2}+x^{\nu _1+1}y^{\mu _1+1}}{(1x)(1y)}}.`$
Since $`\nu _1+1`$ and $`\mu _1+1`$ are both greater than $`\frac{n}{2},`$ equation (11) implies that the coefficient of $`x^{\nu _2}y^{\mu _2}`$ in $`\phi `$ is $`\gamma _{\mu \nu }^\lambda `$. It is convenient to define an auxiliary polynomial by
(12)
$$\zeta =(1xy)(yx)\phi .$$
Let $`\xi `$ be the polynomial obtained by expanding the determinant appearing in (10). Equations (10) and (12) imply
$$\zeta =\frac{\xi }{xy(1x)(1y)}.$$
Let $`\xi _{i,j}`$ be the coefficient of $`x^iy^j`$ in $`\xi `$. ( Then $`\xi _{i,j}`$ is zero if $`i0`$ or $`j0`$, because $`\xi `$ is a polynomial divisible by $`xy`$.) Let $`\zeta _{i,j}`$ be the coefficient of $`x^iy^j`$ in $`\zeta `$. Then
(13) $`{\displaystyle \underset{i,j0}{}}\zeta _{i,j}x^iy^j={\displaystyle \frac{1}{xy(1x)(1y)}}{\displaystyle \underset{i,j0}{}}\xi _{i,j}x^iy^j={\displaystyle \underset{i,j,k,l0}{}}\xi _{ik,jl}x^{i1}y^{j1}.`$
Comparing the coefficient of $`x^iy^j`$ on both sides of equation (13) we obtain that
(14) $`\zeta _{i,j}={\displaystyle \underset{k,l0}{}}\xi _{i+1k,j+1l}={\displaystyle \underset{k=0}{\overset{i}{}}}{\displaystyle \underset{l=0}{\overset{j}{}}}\xi _{k+1,l+1}`$
We compute $`\zeta _{i,j}`$ from (14) by expanding the determinant appearing on (10). We consider two cases.
Case 1. Suppose that $`\lambda _1+\lambda _4>\lambda _2+\lambda _3`$. Then
$$\lambda _1+\lambda _2+4>\lambda _1+\lambda _3+3>\lambda _1+\lambda _4+2\lambda _2+\lambda _3+2>\lambda _2+\lambda _4+1>\lambda _3+\lambda _4.$$
We record the values of $`\xi _{j+1,i+1}`$ in Table 2. We use the convention that $`\xi _{i+1,j+1}`$ is zero whenever the $`(i,j)`$ entry is not in Table 2.
$`i\backslash j`$ $`\lambda _3+\lambda _4`$ $`\lambda _2+\lambda _4+1`$ $`\lambda _2+\lambda _3+2`$ $`\lambda _1+\lambda _4+2`$ $`\lambda _1+\lambda _3+3`$ $`\lambda _1+\lambda _2+4`$ $`\lambda _3+\lambda _4`$ $`0`$ $`1`$ $`+1`$ $`+1`$ $`1`$ $`0`$ $`\lambda _2+\lambda _4+1`$ $`+1`$ $`0`$ $`1`$ $`1`$ $`0`$ $`+1`$ $`\lambda _2+\lambda _3+2`$ $`1`$ $`+1`$ $`0`$ $`0`$ $`+1`$ $`1`$ $`\lambda _1+\lambda _4+2`$ $`1`$ $`+1`$ $`0`$ $`0`$ $`+1`$ $`1`$ $`\lambda _1+\lambda _3+3`$ $`+1`$ $`0`$ $`1`$ $`1`$ $`0`$ $`+1`$ $`\lambda _1+\lambda _2+4`$ $`0`$ $`1`$ $`+1`$ $`+1`$ $`1`$ $`0`$
Table 2
The values of $`\xi _{j+1,i+1}`$ when $`\lambda _1+\lambda _4\lambda _2+\lambda _3`$
Equation (14) shows that the value of $`\zeta _{i,j}`$ can be obtained by adding the entries northwest of the point $`(i,j)`$. In Table 3 we record the values of $`\zeta _{i,j}`$.
$`i\backslash j`$ $`I_1`$ $`I_2`$ $`I_3`$ $`I_4`$ $`I_5`$ $`I_6`$ $`I_7`$ $`I_1`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`I_2`$ $`0`$ $`0`$ $`1`$ $`0`$ $`+1`$ $`0`$ $`0`$ $`I_3`$ $`0`$ $`+1`$ $`0`$ $`0`$ $`0`$ $`1`$ $`0`$ $`I_4`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`I_5`$ $`0`$ $`1`$ $`0`$ $`0`$ $`0`$ $`+1`$ $`0`$ $`I_6`$ $`0`$ $`0`$ $`+1`$ $`0`$ $`1`$ $`0`$ $`0`$ $`I_7`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$ $`0`$
Table 3
The values of $`\zeta _{i,j}`$ when $`\lambda _1+\lambda _4\lambda _2+\lambda _3`$
where
$`I_1`$ $`=[0,\lambda _3+\lambda _4),`$ $`I_2`$ $`=[\lambda _3+\lambda _4,\lambda _2+\lambda _4],`$
$`I_3`$ $`=[\lambda _2+\lambda _4+1,\lambda _2+\lambda _3+1],`$ $`I_4`$ $`=[\lambda _2+\lambda _3+2,\lambda _1+\lambda _4+1],`$
$`I_5`$ $`=[\lambda _1+\lambda _4+2,\lambda _1+\lambda _3+2],`$ $`I_6`$ $`=[\lambda _1+\lambda _3+3,\lambda _1+\lambda _2+3],`$
$`I_7`$ $`=[\lambda _1+\lambda _2+4,\mathrm{}).`$
Case 2. Suppose that $`\lambda _1+\lambda _4\lambda _2+\lambda _3`$. Then
$$\lambda _1+\lambda _2+4>\lambda _1+\lambda _3+3>\lambda _2+\lambda _3+2>\lambda _1+\lambda _4+2>\lambda _2+\lambda _4+1>\lambda _3+\lambda _4.$$
Note that in Table 3, the rows and columns corresponding to $`\lambda _1+\lambda _4+2`$ and $`\lambda _2+\lambda _3+2`$ are the same. Therefore, the values of $`\xi _{i,j}`$ for $`\lambda _1+\lambda _4\lambda _2+\lambda _3`$ are recorded in Table 3, if we set
$`I_3=[\lambda _2+\lambda _4+1,\lambda _1+\lambda _4+1]`$
$`I_4=[\lambda _2+\lambda _4+2,\lambda _2+\lambda _3+1]`$
$`I_5=[\lambda _2+\lambda _3+2,\lambda _1+\lambda _3+2],`$
and define the other intervals as before.
In both cases, let $`\phi _{i,j}`$ be the coefficient of $`x^iy^j`$ in $`\phi `$. Using (12) we obtain that
$`\phi `$ $`={\displaystyle \frac{1}{(1xy)(yx)}}{\displaystyle \underset{i,j0}{}}\zeta _{i,j}x^iy^j`$
(15) $`={\displaystyle \frac{1}{yx}}{\displaystyle \underset{i,j,l0}{}}\zeta _{il,jl}x^iy^j`$
$`={\displaystyle \underset{i,j,k,l0}{}}\zeta _{ikl,j+kl+1}x^iy^j.`$
(Note: We can divide by $`yx`$ because $`\phi =0`$ when $`x=y`$.) Comparing the coefficients of $`x^iy^j`$ on both sides of equation (15), we obtain $`\phi _{i,j}=_{k,l0}\zeta _{ikl,j+kl+1}.`$ Therefore,
(16)
$$\phi _{\nu _2,\mu _2}=\underset{i,j=0}{\overset{\nu _2}{}}\zeta _{\nu _2ij,\mu _2+ij+1}.$$
We have shown that $`\gamma _{\mu \nu }^\lambda =\phi _{\nu _2,\mu _2}`$ can be obtained by adding the entries in Table 3 in all points in $`𝐍^2`$ that can be reached from $`(\nu _2,\mu _2+1)`$. See Table 4.
$``$ $``$ $``$ $``$ $``$ $``$
Table 4
The right-most point in Table 4 has coordinates $`(3,2)`$.
By hypothesis $`\nu _2\mu _2n/2`$. Then, if we start at $`(\nu _2,\mu _2+1)`$ and move as previously described, the only points in $`𝐍^2`$ that we can possibly reach and that are nonzero in Table 3 are those in $`I_2\times I_3`$ or $`I_2\times I_5.`$ Hence, we have that $`\phi _{\nu _2,\mu _2}`$ is the number of points in $`𝐍^2`$ inside $`I_2\times I_3`$ that can be reached from $`(\nu _2,\mu _2+1)`$ minus the ones that can be reached in $`I_2\times I_5.`$
Case 1 The inequality $`\lambda _1+\lambda _4>\lambda _2+\lambda _3`$ implies that $`\lambda _1+\lambda _4+1>\frac{n}{2}`$. Moreover, $`\mu _2\nu _2`$ implies that we are only considering the region of $`𝐍^2`$ given by $`0ij\frac{n}{2}`$. The number of points in $`𝐍^2`$ that can be reached from $`(\nu _2,\mu _2+1)`$ inside $`I_2\times I_3`$ is given by $`\mathrm{\Gamma }(\lambda _3+\lambda _4,\lambda _2\lambda _3,\lambda _2+\lambda _4+1,\lambda _3\lambda _4).`$ Similarly, the number of points in $`𝐍^2`$ that can be reached from $`(\nu _2,\mu _2+1)`$ inside $`I_2\times I_5`$ is given by $`\mathrm{\Gamma }(\lambda _3+\lambda _4,\lambda _2\lambda _3,\lambda _1+\lambda _4+2,\lambda _3\lambda _4).`$
Case 2 The inequality $`\lambda _2+\lambda _3\lambda _1+\lambda _4`$, implies that $`\lambda _1+\lambda _4+1>\frac{n}{2}`$. Moreover, $`\mu _2\nu _2`$ implies that we are only considering the region of $`𝐍^2`$ given by $`0ij\frac{n}{2}`$. The number of points in $`𝐍^2`$ that can be reached from $`(\nu _2,\mu _2+1)`$ inside $`I_2\times I_3`$ is given by $`\mathrm{\Gamma }(\lambda _3+\lambda _4,\lambda _2\lambda _3,\lambda _2+\lambda _4+1,\lambda _1\lambda _2).`$ Similarly, the number of points in $`𝐍^2`$ that can be reached from $`(\nu _2,\mu _2+1)`$ inside $`I_2\times I_5`$ is given by $`\mathrm{\Gamma }(\lambda _3+\lambda _4,\lambda _2\lambda _3,\lambda _2+\lambda _3+2,\lambda _1\lambda _2).`$
###### Corollary 5.
Let $`\mu =(\mu _1,\mu _2)`$, $`\nu =(\nu _1,\nu _2)`$, and $`\lambda =(\lambda _1,\lambda _2)`$ be partitions of $`n`$. Assume that $`\nu _2\mu _2\lambda _2`$. Then
$$\gamma _{\mu \nu }^\lambda =(yx)(yx),$$
where $`x=\mathrm{max}(0,\frac{\mu _2+\nu _2+\lambda _2n}{2})`$ and $`y=\frac{\mu _2+\nu _2\lambda _2+1}{2}`$.
###### Proof.
Set $`\lambda _3=\lambda _4=0`$ in Theorem 4. Then we notice that the second possibility in the definition of $`\mathrm{\Gamma }`$, that is, when $`c<y<c+d`$, never occurs. Note that $`\nu _2+\mu _2\lambda _2\nu _2+\mu _2\lambda _11`$ for all partitions $`\mu ,\nu ,`$ and $`\lambda `$. Therefore,
$$\gamma _{\mu \nu }^\lambda =\sigma _{\lambda _2+1,1}(\nu _2+\mu _2\lambda _2)\sigma _{\lambda _2+1,1}(\nu _2+\mu _2\lambda _11).$$
Suppose that $`\nu _2+\mu _2\lambda _2<0.`$ From the definition of $`\sigma _{\lambda _2+1,1}`$ we obtain that $`\gamma _{\mu \nu }^\lambda =0`$. Therefore, in order to have $`\gamma _{\mu \nu }^\lambda `$ not equal to zero, we should assume that $`\nu _2+\mu _2\lambda _20.`$
If $`0\nu _2+\mu _2\lambda _2<\lambda _2+1`$, then
$$\sigma _{\lambda _2+1,1}(\nu _2+\mu _2\lambda _2)=\frac{\nu _2+\mu _2\lambda _2+1}{2}$$
Similarly, if $`0\nu _2+\mu _2\lambda _2<\lambda _2+1`$, then
$$\sigma _{\lambda _2+1,1}(\nu _2+\mu _2\lambda _11)=\frac{\nu _2+\mu _2+\lambda _2n}{2}$$
It is easy to see that all other cases on the definition of $`\sigma _{k,l}`$ can not occur. Therefore, defining $`x`$ and $`y`$ as above, we obtain the desired result. ∎
###### Example 6.
If $`\mu =\nu =\lambda =(l,l)`$ or $`\mu =\nu =(2l,2l)`$ and $`\lambda =(3l,l)`$, then from the previous corollary, we obtain that
$`\gamma _{\mu \nu }^\lambda ={\displaystyle \frac{l+1}{2}}{\displaystyle \frac{l}{2}}=(l\text{ is even})`$
Note that to apply Corollary 5 to the second family of shapes, we should first use the symmetries of the Kronecker product to set $`\nu =\lambda =(2l,2l)`$ and $`\mu =(3l,l)`$.
###### Corollary 7.
The Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$, where $`\mu `$ and $`\nu `$ are two-row partitions, are unbounded.
###### Proof.
It is enough to construct an unbounded family of Kronecker coefficients. Assume that $`\mu =\nu =\lambda =(3l,l)`$. Then from the previous corollary we obtain that
$`\gamma _{\mu \nu }^\lambda ={\displaystyle \frac{l+1}{2}}`$
## 4. Sergeev’s formula
In this section we state Sergeev’s formula for a Schur function of a difference of two alphabets. See or \[7, section I.3\] for proofs and comments.
###### Definition 8.
Let $`X_m=x_1+\mathrm{}+x_m`$ be a finite alphabet, and let
$`\delta _m=(m1,m2,\mathrm{},1,0)`$. We define $`X_m^{\delta _m}`$ by $`X_m^{\delta _m}=x_1^{m1}\mathrm{}x_{m1}.`$
###### Definition 9.
Let $`i(\alpha )`$ denote the number of inversions of the permutation $`\alpha `$. We define the alternant to be
$$A_m^xP=\underset{\alpha S_m}{}(1)^{i(\alpha )}P(x_{\alpha (1)},\mathrm{},x_{\alpha (m)}),$$
for any polynomial $`P(x_1,\mathrm{},x_n)`$.
###### Definition 10.
Let $`\mathrm{\Delta }`$ be the operation of taking the Vandermonde determinant of an alphabet, i.e.,
$$\mathrm{\Delta }(X_m)=det(x_i^{mj})_{i,j=1}^m.$$
###### Theorem 11 (Sergeev’s Formula).
Let $`X_m=x_1+\mathrm{}+x_m,andY_n=y_1+\mathrm{}+y_n`$ be two alphabets. Then
$$s_\lambda [X_mY_n]=\frac{1}{\mathrm{\Delta }(X_m)\mathrm{\Delta }(Y_n)}A_m^xA_n^yX_m^{\delta _m}Y_n^{\delta _n}\underset{(i,j)\lambda }{}(x_iy_j)$$
The notation $`(i,j)\lambda `$ means that the point $`(i,j)`$ belongs to the diagram of $`\lambda `$. We set $`x_i=0`$ for $`i>m`$ and $`y_j=0`$ for $`j>n`$.
We use Sergeev’s formula as a tool for making some calculations we need for the next two sections.
1. Let $`\mu =(1^{e_1}m_2)`$ be a hook. (We are assuming that $`e_11`$ and $`m_22`$.) Let $`X^1=\{x_1\}`$ and $`X^2=\{x_2\}`$.
(17)
$$s_\mu [x_1x_2]=(1)^{e_1}x_1^{m_21}x_2^{e_1}(x_1x_2).$$
2. Let $`\nu =(\nu _1,\nu _2)`$ be a two-row partition. Let $`Y=\{y_1,y_2\}`$. Then
(18)
$$s_\nu [y_1+y_2]=\frac{(y_1y_2)^{\nu _2}(y_1^{\nu _1\nu _2+1}y_2^{\nu _1\nu _2+1})}{y_1y_2}.$$
3. We say that a partition $`\lambda `$ is a double hook if $`(2,2)\lambda `$ and it has the form $`\lambda =(1^{d_1}2^{d_2}n_3n_4).`$ In particular any two-row shape is considered to be a double hook.
Let $`\lambda `$ be a double hook. Let $`U=\{u_1,u_2\}`$ and $`V=\{v_1,v_2\}`$. Then if $`n_40`$ then $`s_\lambda [u_1+u_2v_1v_2]`$ equals
(19)
$$\begin{array}{c}\frac{(u_1v_1)(u_2v_1)(u_1v_2)(u_2v_2)}{(u_1u_2)(v_1v_2)}(1)^{d_1}(u_1u_2)^{n_32}(v_1v_2)^{d_2}\hfill \\ \hfill \times \left(u_2^{n_4n_3+1}u_1^{n_4n_3+1}\right)\left(v_2^{d_1+1}v_1^{d_1+1}\right).\end{array}$$
On the other hand, if $`n_4=0`$ then to compute $`s_\lambda [u_1+u_2v_1v_2]`$ we should write $`\lambda `$ as $`(1^{d_1}2^{d_21}2n_3).`$ That is, we set $`d_1:=d_1`$, $`d_2:=d_21`$, $`n_3:=2`$, and $`n_4=n_3`$ in (19).
4. Let $`\lambda `$ be a hook shape, $`\lambda =(1^{d_1}n_2).`$ (We are assuming that $`d_11`$ and $`n_22`$.) Let $`U=\{u_1,u_2\}`$ and $`V=\{v_1,v_2\}`$. Then $`s_\lambda [u_1+u_2v_1v_2]`$ equals
(20)
$$\begin{array}{c}(1)^{d_11}\frac{1}{(u_1u_2)}\frac{1}{(v_1v_2)}\times \hfill \\ \hfill \{u_1v_1(u_1v_1)(u_1v_2)(u_2v_1)u_1^{n_22}v_1^{d_11}\\ \hfill u_1v_2(u_1v_2)(u_1v_1)(u_2v_2)u_1^{n_22}v_2^{d_11}\\ \hfill u_2v_1(u_2v_1)(u_2v_2)(u_1v_1)u_2^{n_22}v_1^{d_11}\\ \hfill +u_2v_2(u_2v_2)(u_2v_1)(u_1v_2)u_2^{n_22}v_2^{d_11}.\}\end{array}$$
## 5. The case of two hook shapes
In this section we derive an explicit formula for the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ in the case in which $`\mu =(1^eu)`$, and $`\nu =(1^fv)`$ are both hook shapes. Given a partition $`\lambda `$ the Kronecker coefficient $`\gamma _{\mu \nu }^\lambda `$ tells us whether point $`(u,v)`$ belongs to some regions in $`𝐍^2`$ determined by $`\mu `$, $`\nu `$ and $`\lambda `$.
Recall that we denote the characteristic function by enclosing a proposition $`P`$ with brackets, $`(P)`$.
###### Lemma 12.
Let $`(u,v)𝐍^2`$ and let $`R`$ be the rectangle with vertices $`(a,b)`$, $`(b,a)`$, $`(c,d)`$, and $`(d,c)`$, with $`ab`$, $`cd`$, $`ca`$ and $`db`$. (Sometimes, when $`c=d=e`$, we denote this rectangle as $`(a,b;e)`$.)
Then $`(u,v)R`$ if and only if $`|vu|ab`$ and $`a+bu+vc+d`$
###### Proof.
Let $`L_1`$ be the line of slope $`1`$ passing through $`(u,v)`$, and let $`L_2`$ be the line of slope $`1`$ passing through $`(u,v)`$. Then we have that
$`L_1:y=x+vu`$
$`L_2:y=x+u+v`$
The point $`(u,v)`$ is in $`R`$ if and only if $`L_1`$ is between the lines of slope $`1`$ passing throught $`(a,b)`$ and $`(b,a)`$. That is, $`abvuba`$ and $`L_2`$ is between the lines of slope $`1`$ passing throught $`(a,b)`$ and $`(c,d)`$. That is, $`a+bu+vc+d`$
###### Theorem 13.
Let $`\lambda `$, $`\mu `$ and $`\nu `$ be partitions of $`n`$, where $`\mu =(1^eu)`$ and $`\nu =(1^fv)`$ are hook shapes. Then the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ are given by the following:
1. If $`\lambda `$ is a one-row shape, then $`\gamma _{\mu \nu }^\lambda =\delta _{\mu ,\nu }.`$
2. If $`\lambda `$ is not contained in a double hook shape, then $`\gamma _{\mu \nu }^\lambda =0.`$
3. Let $`\lambda =(1^{d_1}2^{d_2}n_3n_4)`$ be a double hook. Let $`x=2d_2+d_1`$. Then
$$\begin{array}{c}\gamma _{\mu \nu }^\lambda =(n_31\frac{e+fx}{2}n_4)(|fe|d_1)\hfill \\ \hfill +(n_3\frac{e+fx+1}{2}n_4)(|fe|d_1+1).\end{array}$$
Note that if $`n_4=0`$, then we shall rewrite $`\lambda =(1^{d_1}2^{d_21}2n_3)`$ before using the previous formula.
4. Let $`\lambda =(1^dw)`$ be a hook shape. Suppose that $`eu`$, $`fv`$, and $`dw`$. Then
$$\gamma _{\mu \nu }^\lambda =(ed+f)(de+f)(fe+d).$$
###### Proof.
Set $`X=\{1,x\}`$ and $`Y=\{1,y\}`$ in the comultiplication expansion (2) to obtain
(21)
$$s_\lambda [(1x)(1y)]=\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda s_\mu [1x]s_\nu [1y],$$
We use equation (17) to replace $`s_\mu `$ and $`s_\nu `$ in the right hand side of (21). Then we divide the resulting equation by $`(1x)(1y)`$ to get
$$\frac{s_\lambda [1yx+xy]}{(1x)(1y)}=\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda (x)^e(y)^f.$$
Therefore,
$$\gamma _{\mu \nu }^\lambda =\left[(x)^e(y)^f\right]\frac{s_\lambda [1yx+xy]}{(1x)(1y)},$$
when $`\mu `$ and $`\nu `$ are hook shapes.
Case 1. If $`\lambda `$ is not contained in any double hook, then the point $`(3,3)`$ is in $`\lambda `$, and by Sergeev’s formula, $`s_\lambda [1yx+xy]`$ equals zero.
Case 2. Let $`\lambda =(1^{d_1}2^{d_2}n_3n_4)`$ be a double hook. Set $`u_1=1`$, $`u_2=xy`$, $`v_1=x`$, and $`v_2=y`$ in (19). Then we divide by $`(1x)(1y)`$ on both sides of the resulting equation to obtain
(22)
$$\begin{array}{c}\frac{s_\lambda [1yx+xy]}{(1x)(1y)}=(1)^{d_1}(xy)^{n_3+d_21}\hfill \\ \hfill \times (1x)(1y)\left(\frac{1(xy)^{n_4n_3+1}}{1xy}\right)\left(\frac{x^{d_1+1}y^{d_1+1}}{xy}\right).\end{array}$$
Note: If $`n_4=0`$ then we should write $`\lambda =(1^{d_1}2^{d_21}2n_4)`$ in order to use (19).
Let $`p`$ be a point in $`𝐍^2.`$ We say that $`(i,j)`$ can be reached from $`p`$, written $`p(i,j)`$, if $`(i,j)`$ can be reached from $`p`$ by moving any number of steps south-west or north-west. We have defined a weight function by
$$\omega _p(i,j)=\{\begin{array}{cc}x^iy^j,\hfill & \text{ if }p(i,j);\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
Let $`\omega _p(T)=_{(i,j)T}\omega _p(i,j)`$ be the generating function of a region $`T`$ in $`𝐍^2`$. Let $`R`$ be the rectangle with vertices $`(0,d_1)`$, $`(d_1,0)`$, $`(d_1+n_4n_3,n_4n_3)`$ and $`(n_4n_3,d_1+n_4n_3)`$. Then
$`\omega _{(d_1+n_4n_3,n_4n_3)}(R)`$ $`=\left({\displaystyle \frac{1(xy)^{n_4n_3+1}}{1xy}}\right)\left({\displaystyle \frac{x^{d_1+1}y^{d_1+1}}{xy}}\right)`$
$`={\displaystyle \underset{k=0}{\overset{n_4n_3}{}}}{\displaystyle \underset{i+j=d_1}{}}(xy)^kx^iy^j.`$
See Table 5.
$`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$
Table 5.
$`d_1=4`$ and $`n_3n_4=4`$
Recall that we are using matrix coordinates, and that the upper-left corner has coordinates $`(0,0)`$. The coordinates of the four vertices of $`R`$ in Table $`5`$ are $`(0,4)`$, $`(4,0)`$, $`(8,4)`$, and $`(4,8)`$.
We interpret the right-hand side of (22) as the sum of four different generating functions. To be more precise, the right-hand side of (22) can be written as $`_{i=1}^4\omega _{p_i}(r_i)`$ where $`p_1=(n_4+d_21,n_4+d_2+d_11)`$ and $`R_1=\{n_3+d_2+d_11,n_3+d_21;n_4n_3\}`$, $`p_2=(n_4+d_2,n_34+d_2+d_11)`$ and $`R_2=\{n_3+d_2+d_1,n_3+d_21;n_4n_3\}`$, $`p_3=(n_4+d_21,n_4+d_2+d_1)`$ and $`R_3=\{n_3+d_2+d_11,n_3+d_2;n_4n_3\}`$, and $`p_4=(n_4+d_2+d_1,n_4+d_2+d_1)`$ and $`R_4=\{n_3+d_2+d_1,n_3+d_2+d_1;n_4n_3\}`$.
We observe that $`R_1R_2`$ (and $`R_3R_4`$ ) are rectangles in $`𝐍^2`$. Moreover,
(23)
$$\gamma _{\mu \nu }^\lambda =((e,f)R_1R_2)+((e,f)R_3R_4).$$
The vertices of rectangle $`R_1R_4`$ are given (using the notation of 12) by
$`a`$ $`=n_3+d_2+d_11`$ $`b`$ $`=n_3+d_21`$
$`c`$ $`=n_4+d_2+d_1`$ $`d`$ $`=n_4+d_2`$
Similarly, the vertices of rectangle $`R_2R_3`$ are given by
$`a=n_3+d_2+d_1`$ $`b`$ $`=n_3+d_21`$
$`c=n_4+d_2+d_1`$ $`d`$ $`=n_4+d_21`$
Applying Lemma 12 to (23) we obtain
$$\begin{array}{c}\gamma _{\mu \nu }^\lambda =(n_31\frac{e+fx}{2}n_4)(|fe|d_1)\hfill \\ \hfill +(n_3\frac{e+fx+1}{2}n_4)(|fe|d_1+1).\end{array}$$
Case 3. $`\lambda `$ is a hook. Suppose that $`\lambda `$ is a hook, $`\lambda =(1^dw)`$. Set $`u_1=1`$, $`u_2=xy`$, $`v_1=x`$, and $`v_2=y`$ in (20). Then we divide by $`(1x)(1y)`$ on both sides of the resulting equation to obtain
(24)
$$\begin{array}{c}\frac{s_\lambda [1yx+xy]}{(1x)(1y)}=(1)^d\left(\frac{x^{d+1}y^{d+1}}{xy}\right)\left(\frac{1(xy)^w}{1xy}\right)\hfill \\ \hfill +(1)^{d1}xy\left(\frac{x^dy^d}{xy}\right)\left(\frac{1(xy)^{w1}}{1xy}\right).\end{array}$$
We want to interpret this equation as a generating function for a region $`T`$ using the weight $`\omega `$. We proceed as follows:
Let $`R_1`$ be the rectangle with vertices $`(d,0),`$ $`(0,d),`$ $`(d+w1,w1)`$, and $`(w1,d+w1)`$. Then
(25)
$$\omega _{(w1,d+w1)}(R_1)=\left(\frac{1(xy)^w}{1xy}\right)\left(\frac{x^{d+1}y^{d+1}}{xy}\right)=\underset{k=0}{\overset{w1}{}}\underset{i+j=d}{}(xy)^kx^iy^j.$$
(See Table 5.) Similarly, let $`R_2`$ be the rectangle with vertices $`(d,1),`$ $`(1,d),`$ $`(d+w2,w1)`$, and $`(w1,d+w2)`$. Then
(26)
$$\omega _{(w1,d+w2)}(R_2)=xy\left(\frac{1(xy)^{w1}}{1xy}\right)\left(\frac{x^dy^d}{xy}\right)=xy\underset{k=0}{\overset{w2}{}}\underset{i+j=d1}{}(xy)^kx^iy^j.$$
Observe that the points in $`𝐍^2`$ that can be reached from $`(0,d)`$ in $`R_1`$ and the points in $`𝐍^2`$ that can be reached from $`(1,d)`$ in $`R_2`$ are disjoint. Moreover, they completely fill the rectangle $`R_1R_2`$. See Table 6.
$`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$ $`1`$
Table 6
$`d=4,w=6.`$
Note that $`R_2`$ is contained in $`R_1`$. We obtain that
$$\omega _{(w1,d+w1)}(R_1)+\omega _{(w1,d+w2)}(R_2)=|(e,f)R_1|$$
We use apply Lemma 12 to the previous equation to obtain:
$$(|ef|d)(de+fd+2w2).$$
But, by hypothesis, $`eu`$, $`fv`$, and $`dw`$. Therefore, this system is equivalent to $`(de+f)(fe+d)(ed+f)`$, as desired. ∎
###### Corollary 14.
Let $`\lambda `$, $`\mu `$, and $`\nu `$ be partitions of $`n`$, where $`\mu =(1^eu)`$ and $`\nu =(1^fv)`$ are hook shapes and $`\lambda =(\lambda _1,\lambda _2)`$ is a two-row shape. Then the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ are given by
$$\gamma _{\mu \nu }^\lambda =(\lambda _21e\lambda _1)(e=f)+(\lambda _2\frac{e+f+1}{2}\lambda _1)(|ef|1).$$
###### Proof.
In Theorem 13, set $`d_1=d_2=0`$, $`n_3=\lambda _2`$ and $`n_4=\lambda _1.`$
###### Corollary 15.
Let $`\lambda `$, $`\mu `$ and $`\nu `$ be partitions of $`n`$, where $`\mu `$ and $`\nu `$ are hook shapes. Then the Kronecker coefficients are bounded. Moreover, the only possible values for the Kronecker coefficients are $`0`$, $`1`$ or $`2`$.
## 6. The case of a hook shape and a two-row shape
In this section we derive an explicit formula for the Kronecker coefficients in the case $`\mu =(1^{e_1}m_2)`$ is a hook and $`\nu =(\nu _1,\nu _2)`$ is a two-row shape. Given a partition $`\lambda `$, the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ tell us whether the point $`(e_1,\nu _2)`$ belongs to some regions in $`𝐍^2`$ determined by $`\mu `$, $`\nu `$ and $`\lambda `$.
Using the symmetry properties of the Kronecker product, we may assume that if $`\lambda =(1^{d_1}2^{d_2}n_3n_4)`$ then $`n_4n_3d_1`$. (If $`n_4=0`$ then we should rewrite $`\lambda `$ as $`(1^{d_1}2^{d_21}2n_3)`$. Moreover, our hypothesis becomes $`n_32d_1`$.)
Recall that we denote the value of the characteristic function at proposition $`P`$ by $`(P)`$.
###### Theorem 16.
Let $`\lambda `$, $`\mu `$ and $`\nu `$ be partitions of $`n`$, where $`\mu =(1^{e_1}m_2)`$ is a hook and $`\nu =(\nu _1,\nu _2)`$ is a two-row shape. Then the Kronecker coefficients $`\gamma _{\mu \nu }^\lambda `$ are given by the following:
1. If $`\lambda `$ is a one-row shape, then $`\gamma _{\mu \nu }^\lambda =\delta _{\mu ,\nu }.`$
2. If $`\lambda `$ is not contained in any double hook, then $`\gamma _{\mu \nu }^\lambda =0.`$
3. Suppose $`\lambda =(1^{d_1}2^{d_2}n_3n_4)`$ is a double hook. Assume that $`n_4n_3d_1`$.
(If $`n_4=0`$, then we should write $`\lambda =(1^{d_1}2^{d_21}2n_3)`$.) Then
(27)
$$\begin{array}{c}\gamma _{\mu \nu }^\lambda =(n_3\nu _2d_21n_4)(d_1+2d_2<e_1<d_1+2d_2+3)\hfill \\ \hfill +(n_3\nu _2d_2n_4)(d_1+2d_2e_1d_1+2d_2+3)\\ \hfill +(n_3\nu _2d_2+1n_4)(d_1+2d_2<e_1<d_1+2d_2+3)\\ \hfill (n_3+d_2+d_1=\nu _2)(d_1+2d_2+1e_1d_1+2d_2+2).\end{array}$$
4. If $`\lambda `$ is a hook, see Corollary 14.
###### Proof.
Set $`X=1+x`$ and $`Y=1+y`$ in the comultiplication expansion (2) to obtain
(28)
$$s_\lambda [(1x)(1+y)]=\underset{\mu ,\nu }{}\gamma _{\mu \nu }^\lambda s_\mu [1x]s_\nu [1+y].$$
Use (17) and (18) to replace $`s_\mu `$ and $`s_\nu `$ in the right-hand side of (28), and divide by $`(1x)`$ to obtain
(29)
$$\frac{s_\lambda [(1x)(1+y)]}{1x}=\underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}\gamma _{\mu \nu }^\lambda (x)^{e_1}y^{\nu _2}(\frac{1y^{\nu _1\nu _2+1}}{1y}).$$
If $`\lambda `$ is not contained in any double hook, then the point $`(3,3)`$ is in $`\lambda `$, and by Sergeev’s formula, $`s_\lambda [(1x)(1+y)]`$ equals zero.
Since we already computed the Kronecker coefficients when $`\lambda `$ is contained in a hook, we can assume for the rest of this proof that $`\lambda `$ is a double hook. Let $`\lambda =(1^{d_1}2^{d_2}n_3n_4)`$. (Note: If $`n_4=0`$ then we should write $`\lambda =(1^{d_1}2^{d_21}2n_4).`$)
Set $`u_1=1`$, $`u_2=y`$, $`v_1=x`$, and $`v_2=xy`$ in (19), and multiply by $`\frac{1y}{1x}`$ on both sides of the resulting equation.
(30)
$$\begin{array}{c}\underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}\gamma _{\mu \nu }^\lambda (x)^{e_1}y^{\nu _2}\left(1y^{\nu _1\nu _2+1}\right)=(yx)(1xy)(1x)\hfill \\ \hfill \times (x)^{d_1+2d_2}y^{n_3+d_21}\left(\frac{(1y^{n_4n_3+1})(1y^{d_1+1})}{1y}\right).\end{array}$$
We have that $`(yx)(1xy)(1x)=yx(1+y+y^2)+x^2(1+y+y^2)x^3y.`$ Therefore, looking at the coefficient of $`x`$ on both sides of the equation, we see that $`\gamma _{\mu \nu }^\lambda `$ is zero if $`e_1`$ is different from $`d_1+2d_2,d_1+2d_2+1,d_1+2d_2+2,`$ or $`d_1+2d_2+3.`$
Let $`e_1=d_1+2d_2`$ or $`e_1=d_1+2d_2+3`$. Since $`\nu _2n/2`$, we have that
$`\gamma _{\mu \nu }^\lambda `$ $`=[y^{\nu _2}]{\displaystyle \underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda y^{\nu _2}`$
$`=[y^{\nu _2}]{\displaystyle \underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda y^{\nu _2}(1y^{\nu _1\nu _2+1})`$ ($`\nu _1+1>n/2`$)
$`=[y^{\nu _2}]y^{n_3+d_2}(1y^{d_1+1})\left({\displaystyle \frac{1y^{n_4n_3+1}}{1y}}\right)`$ (Eq. 30)
$`=[y^{\nu _2}]y^{n_3+d_2}(1y^{d_1+1}){\displaystyle \underset{k=0}{\overset{n_4n_3}{}}}y^k`$
$`=[y^{\nu _2}]y^{n_3+d_2}{\displaystyle \underset{k=0}{\overset{n_4n_3}{}}}y^k.`$ ($`n_3+d_2+d_1n/2`$)
We have obtained that for $`e_1=d_1+2d_2`$ or $`e_1=d_1+2d_2+3`$
$$\gamma _{\mu \nu }^\lambda =(n_3\nu _2d_2n_4).$$
Let $`e_1=d_1+2d_2+1`$ or $`e_1=d_1+2d_2+2`$. Since $`\nu _2\frac{n}{2}`$ we have that
$`\gamma _{\mu \nu }^\lambda `$ $`=[y^{\nu _2}]{\displaystyle \underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda y^{\nu _2}`$
$`=[y^{\nu _2}]{\displaystyle \underset{\begin{array}{c}\mu =(1^{e_1}m_2)\\ \nu =(\nu _1,\nu _2)\end{array}}{}}\gamma _{\mu \nu }^\lambda (1y^{\nu _1\nu _2+1})`$
$`=[y^{\nu _2}]y^{n_3+d_21}(1+y+y^2)(1y^{d_1+1})\left({\displaystyle \frac{1y^{n_4n_3+1}}{1y}}\right)`$
$`=([y^{\nu _2}]y^{n_3+d_21}(1+y+y^2)\left({\displaystyle \frac{1y^{n_4n_3+1}}{1y}}\right))(n_3+d_2+d_1=\nu _2)`$
$`=\left([y^{\nu _2}]y^{n_3+d_21}(1+y+y^2){\displaystyle \underset{k=0}{\overset{n_4n_3}{}}}y^k\right)(n_3+d_2+d_1=\nu _2)`$
We have obtained that for $`e_1=d_1+2d_2+1`$ or $`e_1=d_1+2d_2+2`$
(31)
$$\begin{array}{c}\gamma _{\mu \nu }^\lambda =(n_3\nu _2d_21n_4)+(n_3\nu _2d_2n_4)\hfill \\ \hfill +(n_3\nu _2d_2+1n_4)(n_3+d_2+d_1=\nu _2).\end{array}$$
###### Corollary 17.
The Kronecker cofficients, $`\gamma _{\mu \nu }^\lambda `$, where $`\mu `$ is a hook and $`\nu `$ is a two-row shape are always $`0`$, $`1`$, $`2`$ or $`3`$.
## 7. Final comments
The inner product of symmetric functions was discovered by J. H. Redfield in 1927, together with the scalar product of symmetric functions. He called them cup and cap products, respectively. D.E. Littlewood reinvented the inner product in 1956.
More recently, I.M. Gessel and A. Lascoux obtained combinatorial interpretations for the Kronecker coefficients in some restricted cases; Lascoux in the case where $`\mu `$ and $`\nu `$ are hooks, and $`\lambda `$ a straight tableaux, and Gessel in the case that $`\mu `$ and $`\nu `$ are zigzag shapes and $`\lambda `$ is an arbitrary skew shape. A. Lascoux interpreted the Kronecker coefficients, when two of the shapes are hooks as counting clases of words under some equivalence relation. We refer to or for a complete statement of his results. The Corollary of Theorem 3 in this paper shows that each class of words, under Lascoux’s equivalence, contains either $`0`$, $`1`$, or $`2`$ different representatives. I. Gessel worked on a more general framework, contemplating the occurrence of skew tableaux. It was shown in that in the case where two of the partitions are hook shapes, and the third one is an arbitrary straight shape, his result is equivalent to Lascoux’s.
In , A.M. Garsia and J.B. Remmel founded a way to relate shuffles of permutations and Kronecker coefficients. From here they obtained a combinatorial interpretation for the Kronecker coefficients when $`\lambda `$ is a product of homogeneous symmetric functions, and $`\mu `$ and $`\nu `$ are arbitrary skew shapes. They also showed how Gessel’s and Lascoux’s results are related.
J.B. Remmel , , and J.B. Remmel and T. Whitehead , obtained formulas for computing the Kronecker coefficients in the same cases considered in this paper. Their approach was mainly combinatorial:
First, they expanded the Kronecker product $`s_\mu s_\nu `$ in terms of Schur functions using the Garsia-Remmel algorithm . The problem of computing the Kronecker coefficients was reduced to computing signed sums of certain products of skew Schur functions.
Then they obtained a description of the coefficients that arise in the expansion of the resulting product of skew Schur functions in terms of counting 3-colored diagrams in , and or 4-colored diagrams in . At this point, they reduced the problem to computing a signed sum of colored diagram.
Finally, they defined involutions on these signed sums to cancel negative terms, and obtained the desired formulas by counting classes of restricted colored diagram.
In general, it is not obvious how to go from the determination of the Kronecker coefficients $`\gamma _{\mu ,\nu }^\lambda `$ when $`\mu `$ and $`\nu `$ are two-row shapes found in this paper, and the one obtained by J.B. Remmel and T. Whitehead . But, in some particular cases this is easy to see. For instance, when $`\lambda `$ is also a two-row shape, both formulas are exactly the same. |
warning/0001/hep-ex0001034.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The neutrino remains as exotic and challenging today as it was seventy years ago when first proposed by Pauli. What is known for certain about neutrinos is minimal indeed. They have spin $`\frac{1}{2}`$, charge 0, helicity -1, and exist in 3 flavors, electron, mu, and tau. Strictly speaking, only 2 flavors are certain: direct observation of the tau neutrino has not yet been achieved, but an experiment, DONUT, at Fermilab is in progress with this objective. Limits on the masses from direct, kinematic experiments (that do not require assumptions about the non-conservation of lepton family number), have been steadily lowered by experiments of ever-increasing sophistication over the years, with the results given in Table 1. As will be discussed below, lower limits on $`\nu _e`$ from tritium beta decay exist , but the data show curious distortions near the endpoint that are not at present understood.
There are strong theoretical and experimental motivations to search for neutrino mass. Presumably created in the early universe in numbers comparable to photons, neutrinos with a mass of only a few eV would contribute a significant fraction of the closure density. A species with a mass of $`94h_0^2`$ eV, where $`h_0`$ is the Hubble constant in units of 100 km s<sup>-1</sup> Mpc<sup>-1</sup>, provides by itself the closure density, but the structure of the universe at small and intermediate scales is incompatible with the dominance of hot dark matter such as neutrinos. On the other hand, evidence has accumulated from surveys of galaxy distributions and of the cosmic microwave background that is consistent with the presence of some hot dark matter.
In gauge theories of the elementary particles, the fermion masses arise from the coupling of left- and right-handed fields. A Dirac mass for neutrinos is expected to be like the quark and charged-lepton masses, but experiment has shown that neutrino masses are tiny. For this reason the minimal Standard Model deprives neutrinos of right-handed fields, forcing them to be always relativistic and massless. A clear demonstration that neutrinos have mass forces a confrontation of our understanding of how mass is generated.
Sensitivity to very small neutrino masses is experimentally accessible if neutrino mass eigenstates are not flavor eigenstates. In that case, a state prepared by the weak interaction (W<sup>+/-</sup> decay) in a specific flavor projection consists of two or more physical mass components propagating slightly differently with time or distance. A remote detector with specific flavor sensitivity will then register ‘a neutrino’ with altered flavor projection as a result of the phase difference that accumulates. Because the neutrinos are in general highly relativistic, the phase difference increases not as the difference in the masses, but as the difference $`\mathrm{\Delta }m^2`$ between the squares of the masses.
If neutrino oscillations occur, the kinematic limits in Table 1 must be understood to refer to appropriately weighted averages of mass eigenstates.
In general, the mass and flavor bases are related by a unitary transformation similar to the Cabibbo-Kobayashi-Maskawa mixing matrix in the quark sector. In the neutrino sector, the matrix is called the Maki-Nakagawa-Sakata (MNS) matrix . The customary 3-flavor version of the MNS matrix is given by Mann .
The CKM matrix is a 3-flavor matrix and one of its most intriguing features is the presence of a free complex phase that provides a natural origin for CP violation within the Standard Model. The MNS matrix has at least this degree of freedom, but in the neutrino sector neutrino-antineutrino mixing can occur as well, giving neutrinos both Dirac and Majorana properties. The MNS matrix is then at least a 6-dimensional one, and can presumably have additional free complex phases. CP violation in neutrinos is potentially a very interesting phenomenon, although experimentally challenging since, as Fisher et al. show , it is necessary to be able to track the oscillatory nature of three flavor components at once. Neutrinoless double beta decay is also a CP microscope, if the mass sensitivity needed can be reached.
There are now 3 experimental signals indicative of neutrino oscillations and mass. The implications, especially if all 3 are correct, are explored below. The reader is referred to the reviews in these proceedings by Suzuki , Mann , and DiLella for a discussion of these epochal experiments.
## 2 Accelerator and Reactor Oscillation Experiments
From the initial experiment discovering the neutrino in 1957 to the present day, reactors and accelerators have been a mainstay of research into the properties of neutrinos. Extensive and modern reviews of the subject are given by DiLella and Fisher et al. .
Certain recent experiments stand out as particularly influential in shaping our present view of the properties of neutrinos:
* The Liquid Scintillator Neutrino Detector (LSND) experiment at Los Alamos National Laboratory has found evidence for $`\overline{\nu }_\mu \overline{\nu }_e`$ with $`0.2\mathrm{\Delta }m^22.0`$ eV<sup>2</sup> and $`0.04\mathrm{sin}^22\theta 0.0015`$. Confirmatory evidence has been obtained in the charge-conjugate channel $`\nu _\mu \nu _e`$.
* The KARMEN Experiment at Rutherford-Appleton Laboratory reports no signal in a region of parameter space overlapping much of that explored by LSND. The small remaining region defines the parameters just given for LSND, plus a small island at $`\mathrm{\Delta }m^2`$= 4.5 eV<sup>2</sup>, $`\mathrm{sin}^22\theta `$= 2.5 x 10<sup>-3</sup>.
* The Brookhaven E776 Experiment excludes (in the $`\nu _\mu \nu _e`$ channel) the small island at 4.5 eV<sup>2</sup>.
* The Chooz long-baseline reactor antineutrino experiment shows that $`\overline{\nu }_e`$ does not transform to anything for $`\mathrm{\Delta }m^2`$$`10^3`$ eV<sup>2</sup>, $`\mathrm{sin}^22\theta `$$`0.1`$. The Palo Verde experiment confirms this, at somewhat lower significance. These negative results are remarkably decisive in ruling out substantial atmospheric $`\nu _\mu \nu _e`$ conversion and also in blocking any possibility of large-$`\mathrm{\Delta }m^2`$, large-angle solutions to the solar neutrino problem.
Together with the results from atmospheric and solar neutrino measurements, these define the scenarios for neutrino mass and mixing that are most likely.
## 3 Atmospheric Neutrinos
The experiments designed to search for proton decay, IMB and Kamiokande , were obliged to deal quantitatively with atmospheric neutrino interactions as the major background to proton instability. It was known that the cosmic ray flux and resulting production rate of neutrinos in the upper atmosphere were uncertain at the level of about a factor of 2, but it gradually became apparent that the ratio of $`\nu _\mu +\overline{\nu _\mu }`$ to $`\nu _e+\overline{\nu _e}`$ was also not well predicted. The latter ratio, naively 2 (from pion and muon decay) and calculable to an accuracy of about 5%, was found to be low by a factor of typically 0.6. As it depends little on the cosmic-ray flux and normalization, the departure from the expected value was somewhat surprising, and neutrino-oscillation solutions were proposed. Towards the end of the operation of Kamiokande, evidence was accumulating for a zenith-angle dependence of the ratio (equivalent to a path-length dependence).
SuperKamiokande has now been in operation since April, 1996, with a fiducial mass of approximately 22,500 tons acquiring events at 10 times the rate of Kamiokande. Other detectors, MACRO and Soudan II, have also been accumulating data. The present situation is summarized in the comprehensive review by Mann . Evidence consistent with neutrino oscillation has emerged in the form of:
* An atmospheric neutrino flavor ratio ($`\nu _\mu /\nu _e`$) 0.68 the expected magnitude,
* A zenith-angle dependence in the $`\nu _\mu `$ flux that departs greatly from the no-oscillation expectation and agrees closely with an oscillation description, and,
* A zenith-angle dependence in the $`\nu _e`$ flux that agrees closely with the no-oscillation expectation.
In addition, the zenith-angle dependence significantly favors (by about 2$`\sigma `$) the $`\nu _\mu \nu _\tau `$ channel over the $`\nu _\mu \nu _s`$ channel when matter effects are taken into account. The best-fit oscillation parameters are,
$`\mathrm{\Delta }m^2`$ $`=`$ $`3.5\times 10^3\mathrm{eV}^2`$
$`\mathrm{sin}^22\theta `$ $`=`$ $`1.0`$
There is no evidence for the subdominant oscillation $`\nu _\mu \nu _e`$ channel at the same mass difference, from which one can set a limit of about 0.05 on the square of the NMS matrix element $`U_{e3}`$ in a 3-flavor description .
Are there any possible loopholes? The following points can be, and in some cases have been, made:
* The “R” value, the ratio of $`\mu `$-flavor to $`e`$-flavor observed compared to that calculated, is smaller than expected for the oscillation parameters constrained by the zenith-angle dependence .
* The experimentally measured value of the rate for inclusive CC reactions of $`\nu _\mu `$ on <sup>12</sup>C is about a factor 2 smaller than calculated, whereas the corresponding $`\nu _e`$ reaction has the expected rate .
* Recoil-order terms in the neutrino-nucleon cross sections, particularly the pseudoscalar form factor, have apparently been neglected. The pseudoscalar contribution has an effect of several percent in the ratio R (because it contains the charged lepton mass), and may have a significantly larger role in the angular distributions where it appears as an interference term. The angular distributions are used to extract the zenith-angle dependence.
* The zenith-angle dependence, a very convincing aspect of the evidence for oscillation because it is so model-independent, does in fact depend on the extent to which pions and muons range out in the earth before decaying, and hence also on the altitude at which they are produced, the primary cosmic-ray spectrum, interaction cross sections, etc.
None of these points is thought to constitute a major concern, and neutrino oscillation appears to be the logical explanation for the results.
## 4 Solar Neutrinos
The sun, it is believed, generates its energy by fusion reactions that can be summarized as
$$4\mathrm{p}+2\mathrm{e}^{}{}_{}{}^{4}\mathrm{He}+2\nu _e+26.731\mathrm{MeV}.$$
Each cycle through the hydrogen-burning process produces 2 electron neutrinos and it follows directly that the neutrino flux at the Earth’s surface is proportional to the thermal energy flux, which is an experimentally measured quantity. The electromagnetic solar constant (irradiance) $`I`$ = 0.1367 W cm<sup>-2</sup>. With a small correction (about 1%) for the energy carried away by the neutrinos themselves, the neutrino flux at the Earth’s surface is 6.44 x 10<sup>10</sup> cm<sup>-2</sup> s<sup>-1</sup>, independent of detailed models of the sun. It is necessary only that the sun be in hydrostatic equilibrium over a period considerably longer than the photon migration time, 10,000 years.
In practice, no detector presently exists that can measure the total flux of solar neutrinos. Detectors have thresholds and strongly energy-dependent sensitivities. In hydrogen burning, a number of pathways lead to <sup>4</sup>He, and a complex spectrum of neutrinos from $`pp`$, $`pep`$, <sup>7</sup>Be, $`hep`$, <sup>13</sup>N, <sup>15</sup>O, <sup>17</sup>F, and <sup>8</sup>B results. The spectral shape of each individual component, whether line or continuum, is determined by laboratory measurement and/or electroweak theory. The relative intensity of each component, on the other hand, depends strongly on the temperature and composition, and therefore on astrophysical models of the sun. The flux component most easily detected, <sup>8</sup>B, is also the most temperature sensitive, varying as the 25th power of the central temperature, and it is moreover a component so minor (0.01%) that it is unconstrained by the sun’s energy output.
Data from 5 experiments (3 different types of experiment) provide information on different combinations of the fluxes. The current results are summarized in Table 2. It is very surprising that, with only 3 independent types of measurement and 8 different neutrino sources in the sun, it is impossible to fit the data (well) without introducing neutrino oscillations or some other non-standard-model physics! This comes about qualitatively as follows:
* The hep flux, as will be discussed, is negligibly small in the flux balance, and is known to be so from the high-energy part of the SK solar neutrino spectrum.
* The pep flux is tied to the $`pp`$ flux in a way that depends very weakly (as the square root of the temperature) on models.
* The CNO and <sup>7</sup>Be neutrinos are detected both by Cl-Ar and by Ga-Ge, but not by Kamiokande and SK. While they cannot be individually disentangled in a model-independent analysis, that is not required to demonstrate the inconsistency of standard-physics solutions with the data.
Consequently, there are really only 3 relevant neutrino sources, namely $`pp`$ \+ $`pep`$, <sup>7</sup>Be + CNO, and <sup>8</sup>B. There are also 3 independent experimental measurements, plus, if one elects to apply it, the luminosity constraint that relates the neutrino flux to solar energy output.
There is no combination of the fluxes with all fluxes non-negative that fits the data. With the luminosity constraint applied to the total flux, the statistical significance of this conclusion is now at about the 3.5$`\sigma `$ level; of course, it depends on there not being large unknown systematic errors in the data or in the detectors’ neutrino cross sections, but it does not depend on models of the sun. Even if the luminosity constraint is abandoned (equivalent to allowing variability of the solar core over times of order 30–10,000 years, or more exotic possibilities), there is no solution at about the 2$`\sigma `$ level.
If the experimental errors have been properly estimated, then, this contradiction means that one of the assumptions made in fitting the data must be incorrect, and there are very few assumptions.
It must be concluded that the shape of the <sup>8</sup>B spectrum is not as expected, containing more strength at high energies and less at low, and/or the neutrino flavor content is not pure electron, which alters the relationship between the water-Čerenkov results and the radiochemical experiments (because elastic scattering, unlike inverse beta decay, can occur via the neutral-current interaction with neutrinos of all active flavors). These features, not permitted in the Minimal Standard Model, are characteristic of neutrino-oscillation solutions. In contrast to the standard-physics solution, such solutions can give an excellent account of all data.
The need for non-standard physics (presumably neutrino oscillations) is model-independent at the roughly 3.5-$`\sigma `$ level, but the derivation of specific oscillation parameters is done in the context of astrophysical solar models and experimental nuclear-physics inputs. Qualitatively, there are three 2-flavor solutions that describe the data reasonably well, and a fourth with a lower probability. These are termed the Large-Mixing-Angle \[LMA\], the Small-Mixing-Angle \[SMA\], the ‘LOW’, and the vacuum solutions; see Fig. 1. Table 3 summarizes the fit results within the framework of a standard solar model .
If neutrino oscillations are indeed the explanation of the solar neutrino problem, independent evidence for them might arise from:
* Spectral distortions in the <sup>8</sup>B flux,
* Day-night variations indicative of MSW regeneration in the Earth,
* Yearly variations beyond those expected from the Earth’s orbital eccentricity, and
* Neutral-current interaction rates larger than expected from measured charged-current rates.
The SuperKamiokande collaboration is continuing a program of careful energy calibration and accumulation of high-statistics data in search of a distortion of the spectrum. The expected effects are quite small at best and, at the upper end of the spectrum, vanishingly small for all but the vacuum solutions. There are, moreover, sources of distortion unrelated to neutrino physics:
* The possible presence of hep neutrinos (<sup>3</sup>He + p $``$ <sup>3</sup>H + e<sup>+</sup> \+ $`\nu _e`$ \+ 18.7 MeV). The hep spectrum is considerably harder than the <sup>8</sup>B one, and contributes extra intensity in the vicinity of the <sup>8</sup>B endpoint at 15 MeV and beyond (to the 19-MeV endpoint of the hep spectrum). Calculation of the rate of the reaction is difficult because the lowest-order Gamow-Teller matrix element is small owing to the near orthogonality of the radial wavefunctions; forbidden terms dominate. Fortunately, comprehensive first-principles calculations of the rate, including the higher partial waves that contribute in the solar plasma, have been recently reported . Horowitz made the first calculation of the continuum <sup>3</sup>P<sub>0</sub> axial-charge transition, finding the S-factor to be
$`S_{0,p}(E)`$ $`=`$ $`1.7\times 10^{17}\mathrm{eV}\mathrm{b},`$
which is almost as large as the ‘standard’ s-wave component in use heretofore ,
$`S_{0,s}(E)`$ $`=`$ $`2.3\times 10^{17}\mathrm{eV}\mathrm{b}.`$
Schiavilla reports:
$`S_{0,s}(E)`$ $`=`$ $`3.9\times 10^{17}\mathrm{eV}\mathrm{b}`$
$`S_{0,p}(E)`$ $`=`$ $`2.4\times 10^{17}\mathrm{eV}\mathrm{b}`$
$`S_{0,d}(E)`$ $``$ $`1\times 10^{18}\mathrm{eV}\mathrm{b},`$
for all the s, p, and d partial waves, respectively. The energy-dependence of $`S_{0,i}`$ is negligible for all partial waves, and thus
$`S_0(E)`$ $`=`$ $`6.3\times 10^{17}\mathrm{eV}\mathrm{b}`$
is an accurate value for this rate for the purposes of neutrino flux calculations. While 3 times larger than the value in previous use, it falls short of the 16.7 times needed to account fully for distortions being seen in SuperKamiokande .
* The shape of the neutrino spectrum from <sup>8</sup>B decay is not directly calculable since the final state (in <sup>8</sup>Be) is broad. The spectrum is inferred from the recoil alpha spectra in laboratory experiments having other objectives . Preliminary results of measurements at Notre Dame University specifically designed to address some possible systematic concerns indicate that the standard spectrum underpredicts the intensity in the endpoint region by a fraction that peaks at about 14% 2 MeV below the endpoint.
* The beta decay of <sup>8</sup>B to the <sup>8</sup>Be ground state is second-forbidden and has not been observed. With similar transitions (e.g. <sup>36</sup>Cl) as a guide, it can be expected to have a branch of order $`10^3`$. Its spectrum would in that case be similar in both magnitude and energy to the hep spectrum.
It is at the moment too soon to draw conclusions concerning neutrino oscillations from the shape of the spectrum at high energies, but the experimental and theoretical uncertainties are rapidly being reduced. One can expect before long to have a useful constraint on oscillation solutions from the spectral shape.
Day-night effects arise from matter regeneration in the varying path through the earth’s core, and are less dependent on details. No certain evidence for time variations beyond statistical fluctuations has shown up to date, but the most recent data from SK yield a two-standard deviation effect,
$`{\displaystyle \frac{ND}{N+D}}=0.065\pm 0.031\pm 0.013.`$
The absence of large day-night effects has already ruled out a large region of parameter space in the range 10$`{}_{}{}^{6}`$ $`\mathrm{\Delta }m^2`$$`10^5`$ eV<sup>2</sup> and 10$`{}_{}{}^{2}`$ $`\mathrm{sin}^22\theta `$. As Suzuki shows (for further discussion, see also Bahcall et al. ), the small but general night enhancement matches better with the LMA solution (the lower part, in the vicinity of $`\mathrm{sin}^22\theta `$=1.0, $`\mathrm{\Delta }m^2`$= 1.9 x 10<sup>-5</sup> eV<sup>2</sup>) than it does the SMA solution where effects are all negligible except when the sun is on the other side of the earth’s core. The latter possibility seems disfavored at almost 3$`\sigma `$. The details of the LMA region are shown in Fig. 2.
The fact that information on solar neutrino solutions was present in atmospheric neutrino data was evidently first noted by Giunti et al. . The solar LMA solution is only valid for mixing with active, not sterile, neutrinos . Hence this solution is in conflict with either the LSND experiment or the indications that the dominant atmospheric signal is $`\nu _\mu \nu _\tau `$. If LMA were nevertheless the correct solution, upward-going electron neutrinos from the atmosphere would convert appreciably and equally to $`\mu `$ and/or $`\tau `$ neutrinos, potentially causing a deficit in $`\nu _e`$ and an increase in $`\nu _\mu `$. The vacuum oscillation length is
$$L_0=2.47\frac{E_\nu }{\mathrm{\Delta }m^2}$$
when distances are in km, energies in GeV, and masses in eV. For the LMA solution, $`\mathrm{\Delta }m^23\times 10^5`$ and the oscillation length becomes one earth diameter (13,000 km) at a neutrino energy of 160 MeV. An inspection of Figs. 2 and 15 in Mann’s paper , shows, if anything, a slight excess in the low-energy $`\nu _e`$ up-down asymmetry. Peres and Smirnov show , however, that when the calculation is done in detail and matter effects are taken into account, the excess is in fact expected for much of the atmospheric and LMA parameter space. The fact that the $`\nu _\mu `$ flux is intrinsically about 2 times the $`\nu _e`$ flux and that both oscillation solutions are near maximal mixing conspire to produce effects that may be of either sign, depending on the specific parameter values.
The implication is that the LMA solution may be favored by the atmospheric and solar neutrino data, good news for the KamLAND experiment, a reactor $`\overline{\nu }_e`$ experiment that reaches the LMA parameter space, as DiLella has described .
Yearly time dependence is the hallmark of vacuum oscillation solutions as the earth’s orbital eccentricity explores different oscillation phases. The eccentricity is small (a total 3.5% distance variation), and for the continuous neutrino spectrum emitted by <sup>8</sup>B and detected with detectors having relatively poor energy resolution, the effects are hard to see. It will be very different when high-statistics detectors primarily sensitive to <sup>7</sup>Be neutrinos (e.g. Borexino ) come on line, because the narrow line width of the source leads to striking time-dependence in the vacuum oscillation signal.
The question of the ratio of charged to neutral currents will be addressed in the Sudbury Neutrino Observatory, a 1000-tonne heavy-water Čerenkov detector now operating 2000 m underground in the INCO Creighton nickel mine near Sudbury, Ontario. SNO will permit observation of both the charged-current (CC) inverse beta decay of the deuteron:
$`\begin{array}{cccc}\hfill \mathrm{d}+\nu _e& & \mathrm{p}+\mathrm{p}+e^{}1.44\mathrm{MeV}\hfill & \end{array}`$
and the neutral-current (NC) neutrino breakup reaction:
$`\begin{array}{cccc}\hfill \mathrm{d}+\nu _x& & \mathrm{p}+\mathrm{n}+\nu _x2.22\mathrm{MeV}\hfill & \end{array}`$
The nuclear-physics uncertainties in the cross sections for these two processes arise mainly from the final states, which are members of the same isospin triplet. As a result, the ratio expected is known to a precision of order 1%, and significant departures (e.g. a factor of 3, as present solar neutrino information would suggest) would point unequivocally to neutrino oscillations to an active species. If the oscillation is to sterile neutrinos, there are nevertheless spectroscopic and time-dependent signatures that may be measurable.
In addition to the NC/CC ratio, SNO will provide good information on the shape of the <sup>8</sup>B spectrum above 5 MeV because in the CC reaction on deuterium the energy of the incident neutrino is transferred largely to the electron. Sensitivity to day-night and yearly effects is similar to that of SuperKamiokande, but the spectroscopic resolution permits the energy-dependence of such effects to be investigated efficiently. The quasi-elastic CC cross section rises quadratically with energy, and the backgrounds at the 2000-m depth of SNO are low, so detection of $`hep`$ neutrinos may be possible.
## 5 Interpreting the Results
The width of the Z<sup>0</sup> permits the existence of 3 light neutrinos and their antineutrinos that couple universally to the weak interaction. If only three different mass eigenstates $`m_i`$, $`i=1,2,3`$, exist, the mass splittings must satisfy
$$\underset{\mathrm{Splittings}}{}\mathrm{\Delta }m_\nu ^2=(m_3^2m_2^2)+(m_2^2m_1^2)+(m_1^2m_3^2)=0,$$
(1)
a trivial condition which is not met by any combination of the independent $`\mathrm{\Delta }m_\nu ^2`$ from experimentally favored neutrino mass differences (Table 4).
How can this difficulty be evaded? One must either assume that an experimental datum is incorrect (or at least misinterpreted), or that a fourth neutrino type exists, one that does not couple to the weak interaction significantly. Specific remedies that have been proposed are,
* A sterile neutrino that mixes with one or more active species.
* Reject the LSND result. Although there is evidence for the effect in both $`\overline{\nu }_\mu \overline{\nu }_e`$ and $`\nu _\mu \nu _e`$, the effect has not been seen in other experiments that explore similar, but not identical, regions . The LSND result is very constraining if correct, which perhaps accounts for the eagerness to reject it.
* Reject the Cl-Ar result or the water-Čerenkov solar-neutrino result. These experiments are both primarily sensitive to the <sup>8</sup>B flux. Together they fragment the allowed oscillation parameter space into islands in which shape distortions of the <sup>8</sup>B spectrum and neutral-current contributions allow the two results to be reconciled. All of those islands lie at small $`\mathrm{\Delta }m^2`$. Now, however, with the Chooz result , rejecting one or the other of those measurements no longer would permit a large-$`\mathrm{\Delta }m^2`$, large-$`\mathrm{sin}^22\theta `$ solution to the solar-neutrino problem.
## 6 Sterile Neutrinos
Heavy sterile neutrinos are a staple ingredient of most extensions to the Standard Model, but light sterile neutrinos have been regarded with distaste. The main reason is that the usual explanation invoked for the lightness of active neutrinos is the see-saw mechanism, realized with the aid of a sterile neutrino of GUT-scale mass. Having ‘used up’ that sterile Majorana neutrino component in producing light active neutrinos, one must then invoke complicated new mechanisms to generate and bring down other light sterile neutrinos.
Light sterile neutrinos are, however, just as natural as light active neutrinos if one does not start from the see-saw.
Our experience with the charged fermions is that they are described by the Dirac equation, to staggering precision. As the neutrino is a neutral fermion, it would be reasonable to suppose the Dirac equation should apply again. Dirac spinors are 4-component objects, with two spin states and distinct neutrinos and antineutrinos.
In deference to the handedness of the weak interaction, it is useful to project the four components, using the R/L and charge conjugation projection operators:
$`\psi _{R/L}=\frac{1}{2}(1\pm \gamma _5)\psi `$
C$`\psi _{R/L}`$C$`{}_{}{}^{1}=\psi _{R/L}^c`$
With 3 active neutrino flavors, mass terms in the Lagrangian have the form
$$_m(x)m_D\overline{\psi }(x)\psi (x)M_D\overline{\mathrm{\Psi }}(x)\mathrm{\Psi }(x)$$
where $`m_D`$ has been replaced by a nondiagonal $`3\times 3`$ matrix $`M_D`$ in flavor space and
$$\mathrm{\Psi }=\left(\begin{array}{c}\psi ^e\\ \psi ^\mu \\ \psi ^\tau \end{array}\right)$$
Mass terms in the Lagrangian must be Lorentz scalars, with no handedness. Following the development of Haxton and Stephenson and Langacker et al. , the resulting mass matrix takes on the form:
$$\begin{array}{c}(\overline{\mathrm{\Psi }}_L^c,\overline{\mathrm{\Psi }}_R,\overline{\mathrm{\Psi }}_L,\overline{\mathrm{\Psi }}_R^c)\\ \end{array}\left(\begin{array}{cccc}0& 0& 0& M_D^T\\ 0& 0& M_D& 0\\ 0& M_D^{}& 0& 0\\ M_D^{}& 0& 0& 0\end{array}\right)\left(\begin{array}{c}\mathrm{\Psi }_L^c\\ \mathrm{\Psi }_R\\ \mathrm{\Psi }_L\\ \mathrm{\Psi }_R^c\end{array}\right).$$
It allows for flavor oscillations if $`M_D`$ is nondiagonal. If CP is conserved, the four mass submatrices are equal .
The mass matrix is comprised of equal-mass active and sterile neutrinos and their antineutrinos, a pair for each generation. This situation is what one would naively expect if neutrinos were exactly like electrons, for example.
The upper left and lower right quadrants of this matrix must be zero because the left- and right-handed projectors annihilate each other. For charged fermions, charge conservation assures the remaining elements other than the ones already specified must be zero. However, for neutral fermions additional terms can be introduced elsewhere if we respect the requirement of hermiticity. Specifically,
$$_m(x)M_D\overline{\mathrm{\Psi }}(x)\mathrm{\Psi }(x)+\overline{\mathrm{\Psi }}_L^cM_L\mathrm{\Psi }_L+\overline{\mathrm{\Psi }}_R^cM_R\mathrm{\Psi }_R$$
so that the mass matrix becomes
$$\begin{array}{c}(\overline{\mathrm{\Psi }}_L^c,\overline{\mathrm{\Psi }}_R,\overline{\mathrm{\Psi }}_L,\overline{\mathrm{\Psi }}_R^c)\\ \end{array}\left(\begin{array}{cccc}0& 0& M_L& M_D^T\\ 0& 0& M_D& M_R^{}\\ M_L^{}& M_D^{}& 0& 0\\ M_D^{}& M_R& 0& 0\end{array}\right)\left(\begin{array}{c}\mathrm{\Psi }_L^c\\ \mathrm{\Psi }_R\\ \mathrm{\Psi }_L\\ \mathrm{\Psi }_R^c\end{array}\right)$$
The new Majorana mass terms break the local gauge invariance associated with a conserved lepton number. It is these nonDirac mass terms that can generate the nonzero $`m_\nu ^{Maj}`$ that must be present if neutrinoless $`\beta \beta `$ decay occurs.
The see-saw arises from setting $`m_L=0`$, $`m_Dm_e,m_\mu ,m_\tau `$, and $`m_RM_{\mathrm{GUT}}`$. After diagonalization, the eigenstates are a pair of light-heavy Majorana neutrinos in each generation. If, however, the Majorana mass terms are small (and so are the Dirac mass terms, although we do not understand why), then a different but equally interesting phenomenology results. The Majorana mass terms introduce a mixing between the active and sterile states, viz,
$$\psi _L\mathrm{cos}\theta \psi _L+\mathrm{sin}\theta \psi _R^c.$$
Such terms can be present within a generation and between generations, and lead to a complex 6-flavor neutrino mass, mixing, and charge-conjugation map. Gelb and Rosen have pursued this with a 4-flavor subset of neutrinos (3 active and one sterile), and show that not only can the observed mass splittings be reproduced, but the mixing angles are natural.
If the positive indications of neutrino oscillation from LSND, from atmospheric neutrinos, and from solar neutrinos are all correct, then either the atmospheric-neutrino or the solar-neutrino mixing must involve a sterile neutrino. There are two contradictory hints about which solution to choose. The zenith-angle distribution for partially-contained atmospheric neutrinos slightly favors active neutrino mixing. On the other hand, the indications of a possible day-night effect slightly favor the LMA solar solution, which only occurs with active neutrinos. In either case, the mass spectrum splits into a pair of doublets with the pair split by the LSND scale ($`0.5`$ eV) (Fig. 3).
Either the standard heirarchical order or the inverted order (solar neutrinos heaviest) are possible, but the sterile partner of the electron neutrino is always the heavier in order to have resonant conversion in the sun. Three sterile neutrinos significantly mixed with the active ones may create conflicts with the pace of evolution of the early universe as determined from the helium abundance, although one such sterile neutrino is probably not ruled out. The same constraint also arises in the ‘degenerate’ scenario in which the mass splittings are as given by oscillation data, but all masses are shifted up (becoming nearly degenerate in the process). Such a scenario is appealing as a source of hot dark matter (HDM), and marginally acceptable in nucleosynthesis if only one sterile neutrino species is mixed.
An important and perhaps unfortunate consequence of this particular sterile neutrino picture is the suppression of neutrinoless double beta decay. Neutrinos retain their Dirac nature by and large, with relatively small Majorana components.
While this structure was forced by the need to avoid conflicts between experimental results, it is important to bear in mind that it may be true even if one or more experiment is presently wrong.
## 7 Direct Methods – Tritium and Double Beta Decay
Oscillation experiments can never yield a value for the mass because such experiments are sensitive only to phase differences that arise from the differences in the squares of the masses.
Only two methods are presently known that have direct mass sensitivity that is at least roughly in the needed range, single beta decay (of tritium especially) and double beta decay. For many years these difficult experiments have been laboriously pursued, but there was always a worry that one was looking at the “wrong” neutrino, because the natural prejudice is to suppose that the $`\nu _e`$, $`\nu _\mu `$, $`\nu _\tau `$ mass heirarchy probably looks like the corresponding charged leptons.
With the discovery of oscillations virtually certain now, this picture has completely changed. The small mass differences that are representative of oscillations, and the links they forge between mass eigenstates, mean that to measure the mass of one eigenstate is to measure them all. The highly sensitive techniques applicable to the electron neutrino bring all the neutrino masses into the laboratory.
There are two tritium beta decay experiments currently in operation, one in Troitsk , the other in Mainz . Both make use of magnetic-electrostatic retarding-field analyzers. The Troitsk analyzer is connected to a gaseous T<sub>2</sub> source, the Mainz one to a solid frozen T<sub>2</sub> source. Steady progress has been made in both laboratories over the years in reducing backgrounds, improving stability and resolution, and checking for systematic effects. The sensitivity of both instruments is now in the range of 2 eV.
Initially both experiments reported large negative values of the parameter $`m^2`$. This parameter, when positive, represents the weighted average of the square of the neutrino mass,
$$m^2=\underset{i}{}|U_{ei}|^2m_i^2$$
but when negative serves as an effective parameter to continue the functional form of the beta spectrum into the non-physical regime (more events near the endpoint instead of fewer) to allow for statistical fluctuations. In fact, all recent tritium experiments have reported $`m^2<0`$, in some cases well beyond the level expected from statistical fluctuations, indicative perhaps of systematic effects. The main effect seen in the Troitsk experiment was traced to electrons trapped in the source region and escaping to the spectrometer only after having suffered energy loss. In the Mainz experiment the main effect was a morphological change in the structure of the tritium film, which increased the energy loss.
Once those problems were eliminated, a curious ‘step’ remained in much of the Troitsk data near the endpoint. (The spectrometer being an integral device, a step would correspond to a spike in the differential spectrum.) This step varied in intensity and position from run to run. In Fig. 4 the position of the step (which was always below the endpoint) and the intensity are shown. The position appeared to show a periodic motion with a period of 0.50 years. The last point, 98.3 is a run taken to test the prediction from the set of earlier data; it cannot be said either to agree or to disagree strongly with it.
In extracting a limit on the mass, the Troitsk group includes a step function in the fit. That is done consistently for all runs, with the step amplitude and position being fit parameters for each run. Unfortunately, this means that approximately half of the runs must be discarded since the step does not show up clearly in them and becomes excessively covariant with other fit parameters. When the step is fit, there remains no (other) non-standard contribution,
$$m_\nu ^2=2.0\pm 3.5\pm 2.1\mathrm{eV}^2$$
and a 95%-CL upper limit (Feldman-Cousins) on the neutrino mass is set at 2.5 eV. Fitting data without the step results in $`15m_\nu ^212`$ eV<sup>2</sup>, unless the fit is restricted to the last 70 eV of the spectrum, in which case $`m_\nu ^2=5\pm 5`$ eV<sup>2</sup>.
With a series of substantial technical improvements to their apparatus, the Mainz group succeeding in increasing the signal-to-background ratio tenfold and are able to take data at a sensitivity competitive with the Troitsk instrument. Out of 4 runs taken in 1997-8, one shows a step 12 eV from the endpoint at the same time as the Troitsk group measured one (shown as “98.2”in Fig. 4). At other times, 97.1, 98.1, and 98.3, no step was seen. The Mainz data fit without a step gives negative values for $`m_\nu ^2`$ very similar to the Troitsk ones ($`15m_\nu ^212`$ eV<sup>2</sup>), but a different prescription for negative $`m_\nu ^2`$ is used, so direct comparison is not possible. When the fit is restricted to the last 15 eV of the spectrum,
$$m_\nu ^2=0.1\pm 3.8\pm 1.8\mathrm{eV}^2$$
and a 95%-CL upper limit (Feldman-Cousins) on the neutrino mass is set at 2.9 eV.
For positive ions, the spectrometers are Penning traps (in which charged particles can be confined by axial magnetic fields and electrostatic potentials) and the Mainz group has begun to explore the possibility that a significant density of ions may accumulate in them. A cloud of such ions would show no kinetic gas pressure but could nevertheless cause inelastic collisions of the electrons being analyzed. One run has been carried out in which an oscillatory clearing electric field was applied at intervals to eject such ions, and in that run no negative $`m_\nu ^2`$ effect was seen. This is promising, but clearly a preliminary result. It must be remembered that the negative $`m_\nu ^2`$ effect was seen in the Mainz data without a step being present, and so they may not be the same phenomenon.
If a step really is present in the integral spectrum, it is a very exciting development. Capture of relic neutrinos produces a spike in the beta spectrum, because the decay energy is transferred entirely to the electron. There are no good laboratory limits on the local density of neutrinos near the earth – perhaps this is an indication that it is very large, of order $`10^{15}`$ cm<sup>-3</sup>. Stephenson et al. present a model in which such a density can arise, and which appears not to conflict with any known facts.
If on the other hand, the step turns out to be instrumental and removable, it is already clear that the present generation of instruments has the capability of reaching a limit in the range of 1-2 eV, a remarkable achievement. Plans are afoot in both groups for next-generation devices that will attack the 1-eV level.
Many nuclei provide an opportunity to search for neutrinoless double beta decay; single beta decay is blocked by energy conservation. Because large high-resolution detectors can be made from Ge (enriched in <sup>76</sup>Ge), the current best limit comes from the Heidelberg-Moscow collaboration :
$$<m_\nu ^{Maj}>=\underset{m}{}\lambda _m|U_{em}|^2m_m0.4\mathrm{eV}^2$$
where $`m`$ is an index summing over Majorana mass terms only and $`\lambda _m`$ is the CP phase.
## 8 The Future
With neutrino oscillations, and therefore mass, becoming more firmly established, the task of the experimentalist is clearer than at any time in the past – not necessarily easier, but clearer. No longer does the parameter ocean extend logarithmically to the horizon in every direction, but instead well-defined islands call out for close exploration.
Table 5, inspired by a similar approach in Fisher et al. , sets forth a list of key experiments and the implications of possible results of those experiments. The questions to be sorted out first include, what neutrino species are involved in the atmospheric and solar oscillation channels, can the LSND result be confirmed, what mixing parameters are responsible for the solar neutrino effects, and what is the magnitude of the mass itself? Answers to these questions will be hard-won but they seem within reach. Even more difficult and subtle issues are the charge-conjugation properties of neutrinos and their transformation under CP. The next decade will be an exciting time in neutrino physics.
Information and help generously given by John Beacom, Wick Haxton, Boris Kayser, Vladimir Lobashev, Ernst Otten, and Peter Rosen are most gratefully acknowledged.
Discussion
Jon Thaler (University of Illinois): Is there any observable effect in electron capture rates?
Robertson: Yes, electron capture decay has been studied as a means for measuring the neutrino mass (as distinct from the antineutrino mass). Capture from higher-lying atomic subshells can be cut off or attenuated if the neutrino has mass. In practice the interpretation of the resulting X-ray spectra is made difficult by rearrangement processes analogous to shakeup and shakeoff. |
warning/0001/astro-ph0001489.html | ar5iv | text | # Colour Change Measurements of Gravitational Microlensing Events by Using the Difference Image Analysis Method
## 1 Introduction
The amplification of source flux caused by gravitational lensing can become theoretically infinite. Points in the source plane on which the amplification of a point source becomes infinite are called caustics. For a point lens, the caustic is a single point behind the lens. For a binary lens, the number of caustics is multiple and they form closed curves. For a real microlensing event, however, the source star has an extended size, and thus the observed amplification deviates from that of a point source event and becomes always finite: extended source effect (Schneider, Ehlers, & Falco 1992; Witt & Mao 1994). The deviation of the light curve caused by the extended source effect becomes important as the source passes very close to the lens caustics.
Detection of the extended source effect is important because one can obtain useful information both about the lens and source star. First, a caustic crossing event provides an opportunity to measure how long it takes for the caustics to transit the face of the source star. By using the source radius crossing time $`t_{}`$, along with an independent determination of the source star size $`\theta _{}`$, one can determine the lens proper motion with respect to the source star by $`\mu =\theta _{}/t_{}`$ (Witt & Mao 1994; Gould 1994; Maoz & Gould 1994; Nemiroff & Wickramasinghe 1994; Loeb & Sasselov 1995; Alcock et al. 1997b, 1997c, 1999a; Afonso et al. 1998; Udalski et al. 1998; Albrow et al. 1999a). Once the value of $`\mu `$ is determined, the mass and location of the lens can be significantly better determined. Second, by analyzing the light curve of an event with the source trajectory passing very close to the caustics, one can investigate the surface structure of the source star such as the surface intensity profile and spots (Valls-Gabaud 1994, 1998; Sasselov 1997; Gaudi & Gould 1999; Heyrovský & Sasselov 2000; Albrow et al. 1999b; Han et al. 2000).
In addition to detecting deviations in the light curve, the extended source effect can also be detected by measuring the colour changes of the source star during the event. The colour changes occur due to the differential amplification over the limb-darkened source star surface (see § 3.1). Throughout this paper, we use a term ‘colour curve’ to refer to the colour changes of an event induced by the extended source effect as a function of time. Detection of the extended source effect from the colour measurements has the following advantages compared to the detection from a single band photometry. First, by detecting colour changes one can increase the chance to detect the extended source effect. If the source star approaches very close to the lens caustics but does not transit the caustics, the amplification induced by the extended source effect simply masquerades as changes in lensing parameters, and thus cannot be detected. By contrast, the colour change cannot be mimicked by the change in lensing parameters, because a point source lensing event should be achromatic (Gould & Welch 1996). Second, by measuring the colour changes one can determine the lens proper motion with relative ease. This is because one can measure $`t_{}`$, and thus $`\mu `$, by simply measuring the turning time of the colour curve (Han, Park, & Jeong 2000, see § 4 for more details). To determine $`t_{}`$ from the light curve, on the other hand, it is required to fit the overall light curve.
However, precise measurements of the colour changes induced by the extended source effect is hampered by blending, which also causes colour changes of the event (Kamionkowski 1995; Buchalter, Kamionkowski, & Rich 1996). To increase the event rate, the current lensing experiments are being conducted towards very dense star fields such as the Galactic bulge and the Magellanic Clouds (Alcock et al. 1993; Aubourg et al. 1993; Udalski et al. 1993; Alard & Guibert 1997). When the brightness of a source star located towards these very dense star fields is measured by using the conventional method based on the extraction of the individual source stars’ point spread functions, it is very likely that the measured source star flux is affected by the unwanted light from nearby unresolved stars (Di Stefano & Esin 1995; Alard 1997; Alcock et al. 1997a; Palanque-Delabrouille et al. 1998; Woźniak & Paczyński 1997; Han 1999; Han, Jeong, & Kim 1998). Since blended stars in general have different colours from the colour of the lensed source star, the measured colour is affected by blending. Han, Park, & Jeong (2000) pointed out the seriousness of blending in colour change measurements by demonstrating that even for a very small fraction ($`<2\%`$) of blended light, the colour change caused by blending can be equivalent to the colour change induced by the extended source effect. Therefore, for the precise measurements of the colour changes induced by the extended source effect, it will be essential to correct for or remove the blending effect.
Although numerous methods have been proposed for the correction of blending (Alard, Mao, & Guibert 1995; Alard 1996; Han 1997, 1998; Goldberg 1998; Goldberg & Woźniak 1998; Han & Kim 1999), most of these methods either have limited applicabilities only to several special cases of blended events or are less practical due to the requirement of space observations.<sup>1</sup><sup>1</sup>1We note, however, that the MACHO group had Hubble Space Telescope followup observations of the source star for one of the events they have detected to clearly identify the lensed source star. From these observations, they could accurately estimate the lens mass by correcting blending (NASA press release #STScI-PR00-03). On the other hand, with the recently developed photometric technique of the “difference image analysis” (DIA) method, one can measure the blending-free light variations of general microlensing events (Tomaney & Crotts 1996; Alard & Lupton 1998; Alard 1999; Alcock et al. 1999b, 1999c). The DIA method detects and measures the variation of the source star flux by subtracting an observed image from a convolved and normalized reference image. Since the blended light is subtracted during the image subtraction process, one can improve the photometric precision by removing the effect of blending. At the moment most microlensing experiments do not use the DIA method in their data analysis. Using the DIA method would require that either the whole microlensing survey process should switch to the DIA method, or that systematic re-processing of the light curves should be performed once events have been detected using the conventional photometric method.
Due to the methodological difference of the DIA photometry from the conventional photometry, the colour change of an event measured by the using the DIA method (DIA colour change) differs from that measured by the conventional method (see § 3). In this paper, we show, however, that the DIA colour curve of an event enables one to obtain the same information about the lens and source star, but with significantly reduced uncertainties due to the absence of blending. We investigate the patterns of the DIA colour curves for both single and binary lens events by constructing colour change maps.
## 2 Gravitational Amplification
### 2.1 Point Source Events
If a microlensing event with a point source is caused by a single lens, the light curve of the event is represented by
$$A_\mathrm{p}=\frac{u^2+2}{u(u^2+4)^{1/2}};u=\left[\left(\frac{tt_0}{t_\mathrm{E}}\right)^2+\beta ^2\right]^{1/2},$$
(1)
where the subscript ‘p’ denotes the event is occurred to a point source, $`u`$ represents the lens-source separation normalized by the angular Einstein ring radius $`\theta _\mathrm{E}`$, and the lensing parameters $`\beta `$, $`t_0`$, and $`t_\mathrm{E}`$ represent the impact parameter of the lens-source encounter, the time of maximum amplification, and the Einstein ring radius crossing time (Einstein time scale), respectively. The angular Einstein ring represents the effective region of gravitational amplification and its size is related to the physical parameters of the lens by
$$\theta _\mathrm{E}=\left(\frac{4GM}{c^2}\frac{D_{ls}}{D_{ol}D_{os}}\right)^{1/2},$$
(2)
where $`M`$ is the mass of the lens and $`D_{ol}`$, $`D_{ls}`$, and $`D_{os}`$ represent the separations between the observer, the lens, and the source, respectively. When the source is located within the Einstein ring, the source star amplification becomes greater than $`3/\sqrt{5}`$.
If an event is caused by a binary lens, on the other hand, the resulting light curve differs from that of a single lens event. When lengths are normalized by the combined angular Einstein ring radius, which represents the angular Einstein ring radius with a lens mass equal to the total mass of the binary system, the lens equation in complex notation for the binary lens system is represented by
$$\zeta =z+\frac{m_1}{\overline{z}_1\overline{z}}+\frac{m_2}{\overline{z}_2\overline{z}},$$
(3)
where $`m_1`$ and $`m_2`$ are the mass fractions of individual lenses (and thus $`m_1+m_2=1`$), $`z_1`$ and $`z_2`$ are the positions of the lenses, $`\zeta =\xi +i\eta `$ and $`z=x+iy`$ are the positions of the source and images, and $`\overline{z}`$ denotes the complex conjugate of $`z`$ (Witt 1990). The amplification of each image, $`A_{\mathrm{p},i}`$, is given by the Jacobian of the transformation (3) evaluated at the images position, i.e.
$$A_{\mathrm{p},i}=\left(\frac{1}{|\mathrm{det}J|}\right)_{z=z_i};\mathrm{det}J=1\frac{\zeta }{\overline{z}}\frac{\overline{\zeta }}{\overline{z}}.$$
(4)
Then the total amplification of the event is given by the sum of the individual amplifications, i.e. $`A_\mathrm{p}=_iA_{\mathrm{p},i}`$. The caustic refers to the source position on which the total amplification becomes infinite, i.e. $`\mathrm{det}J=0`$.
### 2.2 Extended Source Events
If the source of an event has an extended size, the light curve of the event deviates from that of a point source event. For both single and binary lens events, the amplification of an extended source event is the weighted mean of the amplification factor over the source star disk, i.e.
$$A_\nu =\frac{_0^r_{}I_\nu (r)A_\mathrm{p}(\left|𝒓𝒓_L\right|)r𝑑r}{_0^r_{}I_\nu (r)r𝑑r},$$
(5)
where $`I_\nu (r)`$ is the radial surface intensity distribution of the source star with a radius $`r_{}`$ and the vector $`𝒓_L`$ and $`𝒓`$ represent the displacement vector of the source star center with respect to the lens and the orientation vector of a point on the source star surface with respect to the center of the source star, respectively.
## 3 Colour Changes
### 3.1 Measured by the Conventional Photometric Method
In addition to causing deviations in the light curve from that of a point source event, the extended source effect makes the gravitational amplification become wavelength dependent, causing colour changes during the event. The colour changes occur due to the wavelength dependency of the source star radial surface intensity profile caused by limb darkening. If we define $`F_{0,\nu i}=2\pi _0^r_{}I_{\nu i}(r)r𝑑r`$, where $`i=1,2`$, and $`m_{\nu i}`$ as the unamplified source star fluxes and the corresponding magnitudes measured at two different wavelength pass-bands $`\nu 1`$ and $`\nu 2`$, the colour of the un-blended source star measured at a time $`t`$ by using the conventional method is computed by
$$(m_{\nu 2}m_{\nu 1})_0(t)=2.5\mathrm{log}\left[\frac{A_{\nu 2}(t)F_{0,\nu 2}}{A_{\nu 1}(t)F_{0,\nu 1}}\right].$$
(6)
Then the colour curve of the event is represented by
| $`\mathrm{\Delta }(m_{\nu 2}m_{\nu 1})_0(t)`$ | $`=`$ |
| --- | --- |
| | $`2.5\mathrm{log}\left\{\left[{\displaystyle \frac{A_{\nu 2}(t)}{A_{\nu 1}(t)}}\right]\left[{\displaystyle \frac{A_{\nu 2}(t_{\mathrm{ref}})}{A_{\nu 1}(t_{\mathrm{ref}})}}\right]^1\right\},`$ |
(7)
where $`t_{\mathrm{ref}}`$ represents the reference time for the colour change measurements. For a point source event, $`A_{\nu 1}(t)=A_{\nu 2}(t)`$ (i.e. achromatic), and thus $`(m_{\nu 2}m_{\nu 1})_0=2.5\mathrm{log}(F_{0,\nu 2}/F_{0,\nu 1})`$, implying that the colour of the source during the event equals to that of the un-amplified source. For a limb-darkened extended source event, on the other hand, as the source passes close to caustics, different parts of the source star disk with varying surface intensity and spectral energy distribution are amplified by different amount due to the differences in distance to the lens. As a result, the amplification becomes wavelength dependent, i.e. $`A_{\nu 1}(t)A_{\nu 2}(t)`$, and the colour changes during the event. Once the colour curve of the event is constructed, the lens proper motion and the source star surface intensity profile are determined by statistically comparing the observed colour curve to the theoretical curves with various models of limb darkening and source star size.
However, precise measurement of the colour changes induced by the extended source effect is hampered by blending, which is another mechanism causing colour changes of microlensing events. If one defines $`B_{\nu i}`$ as the blended amount of flux in two passbands, the measured colour of a blended event becomes
$`(m_{\nu 2}m_{\nu 1})(t)`$ $`=2.5\mathrm{log}\left[{\displaystyle \frac{A_{\nu 2}(t)F_{0,\nu 2}+B_{\nu 2}}{A_{\nu 1}(t)F_{0,\nu 1}+B_{\nu 1}}}\right]`$ $`=2.5\mathrm{log}\left[{\displaystyle \frac{A_{\nu 2}(t)+f_{\nu 2}}{A_{\nu 1}(t)+f_{\nu 1}}}\right],`$ (8)
where $`f_{\nu i}=B_{\nu i}/F_{0,\nu i}`$ represent the ratios between the blended and the baseline source flux in the individual pass-bands. Then the colour curve of a blended event is represented by
| $`\mathrm{\Delta }(m_{\nu 2}m_{\nu 1})=`$ | |
| --- | --- |
| | $`2.5\mathrm{log}\left\{\left[{\displaystyle \frac{A_{\nu 2}(t)+f_{\nu 2}}{A_{\nu 1}(t)+f_{\nu 1}}}\right]\left[{\displaystyle \frac{A_{\nu 2}(t_{\mathrm{ref}})+f_{\nu 2}}{A_{\nu 1}(t_{\mathrm{ref}})+f_{\nu 1}}}\right]^1\right\}.`$ |
(9)
One finds that equation (9) includes two additional blending parameters of $`f_{\nu 1}`$ and $`f_{\nu 2}`$ compared to the un-blended event colour curve in equation (7). Therefore, to determine the lens proper motion and the source star surface intensity profile from the fit of the blended event colour curve, it is required to include these additional parameters, causing increased uncertainties in the determined quantities.
### 3.2 Measured by the DIA Method
The colour changes induced by the extended source effect can also be measured by using the DIA method. The flux of a source star measured from the subtracted image by using the DIA method is represented by
$$F_\nu =F_{\nu ,\mathrm{obs}}F_{\nu ,\mathrm{ref}}=(A_\nu 1)F_{\nu ,0},$$
(10)
where $`F_{\nu ,\mathrm{obs}}=A_\nu F_{0,\nu }+B_\nu `$ and $`F_{0,\mathrm{ref}}=F_{0,\nu }+B_\nu `$ represent the source star fluxes measured from the images obtained during the progress of the event and from the reference image, respectively. Then the DIA colour curve of the event is represented by
| $`\mathrm{\Delta }(m_{\nu 2}m_{\nu 1})_{\mathrm{DIA}}=`$ | |
| --- | --- |
| | $`2.5\mathrm{log}\left\{\left[{\displaystyle \frac{A_{\nu 2}(t)1}{A_{\nu 1}(t)1}}\right]\left[{\displaystyle \frac{A_{\nu 2}(t_{\mathrm{ref}})1}{A_{\nu 1}(t_{\mathrm{ref}})1}}\right]^1\right\}.`$ |
(11)
From equation (11), one finds that the DIA colour curve does not depends on the blending parameters, and thus it is free from blending. One also finds that although the DIA colour curve has a different form from the curve of the un-blended event constructed by using the conventional method \[cf. equation (7)\], both curves depend on the same parameters of $`A_{\nu 1}`$ and $`A_{\nu 2}`$. Therefore, from the DIA colour curve one can obtain the same information about the lens and source star as that obtained from the colour curve constructed by using the conventional method, but with significantly reduced uncertainties due to the absence of blending.
## 4 Patterns of DIA Colour Change Curves
In this section, we investigate the patterns of the DIA colour curves for both single and binary lens events. To see the patterns, we construct the colour maps, which represent the colour changes as a function of source (lens) position with respect to the position of the lens (source). The constructed map are presented in the upper panel of Figure 1 for the single lens system and Figure 2 for the binary lens system. For the construction of the maps, we assume that the source star has an angular radius (normalized by $`\theta _\mathrm{E}`$) of $`\theta _{}=0.1`$ and it is observed in $`B`$ and $`I`$ bands. The un-blended colour of the source star before amplification is $`(BI)_0=2.15`$ mag, which corresponds to that of a K-type giant (Allen 1973; Schmidt-Kaler 1982; Peletier 1989). For the surface intensity profile, we adopt a linear form of
$$I_\nu (r)=1𝒞_\nu \left[1\sqrt{1(r/r_{})^2}\right],$$
(12)
where the limb-darkening coefficients are $`𝒞_B=0.912`$ and $`𝒞_I=0.053`$, respectively, which correspond to those of a K giant with $`T_{\mathrm{eff}}=4,750`$ K, $`\mathrm{log}g=2.0`$, and a metallicity similar to the sun (Van Hamme 1993). For the binary lens system, we adopt a binary separation (also normalized by $`\theta _\mathrm{E}`$) of $`\mathrm{}=1.0`$ and a mass ratio of $`q=m_1/m_2=1.0`$. The single lens map is presented as a function of lens positions $`(x_L,y_L)`$ with respect the source star (the thick solid circle with its center at the origin), while the binary lens map is presented as a function of source position $`(\xi ,\eta )`$ with respect to the center of the binary system. All lengths in the maps are normalized by $`\theta _\mathrm{E}`$. For both maps, we choose the reference of the colour change measurements to be $`\mathrm{\Delta }(BI)_{\mathrm{DIA}}=0`$ when the source star is not gravitationally amplified. Grey scale is drawn so that it becomes brighter (darker) as the colour of the source star becomes redder (bluer), and the tone of the grey scale changes for every colour change of $`0.02`$ mag. For the single lens map, the grey tone change from the darkest region where $`\mathrm{\Delta }(BI)_{\mathrm{DIA}}0.06`$ mag to the brightest region where $`\mathrm{\Delta }(BI)_{\mathrm{DIA}}>0.04`$ mag. For the binary lens map, the tone changes from the darkest region where $`\mathrm{\Delta }(BI)_{\mathrm{DIA}}0.05`$ mag to the brightest region where $`\mathrm{\Delta }(BI)_{\mathrm{DIA}}>0.07`$ mag. To better show the fine structures of the binary lens event map in the region near the caustics (marked by thick solid curves), the region enclosed by a box (drawn by a short-dashed line) is expanded and presented in the upper panel of Figure 3. The straight lines in the upper panels of Figure 1 and 3 represent various lens (for the single lens case) or source (for the binary lens case) trajectories and the colour curves resulting from the trajectories are presented in the lower panels of the individual figures. The line types of the colour curves are selected so that they match with those of the corresponding lens or source trajectories.
From the colour map of the single lens system and the resulting colour curves, one finds the following patterns. First, the iso-colour-change contours are concentric circles with their center at the center of the source star. Due to the radial symmetry of the map, all the resulting colour curves are symmetric with respect to the time of maximum amplification. Second, the colour does not monotonically change as a function of the separation between the lens and the center of the source star, $`r_L=\sqrt{x_L^2+y_L^2}`$. Outside the source star disk, the colour becomes redder as $`r_L`$ decreases, but within the disk it becomes bluer with the decreasing value of $`r_L`$. As a result, while the colour of the non-source transit event continues to become redder as the lens approaches the source star, the colour curve of a source transit event is characterized by turns at the moment when the lens enters (or leaves) the source star surface. Measurement of the turning time, $`t_{}`$, is important because one can determine the source radius crossing time from the measured value of $`t_{}`$ by
$$t_{}=\left[\beta ^2+\left(\frac{|t_{}t_0|}{t_\mathrm{E}}\right)^2\right]^{1/2}t_\mathrm{E},$$
(13)
where the lensing parameters of $`\beta `$, $`t_0`$, and $`t_\mathrm{E}`$ are determined from the overall shape of the light curve. We note that the described patterns of the DIA colour curves are very similar to those of the un-blended event colour curves constructed by using the conventional method (see Han, Park, & Jeong 2000).
The patterns of the colour map for the binary lens system and the resulting colour curves are as follows. First, compared to the single lens map, the the binary lens map is much more complex. In the regions around fold caustics, the iso-colour-change contours are parallel with the caustics. However, we note that the contours on the left and right sides of the caustics are not symmetric with respect to the caustic line. In the regions around cusps, on the other hand, the contours form separate peaks of blue colour changes. Second, the colour curves of caustic crossing binary lens events are characterized by the turns during caustic crossings. The colour becomes redder as the source approaches the lens caustics. It becomes most reddish at around the time when the edge of the source star touches the caustics. As the source further approaches to the caustics and thus the inner region of the source star disk lies on the caustics, the colour becomes bluer. During each caustic crossing, there are two red peaks which occur when the left and right sides of the source star disk lie on the caustics. By measuring the time separation between the two consecutive red peaks, $`\mathrm{\Delta }t`$, one can estimate the approximate value of the source radius crossing time by
$$t_{}\frac{\mathrm{\Delta }t}{2\mathrm{sin}\varphi },$$
(14)
where $`\varphi `$ is the angle between the source trajectory and the tangential line of the fold caustic at the position where the source crosses the caustics.<sup>2</sup><sup>2</sup>2If the fold caustic is a straight line, the relation in equation (14) is exactly valid. However, for an extended source event the relation is an approximation because the fold caustics can no longer be treated as a straight line. The value of $`\varphi `$ can be determined from the global fit of the binary lens event light curve.
## 5 Uncertainties in DIA Colour Measurements
In this section, we show by example that although the colour changes induced by microlensing are usually small, they can be measured with uncertainties small enough for one to extract useful information about the lens and source stars. We estimate the uncertainties of measured colours, $`\delta [\mathrm{\Delta }(m_{\nu _2}m_{\nu _1})_{\mathrm{DIA}}]`$, in the following ways. Since the uncertainty in the determined source star flux in magnitude is related to the signal-to-noise ratio by
$$\delta m_\nu =\frac{(\delta F_{0,\nu }/F_{0,\nu })}{0.4\mathrm{ln}10}\frac{1.09}{S/N},$$
(15)
the uncertainty in the measured colour is related to the signal-to-noise ratio by
$$\delta [\mathrm{\Delta }(m_{\nu _2}m_{\nu _1})_{\mathrm{DIA}}]\sqrt{2}\delta m_\nu \frac{1.54}{S/N}.$$
(16)
Then, if $`S/N10`$, the uncertainty of the colour measurement will be $`\delta \mathrm{\Delta }(m_{\nu _2}m_{\nu _1})_{\mathrm{DIA}}0.15\mathrm{mag}`$. The signal measured from the subtracted image is proportional to the source flux variation, i.e. $`S(A_\nu 1)F_{0,\nu }t_{\mathrm{exp}}`$, while the noise originates from both the lensed source and blended background stars, i.e. $`N[(F_{0,\nu }A_\nu +B)t_{\mathrm{exp}}]^{1/2}`$. Here $`t_{\mathrm{exp}}`$ is the mean exposure time and $`B`$ represents the average total flux of faint unresolved stars within the effective seeing disk with a radius $`\theta _{\mathrm{see}}`$. Then the signal-to-noise ratio is computed by
$$S/N=F_{0,\nu }(A_\nu 1)\left(\frac{t_{\mathrm{exp}}}{F_{0,\nu }A_\nu +B}\right)^{1/2}.$$
(17)
By using equations (16) and (17), we estimate the uncertainties in the measured colours for several example Galactic bulge single lens events with various impact parameters. We assume the events have a common source size of $`\theta _{}=0.0756`$ and an Einstein time scale of $`t_\mathrm{E}=67.5/2`$ days by adopting the values of the MACHO Alert 95-30 for which the extended source effect was actually detected (Alcock et al. 1997c). The lensed source is a K-type giant with $`I=14.05`$ mag. The observations are assumed to be carried by using a 1 m telescope with a CCD camera that can detect 12 photons/s for $`I=20`$ star. Note that the detection rate is determined considering both the efficiency and gain of the CCD camera. Since the major targets for colour measurements are high amplification events with bright source stars, proper exposure time is important to prevent saturation of images. In addition, since the amplification varies rapidly for these events, we leave $`t_{\mathrm{exp}}`$ as a variable so that the measured signal to be $`40,000`$ photons, which corresponds to the median flux of the linear sensitivity region for typical modern CCD cameras. We estimate $`B`$ by assuming that blended light comes from stars fainter than the crowding limit. Towards the Galactic bulge field, the crowding limit is set when the stellar number density reaches $`10^6`$ stars/$`\mathrm{deg}^2`$. Based on the model luminosity function constructed by combining those of Holtzman et al. (1998) and Gould, Bahcall, & Flynn (1997)<sup>3</sup><sup>3</sup>3For the detailed method of the model LF construction, see Han, Jeong, & Kim (1998)., this number density corresponds to the de-reddened $`I`$-band magnitude of $`I_018.1`$. The background flux is normalized for stars in the seeing disk with $`\theta _{\mathrm{eff}}=2^{\prime \prime }.0`$.
In Figure 4, we present the colour curves of the events along with the estimated uncertainties, which are presented by error bars. From the figure, one finds that it will be difficult to distinguish colour change curves for different non-source transit events ($`\beta >\theta _{}`$) due to large uncertainties. However, the uncertainties for source transit events ($`\beta \theta _{}`$) are small enough for one to distinguish curves with different values of $`\beta `$ and determine the turning times.
Although it is not included in the above uncertainty estimate, another possible source of uncertainty comes from the residual noise due to flat fielding (C. Alard, private communication). Since this noise is constant, it becomes important especially in the wings of the colour change curve at which the signal is very low as $`A_\nu 10`$. However, we find that in the region near the peak amplification where most information from the colour curve is obtained, this noise is not important. For example, let us the residual noise $`\delta F_\nu `$ to be $`5\%`$ of the unlensed source star flux, i.e. $`\delta F_\nu 0.05F_{0,\nu }`$. In the region close to the source star, e.g. $`\beta <0.1`$ (and thus $`A>10`$), the uncertainty in the DIA colour measurements estimated by using equations (15) and (16) will be $`<0.008`$ mag for the same source of the example event in the above error estimate. This uncertainty is substantially (nearly an order) smaller than the colour change of $`>0.05`$ mag for source transit events.
## 6 Summary
We investigate the chrometicity of single and binary microlensing events when their colour changes are measured by using the recently developed DIA photometric method. The findings from this investigation are summarized as follows:
1. The microlensing induced colour change measured by the DIA method differs from that measured by using the conventional photometric method.
2. Despite the difference, from the DIA colour curve one can obtain the same information about the lens and source star as the information obtained from the conventional colour curve. This is because both colour curves depend on the same parameters.
3. However, since the DIA colour curve is not affected by blending, one can determine the lens proper motion and source star surface intensity profile with significantly reduced uncertainties.
4. While the colour of a non-source transit event continues to become redder as the lens approaches the source star, the colour curve of a source transit event is characterized by turns the moments when the lens enters and leaves the source star disk. For the source transit event, one can determine the source radius crossing time by measuring the turning time of the colour curve.
5. The colour curveas of binary lens events are much less symmetric and vary greatly depending on the lens system geometry and the source star trajectory. For a caustic crossing event, one can measure the source radius crossing time by measuring the time separation between the two consecutive red peaks on the colour curve.
This work was supported by the grant (KRF-99-041-D00442) of the Korea Research Foundation. |
warning/0001/math0001163.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Graph theory is a natural instrument for matrix determinants and combined features calculating. It is clear that determinants can be expressed in form of alternating sum of matrix elements products taking by some graphs corresponding to index permutation, each term sign is determined by permutation even. However, the alternation itself is an essential obstruction both in numerical calculating and in analytic research of spectrum properties. It concerns for example cases when matrix elements depend on small (large) parameter and possess on it different orders. Thus under singular perturbed equations research (in particular such as Fokker-Plank equations) necessity appears for spectral analysis of big order matrix having exponentially small elements . The same situation appear under consideration of Markov’s chains connected with diffusion process where the questions of stochastic continuity, possible subprocess kinds and their structure properties run into the necessity of spectral analysis just such kind of matrixes.
Our aim is to obtain the formulas for matrix determinants, algebraic adjunctions, characteristic polynomial coefficients, components of eigen-vectors in the form of signless sums of matrix elements products taking by special graphs. It turns out that one can get such formulas in terms of so called ”tree” structure of some graph corresponding to matrix.
The first step in this direction was made by Kirchhoff who computed the number of connected subgraphs containing all vertices and containing no circuits (spanning trees). This number turned out to be equal to cofactor of any element of so called conductivity or Kirchhoff matrix of non-directed graph. Later this theorem was generalized for all coefficients of characteristic polynomial of this matrix, and also for cases of directed graphs where every arc (ordered pairs of vertices) possesses some quantity called weight . Here we get the same type formulas not for special matrices (as Kirchhoff type) but for arbitrary ones.
## 2 Main definitions and designations
Unification of designations and even terminology proper is not complete yet in graph theory. So firstly we adduce the necessary definitions and notations.
Let $`G`$ be digraph (directed graph). We use $`𝒱G`$ and $`𝒜G`$ to denote the set of vertices and arcs of $`G`$. The subgraph $`H`$ of $`G`$ is called factor if $`𝒱H=𝒱G`$. The outdegree (indegree) of the vertex $`i`$ (the number of arcs going out of (into) $`i`$) we denote $`d^+(i)`$ ($`d^{}(i)`$).
If $`G`$ is digraph in which every arc has its own weight $`g_{ij}`$ (weighted adjacencies digraph), corresponding matrix $`𝐆=\{g_{ij}\}_{i,j=1}^N`$ is called generalized adjacency matrix (the element $`g_{ij}=0`$ if there is no arc $`(i,j)`$ in $`G`$).
The sequence of following each other arcs along their orientations is called a way if all vertices besides possibly the uttermost ones are different. The way connecting the vertices $`m`$ and $`n`$ we denote $`mn`$. The cyclic way is called dicircuit. Linear digraph is digraph every vertex of which has unit in- and out-degree. So it consists of dicircuits.
We associate with any weighed adjacencies digraph $`G`$ the quantity $`\pi _G`$ by the rule
$$\pi _G=\underset{(i,j)𝒜G}{}g_{ij}$$
which is naturally to call by productivity.
Later on forests are the main graph theory object we use. As known there are two forest kinds in digraph situation. Here we call by forest digraph without dicircuits in which every vertex outdegree is equal to 0 or 1 ($`d^+(i)=0,1`$). The only vertex of tree (component of forest) having zero outdegree ($`d^+(i)=0`$) we call root.
Let $`𝒩`$ be finite set and $`𝒲`$ is some subset of $`𝒩`$. By $`_𝒲^k`$, $`k|𝒩||𝒲|`$ we denote the set of forests $`F`$ obeying the following conditions
1) $`𝒱F=𝒩`$;
2) $`F`$ consists of exactly $`k+|𝒲|`$ trees;
3) The set of roots of $`F`$ contain $`𝒲`$ as a subset.
Suppose also $`^k=_{\mathrm{}}^k`$, $`_𝒲=_𝒲^0`$. Note that the set $`^0`$ is empty set and $`^{|𝒩|}`$ consists of the only empty forest having only roots and no arcs.
If some additional condition $`z`$ is put on forests from $`_𝒰^k`$ we denote such set of forests as $`_𝒰^k(z)`$. Other necessary utilized notations we sign as necessary and also we use sometimes the term ”graph” in wide sense designating by it digraphs with weighted adjacencies too.
## 3 Circuit- and tree-like structures
The matrix (graph) spectral analysis can be carried out using some form of characteristic polynomial of matrix itself or some special matrixes (graphs) constructed from it. Thus it is valid known ”theorem on coefficients for digraphs” .
Theorem. Let
$$det(\lambda 𝐈𝐀)=\lambda ^N+a_1\lambda ^{N1}+\mathrm{}+a_N$$
be characteristic polynomial of arbitrary digraph $`A`$ with weighted adjacencies $`a_{ij}`$. Then
$$a_i=\underset{L_i}{}(1)^{p(L)}\pi _L,i=1,2,\mathrm{},N,$$
(1)
where $`_i`$ is the set of all linear directed subgraphs $`L`$ of graph $`A`$ with exactly $`i`$ vertices; $`p(L)`$ means the number of components (dicircuits) of $`L`$.
Coefficients in (1) are expressed in ”circuit” structure of $`A`$, and this theorem is not more than rephrasing from the standard determinant notation $`|\lambda 𝐈𝐀|`$ in the form of matrix elements products sum with sign determined by substitution even (the number of dicircuits) into graph terms.
In terms of ”tree” structure it is known the characteristic polynomial expression not for matrix $`𝐆`$ itself but for its Kirchhoff matrix $`𝐂`$ (or conductivity matrix) determined like
$$𝐂=𝐃𝐆,$$
where $`𝐃`$ is weighted powers matrix
$$𝐃\mathrm{diag}(\underset{j=1}{\overset{N}{}}g_{1j},\underset{j=1}{\overset{N}{}}g_{2j},\mathrm{},\underset{j=1}{\overset{N}{}}g_{Nj}).$$
Corresponding expression has a form :
$$det(\lambda 𝐈𝐂)=(1)^N\underset{k=0}{\overset{N}{}}(\lambda )^k\left[\underset{F^k(G)}{}\pi _F\right],$$
(2)
$$\pi _F=\underset{(i,j)𝒜F}{}f_{ij}=\underset{(i,j)𝒜F}{}g_{ij}.$$
Here the set of forests containing directly $`k`$ trees and being factors of $`G`$ is designated by $`^k(G)`$. Note that as the sum of elements along every line of $`𝐂`$ is equal to zero, so its determinant is equal to zero too and the sum in (refein) one can lead from $`k=1`$. In the following in clear cases we omit indication on graph.
## 4 The characteristic polynomial in tree-like structure terms
Knowing the characteristic polynomial expression of the admittance matrix $`𝐂`$ in terms of tree structure it is not hard to get analogous expression for characteristic polynomial of the matrix $`𝐆`$ itself. For this aim let us construct from $`G`$ some new graph $`G^{}`$ by the next rule. Let $`𝒩^{}=𝒩\{\}`$ be the set of vertices $`𝒩=\{1,2,\mathrm{},N\}`$ of $`G`$ which is supplemented by some new vertex designated $``$. Let also add to $`G`$ the vertex $``$ and lay on out of every vertex $`i𝒩`$ the arc $`(i,)`$ with weight $`g_i=\underset{j=1}{\overset{N}{}}g_{ij}`$, and remove all loops $`(i,i)`$. We denote the obtained graph by $`G^{}`$ ( $`G^{}:𝒱G^{}=𝒩^{}`$, $`𝒜G^{}=\{𝒜G\underset{i𝒩}{}(i,i)\}\underset{i𝒩}{}(i,)`$ , weights of arcs are equal $`g_{ij}`$ , $`i𝒩`$ , $`j𝒩^{}`$). Corresponding generalized adjacencies matrix $`𝐆^{}`$ of $`G^{}`$ has the following form
$$\text{G}\text{}=\left(\begin{array}{cccccc}0& g_{12}& g_{13}& \mathrm{}& g_{1N}& g_1\\ g_{21}& 0& g_{23}& \mathrm{}& g_{2N}& g_2\\ g_{31}& g_{32}& 0& \mathrm{}& g_{3N}& g_3\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ g_{N1}& g_{N2}& g_{N3}& \mathrm{}& 0& g_N\\ 0& 0& 0& \mathrm{}& 0& 0\end{array}\right).$$
Note that, if $`𝐆=𝐏𝐈`$, where $`𝐏`$ is probability matrix setting finite Markov’s chain with killing, so the quantity $`g_i`$ is the probability of killing of the process if it is in state $`i`$. This has a sense of probability of outcoming to the bounder, so the additional vertex $``$ could be interpreted in some sense as a bounder of the finite set $`𝒩=\{1,2,\mathrm{},N\}`$.
It is easy to see, that $`(N+1)\times (N+1)`$ admittance matrix $`𝐂^{}`$ of graph $`G^{}`$ has the form
$$\text{C}\text{}=\left(\begin{array}{cc}\text{G}& \begin{array}{c}g_1\\ g_2\\ \mathrm{}\\ g_N\end{array}\\ & \\ \begin{array}{cccc}0& 0& \mathrm{}& 0\end{array}& 0\end{array}\right),$$
so, using (2) we get the chain of equations
$$det(\lambda 𝐈𝐆)=\frac{1}{\lambda }det(\lambda 𝐈+𝐂^{})=\frac{1}{\lambda }\underset{k=1}{\overset{N+1}{}}\lambda ^k\left[\underset{F^k(G^{})}{}\pi _F\right]=$$
$$=\underset{k=0}{\overset{N}{}}\lambda ^k\left[\underset{F^{k+1}(G^{})}{}\pi _F\right].$$
Note, that since the vertex $``$ in graph $`G^{}`$ has zero outdegree and hence it is a root in every forest $`F^l(G^{})`$, so the sets $`^{k+1}(G^{})`$ and $`_{}^k(G^{})`$ coincide. Thus it is valid
Theorem 1. Characteristic polynomial of an arbitrary $`N\times N`$ matrix $`𝐆`$ can be expressed in the form
$$det(\lambda 𝐈𝐆)=\underset{k=0}{\overset{N}{}}\lambda ^k\left[\underset{F_{}^k}{}\pi _F\right],\pi _F=\underset{(i,j)𝒜F}{}g_{ij},$$
(3)
where $`_{}^k_{}^k(G^{}).`$
Under $`\lambda =0`$ we obtain obvious
Consequence. The determinant of $`N\times N`$ matrix $`𝐆`$ can be expressed in the form
$$det𝐆=(1)^N\underset{F_{}}{}\pi _F,$$
(4)
$$_{}_{}(G^{}).$$
Let $``$ be a subset of $`𝒩`$. Designate by $`𝐆_{}`$ matrix obtained from $`𝐆`$ by striking out the columns and lines with numbers $`i`$. So, $`𝐆_{}`$ is a diagonal minor of $`𝐆`$ of $`(N||)`$-th order. The corresponding to it digraph we denote by $`G_{}`$.
Consequence of consequence. The determinant of minor $`𝐆_{}`$ can be expressed in the form
$$det𝐆_{}=(1)^{N||}\underset{F_{\{\}}}{}\pi _F,$$
(5)
$$_{\{\}}=_{\{\}}(G^{}).$$
Proof. By formula (4)
$$det𝐆_{}=(1)^{N||}\underset{F_{}(G_{}^{})}{}\pi _F.$$
Let us keep in $`𝐆_{}`$ the same numeration of elements as it is in matrix $`𝐆`$. So the elements $`(g_{})_{il}`$ of the corresponding matrix $`𝐆_{}^{}`$ (weights of arcs $`(i,l)`$ of graph $`G_{}^{}`$) are equal:
$`(g_{})_{il}=g_{il},i,l𝒩,`$
$`(g_{})_{il}=0\{i,l\}\mathrm{},`$
$`(g_{})_i=\underset{m\{\}}{}g_{im},i𝒩.`$
So productivity $`\pi _F`$ of any forest $`F_{}(G_{}^{})`$ represents the productivity $`\pi _F^{}`$ of some forest of the set $`_{\{\}}(G^{})`$. The sum of productivities $`\pi _F`$ along all forests $`F_{}(G_{}^{})`$ exhausts the set $`_{\{\}}(G^{})`$, which proves (5).
## 5 Formulas for components of eigenvectors
If it is known the eigenvalue $`\lambda `$ of $`𝐆`$, so to calculate the components $`v_m`$ of corresponding eigenvector $`\stackrel{}{v}`$ it is necessary to decide the standard system $`(\lambda 𝐈𝐆)\stackrel{}{v}=0`$. Let for certainty the $`n`$-th component of eigenvector $`\stackrel{}{v}`$ be not equal to zero. Without loss of generality one can accept it be equal to one: $`v_n=1`$. Then we have for the rest of components following the Kramer rule
$$v_m=\frac{det\mathrm{\Delta }_{nm}^{}(\lambda )}{det\mathrm{\Delta }_{nn}(\lambda )},$$
(6)
where
$$\mathrm{\Delta }_{nn}(\lambda )=\lambda 𝐈𝐆_{nn},$$
$`𝐆_{nn}`$ is algebraic adjunct of the element $`g_{nn}`$ of matrix $`𝐆`$, and $`\mathrm{\Delta }_{nm}^{}(\lambda )`$ is a matrix obtained from $`\mathrm{\Delta }_{nn}(\lambda )`$ by substitution of the $`m`$-th column of $`\lambda 𝐈𝐆`$ by $`n`$-th one with negative sign.
The expression for $`\mathrm{\Delta }_{nn}(\lambda )`$ it is easy to get using already obtained formula (5) concerning diagonal minors $`𝐆_{}`$ of matrix $`𝐆`$. Since the coefficient at $`\lambda ^k`$ at the expression of characteristic polynomial of algebraic adjunct $`𝐆_{nn}`$ is itself a sum of determinants of diagonal minors of matrix $`𝐆`$ of $`(Nk1)`$-th order with sign $`(1)^{Nk1}`$ and not including $`n`$-th column and $`n`$-th line so one can write
$$det(\lambda 𝐈𝐆_{nn})=\underset{k=0}{\overset{N1}{}}\lambda ^k\left[\underset{\genfrac{}{}{0pt}{}{𝒩,n}{\left|\right|=k+1}}{}det(𝐆_{})\right]=$$
$$=\underset{k=0}{\overset{N1}{}}\lambda ^k\left[(1)^{Nk1}\underset{\genfrac{}{}{0pt}{}{𝒩,n}{\left|\right|=k+1}}{}det𝐆_{}\right]=$$
$$=\underset{k=0}{\overset{N1}{}}\lambda ^k[\underset{\genfrac{}{}{0pt}{}{𝒩,n}{\left|\right|=k+1}}{}\{\underset{F_{\{\}}(G^{})}{}\pi _F\left\}\right].$$
The sum at square brackets in the last expression obviously is equal to the sum of productivities $`\pi _F`$ of all forests belonging to the unification
$$\underset{\genfrac{}{}{0pt}{}{𝒩,n}{\left|\right|=k+1}}{}_{\{\}}(G^{}),$$
which in its turn is equal to $`_{\{,n\}}^k(G^{})`$. So, the denominator at (6) has the form
$$det\mathrm{\Delta }_{nn}(\lambda )=det(\lambda 𝐈𝐆_{nn})=\underset{k=0}{\overset{N1}{}}\lambda ^k\left[\underset{F_{\{,n\}}^k}{}\pi _F\right].$$
(7)
$$_{\{,n\}}^k_{\{,n\}}^k(G^{}),$$
It is more difficult to obtain the formula for numerator at (6) having the form
$$det\mathrm{\Delta }_{nm}^{}(\lambda )=\underset{k=0}{\overset{N1}{}}\lambda ^k\left[\underset{F_{\{,n\}}^k(mn)}{}\pi _F\right],$$
(8)
where by $`_{\{,n\}}^k(mn)`$ the subset of $`_{\{,n\}}^k(G^{})`$ having $`mn`$-walk is denoted. The summing over $`k`$ one can lead till $`N2`$, because the set $`_{\{,n\}}^{N1}(mn)`$ is empty.
It is sufficiently to establish (8) under $`\lambda =0`$, because the coefficients at powers of $`\lambda `$ are the diagonal minors of the same matrix and are utterly of analogical to each other form.
Let $`\mathrm{\Delta }_{nm}^{}\mathrm{\Delta }_{nm}^{}(0)`$. It is necessary to prove that
$$det\mathrm{\Delta }_{nm}^{}=\underset{F_{\{,n\}}(mn)}{}\pi _F.$$
(9)
In passing let us remark that using (9) it is easy to obtain the expression for algebraic adjunct $`G_{nm}`$ of the element $`g_{nm}`$. Actually, matrices $`\mathrm{\Delta }_{nm}^{}`$ and $`G_{nm}`$ differ only in such a way. In $`\mathrm{\Delta }_{nm}^{}`$ (if for example $`n<m`$) $`n`$-th column of matrix $`G_{nm}`$ is at the $`(m1)`$-th place and is of opposite sign. Hence
$$G_{nm}=(1)^{N+mn1}\underset{F_{\{,n\}}}{}\pi _F.$$
To prove (9) let us exchange in $`𝐆`$ places of $`n`$-th and $`m`$-th columns. Obtained auxiliary matrix we denote by $`𝐇`$. Nondiagonal elements $`h_{ij}`$ of $`𝐇`$ and corresponding expanded matrix $`𝐇^{}`$ are equal
i) $`h_{ij}=g_{ij},jn,m`$ ,
ii) $`h_{in}=g_{im},h_{im}=g_{in}`$ .
It is easy to see that algebraic adjunct $`𝐇_{nn}`$ of the element $`h_{nn}`$ of $`𝐇`$ differs from matrix $`(\mathrm{\Delta }_{nm}^{})`$ by only the sign of $`m`$-th column. So taking into account (7)
$$det\mathrm{\Delta }_{nm}^{}=(1)^{N1}det𝐇_{nn}=\underset{F_{\{,n\}}(H)}{}\pi _F.$$
(10)
Now in (10) it is necessary to cross from the sum of productivities $`\pi _F`$ of factor forests of $`H^{}`$ to the sum of productivities $`\pi _F^{}`$ of corresponding subgraphs of graph $`G^{}`$. Notice that in $`H^{}`$ weight of arc $`(m,n)`$ is equal to minus sum of weights of arcs $`(m,i)`$ going out of the vertex $`m`$ in graph $`G^{}`$:
$$h_{mn}=g_{mm}=\underset{i𝒩^{}\{m\}}{}g_{mi}.$$
(11)
That is why we divide the set of forests $`_{\{,n\}}(H^{})`$ in to two nonintersecting subsets: the set $`_{\{,n\}}(H^{};(m,n))`$ of forests including the arc $`(m,n)`$ and the set $`_{\{,n\}}(H^{};(m,n))`$ of forests without this arc. Then
$$\underset{F_{\{,n\}}(H^{})}{}\pi _F=\underset{F_{\{,n\}}(H^{};(m,n))}{}\pi _F+\underset{F_{\{,n\}}(H^{};(m,n))}{}\pi _F.$$
(12)
The set of arcs $`𝒜F`$ of the forest $`F_{\{,n\}}(H^{};(m,n))`$ is representable in a form $`𝒜F=(m,n)𝒜P`$ , where $`P`$ is some graph belonging to the set $`_{\{,n,m\}}(H^{})`$. By it as $`h_{ij}=g_{ij}`$, $`jm,n`$ the sets $`_{\{,n,m\}}(H^{})`$ and $`_{\{,n,m\}}(G^{})`$ coincide. Hence, taking into account (11)
$$\underset{F_{\{,n\}}(H^{};(m,n))}{}\pi _F=\underset{i𝒩^{}\{m\}}{}g_{mi}[\underset{F_{\{,n,m\}}(G^{})}{}\pi _F].$$
Note that any forest $`F_{\{,n,m\}}(G^{})`$ consists of exactly of three trees $`T_{},T_n`$ and $`T_m`$ with roots correspondingly $`,n`$ and $`m`$. Therefore the addition to such forest the arc $`(m,i)`$ with weight $`g_{mi}`$ leads either to graph belonging to $`_{\{,n\}}(G^{})`$ (if $`i𝒱T_{}𝒱T_n`$), or (if $`i𝒱T_m`$) to graph, in which outdegree of every vertex is equal to one with the exception of $``$ and $`n`$$`d^+(j)=1j,n`$ , $`d^+()=d^+(n)=0`$, and there is exactly one circuit in graph and this circuit contains the vertex $`m`$. We denote the set of such graphs by $`𝒪_m(G^{})`$. So
$$\underset{F_{\{,n\}}(H^{};(m,n))}{}\pi _F=\underset{F_{\{,n,\}}(G^{})}{}\pi _F\underset{F𝒪_m(G^{})}{}\pi _F.$$
(13)
Let us now consider the last sum at the right part of (12). The productivity $`\pi _F`$ of any forest $`F_{\{,n,\}}(H^{};(m,n))`$ by force of i) ii) one can represent as $`\pi _F=\pi _E`$, where $`E`$ is some subgraph of $`G^{}`$, to obtain which from the forest $`F`$ one must exchange the arcs ending in $`m`$ into arcs ending in $`n`$ and back. The forest $`F_{\{,n,\}}(H^{};(m,n))`$ consists of two trees $`T_{}`$ and $`T_n`$ with the roots $``$ and $`n`$ correspondingly. By this if $`m𝒱T_{}`$ so the graph $`E`$ is still a forest and besides the sequence of arcs from $`m`$ does not lead to $`n`$ (but leads to $``$). In the case of $`mT_n`$ in $`F`$ such an exchanging of arcs results in $`E`$ be graph containing the only circuit and the vertex $`m`$ belongs to this circuit, i.e. $`E𝒪_m(G^{})`$. Thus
$$\underset{F_{\{,n\}}(H^{};(m,n))}{}\pi _F=\underset{F_{\{,n\}}(G^{};mn)}{}\pi _F+\underset{F𝒪_m(G^{})}{}\pi _F.$$
(14)
Here we denote by $`_{\{,n\}}(G^{};mn)`$ the subset of $`_{\{,n\}}(G^{})`$ not containing $`mn`$-way. Combining together (10),(12),(13) and (14) we get, that products $`\pi _F`$ along graphs containing circuits reduce and
$$det\mathrm{\Delta }_{nm}^{}=\underset{F_{\{,n\}}(G^{})}{}\pi _F\underset{F_{\{,n\}}(G^{};mn)}{}\pi _F=\underset{F_{\{,n\}}(G^{};mn)}{}\pi _F,$$
which is the same as (9). Thus it is proved
Theorem 2. Let $`\lambda `$ be the eigenvalue of order 1 of matrix G, $`\stackrel{}{v}`$ — corresponding eigenvector and its $`n`$-th component is not equal to zero, then accurate to constant factor the components of $`\stackrel{}{v}`$ are following
$$v_n=1,v_m=\frac{\underset{k=0}{\overset{N2}{}}\lambda ^k\left[\underset{F_{\{n,\}}^k(mn)}{}\pi _F\right]}{\underset{k=0}{\overset{N1}{}}\lambda ^k\left[\underset{F_{\{n,\}}^k}{}\pi _F\right]}.$$
(15)
where $`_{\{n,\}}^k(mn)`$ is a subset of set $`_{\{n,\}}^k`$, consisting of forests having $`mn`$-way.
Remark. To obtain the expression for $`m`$-th component of eigenvector of the transform matrix $`𝐆^{}`$ it is necessary only to exchange indices $`m,n`$ with each other at the right part of (15).
## 6 Example
Let us consider an easy example demonstrating the application of singless form technique. Let $`𝐌=𝐏𝐈`$, where $`𝐏`$ is a probability matrix, setting the finite Markov’s chain with killing (that is the sum of elements along any line may be less than one, corresponding residual is just a killing probability), and also let the transition probabilities $`M_{ij}=P_{ij}`$, $`ij`$, and the killing ones $`M_i=\underset{j=1}{\overset{N}{}}M_{ij}=1\underset{j=1}{\overset{N}{}}P_{ij}`$ be exponentially small. Namely, let $`N=3`$ and $`M_{ij}=m_{ij}e^{V_{ij}/\epsilon }`$ and $`V_{12}`$=$`V_{13}=4`$, $`V_1=5`$, $`V_{21}=3`$, $`V_{23}=2`$, $`V_2=5`$, $`V_{32}=1`$, $`V_3=4`$, $`V_{31}=3`$; $`\epsilon `$ is a small parameter. Thus if instead of diagonal elements $`M_{ii}`$ we use their expression through the quantities $`M_i`$ ($`M_{ii}=M_i\underset{ji}{\overset{N}{}}M_{ij}`$) the matrix $`𝐌`$ gets the form
$$𝐌=\left(\begin{array}{ccc}M_{12}M_{13}M_1& M_{12}& M_{13}\\ M_{21}& M_{21}M_{23}M_2& M_{23}\\ M_{31}& M_{32}& M_{32}M_3M_{31}\end{array}\right).$$
By virtue of lack of sign and non-negativity of transition (nondiagonal elements of $`𝐌`$) and killing ($`M_i`$) probabilities that are used at tree-like structure formulas, it is ought to keep at the asymptotic of characteristic polynomial coefficients (3) only terms reaching the maximal order on small parameter. Thus the eigenvalues asymptotic one can extract from the equation
$$\lambda ^3+\lambda ^2m_{32}e^{1/\epsilon }+\lambda m_{32}m_{21}e^{5/\epsilon }+m_{32}m_{21}m_1e^{10/\epsilon }=0.$$
The exponential orders $`V_k=\underset{\epsilon 0}{lim}\epsilon \mathrm{ln}a_k`$ of the coefficients $`a_k`$ of characteristic polynomial $`\lambda ^ka_k`$ satisfy convex nonequalities system analogical to : $`V_{k1}V_kV_kV_{k+1}`$, and the eigenvalues asymptotic one can look in the form $`\lambda _k^\epsilon \underset{\epsilon 0}{}\mathrm{\Lambda }_ke^{\frac{V_kV_{k1}}{\epsilon }}`$ . Substituting model $`\lambda `$ in such a form we get the following eigenvalues asymptotics:
$$\lambda _1\underset{\epsilon 0}{}m_1e^{5/\epsilon },\lambda _2\underset{\epsilon 0}{}m_{21}e^{3/\epsilon },\lambda _3\underset{\epsilon 0}{}m_{32}e^{1/\epsilon }.$$
The asymptotic of eigenvectors of the matrix $`𝐌`$ and the transform matrix $`𝐌^{}`$ one can find from (15). Denoting by $`𝐂`$ (by $`𝐂^{}`$) the matrix, composed of ultimate values of vector-lines of $`𝐌`$ (vector-columns of $`𝐌^{}`$) one gets
$$𝐂=\left(\begin{array}{ccc}1& 1& 1\\ 0& 1& 1\\ 0& 0& 1\end{array}\right),𝐂^{}=\left(\begin{array}{ccc}1& 1& 0\\ 0& 1& 1\\ 0& 0& 1\end{array}\right),\mathrm{𝐂𝐂}^{}=𝐈.$$
Notice also that the matrix
$$\stackrel{~}{𝐌}=\left(\begin{array}{ccc}m_1e^{\frac{5}{\epsilon }}& 0& 0\\ m_{21}e^{\frac{3}{\epsilon }}& m_{21}e^{\frac{3}{\epsilon }}& 0\\ 0& m_{32}e^{\frac{1}{\epsilon }}& m_{32}e^{\frac{1}{\epsilon }}\end{array}\right)$$
demonstrates the same in leading order eigenvalues asymptotic and the same ultimate values of eigenvectors’ components.
Although matrices $`𝐌`$ and $`\stackrel{~}{𝐌}`$ outwardly are rather different, however the solutions of the evolution equation
$$\frac{d\stackrel{}{p}}{dt}=𝐋\stackrel{}{p}$$
with guiding matrix $`𝐋`$ equal to $`𝐌^{}`$ or to $`\stackrel{~}{𝐌}^{}`$ are the same under exponentially large times $`t=\tau e^{V/\epsilon }`$, $`V>0`$, and in ”slow” time” $`\tau `$ with exponential scale $`V`$ the corresponding evolution $`\stackrel{}{p}_V(\tau )=\underset{\epsilon 0}{lim}\stackrel{}{p}(\tau e^{V/\epsilon })`$ has the form
$$\stackrel{}{p}_V(\tau )=𝐂^{}\left[\underset{\epsilon 0}{lim}\mathrm{diag}\{e^{\mathrm{\Lambda }_1\tau \mathrm{exp}(\frac{V5}{\epsilon })},e^{\mathrm{\Lambda }_2\tau \mathrm{exp}(\frac{V3}{\epsilon })},e^{\mathrm{\Lambda }_1\tau \mathrm{exp}(\frac{V1}{\epsilon })}\}\right]𝐂\stackrel{}{p}_V(0).$$
Under exponential scale $`V=1`$ part of the initial distribution concentrated at the third state crosses to the second one. Such crossing rapidness at ”slow” time $`\tau `$ is determined by the quantity $`\mathrm{\Lambda }_3=m_{32}`$. Under exponential scale $`V=3`$ if $`\tau \mathrm{}`$ entire distribution turns out to be at the first state, corresponding rapidness is determined by the quantity $`\mathrm{\Lambda }_2=m_{21}`$. Finally at the scale $`V=5`$ killing of the process takes place and it is governed by the quantity $`\mathrm{\Lambda }_1=m_1`$. Note, that in spite of the killing probability $`M_3=m_3e^{4/\epsilon }`$ at state 3 is greater than the killing probability $`M_1=m_1e^{5/\epsilon }`$ at state 1, but the part of distribution concentrated at state 3 crosses to the second state (and then to state 1) before killing at state 3 could take place.
Considered example shows from one side small matrix elements are not negligible comparatively with the large ones. From another side nevertheless one can neglect by some of the elements, may be not small ( here these are elements $`M_{12}`$, $`M_2`$, $`M_{23}`$, $`M_{31}`$, $`M_3`$). Which of elements do not affect on the spectrum and on the evolution is determined by signless formulas (3) and (15).
## 7 Discussion
Obtained expressions for characteristic polynomial (3) and components of eigenvectors (15) at the situation of arbitrary matrices give few substantial for practical computation in comparison with usual methods. Moreover, the formula (1) being only reformulation of standard characteristic polynomial expression is more preferable than (3), because the volume of calculations needed for determination of coefficients at powers of $`\lambda `$ using (1) is less than using (3). The reason is that $`G`$ has less subgraphs being circuits than $`G^{}`$ has forests. In such case signless formulas obtained are not of interest for practical calculations.
However, the situation strongly changes in the case of large matrices with non-negative elements which coefficients depend on small (large) parameter and may possess on it different orders (such matrices, as it was noted before, often appear in singular and stochastic problems). Usual methods in such a situation are practically unsuitable due to large amount of calculations and low accuracy. On the contrary, formulas (3) and (15) due to sign absence at sums become extremely effective and allow to keep an eye on only higher order in asymptotic. This circumstance permits not only to carry out practical calculations, but also (which is highly essentially from theoretical point oh view) to determine singular limits and analyze structure spectrum properties. Here it is necessary to note that though for matrix $`𝐆`$ with non-negative elements introducing quantities $`g_i=\underset{j=1}{\overset{N}{}}g_{ij}`$ are not positive however at the expense of spectrum shift it is always possible to make them non-negative ones. The problem of picking out at characteristic polynomial terms having highly order is reducing to determination forests $`F`$ with maximal productivity $`\pi _F`$ and is a separate difficult one. This problem is solved, the effective technique for determination of extreme forests is elaborated . Analogical technique can be used for analysis of large very sparse matrices.
This work was supported RFBR, grants N-99-01-00696 and N-98-01-01063. |
warning/0001/hep-th0001220.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The path integral formulation of the quantum mechanics of fermionic systems is usually associated with the introduction of Grassmann variables as pseudoclassical configuration variables (Nevertheless see Ref. for an alternative approach in the relativistic case). The path integral treatment of a single fermionic degree of freedom is very well understood but, surprisingly, the extension of the formalism to a space time description of relativistic and non relativistic spinning particles or to the solution of potential problems has not been developed yet. Several approaches using different sets of Grassmannian non-commuting variables have been proposed in the literature allowing to construct a pseudo-classical description of the dynamics of the spinning particle, (or pseudomechanics ) for the non-relativistic case, and for the Dirac electron but this is not enough for a quantum description. With the pseudo-classical action identified, in order to write down a path integral expression for the fermionic propagator, it is necessary to construct an explicit representation of the spin observables and the polarized states of the particle in the Hilbert space associated to the Grassmann variables. This will we done in what follows. On the other hand, in the usual approach, this obstacle is bypassed by showing, instead that the constraint which emerge from imposing a variational principle to the action functional is equivalent in the operatorial formalism to the wave equation.Then the form of the propagator is borrowed from this formalism. This strategy, although enlightening from the conceptual point of view is not useful for computational purposes. Another path integral formalism, which is based on the use of Grassmannian coherent states , has been devised for the description of fermionic systems. It allows the computation of the propagator and bound state energies but the relation of this formalism with the pseudo-classical description is not completely clear and in particular does not provide a direct interpretation of the path integral as a sum over histories in configuration space. In what follows we also discuss how we can get this interpretation for spinning particles. First, we show that with an explicit realization of the spin observables one can represent the spin polarized states in the Grassmannian sector of the superspace. Then we derive the path integral formulation of the non-relativistic electron as a sum over histories directly from the pseudo-classical description. Finally, we show that being careful with the boundary conditions of the Grassmann functions one is able to compute the probability amplitudes using a semiclassical expansion.
## 2 Wave functions and spin observables
Consider a real Grassmannian vector $`\stackrel{}{ϵ}`$ satisfying the anti-conmutation relations,
$$ϵ_iϵ_j=ϵ_jϵ_i$$
and the super-configuration space of coordinates ($`\stackrel{}{x}`$,$`\stackrel{}{ϵ}`$). Let us consider wave functions of both $`\stackrel{}{x}`$ and $`\stackrel{}{ϵ}`$ with the general expansion ,
$$\varphi (x,\stackrel{}{ϵ})=\varphi (x)+\varphi _i(x)ϵ_i+\varphi _{ij}(x)ϵ_iϵ_j+\varphi _{ijk}(x)ϵ_iϵ_jϵ_k.$$
(1)
The functions $`\varphi _{ij}(x)`$ and $`\varphi _{ijk}`$ are anti-symmetric. The wave functions belong to a $`8`$-dimensional complex vectorial space, $`(C/^8)`$. The internal product in this space may be defined by,
$$\varphi _1|\varphi _2=𝑑ϵ_3𝑑ϵ_2𝑑ϵ_1(\varphi ^{})I_ϵ\varphi _2.$$
(2)
where $`dϵ_3dϵ_2dϵ_1`$ is the Berezin integration measure. The operator $`I_ϵ`$ is,
$$I_ϵ=\frac{ϵ_{ijk}}{3!}\left(ϵ_i+_i\right)\left(ϵ_j+_j\right)\left(ϵ_k+_k\right).$$
(3)
and $`_m=\frac{}{ϵ_m}`$ are the Grassmannian right derivatives which satisfy, $`_mϵ_k=\delta _{mk}`$ and $`\text{ }_mϵ_kϵ_j=\delta _{mk}ϵ_j\delta _{mj}ϵ_k.`$ If one considers the eight independent functions in $`\varphi `$ as the components of a vector in $`(C/^8)`$, the internal product defined in (2) corresponds to the usual product in $`(C/^8)`$. The $`\delta `$ function in the odd sector is given by
$$\delta (\stackrel{}{ϵ}^{}\stackrel{}{ϵ})=\left(ϵ_1^{}ϵ_1\right)\left(ϵ_2^{}ϵ_2\right)\left(ϵ_3^{}ϵ_3\right).$$
(4)
Introduce the non-Hermitian position operator $`E`$ in the odd sector of the configuration space and the continuous set of eigenvectors $`|\stackrel{}{ϵ}`$ of $`E`$. Similarly as in the coherent state representation we have the relations
$$\stackrel{}{ϵ}|\stackrel{}{ϵ}^{}=e^{\stackrel{}{ϵ}^{}.\stackrel{}{ϵ}}.$$
(5)
An arbitrary wave function $`\varphi (x,ϵ_i)`$ is represented in Dirac notation in the form, $`\varphi (x,ϵ_i)=\stackrel{}{x},\stackrel{}{ϵ}|\varphi `$. The identity operator in the odd sector is
$$1\text{l}=𝑑ϵ_3𝑑ϵ_2𝑑ϵ_1|\stackrel{}{ϵ}^{}I_ϵ\stackrel{}{ϵ}|,$$
(6)
and we note also that, $`I_ϵe^{ϵ^{}ϵ}=\delta (ϵϵ^{})`$. The physical sector of this space should be expanded by the spin polarized states. To represent the spin observables $`\stackrel{}{S}`$ we introduce the differential operators,
$$S_i=\frac{i}{4}ϵ_{ijk}\left(ϵ_j+_j\right)\left(ϵ_k+_k\right)$$
(7)
which satisfy the angular momentum algebra, $`[S_i,S_j]=iϵ_{ijk}S_k`$. This representation of the spin observables is the natural generalization of the usual representation of the fermionic position-momentum algebra which led to the coherent state formulation of the path integral. It has been discussed also in Ref. . Looking at things from another point of view, the Hilbert space of states $`|\stackrel{}{ϵ}`$ and the operators (7) provide a fermionic coherent state representation of the $`SU(2)`$ algebra alternative to the bosonic approach . This construction may be generalized to other groups. The complete set of eigenfunctions of $`S_3`$ is given in the following table.
| $`f_\lambda ^n`$ | $`\lambda `$ | $`\varphi (\stackrel{}{ϵ})`$ |
| --- | --- | --- |
| $`f_+^1`$ | $`\frac{1}{2}`$ | $`1iϵ_1ϵ_2`$ |
| $`f_+^2`$ | $`\frac{1}{2}`$ | $`ϵ_3iϵ_1ϵ_2ϵ_3`$ |
| $`f_+^3`$ | $`\frac{1}{2}`$ | $`ϵ_1+iϵ_2`$ |
| $`f_+^4`$ | $`\frac{1}{2}`$ | $`ϵ_1ϵ_3iϵ_2ϵ_3`$ |
| $`f_{}^4`$ | $`\frac{1}{2}`$ | $`1+iϵ_1ϵ_2`$ |
| $`f_{}^3`$ | $`\frac{1}{2}`$ | $`ϵ_3+iϵ_1ϵ_2ϵ_3`$ |
| $`f_{}^2`$ | $`\frac{1}{2}`$ | $`ϵ_1iϵ_2`$ |
| $`f_{}^1`$ | $`\frac{1}{2}`$ | $`ϵ_1ϵ_3+iϵ_2ϵ_3`$ |
The eigenvalues of $`S_3`$, denoted by $`\lambda `$ are which, as can be seen, degenerate. To construct a particular base of states, we take a linear combination of them in such a way that the action of the up and down operators $`S_+=S_1+iS_2`$ and $`S_{}=S_1iS_2`$ is well defined. A possible choice of the eigenfunctions which represent the polarized states is
$`\stackrel{}{ϵ}|+`$ $`=`$ $`(1+ϵ_3)(1iϵ_1ϵ_2),`$ (8)
$`\stackrel{}{ϵ}|`$ $`=`$ $`(1ϵ_3)(ϵ_1iϵ_2).`$
## 3 The action functional and the path integral
Let us consider now a spinning particle whose dynamics is determined by the Hamiltonian operator $`H`$. We want to discuss a discretization procedure in the trajectories of the particle in the configuration space which allows to represent the Green functions in terms of a path integral. The evolution operator is given by,
$$U(t_ft_i)=e^{iH(t_ft_i)}.$$
(9)
with matrix elements, $`U(x,\stackrel{}{ϵ},t;x^{},\stackrel{}{ϵ}^{},t^{})=x^{},\stackrel{}{ϵ}^{}|U(t,t^{})|x,\stackrel{}{ϵ}.`$ For the spin polarized states ($`k=+,`$), we define the physical propagator, $`K(k,t_f;j,t_i=x_f,k|U|x_i,j`$ which may be projected in the form,
$$K(k,t_f;j,t_i)=𝑑ϵ_i𝑑ϵ_fk|ϵ_fI_{ϵ_f}ϵ_f|U(x_f,ϵ_f,t_f;x_i,ϵ_i,t_i)|ϵ_iI_{ϵ_i}ϵ_i|j.$$
(10)
(we drop the arrow on the Grassmann coordinates).Now consider a discretization $`\{t_1,t_2,\mathrm{}.,t_{2N}\}`$ with $`\delta =t_kt_{k1}`$ of the time interval, and using the resolution of unity let us compute the matrix element of the evolution operator. We get,
$`x_f,ϵ_f|e^{iH(t_ft_i)}|x_i,ϵ_i=`$
$`\underset{2N\mathrm{}}{lim}{\displaystyle \mathrm{\Pi }_{k=1}^{\mathrm{}}\frac{dp_k}{2\pi }\mathrm{\Pi }_{j=1}^{\mathrm{}}𝑑x_j𝑑ϵ_1^{}𝑑ϵ_1𝑑ϵ_2^{}𝑑ϵ_2\mathrm{}𝑑ϵ_{2N}^{}𝑑ϵ_{2N}}`$
$`x_f,ϵ_f||p_{2N},ϵ_{2N}I_{ϵ_{2N}}p_{2N},ϵ_{2N}|e^{iH\delta }|x_{2N},ϵ_{2N}^{}I_{ϵ_{2N}^{}}x_{2N},ϵ_{2N}^{}|\mathrm{}`$
$`\mathrm{}ϵ_k,p_k|e^{iH\delta }|x_k,ϵ_k^{}I_{ϵ_k^{}}x_k,ϵ_k^{}||p_{k1},ϵ_{k1}I_{ϵ_{k1}}p_{k1},ϵ_{k1}|\mathrm{}`$
$`\mathrm{}ϵ_1,p_1|e^{iH\delta }|x_1,ϵ_1^{}I_{ϵ_1^{}}x_1,ϵ_1^{}|x_i,ϵ_i.`$
The general term can be expanded in the form,
$`I_{ϵ_k}ϵ_k,p_k|e^{iH\delta }|x_k,ϵ_k^{}I_{ϵ_k^{}}x_k,ϵ_k^{}|p_{k1},ϵ_{k1}=`$
$`I_{ϵ_k^{}}e^{iH(p_k,x_k,ϵ_k)\delta }e^{ip_kx_k}e^{ϵ_k^{}ϵ_k}I_{ϵ_k^{}}e^{ip_{k1}x_{k1}}e^{ϵ_{k1}ϵ_k^{}}=`$
$`e^{iH_0(p_k,x_k)\delta }e^{ip_kx_k}e^{ip_{k1}x_{k1}}I_{ϵ_k^{}}e^{iH_1(x_k,ϵ_k)\delta }e^{ϵ_k^{}ϵ_k}I_{ϵ_k^{}}e^{ϵ_{k1}ϵ_k^{}}`$
were we use that the hamiltonian function satisfies,
$$H(p_k,x_k,ϵ_k)=H_0(p_k,x_k)+H_1(x_k,ϵ_k).$$
(11)
This is obvious if $`H`$ does not depend on the $`ϵ`$ derivatives, but it is also true in the general case due to the external integrals. Let us focus in the Grassmann sector alone and note that under the integral sign we have,
$`ϵ_f|ϵ_{2N}I_{ϵ_{2N}}e^{iH_1(x_{2N},ϵ_{2N)}\delta }e^{ϵ_{2N}^{}ϵ_{2N}}I_{ϵ_{2N}^{}}e^{ϵ_{2N1}ϵ_{2N}^{}}\mathrm{}I_{ϵ_1}e^{iH_1(x_1,ϵ_1)\delta }e^{ϵ_iϵ_1}=`$
$`e^{\frac{1}{2}ϵ_{2N}ϵ_f}e^{iH(p_{2N},x_{2N},ϵ_{2N})\delta +\frac{1}{2}ϵ_{2N}\left(ϵ_fϵ_{2N1}\right)}e^{iH(p_{2N1},x_{2N1},ϵ_{2N1})\delta +\frac{1}{2}ϵ_{2N1}\left(ϵ_{2N}ϵ_{2N2}\right)}\mathrm{}`$
$`e^{iH(p_1,x_1,ϵ_1)\delta +\frac{1}{2}ϵ_1\left(ϵ_2ϵ_i\right)}e^{\frac{1}{2}ϵ_iϵ_1}.`$
The terms at the end of the interval are of the form,
$`\stackrel{}{ϵ}_{2N}\stackrel{}{ϵ}_f`$ $``$ $`(\stackrel{}{ϵ}_f\delta \dot{\stackrel{}{ϵ}_f})\stackrel{}{ϵ}_f=\delta \stackrel{}{ϵ}_f\dot{\stackrel{}{ϵ}_f},`$ (12)
$`\stackrel{}{ϵ}_i\stackrel{}{ϵ}_1`$ $``$ $`\stackrel{}{ϵ}_i(\stackrel{}{ϵ}_i+\delta \dot{\stackrel{}{ϵ}_i})=\delta \stackrel{}{ϵ}_i\dot{\stackrel{}{ϵ}_i}.`$ (13)
In the limit $`2N\mathrm{}`$ ($`\delta 0`$), they reduce to boundary terms which depend only of initial and final values
$$g(\stackrel{}{ϵ}_i,\stackrel{}{ϵ}_f)=\underset{\delta 0}{lim}\frac{1}{2}\left\{_{t_i}^{t_i+\delta }\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}𝑑t+_{t_f\delta }^{t_f}\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}𝑑t\right\}.$$
(14)
Incorporating the bosonic sector we are left with,
$$U(\stackrel{}{ϵ}_f,t_f;\stackrel{}{ϵ}_i,t_i)=D[ϵ]D[x]D[p]e^{g(\stackrel{}{ϵ}_i,\stackrel{}{ϵ}_f)+i{\scriptscriptstyle \left\{\dot{\stackrel{}{x}}\stackrel{}{p}{\scriptscriptstyle \frac{i\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}}{2}}H(x,p,ϵ)\right\}𝑑t}}.$$
(15)
The action functional recovered in the measure of the path integral is the one that appears in the pseudoclassical description of the spinning particle . To compute the quantum ampitude between physical states, one introduces (15) in (10).
## 4 The semiclassical approximation and the variational principle
Bosonic path integrals with quadratic potentials may be computed using a semiclassical approximation . In this section we show that a similar result holds also in the case under consideration if proper care is given to the boundary terms. The point here is that, since the equations for $`\stackrel{}{ϵ}`$ are first order, it is not possible in general to find trajectories $`x(t)`$ and $`\stackrel{}{ϵ}(t)`$, extremals of $`S`$ in the time interval $`t_ft_i`$, with $`x(t_i)=x_i`$, $`x(t_f)=x_f`$, $`ϵ(t_i)=ϵ_i`$ and $`ϵ(t_f)=ϵ_f`$. So we introduce two Lagrange multipliers $`\pi _i,\pi _f`$ and consider instead an extended action
$$S^{}[ϵ(t),\pi _i,\pi _f]=S[ϵ(t)]+\pi _i(ϵ(t_i)ϵ_i)\pi _f(ϵ(t_f)ϵ_f).$$
(16)
The equations of motion are
$`2\left(\left({\displaystyle \frac{L}{\dot{ϵ}}}\right)(t_f)\pi _f\right)\delta (tt_f)2\left(\left({\displaystyle \frac{L}{\dot{ϵ}}}\right)(t_i)\pi _i\right)\delta (tt_i)=0,`$ (17)
$`x(t_i)=x_i,x(t_f)=x_f.`$
Now we can fix the values of the Lagrange multipliers to guarantee that the boundary conditions, which here appear as independent equations, are satisfied. In fact, we still have the freedom to fix $`\pi _i`$ to zero. Then the solution to the equations may be written in the form
$$\stackrel{}{ϵ}_{class}=\stackrel{}{ϵ_0}+\stackrel{}{\delta _ϵ}\mathrm{\Theta }(tt_f)$$
(18)
where $`\stackrel{}{ϵ_0}`$ satisfies
$$\dot{\stackrel{}{ϵ}_0}=i\frac{H}{\stackrel{}{ϵ}_0}.$$
(19)
and $`\delta _ϵ`$ is a jump at the end of the trajectory. To perform the semiclassical expansion let us consider first the free case with the action given by,
$$iS=g(\stackrel{}{ϵ}_i,\stackrel{}{ϵ}_f)+i\left[\frac{i\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}}{2}\right]𝑑t.$$
(20)
The solution to the equations of motion which satisfies the boundary condition is simply
$$\stackrel{}{ϵ}(t)=\stackrel{}{ϵ}_i+2(\stackrel{}{ϵ}_f\stackrel{}{ϵ}_i)\mathrm{\Theta }(tt_f).$$
(21)
Consider the path integral (15) computed in the previous section with the boundary term (14), and let us perform an expansion around $`\stackrel{}{ϵ}_{class}`$,
$$\stackrel{}{ϵ}(t)=\stackrel{}{ϵ}_{class}(t)+\stackrel{}{\xi }(t).$$
(22)
Substituting (22) and (21) in (15) we get the expected result,
$$U(\stackrel{}{ϵ}_f,t_f;\stackrel{}{ϵ}_i,t_i)=Ne^{ϵ_iϵ_f}.$$
(23)
Let us turn out our attention to a more general case. The most general even Hamiltonian function has the form,
$$H(x,p,ϵ)=H_0(x,p)+H_{ij}(x)ϵ_iϵ_j.$$
(24)
Then, the equation of motion is linear in $`ϵ`$. Using the equation of motion for $`ϵ_0`$, and the linearity of the equation of motion it is readily seen that the boundary term takes the form,
$$g(\stackrel{}{ϵ}_i,\stackrel{}{ϵ}_f)=\frac{\stackrel{}{ϵ}_0(t_f)\stackrel{}{ϵ}_f}{2}.$$
(25)
This result generalizes for the interacting case the expression obtained by Galvao and Teitelboim . Consider again expressions of the form (22) and (22). Substitution in (15) leads us to the expression,
$`g(\stackrel{}{ϵ_i},\stackrel{}{ϵ_f})+iS=\stackrel{}{ϵ_0}(t_f)\stackrel{}{ϵ_f}+`$ (26)
$`i{\displaystyle _{t_i}^{t_f}}𝑑t\left\{\left[i{\displaystyle \frac{\stackrel{}{ϵ_0}\dot{\stackrel{}{ϵ_0}}}{2}}+H_0(x,p)+H_{ij}ϵ_{0i}ϵ_{0j}\right]+\left[{\displaystyle \frac{i}{2}}\stackrel{}{\xi }\dot{\stackrel{}{\xi }}+H_{ij}\xi _i\xi _j\right]\mathrm{}\right\}`$
for the exponent. (The dots appear to denote possible bosonic contributions). Then, in the case when the spin degrees of freedom factorize, we get the following simple expression for the matrix elements
$$U(\stackrel{}{ϵ}_f,t_f;\stackrel{}{ϵ}_i,t_i)=Ne^{2g(\stackrel{}{ϵ}_f,\stackrel{}{ϵ}_i)}.$$
(27)
where $`N`$ is a normalizacion constant.
## 5 Spin precession
Let us recover the known results for particle in a uniform magnetic field for example. The action is given by
$$S=g(\stackrel{}{ϵ}_i,\stackrel{}{ϵ}_f)+\left(\dot{\stackrel{}{x}}\stackrel{}{p}\frac{i}{2}\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}+\frac{p^2}{2m}+\frac{q}{2m}\left(\frac{ϵ_{ijk}}{2}iϵ_jϵ_k\right)B_i\right)𝑑t.$$
(28)
Defining,
$$M_{jk}=\frac{q}{2m}\frac{ϵ_{ijk}}{2}B_i,$$
(29)
we are left with the Lagrangian,
$$L=\frac{i\stackrel{}{ϵ}\dot{\stackrel{}{ϵ}}}{2}+i\frac{\stackrel{}{ϵ}^TM\stackrel{}{ϵ}}{2}.$$
(30)
The equation of motion is simply, $`\dot{\stackrel{}{ϵ}}_0(t)=M\stackrel{}{ϵ}_0(t)`$, and the classical trajectory with arbitrary boundary conditions is given by,
$`\stackrel{}{ϵ}_{class}(t)`$ $`=`$ $`e^{M(tt_i)}\stackrel{}{ϵ}_i+\stackrel{}{\delta }_ϵ\mathrm{\Theta }(tt_f),`$ (31)
$`\dot{\stackrel{}{ϵ}}_{class}(t)`$ $`=`$ $`Me^{M(tt_i)}\stackrel{}{ϵ}_i+\stackrel{}{\delta }_ϵ\delta (tt_f),`$
Defining $`\omega =\frac{qB}{m}`$ we have,
$$e^{M(tt_i)}=\left(\begin{array}{ccc}cos(\omega (tt_i))& sen(\omega (tt_i))& 0\\ sen(\omega (tt_i))& cos(\omega (tt_i))& 0\\ 0& 0& 1\end{array}\right).$$
(32)
In this case the boundary term $`g(\stackrel{}{ϵ}_f,\stackrel{}{ϵ}_i)`$ is nontrivial and takes the form
$$g(\stackrel{}{ϵ}_f,\stackrel{}{ϵ}_i)=\frac{\stackrel{}{ϵ}_i^te^{M(t_ft_i)}\stackrel{}{ϵ}_f}{2}.$$
(33)
According with the discussion of the previous section the we have now,
$$U(\stackrel{}{ϵ}_f,t_f;\stackrel{}{ϵ}_i,t_i)=e^{\stackrel{}{ϵ}_i^te^{M(t_ft_i)}\stackrel{}{ϵ}_f}.$$
(34)
To recover the standard result we compute the time evolution of an arbitrary wave function
$$\varphi (\stackrel{}{ϵ}_f,t_f)=𝑑ϵ_iI_{\stackrel{}{ϵ}_f}e^{\stackrel{}{ϵ}_i^te^{M(t_ft_i)}\stackrel{}{ϵ}_f}\varphi (\stackrel{}{ϵ}_i,t_i)=\varphi (e^{M(t_ft_i)}\stackrel{}{ϵ}_f,t_i).$$
(35)
With the initial state, $`|\varphi _i=cos(\frac{\theta }{2})e^{\frac{i\phi }{2}}|++sen(\frac{\theta }{2})e^{\frac{i\phi }{2}}|`$ and $`\stackrel{}{B}`$ directed in the $`x_3`$ direction, we get
$`\varphi (t)|S_3|\varphi (t)`$ $`=`$ $`cos(\theta ),`$ (36)
$`\varphi (t)|S_1|\varphi (t)`$ $`=`$ $`sen(\theta )cos(\phi +\omega t),`$
$`\varphi (t)|S_2|\varphi (t)`$ $`=`$ $`sen(\theta )sen(\phi +\omega t).`$
## 6 Conclusion
In this paper we have assembled many sparse elements of the the theory of spinning particle already found in the literature, and developed a little some of them, to construct a path integral representation of the quantum amplitudes of a non-relativistic electron in an external electromagnetic field. This fermionic path integral shares the interpretation of a sum over (pseudo) classical histories with its bosonic counterpart. The clue in this approach is to build up the path integral from the explicit realization of the spin operators. The main technical point in the computations concerns the correct handling of the boundary contributions. There are various natural ways to develop further the work presented in this paper. First, one can extend the computational techniques to cases where the spin and the translational degrees of freedom are mixed by the interaction (For example in Ref.). One can also generalize this approach to the relativistic Dirac particle as we discuss elsewhere . Finally the relation between the Grassmannian representation of the spin observables and the fermionic $`SU(2)`$ coherent states may be generalized for other groups.
## 7 Acknowledgments
It is a pleasure to thank the Organizing Committee and all the people involved in the organization of the VI Wigner Symposium for the high level scientific atmosphere and pleasant environment that they created even in adverse conditions. J.S wish also thank J.D.Vergara for very useful discussions on the topics treated in this work |
warning/0001/astro-ph0001414.html | ar5iv | text | # Predicting Future Duration from Present Age: A Critical Assessment11footnote 1This work was supported in part by the U.S. Office of Naval Research (Grant No. N00014-93-1-0116).
(21 January 2000)
## Abstract
Using a temporal version of the Copernican principle, Gott has proposed a statistical predictor of future longevity based on present age \[J. R. Gott III, Nature 363, 315 (1993)\] and applied the predictor to a variety of examples, including the longevity of the human species. Though Gott’s proposal contains a grain of truth, it does not have the universal predictive power that he attributes to it.
Returning from a five-week residence at the Isaac Newton Institute this past summer, I found on my desk the July 21 issue of The New Yorker, containing a provocative story by the well known science writer Timothy Ferris . The story, entitled “How to Predict Everything,” describes how J. Richard Gott, a Princeton astrophysicist, makes universal probabilistic predictions for a phenomenon’s future duration based on knowing how long the phenomenon has lasted. The justification for Gott’s rule is said to be a temporal version of the Copernican principle: when you observe a phenomenon in progress, your observation does not occur at a special time.
Here is Gott’s account, as related to Ferris, of how he conceived his rule while contemplating the Berlin Wall.
> Standing at the Wall in 1969, I made the following argument, using the Copernican principle. I said, Well, there’s nothing special about the timing of my visit. I’m just travelling—you know, Europe on five dollars a day—and I’m observing the Wall because it happens to be here. My visit is random in time. So if I divide the Wall’s total history, from the beginning to the end, into four quarters, and I’m located randomly somewhere in there, there’s a fifty-per-cent chance that I’m in the middle two quarters—that means, not in the first quarter and not in the fourth quarter.
>
> Let’s suppose that I’m at the beginning of that middle fifty per cent. In that case, one quarter of the Wall’s ultimate history has passed and there are three quarters left in the future. In that case, the future’s three times as long as the past. On the other hand, if I’m at the other end, then three quarters have happened already, and there’s one quarter left in the future. In that case, the future is one-third as long as the past.
>
> (The Wall was) eight years (old in 1969). So I said to a friend, “There’s a fifty-per-cent chance that the Wall’s future duration will be between (two and) two-thirds of a year and twenty-four years.” Twenty years later, in 1989, the Wall came down, within those two limits that I had predicted. I thought, Well, you know, maybe I should write this up.
Ferris goes on to recount how Gott applies his method to the longevity of the human species.
> The question that Gott has been asking lately is how long the human species is going to last. Since scientists generally make predictions at the ninety-five-per-cent confidence level, Gott begins with the assumption that you and I, having no reason to think we’ve been born in a special time, are probably living during the middle ninety-five per cent of the ultimate duration of our species. In other words, we’re probably living neither during the first two and a half per cent nor during the last two and a half per cent of all the time that human beings will have existed.
>
> Homo sapiens has been around for two hundred thousand years,” Gott said … . “That’s how long our past is. Two and half per cent is equal to one-fortieth, so the future is probably at least one-thirty-ninth as long as the past but not more than thirty-nine times the past. If we divide two hundred thousand years by thirty-nine, we get about fifty-one hundred years. If we multiply it by thirty-nine, we get 7.8 million years. So if our location in human history is not special, there’s a ninety-five-per-cent chance we’re in the middle ninety-five per cent of it. Therefore the human future is probably going to last longer than fifty-one hundred years but less than 7.8 million years.
>
> “Now, those numbers are interesting, because they give us a total longevity that’s comparable to that of other species.”
These glib predictions astonished me, not because Gott concludes from them that homo sapiens is unlikely to last longer than other species—that is a legitimate subject for inquiry and debate—but because they are put forward as a universal rule, applicable no matter what other information one has about the phenomenon in question. In making statistical predictions of future longevity, Gott dismisses the entire process of assembling and organizing information about a phenomenon, evaluating that information critically, and if possible, formulating laws that describe the phenomenon. Put succinctly, he rejects as irrelevant the process of rational, scientific inquiry, replacing it with a single, universal statistical rule. That has to be wrong.
I decided it was important to find the flaws in Gott’s reasoning: flawed thinking is an inevitable, even necessary part of the scientific enterprise, but when it makes its way into The New Yorker, the time has come to find the flaws and draw attention to them. I began by requesting from the UNM Library a copy of the Nature article where Gott proposes his rule and applies it to the above examples and others. A citation search turned up two other pieces in which Gott adds to the content of his Nature article: a Letter to Nature responding to letters criticizing the original article and a chapter in the proceedings of an Astronomical Society of the Pacific (ASP) Symposium . The present paper analyzes what I found in Gott’s papers and reports my conclusions.
Gott’s delta-t argument
Gott justifies his probabilistic predictions by making what he calls the delta-$`t`$ argument . Suppose there is a phenomenon that has a beginning, or birth, at time $`t_0`$ and an end, or death, at time $`t_0+T`$, $`T`$ being the duration of the phenomenon. You observe the event at a time $`t`$ between the beginning and the end, corresponding to a present age, $`t_p=tt_0`$, and a future duration, $`t_f=Tt_p`$. If there is nothing special about the observation time—this is the content of the temporal Copernican principle—Gott reasons that $`t_p=tt_0`$ is a random variable uniformly distributed between 0 and $`T`$. This means that $`t_p`$ lies between $`aT`$ and $`bT`$, $`0ab1`$, with probability $`baf`$; in symbols, we write
$$P(aT<t_p<bT)=ba=f.$$
(1)
Gott’s next step is to infer from Eq. (1) that the duration $`T`$ lies between the corresponding bounds, $`t_p/b`$ and $`t_p/a`$, with the same probability $`f`$. Translated to future duration, this says that $`t_f`$ lies between $`(b^11)t_p`$ and $`(a^11)t_p`$ with probability $`f`$, i.e.,
$$G\left(\frac{1b}{b}t_p<t_f<\frac{1a}{a}t_p\right)=ba=f.$$
(2)
All of Gott’s predictions flow from this probability rule. I use the letter $`G`$ to distinguish probabilities based on this rule.
Gott phrases his predictions in terms of particular $`f\times 100\%`$ confidence levels, which he obtains by letting $`a`$ and $`b`$ be equidistant from 0 and 1, i.e., $`a=1b`$. The resulting choices, $`a=\frac{1}{2}(1f)`$ and $`b=\frac{1}{2}(1+f)`$, lead to Gott’s confidence-level prediction:
$$\begin{array}{c}\frac{1f}{1+f}t_p<t_f<\frac{1+f}{1f}t_p\\ \text{(}f\times 100\%\text{ confidence level)}\end{array}.$$
(3)
For example, in his encounter with the then $`(t_p\mathrm{=})`$8-year-old Berlin Wall, Gott used $`f=1/2`$, with $`a=1/4`$ and $`b=3/4`$, which led him to predict with 50% confidence that the total duration of the Wall would lie between $`4t_p/3=10\frac{2}{3}`$yr and $`4t_p=32`$yr or, equivalently, that the future duration would lie between $`t_p/3=2\frac{2}{3}`$yr and $`3t_p=24`$yr. In most of his work, Gott uses a 95% confidence level, corresponding to $`f=0.95`$.
Another form of Gott’s rule arises from letting $`b=1`$ and $`a=(1+Y)^1`$. Inserting these choices into Eq. (2), one finds that $`t_f<Yt_p`$ with probability $`Y/(1+Y)`$; equivalently, the probability that the future duration is not less than $`Yt_p`$ is $`(1+Y)^1`$, i.e.,
$$G(t_fYt_p)=\frac{1}{1+Y}.$$
(4)
In his Nature article, Gott derives Eq. (4) independently of the delta-$`t`$ argument by assuming that the phenomenon of interest is an exponential decay \[see Gott’s Eq. (6) and preceding discussion\]. There being no hint in the delta-$`t`$ argument that Gott restricts his method to exponential decays, this derivation must be intended as an example of his method. I defer discussion of this derivation, since its status can be appreciated only after exposing and correcting the flaws in Gott’s reasoning.
The delta-$`t`$ argument implies that Gott’s rule provides a universal method for predicting the future duration of any phenomenon, the only assumption being that the observation time is not special. Moreover, it is clear from the variety of phenomena to which Gott applies his rule—durations of the Berlin Wall, Stonehenge, and the Soviet Union, the publication lifetime of Nature, longevity of the human species, and in his ASP contribution and in his conversations with Ferris, running times of plays in New York—that he places no restrictions on the applicability of his rule.
It is not hard to find an error in the delta-$`t`$ argument: the step from Eq. (1) to Gott’s rule (2) has no justification in probability theory. This error that has been pointed out by Buch, in a Letter to Nature criticizing Gott’s method . The total duration $`T`$ (or the future duration $`t_f`$) is unknown and thus must be treated as a random variable described by a prior probability distribution. This prior distribution expresses whatever information one possesses that can be used to make probabilistic statements about the phenomenon’s duration. After collecting the data that the phenomenon’s present age is $`t_p`$, the only procedure authorized by probability theory is to update the prior distribution to a new, posterior distribution for $`T`$ (or $`t_f`$), which reflects both the prior information and the present age. The formal procedure for this updating is called Bayes’s theorem .
The error just identified is sufficient to invalidate the delta-$`t`$ argument. To correct it requires an analysis that uses Bayes’s theorem to update probabilities. Indeed, Gott has endorsed a Bayesian analysis suggested by Buch ; this Bayesian analysis, said to be based on the temporal Copernican principle, leads to Gott’s rule, provided one uses a particular prior distribution, $`dT/T`$, called the Jeffreys prior . The reader should be aware, however, that the Bayesian analysis suggested by Buch and endorsed by Gott is also flawed. In considering the Buch-Gott Bayesian analysis below, we will uncover this flaw, thus revealing a second error in the delta-$`t`$ argument, just as serious as the first, but more insidious because it is more subtle: Equation (1) is an incorrect mathematical formulation of the temporal Copernican principle. The pay-off for identifying this second flaw is that it clarifies the meaning and status of the temporal Copernican principle. In developing a proper Bayesian analysis based on the temporal Copernican principle, we will discover that Gott’s rule is a universal consequence of the Copernican principle, in the situation where one knows the phenomenon to be in progress, but does not know its present age. Not knowing the present age, one cannot make Gott’s predictions of future duration.
Before turning to the Bayesian analysis, however, I introduce a few examples that show that Gott’s rule cannot be a universal predictor and also serve to put some flesh on the dry bones of the subsequent Bayesian analysis.
Examples of using Gott’s rule
I advise my students to test the solution to a homework problem by considering special cases where the solution is already known. This common-sense technique, a good rule in scientific thinking and in everyday life, provides compelling evidence that Gott’s predictions cannot have the universal validity that he attributes to them.
* Exponential decay. Consider an atom that is excited to a metastable energy level at some unknown time and then decays exponentially to the ground state with a decay constant $`\tau ^1=(20\mathrm{min})^1`$. You come along at time $`t`$ and are told that the atom is in the metastable level, having been excited a time $`t_p=15`$min ago. According to Gott, you can predict with 95% confidence that the decay will occur between $`t_f=t_p/39=23.1`$s and $`t_f=39t_p=9.75`$hr into the future; more telling is that Eq. (4) predicts that $`t_f4t_p=60.0`$min with probability 1/5. These predictions contradict the defining property of an exponential decay: being informed that the atom is still in the excited state at time $`t`$ simply resets the clock so that your expectations for its future decay are the same as though it had been initially excited at time $`t`$. Specifically, you predict that it will survive a further time $`t_f`$ without decaying with probability $`e^{t_f/\tau }`$, corresponding to 95% confidence of decay between $`t_f=0`$ and $`t_f=\tau \mathrm{ln}20=3.00\tau =59.9`$min. Though the numerical discrepancies between Gott’s predictions and the predictions of an exponential decay are important, they are only a symptom of the real problem: Gott’s rule, by including present age in the prediction of future duration, is inconsistent with the very notion of an exponential decay.
Buch has pointed out that Gott’s rule is inconsistent with the properties of an exponential decay. In his reply to Buch and in his ASP contribution , Gott admits that his method doesn’t apply to an exponential decay whose decay constant is known. Instead, he says that it applies to an exponential decay whose decay constant is unknown and distributed according to the Jeffreys prior $`d\tau /\tau `$; this leads to the Jeffreys prior $`dT/T`$ for total duration $`T`$ and is not an exponential decay at all. Gott also reasserts his exponential-decay derivation of Eq. (4), to be discussed below. All this leaves one thoroughly confused—does Gott regard his rule as universal or not?—but his subsequent conversations with Ferris make clear that he does not acknowledge any restrictions on the use of his rule.
* Longevity of an individual. Suppose you are going to a meeting of your book club, to be held at a member’s house that you’ve never been to before. You find the right street, but having forgotten the street address, you choose between two houses where there is evident activity. Knocking at one, you are told that the activity within is a birthday party, not a book-club meeting. Your friendly enquiry about the age of the celebrant elicits the reply that she is celebrating her $`(t_p\mathrm{=})`$50th birthday. According to Gott, you can predict with 95% confidence that the woman will survive between $`t_p/39=1.28`$years and $`39t_p=1,950`$years into the future. Since the wide range encompasses reasonable expectations regarding the woman’s survival, it might not seem so bad, till one realizes that Eq. (4) predicts that with probability 1/2 the woman will survive beyond 100 years old and with probability 1/3 beyond 150. Few of us would want to bet on the woman’s survival using Gott’s rule.
One might object at this point that Gott probably didn’t intend his rule to apply to an individual’s longevity, but in his ASP contribution , Gott applies the rule to himself: “At the time my (Nature) paper was published on May 27, 1993, I was 46.3 years old, so the 95% delta-$`t`$ argument predicted that my future longevity would be at least 1.2 years but less than 1,806 years. I have survived past the lower limit already and so if I don’t make it past the upper limit, then that prediction will indeed prove correct for me!”
* Deterministic phenomena. The best testing ground for ideas comes from extreme cases, and here the most extreme case is a deterministic phenomenon. Putting the example in a dramatic context, suppose you are captured by terrorists, who confine you to a small room. You are told that at some time in the next 24 hours, a timer will be set and that after it has ticked for 30 minutes, poison gas will fill the room, killing you. You are then drugged and wake up to find the timer ticking and reading 20 minutes since being set. According to Gott, you can predict with 95% confidence that the time to release of the gas lies between $`t_f=20\mathrm{min}/39=30.8`$s and $`t_f=39\times 20\mathrm{min}=13.0`$hr; even worse, Eq. (4) predicts that with probability 2/3 the time to release is $`10`$min or more. These reassuring predictions provide scant comfort, since you know you have $`10`$min to live.
These examples demonstrate that Gott’s rule cannot be a universal method for predicting future durations. If the rule has any validity, it must involve other information than the present age of a phenomenon. As E. T. Jaynes taught us , when probabilistic predictions violate one’s intuition, the proper response is neither to accept the nonintuitive predictions without question nor to dismiss them out of hand, but rather to identify the information underlying the prediction. You will either find the information inapplicable to the situation at hand, thereby confirming your intuition and allowing you to discard the predictions, or you will sharpen your intuition.
The objective of this paper is to identify the prior information that underlies Gott’s rule. The tool is Bayesian analysis. We will discover that the temporal Copernican principle contains a grain of truth, but that grain of truth does not include Gott’s predictions of future duration.
A flawed, but instructive Bayesian analysis
Return to the general situation introduced above, that of a phenomenon with a birth time $`t_0`$ and a duration $`T`$. You observe the phenomenon at time $`t`$. It is often useful to replace one or both quantities, $`t_0`$ and $`T`$, by the present age, $`t_p=tt_0`$, and the future duration, $`t_f=Tt_p`$. In developing the Bayesian analysis, I first formulate and analyze a flawed approach, advanced by Buch and endorsed by Gott , which is modeled on Gott’s delta-$`t`$ argument. For this purpose, it is most convenient to use $`t_p`$ and $`T`$ as the primary variables. The reason for going through this flawed analysis is that it turns up the second error in Gott’s delta-$`t`$ argument.
Your prior information about the phenomenon is expressed in a prior probability $`p(t_p,T)dt_pdT`$, the joint probability that the phenomenon has lasted a time between $`t_p`$ and $`t_p+dt_p`$ at the time of observation and that the phenomenon will last a total time between $`T`$ and $`T+dT`$. The joint probability density can be written as $`p(t_p,T)=p(t_p|T)w(T)`$, where $`p(t_p|T)`$ is the conditional probability density for the present age, given a total duration $`T`$, and $`w(T)`$ is your prior probability density for the total duration. Throughout I use upper-case letters for probabilities and lower-case letters for probability densities.
Before going further, it is useful to introduce two quantities related to $`w(T)`$: $`\lambda (T)`$ is the death rate—i.e., $`\lambda (T)dT`$ is the probability that the phenomenon, having lasted a time $`T`$, ends in the next $`dT`$—and $`Q(T)`$ is the survival probability—the probability that the phenomenon lasts at least a time $`T`$. These quantities are related by
$$w(T)=\frac{dQ}{dT}=Q(T)\lambda (T)$$
(5)
or, equivalently, by
$$Q(T)=_T^{\mathrm{}}𝑑T^{}w(T^{})=\mathrm{exp}\left(_0^T𝑑T^{}\lambda (T^{})\right).$$
(6)
An exponential decay is characterized by a constant death rate, $`\lambda (T)=\lambda _0`$, in which case $`Q(T)=e^{\lambda _0T}`$ and $`w(T)=\lambda _0e^{\lambda _0T}`$.
Gott’s formulation of the temporal Copernican principle is the following: if there is nothing special about the observation time, the present age is a random variable uniformly distributed between 0 and $`T`$, i.e.,
$$p(t_p|T)dt_p=dt_p/T,0t_pT.$$
(7)
This is the probability-density version of Eq. (1). We now use Bayes’s theorem,
$$P(X|Y)P(Y)=P(X,Y)=P(Y|X)P(X),$$
(8)
to find your posterior probability density for the total duration, given the present age:
$$p(T|t_p)=\frac{p(t_p|T)w(T)}{p(t_p)}=\{\begin{array}{cc}0,\hfill & T<t_p\text{,}\hfill \\ w(T)/Tp(t_p),\hfill & Tt_p\text{.}\hfill \end{array}$$
(9)
The unconditional probability density $`p(t_p)`$ for the present age, which is a normalization constant in this expression, is given by
$$p(t_p)=_{t_p}^{\mathrm{}}𝑑T\frac{w(T)}{T}.$$
(10)
One can easily verify that Gott’s rule, embodied in Eqs. (2)–(4), is equivalent to a posterior density
$$p(T|t_p)=\{\begin{array}{cc}0,\hfill & T<t_p\text{,}\hfill \\ t_p/T^2,\hfill & Tt_p\text{.}\hfill \end{array}$$
(11)
To get this posterior from the present analysis, one must assume the (unnormalizable) Jeffreys prior,
$$w(T)=\frac{1}{T}.$$
(12)
Buch concludes that Gott’s rule is unreasonable because it corresponds to an unnormalizable prior density. Gott replies (correctly, I think) that there is nothing wrong with an unnormalizable prior, since the posterior density for $`T`$ can be normalized. He defends the Jeffreys prior as being the appropriate “vague Bayesian prior” to use in a situation where one initially knows nothing about the magnitude of the duration .
Jaynes has delineated the conditions for using the Jeffreys prior, showing that it should be used when one’s prior information is unchanged by a rescaling of the total duration, $`T^{}=\alpha T`$ ($`\alpha >0`$). If one’s prior information is unchanged by the rescaling, then the density for $`T^{}`$, $`w^{}(T^{})=w(T)dT/dT^{}=w(T)/\alpha `$, should have the same functional form as the original density, i.e., $`w^{}(T^{})=w(T^{})`$. This gives $`w(\alpha T)=w(T)/\alpha `$, which implies that $`w(T)1/T`$. This might seem to be progress in identifying the information that underlies Gott’s rule—use it when one has no prior information about time scales associated with the phenomenon—but it turns out not to be, because the present Bayesian analysis is wrong. The reason for presenting it is not to consider its consequences, but to identify where it goes wrong.
A straightforward Bayesian analysis
That something is wrong is made apparent by a different analysis of the same situation, this time a straightforward Bayesian analysis that does not invoke the Copernican principle. Your prior information about the total duration is expressed in the prior density $`w(T)`$. You observe the phenomenon still to be in progress a time $`t_p`$ after its beginning. The conditional probability for this observation, given a total duration $`T`$, is 0 if $`t_p>T`$ and 1 if $`t_pT`$. Thus Bayes’s theorem implies, with $`O`$ denoting the observation,
$$p(T|O)=\frac{P(O|T)w(T)}{P(O)}=\{\begin{array}{cc}0,\hfill & T<t_p\text{,}\hfill \\ w(T)/Q(t_p),\hfill & Tt_p\text{.}\hfill \end{array}$$
(13)
Here the normalization constant is the survival probability, i.e., $`P(O)=Q(t_p)`$.
The posterior density (13) is so eminently reasonable that one could have written it down without using the formal apparatus of Bayes’s theorem. It says that the effect of discovering the present age is to rule out durations shorter than the present age; your posterior expectations for durations longer than the present age are the same as your prior expectations, with appropriate renormalization. Notice that this inference updates sensibly: subsequent observations that find the phenomenon still in progress simply exclude a wider interval of durations. Yet putting this simple inference in the context of the Copernican principle apparently yields a different posterior density (9) for the total duration. How can that be? There’s nothing wrong with the Bayesian inference in either analysis, so the culprit must be Gott’s formulation of the temporal Copernican principle. Thus we arrive at the second error in the delta-$`t`$ argument: the uniform density (7) for $`t_p`$—and, by extension, Eq. (1)—is not the correct mathematical formulation of the temporal Copernican principle.
Where the uniform density goes wrong is in assuming that your observation occurs while the phenomenon is in progress. If your observation does not occur at a special time, then it is very likely that it occurs before the phenomenon begins or after it has ended. Including these other possibilities leads to a proper Bayesian formulation of the temporal Copernican principle, which is consistent with the inference expressed in Eq. (13).
A proper Bayesian analysis of the temporal Copernican principle
In formulating a proper Bayesian analysis, it is convenient to choose the birth time $`t_0`$ and the total duration $`T`$ as the primary variables. Your prior knowledge about these two quantities is incorporated in two probability densities: (i) $`\gamma (t_0)`$ gives the probability $`\gamma (t_0)dt_0`$ that the phenomenon begins between times $`t_0`$ and $`t_0+dt_0`$; (ii) $`p(T|t_0)`$ gives the probability $`p(T|t_0)dT`$ that the phenomenon lasts a time between $`T`$ and $`T+dT`$, given that it began at time $`t_0`$. The corresponding joint probability density is $`p(t_0,T)=p(T|t_0)\gamma (t_0)`$.
The temporal Copernican principle—that your observation does not take place at a special time—is a time-translation symmetry that restricts the form of the prior densities . To say that your observation time is not special is to say that your prior information is unchanged if the entire phenomenon is displaced in time while your observation time remains fixed. To be consistent with this translation symmetry, your prior probability density should be unchanged by such a time translation; i.e., $`p(t_0,T)`$ should be independent of the birth time $`t_0`$. Thus the temporal Copernican principle can be captured precisely in the following two statements:
1. The phenomenon is equally likely to begin at any time. This means that $`\gamma (t_0)`$ is a constant. In order to work with normalizable probabilities, I replace the exact symmetry with the approximate one that $`\gamma (t_0)`$ has a constant value, $`1/\mathrm{\Delta }`$, at all times within a very long time interval. The duration $`\mathrm{\Delta }`$ of this very long time interval exceeds all other times relevant to the problem, particularly typical durations.
2. Probabilities for total duration are independent of birth time. This means that the conditional probability density $`p(T|t_0)`$ does not depend on $`t_0`$ and can be written as $`p(T|t_0)=w(T)`$, where $`w(T)`$ is the probability density introduced above.
Should you be dissatisfied with these restrictions on the prior probabilities, it means that you do not accept the temporal Copernican principle as applying to your prior information. Dissatisfaction should not be surprising, for one would not expect the Copernican principle to apply to all situations. The three examples introduced above illustrate considerations that arise in using the temporal Copernican principle. In all three examples, it is easy to accept that ignorance of the birth time is described by the time-translation symmetry of the temporal Copernican principle: the atom can be excited at any time during an interval much longer than the decay time; for the woman at the birthday party, the situation could be phrased in terms of an individual whose birth could occur at any time over a period much longer than a typical human lifetime; the timer can be set at any time within a 24-hour period, a period somewhat longer than the 30 minutes that the timer ticks. Moreover, in the cases of the atom and the poison gas, duration probabilities are independent of the birth time. In contrast, in the case of the longevity of an individual, the prior conditional probability for the individual’s lifetime would depend on the time of birth. Your prior expectation for the longevity of an individual born, say in Britain, would depend on whether the individual was born in the second half of the 20th Century, at the beginning of the 19th Century, or 10,000 years ago, at the end of the last Ice Age.
At time $`t`$ you make your observation. In Gott’s formulation the observation yields the present age, but we now understand that getting the present age presupposes that your observation finds the phenomenon in progress. The first result of the observation is simply to determine whether the phenomenon has not yet begun, is already over, or is in progress. Only the last of these possibilities, denoted by $`I`$ for “in progress,” is of interest to us. The conditional probability to find the phenomenon in progress, given a birth time $`t_0`$ and a duration $`T`$, is
$$P(I|t_0,T)=\{\begin{array}{cc}1,\hfill & t_0tt_0+T\text{,}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
(14)
The unconditional probability to find the phenomenon in progress is given by
$`P(I)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t_0{\displaystyle _0^{\mathrm{}}}𝑑TP(I|t_0,T)\gamma (t_0)w(T)`$
$`=`$ $`{\displaystyle _{\mathrm{}}^t}𝑑t_0\gamma (t_0)\underset{=Q\left(tt_0\right)}{\underset{}{{\displaystyle _{tt_0}^{\mathrm{}}}𝑑Tw(T)}}`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑t^{}\gamma (tt^{})Q(t^{}).`$
The assumption that $`\gamma (t_0)`$ is constant for all times of interest means that $`\gamma (tt^{})=1/\mathrm{\Delta }`$ for all times $`t^{}`$ such that the survival probability $`Q(t^{})`$ is significantly different from zero. This allows us to put $`P(I)`$ in the form
$$P(I)=\overline{T}/\mathrm{\Delta },$$
(16)
where
$$\overline{T}=_0^{\mathrm{}}𝑑TTw(T)=_0^{\mathrm{}}𝑑TQ(T)$$
(17)
is the mean total duration with respect to the prior density $`w(T)`$. The present analysis assumes that $`\overline{T}`$ is finite, which requires, for large durations $`T`$, that $`Q(T)`$ go to zero faster than $`1/T`$ or, equivalently, that $`w(T)`$ go to zero faster than $`1/T^2`$. For an exponential decay, $`\overline{T}^1=\lambda _0`$ is the decay constant. Notice that the probability to find the phenomenon in progress is very small.
Bayes’s theorem gives the posterior probability density for $`t_0`$ and $`T`$, given that the phenomenon is occurring:
$$p(t_0,T|I)=\frac{P(I|t_0,T)\gamma (t_0)w(T)}{P(I)}=\{\begin{array}{cc}\gamma (t_0)w(T)/P(I),\hfill & t_0tt_0+T\text{,}\hfill \\ 0,\hfill & t_0>t\text{ or }t_0+T<t\text{.}\hfill \end{array}$$
(18)
If $`tt_0`$ is large enough in this expression that $`\gamma (t_0)`$ does not have its constant value, then $`Ttt_0`$ is so large that $`w(T)`$ is negligible. Thus we can again replace $`\gamma (t_0)`$ by the constant value $`1/\mathrm{\Delta }`$, leaving
$$p(t_0,T|I)=\{\begin{array}{cc}w(T)/\overline{T},\hfill & t_0tt_0+T\text{,}\hfill \\ 0,\hfill & t_0>t\text{ or }t_0+T<t\text{.}\hfill \end{array}$$
(19)
It is instructive to consider Eq. (19) from a variety of perspectives. A first question asks how the probability density for total duration changes on learning that the phenomenon is occurring:
$$p(T|I)=_{\mathrm{}}^{\mathrm{}}𝑑t_0p(t_0,T|I)=\frac{Tw(T)}{\overline{T}}.$$
(20)
Notice that $`p(T|I)`$ is biased toward longer durations than the prior density $`w(T)`$. This is because the phenomenon is very unlikely to be in progress at a random time selected from the long time interval $`\mathrm{\Delta }`$, so finding it in progress prejudices you to think that it has a longer duration than your original expectations.
A useful, equivalent form for Eq. (19) comes from changing variables to present age and future duration. The Jacobian of the transformation from $`(t_0,T)`$ to $`(t_p,t_f)`$ is $`1`$, which implies that $`dt_0dT=dt_pdt_f`$. Hence, the probability density for present age and future duration, given that the phenomenon is occurring, is
$$p(t_p,t_f|I)=p(t_0,T|I)=\{\begin{array}{cc}w(t_p+t_f)/\overline{T},\hfill & t_p0\text{ and }t_f0\text{,}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}$$
(21)
Knowing the phenomenon is in progress is equivalent to saying that both the present age and future duration are nonnegative, so we can regard that condition as implicit and omit it from subsequent expressions. The content of Eq. (21) is the following: if you know the phenomenon is in progress, but don’t know its present age, you treat uniformly the split of total duration into past and future; more precisely, you assign the same probability, governed by $`w(T)`$, to all ways of splitting $`T`$ into past and future. That’s the temporal Copernican principle. Indeed, Eq. (21) is the mathematical embodiment of the temporal Copernican principle for phenomena known to be in progress.
Equation (21) has three immediate consequences that highlight the connection between the Copernican principle and Gott’s rule. We proceed by noting that once the phenomenon is known to be in progress, the total duration is the sum of the present age and the future duration, i.e., $`p(T|t_p,t_f,I)=\delta (Tt_pt_f)`$. Another application of Bayes’s theorem then gives
$$p(t_p,t_f|T,I)=\frac{p(T|t_p,t_f,I)p(t_p,t_f|I)}{p(T|I)}=\frac{1}{T}\delta (Tt_pt_f).$$
(22)
This is a conditional version of Eq. (21), with the same content.
An obvious consequence is that if you know the phenomenon is in progress and also know its total duration, then you conclude that the present age is uniformly distributed between 0 and $`T`$:
$$p(t_p|T,I)=_0^{\mathrm{}}𝑑t_fp(t_p,t_f|T,I)=\{\begin{array}{cc}1/T,\hfill & t_pT\text{,}\hfill \\ 0,\hfill & t_p>T\text{.}\hfill \end{array}$$
(23)
This is the precise statement of what Gott is trying to capture in his initial assumption (1) about the present age. The starting point (7) of the flawed Bayesian analysis also asserts that $`t_p`$ is uniformly distributed between 0 and $`T`$, but it is different from Eq. (23) in a subtle, but crucial way: because $`p(t_p|T,I)`$ is conditioned on knowing the phenomenon is occurring, further statistical inference uses the conditional density $`p(T|I)=Tw(T)/\overline{T}`$, instead of the prior density $`w(T)`$; we find below \[see Eq. (28)\] that this is how the present Bayesian analysis comes into agreement with the straightforward inference of the preceding section.
A second obvious, but important consequence of Eq. (21) is that
$$P\left(\frac{1b}{b}t_p<t_f<\frac{1a}{a}t_p|T,I\right)=_0^T𝑑t_p_{(b^11)t_p}^{(a^11)t_p}𝑑t_fp(t_p,t_f|T,I)=_{aT}^{bT}\frac{dt_p}{T}=ba.$$
(24)
Since the condition on future duration in the probability on the left is equivalent to $`aT<t_p<bT`$, this is just the statement that in dividing the total duration into past and future, possibilities satisfying the condition are a fraction $`ba`$ of all the possibilities. Furthermore, since the conditional probability (24) is independent of $`T`$, the same result holds no matter what the prior density for $`T`$:
$$P\left(\frac{1b}{b}t_p<t_f<\frac{1a}{a}t_p|I\right)=_0^{\mathrm{}}𝑑TP\left(\frac{1b}{b}t_p<t_f<\frac{1a}{a}t_p|T,I\right)p(T|I)=ba.$$
(25)
Setting $`b=1`$ and $`a=(1+Y)^1`$ in this result yields the third consequence of Eq. (21):
$$P(t_fYt_p|I)=\frac{1}{1+Y}.$$
(26)
Equations (25) and (26) are precise statements of Gott’s rule in the forms (2) and (4). Indeed, they look just like Gott’s rule, with the crucial difference that they are conditioned on knowing the phenomenon is in progress, without knowing its present age.
We are now in the curious position of affirming that for a phenomenon known to be in progress, but whose present age is unknown, the temporal Copernican principle leads to universal statistical predictions, which are described by Gott’s rule. Indeed, all the manipulations in Gott’s delta-$`t`$ argument are valid in this situation. The down side for Gott is that this conclusion does not authorize his predictions: in these circumstances, Gott’s rule has no power to predict future durations from present ages, for the simple reason that the present age is unknown.
The results of the present Bayesian analysis make perfect sense in the three examples introduced above. Oddly enough, the deterministic example is the simplest: if you find the timer ticking, but there is nothing to indicate how long it has been ticking, it is reasonable to assign the same probability to all ways of dividing the 30-minute interval into past and future.
The case of the women’s longevity is a bit more complicated. If you encounter an individual, but are given no clue as to the individual’s age, a first cut might treat the past-future split uniformly. For a person born in Britain, a more careful analysis would give greater weight to the future than to the past, because of the increase in life expectancy in this century. I have already indicated that the case of an individual’s longevity does not fit into the temporal Copernican principle for just this reason. It should be emphasized that the problem is not the use of Bayesian analysis: the increase in life expectancy could be incorporated into a more complicated Bayesian analysis, which would automatically produce a bias toward the future.
The case of the atom is particularly interesting because of Gott’s exponential-decay derivation of Eq. (4). If you find the atom in the excited state, but you are not told when it was excited, it is reasonable to assign the same probability to all ways of splitting a particular duration $`T`$ into past and future and to weight the result by an exponential $`e^{\lambda _0T}=e^{\lambda _0(t_p+t_f)}`$, which expresses the probability for duration $`T`$. The final result, properly normalized, is the probability density (21) specialized to an exponential decay:
$$p(t_p,t_f|I)=\lambda _0^2e^{\lambda _0(t_p+t_f)}.$$
(27)
This allows us to understand Gott’s exponential-decay derivation of Eq. (4) \[see Gott’s Eq. (6)\]: he starts with Eq. (27) \[see Gott’s Eqs. (3) and (4)\], from which he immediately derives Eq. (26), all without realizing that Eq. (27) applies to an exponential decay whose present age is unknown. Having gone through a proper Bayesian analysis, we now understand that Eq. (26) does not depend at all on assuming an exponential decay, but rather is a universal consequence of the temporal Copernican principle, valid no matter what the prior density $`w(T)`$, provided the present age is unknown.
The next task is to find out what happens if you do discover the present age. When you determine the present age of the phenomenon, your Bayesian posterior for the total duration is given by
$$p(T|t_p,I)=\frac{p(t_p|T,I)p(T|I)}{p(t_p|I)}=\{\begin{array}{cc}0,\hfill & T<t_p\text{,}\hfill \\ w(T)/Q(t_p),\hfill & Tt_p\text{,}\hfill \end{array}$$
(28)
where
$$p(t_p|I)=_0^{\mathrm{}}𝑑Tp(t_p|T,I)p(T|I)=Q(t_p)/\overline{T}.$$
(29)
This posterior density is identical to the one that emerged from the straightforward Bayesian analysis that wholly ignored the Copernican principle. This is as it should be, because in the language of this section, the straightforward Bayesian inference corresponds to first learning the birth time $`t_0`$ and then discovering that the phenomenon has survived a time $`t_p`$, a situation that is equivalent to first learning that the phenomenon is in progress and then discovering its present age. Once you are informed of the present age or, equivalently, of the birth time, you are at a special time, the time $`t_p`$ since the phenomenon began. The temporal Copernican principle becomes irrelevant. It just gets in the way of the obvious inference expressed in Eq. (13).
At this point it is profitable to re-read Gott’s account of his 1969 encounter with the Berlin Wall. If Gott had not known when the Wall was built, the logic of the first two paragraphs of his account would be impeccable. Under those circumstances, it would be reasonable to assign probability 1/2 to the encounter’s occurring during the middle two quarters of the Wall’s total history. Since he did know that the Wall was built in 1961, however, his encounter did occur at a special time, the time eight years after the Wall’s construction. The predictions made in the third paragraph of his account do not follow from the argument in the first two paragraphs. Indeed, his posterior expectations for the Wall’s duration should have been a renormalized version of his prior expectations, whatever those were, with durations up to eight years excluded.
We can now give a succinct account of how Gott’s delta-$`t`$ argument goes awry: the first two steps are wrong. The step from Eq. (1) to Gott’s rule (2) is a non-Bayesian inference having no justification in probability theory; just as important, Eq. (1) is itself an incorrect expression of the temporal Copernican principle, because it assumes that an observation at a random time will find the very unlikely result that the phenomenon is in progress. In repairing these errors, we discovered that Gott’s rule for relating future duration to present age is indeed a universal consequence of the temporal Copernican principle, but only in a situation—not knowing the present age—which leaves the rule shorn of predictive power. Gott’s predictions require knowing how long a phenomenon has lasted, but once you obtain this information, the temporal Copernican principle no longer has any impact, because you are at a special time within the lifetime of the phenomenon.
Gott’s rule as a predictor
All of Gott’s predictions—from the future duration of the Berlin Wall to the longevity of the human species—are now detached from their original mooring in the temporal Copernican principle and left to float free of justification. Yet a flawed analysis might lead to reasonable predictions. There might be some justification for Gott’s predictions other than the Copernican principle. Both the straightforward Bayesian analysis and the analysis based on the Copernican principle culminate in the same inference \[Eqs. (13) and (28)\]: once you know the present age, your expectations about total duration are the same as your prior expectations, except that durations shorter than the present age are excluded. Thus all questions about the applicability of Gott’s predictions reduce to determining what prior density underlies his predictions.
As noted above, Gott’s rule follows from a posterior density
$$p(T|O)=\{\begin{array}{cc}0,\hfill & T<t_p\text{,}\hfill \\ t_p/T^2,\hfill & Tt_p\text{.}\hfill \end{array}$$
(30)
Within the correct Bayesian analysis, this posterior comes from an unnormalizable prior density
$$w_g(T)=\frac{1}{T^2}.$$
(31)
This prior density, distinguished by a subscript $`g`$, corresponds to a survival probability $`Q_g(T)=1/T`$ and to a death rate $`\lambda _g(T)=1/T`$. One way to characterize $`w_g(T)`$ is that the characteristic time associated with the death rate, $`\lambda _g^1(T)`$, is always the same as the age $`T`$.
The prior density $`w_g(T)`$ is different from the Jeffreys prior that Gott identifies with his predictions, the reason being that Gott uses the flawed Bayesian analysis given above. Yet within the Bayesian analysis using the temporal Copernican principle, $`w_g(T)`$ has a scale-free status similar to that found by Jaynes for the Jeffreys prior. Suppose that once you know the phenomenon is in progress, anything else you know, coming from the prior information about $`T`$, is unchanged by a simultaneous change in the scale of the past and the future. Under such a scale change, $`t_p^{}=\alpha t_p`$ and $`t_f^{}=\alpha t_f`$, the new and old probability densities are related by
$$p^{}(t_p^{},t_f^{}|I)=p(t_p,t_f|I)dt_pdt_f/dt_p^{}dt_f^{}=p(t_p,t_f|I)/\alpha ^2.$$
(32)
To say that all your information is unchanged by this scale change is to say that the old and new densities should have the same functional form, i.e., $`p^{}(t_p^{},t_f^{})=p(t_p^{},t_f^{})`$, which implies that $`p(\alpha t_p,\alpha t_f|I)=p(t_p,t_f|I)/\alpha ^2`$. Using Eq. (21) to write this in terms of the prior density, one finds that $`w(\alpha T)=w(T)/\alpha ^2`$, which implies that the prior density has the form (31).
As discussed above, the Jeffreys prior applies when your prior information about the duration, before any observation, is scale-invariant. Once you know the phenomenon is in progress, however, $`w_g(T)`$ captures the notion of scale invariance, because it corresponds to invariance of $`p(t_p,t_f|I)`$ under simultaneous rescaling of the past and future. In contrast, the Jeffreys prior corresponds to invariance of $`p(t_p,t_f|I)`$ under rescaling of $`t_p`$ or $`t_f`$, but not both simultaneously.
We have now uncovered the prior information that underlies the use of Gott’s rule as a predictor of future duration; namely, knowing that a phenomenon is in progress, you cannot identify any time scales associated with the phenomenon either into the past or into the future. One way of thinking about this is that for a phenomenon that has no time scales, discovering the present age does not put you at a special time in the phenomenon’s history, so some consequences of the temporal Copernican principle survive. Whether the scale-free prior information is appropriate must be judged case by case; it is not a universal rule. The scale-free prior certainly does not apply to the three examples introduced in this article, each of which has an obvious time scale: for the atom, the scale is the decay time; for an individual, the scale is a typical human lifetime; for the deterministic phenomenon, the scale is the 30 minutes that the timer ticks. Ignoring these time scales is the reason that Gott’s rule leads to absurd predictions for these examples.
The examples Gott discusses at the beginning of his Nature article all have readily identifiable time scales that make application of Gott’s rule problematic. The survival of a human institution—a political institution such as the government of the former Soviet Union or a cultural institution such as a periodical like Nature—is influenced by the 30-year time scale of a generation or by a typical human lifetime, since loyalty to and management of such institutions change on these time scales. Physical manifestations of human institutions, such as the Berlin Wall or Stonehenge, are influenced by these same human time scales and, in addition, by the time scale over which erosion leads to disintegration.
The success of Gott’s rule
Even though there is little reason to adopt Gott’s rule, he portrays his predictions as successful . Consider, for example, his 95%-confidence prediction that Nature, given its 123-year history of publication in 1993, would continue to publish for a period between 3.15 years and 4,800 years. Gott would consider this prediction successful because Nature has already surpassed the lower bound and is very unlikely to exceed the upper bound. Yet there’s the hitch: the upper bound is far too large; without doing any analysis, anyone could have written down a similar very large 95% confidence interval and achieved the same “success.” To assess Gott’s rule, one should direct attention not at the the 95% confidence predictions, but at the high probabilities the rule assigns to very long future durations. Gott’s rule in the form (4) predicts that with probability 1/2, Nature will continue to publish for more than 123 years after 1993, with probability 1/5 for more than 492 years, with probability 1/10 for more than 1,107 years, and with probability 1/20 for more than 2,337 years. These probabilities posit a great deal of faith in the durability of human institutions.
To make this point more quantitatively, it is useful to consider a particular form of the probability that future duration exceeds some multiple of present age:
$`P(t_fYt_p|O)`$ $`=`$ $`P(T(1+Y)t_p|O)`$ (33)
$`=`$ $`{\displaystyle _{(1+Y)t_p}^{\mathrm{}}}𝑑Tp(T|O)={\displaystyle \frac{Q[(1+Y)t_p]}{Q(t_p)}}=\mathrm{exp}\left({\displaystyle _{t_p}^{(1+Y)t_p}}𝑑T\lambda (T)\right).`$
This form makes clear that $`P(t_fYt_p|O)`$ depends only on the death rate during the interval between the present age and the lower bound for longevity. For a death rate $`\lambda _g(T)=1/T`$, one gets Gott’s rule.
Now let’s apply this to the example of a periodical like Nature. At start-up a new publication confronts a variety of short-term, rapid-death scenarios. Should it survive these initial hazards and become established like Nature, the next time scale it faces might be roughly a human lifetime. If this time scale is modeled by a constant death rate $`\tau ^1=(60\mathrm{yr})^1`$, then one finds from Eq. (33) that $`P(t_fYt_p|O)=e^{Yt_p/\tau }`$. For Nature, this gives predictions quite different from Gott’s: for example, a probability $`0.129`$ to continue publishing for more than 123 years beyond 1993 and a probability $`2.75\times 10^4`$ to continue publishing for more than 492 years.
Should these predictions seem unduly pessimistic, it is because the constant decay rate does not recognize a long publication record as providing evidence for continued success. A prejudice that success begets success can be incorporated, without discarding the time scale, by choosing, for example,
$$\lambda (T)=\frac{1}{T}\left(\frac{T}{\tau }\right)^\beta ,$$
(34)
where $`\beta 0`$. For $`\beta =1`$, this gives a constant death rate $`\tau ^1`$, and for $`\beta =0`$, it gives Gott’s rule. For intermediate values, it gives a death rate that decreases with age, but with the time scale $`\tau `$ still having an effect. The resulting probability (33) is
$$P(t_fYt_p|O)=\mathrm{exp}\left[\left(\frac{t_p}{\tau }\right)^\beta \frac{(1+Y)^\beta 1}{\beta }\right],$$
(35)
For Nature this gives, assuming $`\beta =1/2`$, a probability $`0.305`$ to continue publishing for more than 123 years beyond 1993, a probability $`2.90\times 10^2`$ for more than 492 years, and a probability $`2.05\times 10^3`$ for more than 1,107 years. The point here is not the particular values nor even the death-rate model, but rather that there is one or more time scales, which can and should be incorporated in the prior distribution.
Gott stresses the success of his predictions for the 44 Broadway and off-Broadway plays listed in The New Yorker on 27 May 1993, the day his original Nature article was published. For example, Gott’s 95%-confidence rule predicted that Cats, having played for 3,885 days, would continue to play for a period between 100 days and 415 years. Gott regards this prediction as a success because the production continues today, thereby surpassing the lower bound, and is unlikely to exceed the upper bound . Yet since Cats had run 6,263 days through 30 November 1999, when I determined that it was still running, the same rule predicts that with probability 1/5, it will continue to run for at least another 68.6 years, with probability 1/10 for at least another 154 years, and with probability 1/20 for another 326 years. Such predictions ignore obvious time scales. A new production faces a variety of short- to medium-term scales, including the time to the first reviews, the time over which a producer is willing to back a losing production, the annual cycle of openings and closings, and the time over which a star performer tires of a particular part and moves on to other challenges. An established production like Cats, having survived these initial hurdles, must deal with the decade- to generation-long scale over which taste and fashion change substantially and the production experiences a nearly 100% turnover of personnel. Including this long-term scale would temper Gott’s predictions for extraordinarily long runs.
The problems with Gott’s long-term predictions show up more dramatically in phenomena, such as the longevity of an individual, where an initial period of low death rate is followed by relatively rapid extinction. We don’t need a detailed model to tell us whether we should believe Gott’s prediction, based on his age of 46.3 years on 27 May 1993, that he has a 1/3 chance to survive to more than 139 years old.
There are two reasons, in my view, why Gott is able to get away with making his scale-free predictions for the survival of governments and plays and periodicals. First, statistical models for the longevity of these phenomena are not well developed, so Gott is protected from the absurdities that arise immediately in the three examples used in this article. Although there are readily identifiable time scales associated with the phenomena Gott considers, how to incorporate them into prior probabilities for duration has not been much investigated. There’s a good reason for this: to assess the viability of an established government or play or periodical, readily available current data about the particular phenomenon in question—data such as the popularity of the government, the balance sheet of the play or periodical, trends in attendance at the play or the number of subscribers of the periodical—are far more cogent than prior information about longevity together with the present age. Second, the intervals that Gott finds for survival times are so wide that he is likely to be right, till he is forced to place bets based on the high probabilities he assigns to long survival times. A negative feature of such bets, however, is that the bettors might not survive till the bets are settled. Even for the case of human longevity, where one could easily formulate bets that Gott would almost certainly lose, the time scales are long enough that one might not get much personal satisfaction from winning.
A way to overcome this difficulty is to bet on the survival of creatures with a shorter lifetime than humans, but for which data on present age and future survival are readily obtainable. For this purpose, I sent an e-mail on 21 October 1999 and again on 2 December 1999, to my department’s most comprehensive e-mail alias, which includes faculty, staff, and graduate students, requesting information on pet dogs. The responses were compiled and checked for accuracy on 6 December; a notarized list of the 24 dogs, including each dog’s name, date of birth, and breed, and the caretaker’s name, was deposited in my departmental personnel file on 21 December 1999. Gott’s rule predicts that each dog will survive to twice its present age with probability 1/2. For each of the 6 dogs above 10 years old on the list, I am offering to bet Gott $1,000 US, at odds of 2:1 in his favor, that the dog will not survive to twice its age on 3 December 1999. The reason for weighting the odds in Gott’s favor is to test his belief in his own predictions: given the odds, his rule says that his expected gain, at $1,000 per bet, is $6,000; moreover, the probability that he will be a net loser (by losing five or more of the six bets) is 7/64=0.109.
Discussion
The stated objective of this article is to determine what prior information underlies Gott’s rule. Gott proposed his rule as a predictor of future duration based on knowing the present age and nothing else. What we have discovered is that the actual prior information underlying Gott’s rule is both less and more than he thought. On the one hand, Gott’s rule is a consequence of the temporal Copernican principle for a phenomenon whose age is unknown, but this universal form of Gott’s rule has no predictive power for future durations. On the other hand, Gott’s rule as a predictor of future durations is a consequence of discovering the present age of a phenomenon that has no identifiable time scales in the past or future.
What about the focus of Gott’s Nature article, the longevity of the human species? A species’s survival depends on its ability to adapt to short- and long-term environmental changes produced by other species in its ecosystem and by climatological and geological processes. The adaptations are made possible by existing genetic variability in the gene pool and by random mutation. How homo sapiens fits into this picture is a complicated question, certainly not amenable to a universal statistical rule. As Ferris puts it, “ … in my experience most people either think we’re going to hell in a handbasket or assume that we’re going to be around for a very long time.” Both views are a reflection of advancing technology. The first comes from alarm at technology’s increasing impact—changes might be so rapid that we (and certainly other) species could not adapt. The second comes from a belief that technology can save us—by controlling the environment or by making possible remarkable adaptations such as escaping our earthly environment or changing our genetic constitution.
Gott dismisses all such thinking as the illusions of those who don’t appreciate the power of the Copernican principle. He contends that everything relevant to assessing our future prospects is contained in the statement that we are not at a special time. This article shows that the Copernican principle is irrelevant to considerations of the longevity of our species. Perhaps we are still subject to the factors that determine the survival of other species. More likely, our survival—and the survival of many other species along with us—depends on what we do now and in the future. We better think hard about it.
> Carlton M. Caves has been Professor of Physics and Astronomy at the University of New Mexico since 1992. He received a Ph.D. in Physics from the California Institute of Technology in 1979, working mainly on the theory of gravitational-wave detection. His subsequent career has proceeded from relativity theory through quantum optics to quantum information science, the unifying strand being an interest in how quantum mechanics impacts what we can measure and what we can know. His interest in probability theory springs from the conviction that to have any hope of understanding the weird features of quantum mechanics, one must first have a clear idea of what probability means. He maintains that the only consistent way of thinking about probabilities is the Bayesian interpretation, which holds that probabilities, far from being physical properties, are a measure of credible belief based on what one knows. His main accomplishment is to have had a series of excellent Ph.D. students and postdocs, who continue to teach him most of what he knows. |
warning/0001/cond-mat0001360.html | ar5iv | text | # Motion of dark solitons in trapped Bose-Einstein condensates
## Abstract
We use a multiple time scale boundary layer theory to derive the equation of motion for a dark (or ‘grey’) soliton propagating through an effectively one-dimensional cloud of Bose-Einstein condensate, assuming only that the background density and velocity vary slowly on the soliton scale. We show that solitons can exhibit viscous or radiative acceleration (anti-damping), which we estimate as slow but observable on experimental time scales.
The success of the Gross-Pitaevski mean field theory in describing experimentally observed dilute Bose condensates shows that one really can persuade a large number of particles to behave as a field. There is thus a pleasant circularity in investigating situations where this field in turn behaves in a particle-like manner, in that it contains topological defects or solitons. In this paper we discuss one particular particle-like configuration of the Gross-Pitaevski mean field, namely the one-dimensional dark soliton. Quasi-one-dimensional traps are realistic prospects in the relatively near future , and dark solitons are expected to emerge in them from generic violent collisions between condensates . A controlled method for creating them by adiabatic state engineering with lasers has also recently been proposed . And they are expected to play a crucial role in the eventual decay of superfluid currents in tight toroidal traps , which would be a valuable analog of the thin superconducting wires whose resistivity is one of the triumphs of non-equilibrium statistical mechanics . Although dark solitons have been studied extensively in nonlinear optics , optical fibres are spatially homogeneous on the relevant scale. In this Letter we extend or correct previous treatments of dark soliton motion in Bose condensates , by using multiple scale analysis to derive equations of motion for a dark soliton moving through a background condensate which changes slowly in both space and time, and is subject to a generic slowly-varying potential (not necessarily harmonic). This powerful analytical method may also be useful for other structures.
The Gross-Pitaevski equation (GPE) governs the evolution of the c-number ‘macroscopic wave function’ $`\psi (\stackrel{}{x},t)`$ of a Bose-Einstein condensate. (This is of course a mean field approximation to the full quantum field theory; we will consider dissipation from quasi-particle interactions very briefly below.) Incorporating a chemical potential by extracting a factor $`e^{i\mu t/\mathrm{}}`$, and then appropriately scaling the wave function, space, and time, one can write this equation in the convenient form
$$i_t\psi =\frac{1}{2}^2\psi +(|\psi |^2+V(\stackrel{}{x})1)\psi .$$
(1)
We assume here a positive scattering length; and we do not restrict the normalization constant $`U𝑑x|\psi |^2`$, which is the number of particles rescaled by the strength of their mutual repulsion. Crucially, we assume a trap so thin that one can apply the GPE in one dimension. The approach to this limit from three dimensions has recently been discussed . The essential requirement is that the transverse thickness of the trap be less than the healing length, to stabilize against buckling modes in the GPE. Making transverse confinement stronger than the temperature will make even the quantum field theory effectively one dimensional. Experimental capability is already approaching both these limits.
Eqn. (1) in one dimension with constant $`V`$ has been extensively studied in nonlinear optics , and a solution with a localized structure has long been known: $`\psi _{DS}=\mathrm{tanh}\sqrt{1V}x`$. This time-independent solution is known as a dark soliton, because it describes a small dark spot in a light pulse; in our case this becomes a small ‘bubble’ of low condensate density in the dilute Bose gas. If $`\psi (x)`$ were restricted to be real, the dark soliton would be topologically stable, like other ‘kink’ solitons; but by taking $`\psi `$ into the complex plane one can deform it into a configuration with constant density and phase, eliminating the ‘bubble’. So unlike two-dimensional vortices, dark solitons are not topologically stable. Before considering their motion, therefore, one should first examine their stability; but in fact the two problems are closely connected, because the complex deformations of $`\psi _{DS}`$ include the larger family of dark solitons moving with arbitrary (sub-critical) velocities . These exact non-stationary solutions to (1), moving with constant velocity $`p\dot{q}`$, are
$$\psi _{GS}=ip+\sqrt{v_c^2p^2}\mathrm{tanh}\sqrt{v_c^2p^2}[xq(t)],$$
(2)
where $`v_c\sqrt{1V}`$ is the Landau critical velocity (which the soliton cannot exceed). For $`p0`$ we recover the motionless dark soliton at position $`x=q`$. Since for non-zero $`p`$ the condensate density $`|\psi |^2`$ never vanishes, moving dark solitons are also called grey solitons.
For moving solitons the difference between maximum and minimum densities is $`v_c^2p^2`$, and the phase slip across the soliton is $`\pi +2\mathrm{arctan}(p/\sqrt{v_c^2p^2})`$. This means that in the limit $`p\pm v_c`$, the soliton becomes identical with motionless condensate. Thus the soliton with maximum speed is the ground state; the energy of slower solitons is higher! In this sense one may say that dark solitons have negative kinetic energy. To be precise: The solutions given in (2) have fixed chemical potential (set to one), and their free energy $`GEU`$ is
$$G=G_0+\frac{4}{3}(1Vp^2)^{\frac{3}{2}},$$
(3)
where $`E\frac{1}{2}𝑑x[|\psi ^{}|^2+|\psi |^4+2V|\psi |^2]`$, and $`G_0`$ is the free energy of the ground state $`\psi =\sqrt{1V}`$. (Alternatively we could change the chemical potential with $`p`$ so as to keep the particle number constant; this slightly different family of grey soliton solutions with constant $`U`$ has energy $`E=E_0+(4/3)(1Vp^2)^{3/2}`$.) Thus, dark solitons are energetically as well as topologically unstable; but their instability is to acceleration, not to filling or collapse. Bogoliubov theory shows that acceleration is indeed their only instability, and as we discuss below, the ‘anti-damping’ time scale should be quite long.
So apart from slow anti-damping, dark solitons in bulk behave as robust free particles, obeying $`\ddot{q}=\dot{p}=0`$. We now consider a dark soliton in a slowly varying medium, where we will be able to derive a more complicated equation of motion, if we interpret the ‘slow variation’ of $`V(x)`$ as implying that there exists a length scale $`\mathrm{\Lambda }`$ which is both large compared to the soliton scale and small compared to the trap scale. Precisely: there is a small dimensionless $`ϵ`$, such that $`\mathrm{exp}(\sqrt{1V(q)p^2}\mathrm{\Lambda })<<ϵ`$ for all phase space points $`(q,p)`$ through which the soliton will actually pass, but $`V(x)=V(q)+V^{}(q)(xq)+𝒪(ϵ^2)`$ as long as $`|xq|<\mathrm{\Lambda }`$ (with $`V^{}(q)|\mathrm{\Lambda }|`$ being of order $`ϵ`$). We will then examine an interval $`|xq|<\mathrm{\Lambda }`$ around a grey soliton in a trap, the interval moving with the soliton, and smoothly patch this interval into a background condensate cloud in the hydrodynamic limit. Applying a simple form of multiple time scale analysis will then yield the equation of motion. This involved procedure (‘boundary layer theory’) is indeed necessary: merely treating $`V(x)`$ as a perturbation is only valid if the potential is everywhere small, whereas we are interested in cases where, over large enough distances, it can change greatly. And ordinary perturbation theory will be valid only for a short time, but we are interested in large changes over longer times (such as the reflection of the soliton from a barrier).
We begin with the simplest step of considering the background cloud. We will assume that the background cloud consists of condensate varying slowly on the healing length scale and its associated time scale, except possibly for small high frequency perturbations. For the dominant low-frequency component, we define $`\psi =\sqrt{\rho }e^{i\theta }`$ for real $`\rho ,\theta `$, and stipulate that spatial and temporal derivatives of $`\rho `$ and $`v_x\theta `$ are of order $`ϵ`$. We may therefore neglect $`\sqrt{\rho }^{\prime \prime }/\sqrt{\rho }`$ in the GPE to obtain the hydrodynamic equations
$`_t\rho `$ $`=`$ $`_x(\rho _x\theta )_x(\rho v)`$ (4)
$`_t\theta `$ $`=`$ $`1\rho Vv^2/2.`$ (5)
We now patch our family of solitons into this background condensate: within $`|xq|<\mathrm{\Lambda }`$ we write $`\psi =e^{i\overline{\theta }+i\overline{v}(xq)}[\psi _0+ϵ\psi _1(xq,t)]+𝒪(ϵ^2)`$, for
$$\psi _0=i(p\overline{v})+\kappa \mathrm{tanh}\kappa (xq),$$
(6)
where $`\overline{\theta }(t)[\theta (q\mathrm{\Lambda },t)+\theta (q+\mathrm{\Lambda },t)]/2`$, $`\overline{\rho }`$ and $`\overline{v}`$ are similarly defined, and $`\kappa ^2\overline{\rho }(p\overline{v})^2`$. Since $`p`$ would be constant if $`ϵ0`$, we conclude that $`\dot{p}`$ is order $`ϵ`$; in fact, $`p`$, $`\overline{v}`$, and $`\kappa `$ may be taken as functions of the ‘slow time’ $`ϵt`$.
We then expand the Gross-Pitaevski equation to order $`ϵ`$ within $`|xq|<\mathrm{\Lambda }`$, keeping in mind that $`\dot{p},V^{},\dot{\kappa }`$ etc., are all $`𝒪(ϵ)`$. First using (4) to establish
$`_t\overline{\theta }`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\pm }{}}\left[(_t+p_x)\theta (x,t)\right]_{x=q\pm \mathrm{\Lambda }}`$ (7)
$`=`$ $`p\overline{v}+1\overline{\rho }V(q)\overline{v}^2/2+𝒪(ϵ^2),`$ (8)
we find from the zeroth order terms $`\dot{q}=p`$ as always, plus the following at first order in $`ϵ`$:
$`[V^{}(q)+\dot{\overline{v}}](xq)\psi _0+[\dot{p}\dot{\overline{v}}i\dot{\kappa }{\displaystyle \frac{}{\kappa }}\kappa \mathrm{tanh}\kappa (xq)]`$ (9)
$`=ϵ\left[i_t\psi _1|_{xq}i(p\overline{v})\psi _1^{}+{\displaystyle \frac{1}{2}}\psi _1^{\prime \prime }\psi _0^2\psi _1^{}(2|\psi _0|^2\overline{\rho })\psi _1\right].`$ (10)
We will abbreviate Eqn. (9) as $`i_t\psi _1|_{xq}+(\psi _1,\psi _1^{})=𝒥(x,ϵt)`$. (Note that it is a straightforward but very important step in obtaining (9) to distinguish $`_t`$, which is, as usual, differentiation with respect to $`t`$ with $`x`$ fixed, from differentiation with respect to $`t`$ with $`xq`$ fixed: $`_tf(xq,t)|_x=_tf(xq,t)|_{xq}p_xf(xq,t)|_t`$ for any function $`f`$.)
We could then proceed to solve Eqn. (9) using the Green’s function for the homogeneous part. To construct this we would need all the independent solutions to the homogeneous equation; but in fact for our purpose we will require only the four independent solutions $`u_1,\mathrm{},u_4`$ to the time-independent equation $`(u_j,u_j^{})=0`$. Distinguishing the fast and slow parts of $`\psi _1`$ by defining $`\psi _1\varphi (xq,ϵt)+\chi (xq,t)`$, we can use (9) to show that the real parts of certain integrals are constrained to vanish:
$`\mathrm{Re}{\displaystyle _{q\mathrm{\Lambda }}^{q+\mathrm{\Lambda }}}𝑑x\left(2iu_j^{}(x)_t\chi |_{xq}+_x[\chi ^{}u_j^{}u_j^{^{}}\chi 2i(p\overline{v})u_j^{}\chi ]\right)`$ (11)
$`=`$ $`\mathrm{Re}{\displaystyle _{q\mathrm{\Lambda }}^{q+\mathrm{\Lambda }}}𝑑x\left(2u_j^{}(x)𝒥(x,ϵt)_x[\varphi ^{}u_j^{}u_j^{^{}}\varphi 2i(p\overline{v})u_j^{}\varphi ]\right)=0.`$ (12)
Here the crucial final equality follows from the fact that the two sides of the preceding equation vary on different time scales, and so must separately equal zero. (Since the first line is linear in $`\chi `$, which must be fast, a non-zero constant is not allowed.) This is the great strength of the combination of boundary layer and multiple time scale analysis, that it allows us to obtain the motion of a short-scale defect in a long-scale background, by solving only time-independent equations.
Eqn. (11) gives us four constraints, which since all four $`u_j(x)`$ may be obtained explicitly, can be evaluated. In addition we require that our soliton $`\psi `$ match smoothly into the background flow as $`|xq|\mathrm{\Lambda }`$, and this introduces constraints from (4) as well. Together these constraints fix the hitherto unknown $`\dot{p}`$, and also relate $`\rho ,\theta ,v`$ at $`x=q\mathrm{\Lambda }`$ to their values at $`x=q+\mathrm{\Lambda }`$. We illustrate the procedure with the simplest but most important constraint, the one involving $`u_1(x)=\mathrm{sech}^2\kappa (xq)`$. Since $`u_1(\pm \mathrm{\Lambda })`$ and $`u_1^{}(\pm \mathrm{\Lambda })`$ are exponentially negligible, we discard terms of this order in (11). We can then extend the limits of integration to infinity and shift the integration dummy variable $`xqy`$, to obtain
$`\kappa {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑y[\kappa (V^{}(q)+\dot{\overline{v}})y\mathrm{tanh}\kappa y+\dot{p}\dot{\overline{v}}]\mathrm{sech}^2\kappa y=2\dot{p}+V^{}(q)\dot{\overline{v}}=0.`$ (13)
This is the equation of motion, accurate to $`𝒪(ϵ)`$, for a dark soliton in an otherwise hydrodynamic condensate in an inhomogeneous potential. We will examine it in some simple limits, before discussing the conditions obtained from the other $`u_j`$, and from requiring (4) as $`|xq|\mathrm{\Lambda }`$.
With $`v=0`$, Eqn. (13) implies
$$\ddot{q}=\frac{1}{2}V^{}(q).$$
(14)
In a harmonic trap, this implies oscillation of the soliton with frequency $`1/\sqrt{2}`$ times that of the dipole mode of the condensate (the trap frequency). This result can also be obtained for small oscillations by solving the Bogoliubov equations for a motionless soliton in a trap, using a simpler, time-independent version of the ‘boundary layer’ approach that led to (13) . We have confirmed this frequency to rather more than the expected accuracy in numerical simulations of harmonic traps over a wide range of condensate densities and oscillation amplitudes; we have also confirmed that the center of mass is decoupled and oscillates at the trap frequency. Eqn. (14) also holds for arbitrary potentials, however, as long as they vary slowly on the healing length scale. We have therefore further confirmed the good accuracy of our equation of motion by solving Eqn. (1) numerically over a wide range of parameters and for various potentials; a generic example is shown in Fig. 1. Since with lasers one can generate micro-wells or barriers in a trap, it should be possible to realize similar potentials experimentally.
We now consider a stationary background flow, such as in an inhomogeneous toroidal trap holding a persistent current. In general the system is quite complicated; but in the limit where both the inhomogeneous potential $`V`$ and the average kinetic energy $`v_0^2`$ are small compared to the chemical potential, we have $`\rho 1V`$, $`vv_0[1+V]`$, which with $`_tv=0`$ implies the easily solvable equation
$$\ddot{q}=V^{}(q)[v_0\dot{q}1]/2.$$
(15)
Despite the $`\dot{q}`$ term, Eqn. (15) is not dissipative: it may be derived variationally from the Lagrangian $`(2/v_0^2)(1v_0\dot{q})[\mathrm{ln}(1v_0\dot{q})1]V`$, and the energy $`\dot{q}\frac{L}{\dot{q}}L`$ is conserved.
A simple example of the generally still more complex case where $`\rho `$ and $`v`$ are time-dependent is a soliton moving in a harmonic trap of frequency $`\mathrm{\Omega }`$ in which the collective dipole mode has also been excited:
$$\ddot{q}=\frac{\mathrm{\Omega }^2}{2}[q+Q\mathrm{cos}\mathrm{\Omega }(tt_0)],$$
(16)
where $`Q`$ is the dipole amplitude. As required by the Ehrenfest theorem for a condensate in a harmonic trap, the rigid dipole oscillation of background and soliton together, $`q=Q\mathrm{cos}\mathrm{\Omega }(tt_0)`$, is a solution to (16).
Since this Ehrenfest theorem states that the centre of mass of the condensate must oscillate at the trap frequency $`\mathrm{\Omega }`$, but (14) makes the small soliton ‘bubble’ oscillate at $`\mathrm{\Omega }/\sqrt{2}`$, it is clear that the background condensate must be perturbed by the soliton moving through it. This brings us back to the constraints we have not yet examined, which turn out to imply discontinuities of $`𝒪(ϵ)`$ in both $`\rho `$ and $`v`$ between $`xq=\pm \mathrm{\Lambda }`$. These are in addition to the trivial discontinuities $`\mathrm{\Lambda }`$ due to background gradients. It is both convenient, and consistent with our $`𝒪(ϵ)`$, ‘boundary layer’ approach, to consider the entire interval $`|xq|<\mathrm{\Lambda }`$ to be pointlike as far as the background condensate is concerned; so, formally letting $`\mathrm{\Lambda }0`$ after obtaining all our results so far, the discontinuities across the soliton become abrupt. The requirement for them can then be expressed as delta function sources, at $`x=q(t)`$, which must be added to the hydrodynamic equations. The result can be shown to be
$`_t\rho `$ $`=`$ $`_x(\rho v)+2\dot{\kappa }\delta (xq)`$ (17)
$`_tv`$ $`=`$ $`_x(v^2/2+\rho +V)+\rho ^1\delta (xq)[\kappa (V^{}+\dot{v})+2(pv)\dot{\kappa }].`$ (18)
In most cases indeed these delta function sources are unimportant, since the soliton couples only to the smooth part $`\overline{v}=[v(q+,t)+v(q,t)]/2`$, and the sources generate only discontinuities. The effect of these on $`\overline{v}`$ depends on the boundary conditions for the entire condensate, and solving (13) and (17) together to determine this effect is generally not much easier than numerically solving the GPE with the dark soliton. There are nevertheless some important points that can be learned from the source terms. For instance, they preserve the Ehrenfest theorem in a harmonic trap, as may be checked straightforwardly by evolving $`X=2\kappa q+𝑑xx\rho `$ under (13) and (17). And because of its coupling to the background fluid, one can deduce that a dark soliton oscillating in a small well within a large sample of bulk condensate will generate sound waves, and so exhibit radiative anti-damping. Numerical integration of the GPE confirms this prediction: the soliton eventually escapes from the micro-well, the radiation ceasing as it enters the region of constant potential .
In a finite trap, however, coupling to the background condensate modes does not provide dissipation. In this case dissipation can only come from corrections to mean field theory; in particular, from collisions with uncondensed atoms of the thermal cloud. A simple estimate of the anti-damping time scale is provided by the rate at which the soliton encounters particles, divided by the number $`2\kappa `$ of particles ‘in’ the soliton (for the ‘soliton mass’). At current experimental temperatures and densities, with 99% of the particles in the condensate, this time is on the order of one second; which agrees with the calculation in Ref. of the dark soliton decay time. It is clear therefore that the instability of dark solitons is by no means fast enough to prevent their observation.
Acknowledgements
We are happy to acknowledge valuable discussions with J.I. Cirac, V. Perez-Garcia, and P. Zoller. This work was supported by the European Union under the TMR Network ERBFMRX-CT96-0002 and by the Austrian FWF. |
warning/0001/math0001062.html | ar5iv | text | # The Bianchi–Darboux transform of 𝐿-isothermic surfaces
## 1. Introduction
Certain types of integrable non-linear PDEs (soliton equations) arise in differential geometry as compatibility conditions for the linear equations obeyed by frames adapted to surfaces in higher dimensional manifolds. In a number of situations, the construction of new solutions of the arising PDE relies on the existence of Bäcklund type transformations for the surfaces and on their permutability properties. Well-known examples include pseudo-spherical surfaces, affine minimal surfaces, and isothermic surfaces in Möbius geometry \[CT\],\[BHPP\],\[CGS\],\[TT\]. The loop group approach to soliton theory, recently developed by Terng–Uhlenbeck and others \[TU\], explains the uniformity of results in several of these examples.
This paper concentrates on $`L`$-isothermic surfaces. An immersion $`f:M^3`$ is called $`L`$-isothermic if about each point of $`M`$ there exist curvature line coordinates $`(x,y)`$ which are conformal with respect to the third fundamental form of $`f`$. $`L`$-isothermic surfaces are invariant under a group of contact transformations, the Laguerre group, which is isomorphic to the identity component of the isometry group of Minkowski 4-space. This group acts transitively on the set of oriented 2-spheres and points as well as on the set of oriented planes of $`^3`$. We first became interested in $`L`$-isothermic surfaces when studying the (infinitesimal) deformation problem for surfaces under the Laguerre group. Namely, $`L`$-isothermic surfaces are the only surfaces allowing 1-parameter families of second order deformations \[MN1\],\[MN2\]. Such families correspond to the solutions of the non-linear fourth-order equation $`\mathrm{\Delta }(w_{xy}w^1)=0`$ which in turn is equivalent to the Gauss-Codazzi equations and arises as the integrability condition of a linear differential system containing a free parameter \[MN3\].
The theory of $`L`$-isothermic surfaces presents many analogies with that of isothermic surfaces in Möbius geometry. The latter has received much attention in recent studies \[Bu\],\[BHPP\],\[He\],\[HP\],\[HMN\], especially in relation with the general theory of curved flats in pseudo-Riemannian symmetric spaces \[FP\]. Curved flats arise in 1-parameter families as solutions to a certain integrable system expressed in the form of a zero-curvature equation with spectral parameter. This approach to the study of isothermic surfaces provides a uniform framework for understanding their theory of transformations and deformation. The notion of isothermic submanifolds have also been defined in other contexts; for instance, a general theory of isothermic submanifolds in symmetric $`R`$-spaces have been recently developed in \[BPP\].
In this paper, we study a geometric transformation for $`L`$-isothermic surfaces by realizing them as enveloping surfaces of a 2-sphere congruence. The Laguerre invariant conditions that the congruence preserve curvature lines and that the third fundamental forms of the two envelopes at corresponding points be conformal yield that both surfaces are $`L`$-isothermic. This is the content of Theorem 1. The resulting congruence is an analogue of the Darboux congruence occurring in Möbius geometry and its two envelopes are said to form a Bianchi–Darboux pair. A Bianchi–Darboux pair gives rise to a curved flat in the Grassmannian $`\stackrel{~}{G}_{1,1}(_1^4)`$ of oriented 2-planes of signature $`(1,1)`$ in $`_1^4`$; conversely, a curved flat in this Grassmannian only determines the normals of a Bianchi–Darboux pair. Moreover, the spectral deformation of curved flats amounts to second order deformation as in the conformal situation. In Section 5, we prove the existence of Bianchi–Darboux transforms by explicitly constructing new $`L`$-isothermic surfaces from a given one; the construction requires solving an integrable linear differential system. This furnishes a geometrical method of deriving new solutions of the defining PDE $`\mathrm{\Delta }(w_{xy}w^1)=0`$ from any given one. Finally, in Section 6, we establish a permutability theorem for the Bianchi–Darboux transformation.
###### Remark Acknowledgements
It is a pleasure to thank Fran Burstall and Udo Hertrich-Jeromin for many interesting and instructive discussions concerning isothermic surfaces and their transformation theory. Thanks also to the referee for his comments and questions.
## 2. $`L`$-isothermic surfaces
Let $`f:M^3`$ be an immersion of a surface $`M`$ in Euclidean space with no parabolic points oriented by a field of unit normals $`n:MS^2`$. Consider on $`M`$ the unique complex structure compatible with the given orientation and the conformal structure defined by the third fundamental form $`\text{rm III}=\text{d}n\text{d}n`$. (Here $``$ denotes the Euclidean inner product.) Accordingly, the second fundamental form rm II decomposes into bidegrees:
$$\text{rm II}=\text{rm II}^{(2,0)}+\text{rm II}^{(1,1)}+\text{rm II}^{(0,2)}.$$
$`\text{rm II}^{(2,0)}`$ is a globally defined $`(2,0)`$ symmetric bilinear form on $`M`$ which plays the role of the usual Hopf differential for the pair of quadratic forms rm III, rm II. We refer to $`\text{rm II}^{(2,0)}`$ simply as the Hopf differential.
###### Definition Definition
A Riemann surface $`M`$ equipped with a holomorphic quadratic differential $`q`$ is called a polarized surface. An immersion $`f:(M,q)^3`$ is called $`L`$-isothermic if $`\text{rm II}^{(2,0)}=\mu q`$, for a real-valued smooth function $`\mu :M`$.
Near any point $`pM`$ where $`q_{|p}0`$ there exists local complex coordinate $`z=x+iy:\mathrm{\Omega }M`$ such that $`q_{|\mathrm{\Omega }}=\text{d}z^2`$ and $`\text{rm III}_{|\mathrm{\Omega }}=e^{2u}\text{d}z\text{d}\overline{z}`$. Then $`(x,y)`$ are curvature line coordinates which are conformal for the third fundamental form of $`f`$. $`(x,y)`$ will be called conformal principal coordinates for $`f`$ <sup>1</sup><sup>1</sup>1The umbilic points of the immersion are precisely the zeros of the Hopf diffrential. The idea here is that not all umbilic points prohibit conformal principal coordinates. If an umbilic is caused by a zero in the polarization, then it is a bad umbilic, where no conformal principal coordinates may be found; instead, umbilic points caused by zeros in $`\mu `$ will not cause problems..
Examples of $`L`$-isothermic surfaces include surfaces of revolution, molding surfaces \[BCG\], surfaces with plane lines of curvature in both systems \[MN4\], and the class of Weingarten surfaces on which $`aH+bK=0`$, for constants $`a,b`$ with $`a0`$, where $`H`$ and $`K0`$ denote the mean and Gaussian curvatures, respectively. The last example follows as an application of Hopf’s classical argument: $`H/K`$ which is $`H(\text{rm III},\text{rm II}):=(1/2)\text{tr}_{\text{rm III}}\text{rm II}`$ is constant if and only if $`\text{rm II}^{(2,0)}`$ is holomorphic \[H\].
## 3. The geometry of $`L`$-isothermic surfaces
A pair of real quadratic forms rm III and rm II on $`M`$ such that: rm III is positive definite, the intrinsic curvature $`K(\text{rm III})1`$, and such that rm II satisfies the Codazzi equations with respect to the metric rm III, defines, up to contact tansformations, an immersion $`F=(f,n):M^3\times S^2`$ in the space of contact elements of $`^3`$ satisfying the contact condition $`\text{d}fn=0`$. A map satisfying this condition is called a Legendre immersion. The Euclidean projection $`f:M^3`$ need not be an immersion nor even have constant rank.
Let $`\mathrm{\Lambda }=^3\times S^2`$ denote the space of contact elementsof $`^3`$. Now each contact element $`(x,n)\mathrm{\Lambda }`$ corresponds to a null line in Minkowski 4-space $`_1^4`$ with its Lorentz scalar product $`,=(u^1)^2+(u^2)^2+(u^3)^2(u^4)^2`$ via
$$(x,n)[x,n]:=\{{}_{}{}^{t}(xtn,t),t\}_1^4.$$
Under this identification, the identity component $`L=^4SO_0(3,1)`$ of the isometry group of $`_1^4`$ acts transitively on $`\mathrm{\Lambda }`$ and preserves the contact condition. The action of $`L`$ permutes Legendre immersions: let $`F=(f,n)`$ be a Legendre immersion and, for $`pM`$, consider the null line $`[f(p),n(p)]`$. For each $`AL`$, $`A[f(p),n(p)]\mathrm{\Lambda }`$ and intersects $`^3=\{{}_{}{}^{t}(x,t)_1^4:t=0\}`$ in the unique point $`f^{}(p)`$. If $`{}_{}{}^{t}(n^{},1)`$ denotes the null direction of $`A[f(p),n(p)]`$, then $`AF:=(f^{},n^{})`$ defines another Legendre immersion. Notice that the action does not preserve the Euclidean fibration $`\mathrm{\Lambda }^3`$.
###### Remark Remark
A standard model for Laguerre geometry is obtained by identifying $`_1^4`$ with the set $`𝒮`$ of oriented spheres (including point spheres) of $`^3`$ by
$$\mathrm{\Sigma }:𝒮_1^4,\sigma _r(p){}_{}{}^{t}(p,r),$$
where $`\sigma _r(p)`$ denotes a sphere with center $`p`$ and signed radius $`r`$. Note that if $`r>0`$ (resp. $`r<0`$) then $`{}_{}{}^{t}(p,r)`$ is the vertex of the backward (resp. forward) pointing light-cone which intersects $`^3`$ in exactly the sphere $`\sigma _r(p)`$. Two spheres $`\sigma _r(p)`$ and $`\sigma _r^{}(p)`$ are in oriented contact if and only if $`{}_{}{}^{t}(pp^{},rr^{})`$ is a null vector. Thus each contact element $`(x,n)`$ determines a null line of 2-spheres all of which are in oriented contact at $`x`$ with normal $`n`$. Following the classical terminology, a Laguerre transformation is a contact transformation of $`\mathrm{\Lambda }`$ induced by an element in the group $`L`$. In terms of $`^3`$, a Laguerre transformation takes oriented planes in $`^3`$ to oriented planes, and oriented spheres to oriented spheres. In this context, two immersions $`f,f^{}:M^3`$ are said to be Laguerre equivalent if there exists $`AL`$ such that their respective Legendrian lifts $`F,F^{}`$ satisfy $`AF=F^{}`$. (For more details about Laguerre geometry we refer to \[Bl\],\[C\],\[MN1\].)
The theory of $`L`$-isothermic surfaces belongs in the Laguerre geometry:
###### (Laguerre invariance)
Let $`F=(f,n)`$ be the Legendrian lift of an $`L`$-isothermic surfaces $`f`$ and $`A:\mathrm{\Lambda }\mathrm{\Lambda }`$ be a Laguerre contact transformation. Then the Euclidean projection of $`AF=(f^{},n^{})`$ is $`L`$-isothermic also.
###### Demonstration Proof
Let $`^+`$ be the positive light-cone in $`_1^4`$. The projective light-cone $`(^+)`$ identifies with the conformal 2-sphere $`S^2`$ and the projection $`^+(^+)`$ is a principal $`^+`$–bundle which is trivial. For each $`AL`$, $`\text{d}ASO_0(3,1)`$, which preserves the light-cone $`^+`$ and descends to an orientation-preserving conformal diffeomorphism $`\stackrel{~}{A}`$ of $`S^2`$. This implies that $`n^{}=\stackrel{~}{A}(n)`$. Thus, the conformal class of rm III and then the conformal class of rm II (mod rm III) are Laguerre invariant. ∎
###### Remark 3.1 Notations: moving frames in Laguerre geometry
Here and in the following we shall consider $`_1^4`$ with linear coordinates $`x^1,\mathrm{},x^4`$ such that $`,=2x^1x^4+(x^2)^2+(x^3)^2`$; an orientation for which $`\text{d}x^1\mathrm{}\text{d}x^40`$; and a time-orientation given by $`x^1+x^4>0`$. By a Laguerre frame is meant a position vector $`a_0_1^4`$ and an oriented basis $`a_1,a_2,a_3,a_4`$ of $`_1^4`$ such that $`(a_0;a_1,a_2,a_3,a_4)L`$. Up to the choice of a reference frame, the manifold of Laguerre frames may be identified with the group $`L`$. For any $`A=(a_0,a)L`$, $`a_i`$, $`i=1,2,3,4`$, denote the column vectors of the matrix $`a`$. Geometrically, the null directions $`a_1`$, $`a_4`$ represent oriented planes which are in oriented contact with the oriented sphere represented by $`a_0`$. By $`[a_0,a_1],[a_0,a_4]`$ we denote the lines in $`_1^4`$ through $`a_0`$ with null directions $`a_1`$ and $`a_4`$, respectively.
The transitive action of $`L`$ on $`\mathrm{\Lambda }`$ defines a principal $`L_0`$-bundle
$$\pi _L:L\mathrm{\Lambda }=L/L_0,A[a_0,a_1].$$
###### Definition Definition
A Laguerre frame field in $`\mathrm{\Lambda }`$ is a local smooth section $`A=(a_0,a)`$ of $`\pi _L`$.
A (local) Laguerre framing along a Legendre immersion $`F:M\mathrm{\Lambda }`$ is a smooth map $`A:𝒰L`$ defined on an open subset $`𝒰M`$, such that $`\pi _LA=[a_0,a_1]=F`$.
###### Remark Remark (Principal frames)
Let $`F:(M,q)\mathrm{\Lambda }`$ be a Legendrian immersion. Then $`F`$ can be equipped with a Laguerre framing $`A=(a_0;a_1,a_2,a_3,a_4):ML`$ such that $`F(p)=[a_0,a_1](p)`$ and span$`\{a_2,a_3\}_{|p}=\text{d}F_{|p}(T_pM)`$, for any $`pM`$. Let $`\alpha =A^1\text{d}A`$ be its connection form. Then $`\alpha _0^4=0,\alpha _1^2\alpha _1^30`$. Differentiating $`\alpha _0^4=0`$, it follows that the quadratic form $`\alpha _0^2\alpha _1^2+\alpha _0^3\alpha _1^3`$ is symmetric and then diagonalizable. So we may assume that $`\alpha _0^2\alpha _1^2=\alpha _0^3\alpha _1^3=0`$. We call such an $`A`$ a principal framing along $`F`$. An easy calculation using the Maurer–Cartan equation of $`L`$ yields
$$\text{d}(\alpha _1^2+i\alpha _1^3)=(\alpha _1^1\alpha _2^3)(\alpha _1^2+i\alpha _1^3),$$
from which follows that $`F`$ is $`L`$-isothermic if and only if $`\text{d}(\alpha _1^1\alpha _2^3)=0`$.
## 4. Sphere congruences and the Bianchi–Darboux transformation
In classical surface theory, a sphere congruence is an immersion $`S:M_1^4`$ of a surface $`M`$ into the space $`_1^4`$ of all oriented 2-spheres (including points) of $`^3`$. A Legendre immersion $`F=(f,n)`$ is said to envelope the sphere congruence $`S`$ if for each $`pM`$, the sphere $`\mathrm{\Sigma }^1(S(p))`$ is in oriented contact with the Legendre surface at $`F(p)`$. If $`S`$ is a space-like immersion, i.e., the induced metric on $`M`$ is positive definite, then there exist two enveloping surfaces \[Bl\]. It follows that the spheres of the congruence are the common tangent spheres of the two envelopes. Accordingly, there results a mapping between the enveloping surfaces such that the sphere congruence consists of the spheres passing through the points on the envelopes.
###### Definition Definition
A space-like sphere congruence is called Ribaucour if the resulting mapping between the two envelopes preserves curvature lines. A Ribaucour sphere congruence is called Darboux <sup>2</sup><sup>2</sup>2In the context of Möbius geometry a sphere congruence is an immersion of a surface into the Lorentzian sphere $`S_1^4`$ interpreted as the space of all oriented spheres (excluding points) and oriented planes in $`^3`$. A sphere congruence is then called Ribaucour if the curvature lines on the two envelopes correspond; a Ribaucour sphere congruence is Darboux if the correspondence between the two envelopes is conformal with respect to the first fundamental forms of the envelopes, see for example \[BHPP\] and the literature therein. if the resulting mapping is conformal with respect to the third fundamental forms induced on $`M`$ by the two envelopes $`F`$, $`\widehat{F}`$.
Let $`S:M_1^4`$ be a space-like congruence enveloping $`F`$ and $`\widehat{F}`$; we can consider an adapted Laguerre frame $`A=(S;a_1,\mathrm{},a_4):ML`$ such that $`F=[S,a_1]`$, $`\widehat{F}=[S,a_4]`$, $`a_2,a_3`$ are tangential, and $`a_1,a_4`$ are normal null directions. The connection form $`\alpha =(\alpha _0,\alpha ^{})=A^1\text{d}A`$ of $`A`$ is then
$$[\left(\begin{array}{c}0\\ \alpha _0^2\\ \alpha _0^3\\ 0\end{array}\right),\left(\begin{array}{cccc}\alpha _1^1& \alpha _2^1& \alpha _3^1& 0\\ \alpha _1^2& 0& \alpha _2^3& \alpha _2^1\\ \alpha _1^3& \alpha _2^3& 0& \alpha _3^1\\ 0& \alpha _1^2& \alpha _1^3& \alpha _1^1\end{array}\right)].$$
Observe that $`\alpha _0^2\alpha _0^30`$ on $`M`$. The Maurer–Cartan equation $`\text{d}\alpha +\alpha \alpha =0`$ yields the Ricci equations:
$`0=\alpha _0^2\alpha _2^1+\alpha _0^3\alpha _3^1,`$ $`1`$$`2`$$`3`$
$`0=\alpha _0^2\alpha _1^2+\alpha _0^3\alpha _1^3,`$
$`d\alpha _1^1+\alpha _2^1\alpha _1^2+\alpha _3^1\alpha _1^3=0,`$
the Gauss equation:
$$\text{d}\alpha _2^3=\alpha _2^1\alpha _1^3+\alpha _1^2\alpha _3^1,$$
$`4`$
and the Codazzi equations:
$`\text{d}\alpha _0^2=\alpha _2^3\alpha _0^2,\text{d}\alpha _0^3=\alpha _2^3\alpha _0^2`$ $`5`$$`6`$$`7`$
$`\text{d}\alpha _1^2=\alpha _1^1\alpha _1^2+\alpha _1^2\alpha _2^3,\text{d}\alpha _1^3=\alpha _1^1\alpha _1^3+\alpha _1^2\alpha _2^3`$
$`\text{d}\alpha _2^1=\alpha _2^1\alpha _1^1+\alpha _2^3\alpha _3^1,\text{d}\alpha _3^1=\alpha _3^1\alpha _1^1\alpha _2^3\alpha _2^1.`$
We thus can state:
###### Theorem 1
Let $`S:M_1^4`$ be a flat space-like immersion with flat normal bundle. Then $`S`$ is a Darboux sphere congruence and both its enveloping surfaces – which have opposite orientations – are $`L`$-isothermic.
###### Demonstration Proof
Assume $`S`$ induces a flat space-like metric on $`M`$, i.e., $`\text{d}\alpha _1^1=0`$, and has flat normal bundle, i.e., $`\text{d}\alpha _2^3=0`$. According to (2), the second fundamental form $`\alpha _0^2\alpha _1^2+\alpha _0^3\alpha _1^3`$ of $`S`$ in the normal direction $`a_1`$ is diagonalizable. So we may choose $`A`$ such that
$$\alpha _0^2=h_2\alpha _1^2,\alpha _0^3=h_3\alpha _1^3,h_2h_30.$$
$`8`$
In particular, $`A`$ becomes a principal framing along $`F`$. We can now write $`\alpha _2^1=a_{11}\alpha _1^2+a_{12}\alpha _1^3`$, $`\alpha _3^1=a_{21}\alpha _1^2+a_{22}\alpha _1^3`$ for smooth functions $`a_{ij}`$. From equations (3) and (4) we obtain $`a_{12}=a_{21}`$, $`a_{11}=a_{22}`$. Further substituting into (1) and (2) yields
$$a_{12}(h_2h_3)=0,$$
and hence $`a_{12}=0`$. Therefore $`\alpha _2^1=a_{11}\alpha _1^2`$ and $`\alpha _3^1=a_{11}\alpha _1^3`$. This implies, in particular, that the correspondence induced by $`S`$ on the two envelopes preserves curvature lines and that $`\text{d}a_1,\text{d}a_1\text{d}a_4,\text{d}a_4`$, that is, $`S`$ is a Darboux sphere congruence. Also, this tells us that the two envelopes have opposite orientations.
As for the isothermic property of the envelopes, by the Codazzi equations (7),
$`\text{d}a_{11}\alpha _1^2`$ $`=2a_{11}(\alpha _1^1\alpha _1^2+\alpha _2^3\alpha _1^3)`$
$`\text{d}a_{11}\alpha _1^3`$ $`=2a_{11}(\alpha _1^1\alpha _1^3\alpha _2^3\alpha _1^2),`$
and from these $`2(\alpha _1^1\alpha _2^3)=\text{d}\mathrm{log}|a_{11}|`$, which is the condition for $`F`$ being $`L`$-isothermic according to the remark in the previous section.
Concerning the second envelope $`\widehat{F}`$, $`B=(b_0;b_1,b_2,b_3,b_4):=`$$`(a_0;a_4,a_2,a_3,a_1)`$ defines a frame along $`\widehat{F}`$. Its connection form is computed to be
$$\beta =B^1\text{d}B=[\left(\begin{array}{c}0\\ \alpha _0^2\\ \alpha _0^3\\ 0\end{array}\right),\left(\begin{array}{cccc}\alpha _1^1& \alpha _1^2& \alpha _1^3& 0\\ \alpha _2^1& 0& \alpha _2^3& \alpha _1^2\\ \alpha _3^1& \alpha _2^3& 0& \alpha _1^3\\ 0& \alpha _2^1& \alpha _3^1& \alpha _1^1\end{array}\right)].$$
$`B`$ is then a principal framing along $`\widehat{F}`$, and the one form $`\beta _1^1\beta _2^3`$ is closed. So, also $`\widehat{F}`$ is $`L`$-isothermic. ∎
If $`S`$ is a Darboux congruence, then $`S`$ is either as in Theorem 1, or the normals $`a_1`$ and $`a_4`$ of the two envelopes are Möbius equivalent in $`S^2`$, that is $`F`$ and $`\widehat{F}`$ are Laguerre equivalent and have the same orientations. In the latter case the congruence $`S`$ belongs to a fixed sphere complex \[Bl\].
Up to this degenerate situations, we can state:
###### Theorem 2
A map $`S:M_1^4`$ defines a nondegenerate Darboux sphere congruence if and only if it is a flat space-like immersion with flat normal bundle.
###### Definition Definition
The $`L`$-isothermic surfaces enveloping a nondegenerate Darboux congruence are called Bianchi–Darboux transforms of each other. They form a Bianchi–Darboux pair.
###### Remark Remark
Let $`F=(f,n)`$, $`\widehat{F}=(\widehat{f},\widehat{n})`$ be a Bianchi–Darboux pair enveloped by $`S:M_1^4`$. Next, let $`𝔰𝔬(3,1)=𝔨𝔪`$, $`𝔨=𝔰𝔬(2)\times 𝔰𝔬(1,1)`$ be a symmetric decomposition of the Lie algebra of $`SO_0(3,1)`$. Then, according to (3) and (4), the flatness of both the induced metric and the normal bundle of $`S`$ is expressed by
$$\alpha _𝔪^{}\alpha _𝔪^{}=0,$$
$`9`$
where $`\alpha _𝔪^{}`$ denotes the $`𝔪`$-part of $`\alpha ^{}`$. This condition expresses the fact that the map $`(n,\widehat{n}):MN=S^2\times S^2\mathrm{}\mathrm{\Delta }`$ into the manifold of ordered pairs of distinct points of $`S^2`$ is a curved flat, see \[FP\],\[BHPP\]. $`N`$ may be viewed as the Grassmannian $`\stackrel{~}{G}_{1,1}(_1^4)`$ of oriented 2-planes of signature $`(1,1)`$ in $`_1^4`$ — a pseudo-Riemannian symmetric space with invariant metric of signature $`(2,2)`$. Note that the converse is not true, that is, from a curved flat in $`N`$ we can only construct the normals of a Bianchi–Darboux pair.
###### Definition Definition
A map $`A=(a_0,a):ML`$ such that $`\alpha =(\alpha _0,\alpha ^{})=A^1dA`$ satisfies (9) and $`\alpha _0^1=\alpha _0^4=0`$ is called a curved flat framing and $`a_0:M_1^4`$ the associated sphere congruence. Note that, up to a gauge change, such a frame can be chosen so that $`\alpha _1^1=0`$.
The above discussion yields:
###### Corollary 1
The sphere congruence $`a_0:M_1^4`$ of a curved flat framing $`A=(a_0,a):ML`$ is a Darboux sphere congruence enveloped by the two $`L`$-isothermic immersions $`F=[a_0,a_1]`$ and $`\widehat{F}=[a_0,a_4]`$. Conversely, a (nondegenerate) Darboux congruence $`S`$ defines a curved flat framing $`(S;a_1,a_2,a_3,a_4):ML`$, where $`\text{d}S(TM)=\text{span}\{a_2,a_3\}`$ and $`a_1`$, $`a_4`$ generate the null subbundles of the normal bundle of $`S`$.
In the next section we shall discuss the existence of Bianchi–Darboux transforms.
## 5. Integrability: construction of the Bianchi–Darboux transform
Let $`F:M\mathrm{\Lambda }`$ be $`L`$-isothermic, and consider a principal frame $`A=`$$`(a_0;a_1,a_2,a_3,a_4)`$ along $`F`$, $`F=[a_0,a_1]`$. Let $`\alpha `$ denote its connection form and $`\rho `$ be a smooth positive function such that
$$2(\alpha _1^1\alpha _2^3)=\text{d}\mathrm{log}\rho .$$
Define
$$\alpha _\rho =(\alpha _0,\alpha _\rho ^{}):=[\left(\begin{array}{c}\alpha _0^1\\ \alpha _0^2\\ \alpha _0^3\\ 0\end{array}\right),\left(\begin{array}{cccc}\alpha _1^1& \alpha _2^1+\rho \alpha _1^2& \alpha _3^1\rho \alpha _1^3& 0\\ \alpha _1^2& 0& \alpha _2^3& \alpha _2^1+\rho \alpha _1^2\\ \alpha _1^3& \alpha _2^3& 0& \alpha _3^1\rho \alpha _1^3\\ 0& \alpha _1^2& \alpha _1^3& \alpha _1^1\end{array}\right)].$$
Note that $`0=\text{d}\alpha _\rho +\alpha _\rho \alpha _\rho `$.
Consider a solution $`v={}_{}{}^{t}(v^1,\mathrm{},v^4):M^+_1^4`$ of the completely integrable linear system
$$\text{d}v=\alpha _\rho ^{}v.$$
By a direct computation:
###### Lemma
The one form $`\gamma :=v^2\alpha _0^2+v^3\alpha _0^3v^4\alpha _0^1`$ is closed.
Next put $`s:=\frac{r}{v_4}`$, where <sup>3</sup><sup>3</sup>3As it will be clear from the following discussion, the function $`r`$, which locally integrates $`\gamma `$, gives essentially the signed radius of the Darboux sphere congruence we are going to construct (cf. Section 5.1). $`dr=\gamma `$, and consider the gauged frame $`\stackrel{~}{A}=Ag(s,v)`$, where <sup>4</sup><sup>4</sup>4Observe that the group of all elements $`g(s,v)`$, $`s`$,$`v^+`$, is diffeomorphic to the structure group of the fibration $`𝒫(F)M`$ of principal frames along $`F`$.
$$g(s,v)=[\left(\begin{array}{c}s\\ 0\\ 0\\ 0\end{array}\right),\left(\begin{array}{cccc}1/v^4& v^2/v^4& v^3/v^4& v^1\\ 0& 1& 0& v^2\\ 0& 0& 1& v^3\\ 0& 0& 0& v^4\end{array}\right)].$$
Observe that $`[\stackrel{~}{a_0},\stackrel{~}{a_1}]=[a_0,a_1]=F`$. The connection form $`\stackrel{~}{A}^1\text{d}\stackrel{~}{A}`$ takes the form
$$[\left(\begin{array}{c}0\\ \alpha _0^2+s\alpha _1^2\\ \alpha _0^3+s\alpha _1^3\\ 0\end{array}\right),\left(\begin{array}{cccc}0& \rho v^4\alpha _1^2& \rho v^4\alpha _1^3& 0\\ \frac{1}{v^4}\alpha _1^2& 0& \alpha _2^3+\frac{1}{v^4}(v^3\alpha _1^2v^2\alpha _1^3)& \rho v^4\alpha _1^2\\ \frac{1}{v^4}\alpha _1^3& \alpha _2^3\frac{1}{v^4}(v^3\alpha _1^2v^2\alpha _1^3)& 0& \rho v^4\alpha _1^3\\ 0& \alpha _1^2& \alpha _1^3& 0\end{array}\right)].$$
Thus $`\stackrel{~}{A}`$ is a curved flat framing and $`\stackrel{~}{a}_0:M_1^4`$ is a Darboux congruence which envelopes the Bianchi–Darboux pair $`F`$ and $`\widehat{F}=[\stackrel{~}{a_0},\stackrel{~}{a_4}]`$.
We thus have proved:
###### Theorem 3
Let $`F:M\mathrm{\Lambda }`$ be an $`L`$-isothermic immersion and $`A`$ be any principal frame along $`F`$. Then any solution $`v:M^+`$ of the linear system
$$\text{d}v=\alpha _\rho ^{}v$$
defines an $`L`$-isothermic immersion
$$\widehat{F}:=[a_0+sa_1,v^1a_1+v^2a_2+v^3a_3+v^4a_4]:M\mathrm{\Lambda }$$
which is a Bianchi-Darboux transform of $`F`$.
###### Remark Remark
The space of principal frames $`𝒫(F)`$ along a Legendrian immersion $`F`$ can be viewed as a 6-dimensional integral submanifold of the exterior differential system $`\omega _0^4=0`$, $`\omega _0^2\omega _1^2=\omega _0^3\omega _1^3=0`$ on $`L`$ with independence condition $`\omega _1^2\omega _1^30`$, where $`\omega `$ denotes the Maurer–Cartan form of $`L`$. Proving the existence of Bianchi–Darboux transforms for an $`L`$-isothermic $`F`$ can be reduced to checking the Frobenius integrability condition for the Pfaffian system on $`𝒫(F)\times `$ given by:
$$\omega _0^1=0,\omega _2^1a\omega _1^2=0,\omega _3^1+a\omega _1^3=0,\text{d}a+2a(\omega _1^1\omega _2^3)=0.$$
###### Remark 5.1 The Bianchi-Darboux transform in Euclidean terms
Let $`f:(M,q)^3`$ be an $`L`$-isothermic immersion with normal $`n`$ and conformal principal coordinate $`z=x+iy`$. If we identify $`𝔼(3)`$ with the subgroup of $`L`$ consisting of all $`AL`$ fixing the time-like vector $`ϵ_1+ϵ_4`$ ($`ϵ_1,\mathrm{},ϵ_4`$ the canonical basis of $`_1^4`$), the Euclidean framing $`(f;t_1,t_2,t_3):M𝔼(3)`$ defined by $`t_1=n`$, $`t_2=\frac{f_x}{f_x}`$, $`t_3=\frac{f_y}{f_y}`$ corresponds to the Laguerre framing
$$E=(e_0;e)=(\left(\begin{array}{c}\frac{f^1}{\sqrt{2}}\hfill \\ f^2\\ f^3\\ \frac{f^1}{\sqrt{2}}\hfill \end{array}\right);\left(\begin{array}{cccc}\frac{1+t_1^1}{2}\hfill & \frac{t_2^1}{\sqrt{2}}\hfill & \frac{t_3^1}{\sqrt{2}}\hfill & \frac{1t_1^1}{2}\hfill \\ \frac{t_1^2}{\sqrt{2}}\hfill & t_2^2& t_3^2& \frac{t_1^2}{\sqrt{2}}\hfill \\ \frac{t_1^3}{\sqrt{2}}\hfill & t_2^3& t_3^3& \frac{t_1^3}{\sqrt{2}}\hfill \\ \frac{1t_1^1}{2}\hfill & \frac{t_2^1}{\sqrt{2}}\hfill & \frac{t_3^1}{\sqrt{2}}\hfill & \frac{1+t_1^1}{2}\hfill \end{array}\right)),$$
whose connection form can be written as
$$[\left(\begin{array}{c}0\\ h_1\text{d}x\\ h_2\text{d}y\\ 0\end{array}\right),\left(\begin{array}{cccc}0& \frac{e^u}{\sqrt{2}}\text{d}x& \frac{e^u}{\sqrt{2}}\text{d}y& 0\\ \frac{e^u}{\sqrt{2}}\text{d}x& 0& u_y\text{d}xu_x\text{d}y& \frac{e^u}{\sqrt{2}}\text{d}x\\ \frac{e^u}{\sqrt{2}}\text{d}y& u_y\text{d}x+u_x\text{d}y& 0& \frac{e^u}{\sqrt{2}}\text{d}y\\ 0& \frac{e^u}{\sqrt{2}}\text{d}x& \frac{e^u}{\sqrt{2}}\text{d}y& 0\end{array}\right)]$$
for $`u,h_1`$, and $`h_2`$ smooth functions with $`h_1h_20`$ at each point.
In this setting $`\rho =me^{2u}`$, for a constant $`m`$, and the linear system becomes
$$\{\begin{array}{c}\text{d}v^1=\frac{e^ume^u}{\sqrt{2}}\text{d}xv^2+\frac{e^u+me^u}{\sqrt{2}}\text{d}yv^3\\ \\ \text{d}v^2=\frac{e^u}{\sqrt{2}}\text{d}xv^1(u_y\text{d}xu_x\text{d}y)v^3+\frac{e^ume^u}{\sqrt{2}}\text{d}xv^4\\ \\ \text{d}v^3=\frac{e^u}{\sqrt{2}}\text{d}yv^1+(u_y\text{d}xu_x\text{d}y)v^2+\frac{e^u+me^u}{\sqrt{2}}\text{d}yv^4\\ \\ \text{d}v^4=\frac{e^u}{\sqrt{2}}\text{d}xv^1\frac{e^u}{\sqrt{2}}\text{d}yv^2\end{array}.$$
$`10`$
Let $`{}_{}{}^{t}(v^1,v^2,v^3,v^4):M^+`$ be a solution to (10) and let $`r`$ be a smooth function that locally integrates the closed 1-form $`\gamma =v^2h_1\text{d}x+v^3h_2\text{d}y`$; we will refer to the functions $`r,v^1,v^2,v^3,v^4`$ as transforming functions. We are in a position to state:
###### Theorem 4
Let $`f:M^3`$ be an $`L`$-isothermic immersion and let $`r,v^1,v^2,v^3,v^4`$ be a set of transforming functions. Then
$$\widehat{f}=f+\frac{\sqrt{2}r}{v^1+v^4}\frac{f_x\times f_y}{f_x\times f_y}\frac{rv^2}{v^4(v^1+v^4)}\frac{f_x}{f_x}\frac{rv^3}{v^4(v^1+v^4)}\frac{f_y}{f_y}.$$
is a Bianchi–Darboux transform of $`f`$. All Bianchi–Darboux transforms of $`f`$ arise this way.
###### Demonstration Proof
Let $`F=[e_0,e_1]`$ be the Legendrian lift of $`f`$. According to Theorem 3, its Bianchi–Darboux transform $`\widehat{F}=[e_0+\frac{r}{v^4}e_1,v^1e_1+v^2e_2+v^3e_3+v^4e_4]`$, that is
$$\widehat{F}=e_0+\frac{r}{v^4}e_1+\mu (v^1e_1+v^2e_2+v^3e_3+v^4e_4),$$
for a smooth function $`\mu :M`$. Now $`\widehat{F}`$ takes values in $`𝔼(3)`$ if and only if $`\widehat{F},ϵ_1+ϵ_4=0`$ if and only if $`\mu =r/v^4(v^1+v^4)`$. Substituting and using the above realization of $`^3`$ in $`_1^4`$, we obtain the required expression for $`\widehat{f}`$. ∎
## 6. Superposition and permutability
Let $`A^{(1)},A^{(2)}:ML`$ be two curved flat framings having the same $`L`$-isothermic map $`F=[a_0^{(1)},a_1^{(1)}]=[a_0^{(2)},a_1^{(2)}]:M\mathrm{\Lambda }`$ as first envelope, and with second envelopes $`F^{(1)}`$ and $`F^{(2)}`$, respectively. Let $`\alpha ^{(1)}`$, $`\alpha ^{(2)}`$ be the corresponding connection forms.
The deformed forms associated with $`\alpha ^{(1)}`$ are given by
$$\alpha _\lambda ^{(1)}=[\left(\begin{array}{c}0\\ \alpha _0^2\\ \alpha _0^3\\ 0\end{array}\right),\left(\begin{array}{cccc}0& \lambda \alpha _2^1& \lambda \alpha _3^1& 0\\ \alpha _1^2& 0& \alpha _2^3& \lambda \alpha _2^1\\ \alpha _1^3& \alpha _2^3& 0& \lambda \alpha _3^1\\ 0& \alpha _1^2& \alpha _1^3& 0\end{array}\right)].$$
for some constant $`\lambda `$. Similarly for $`\alpha ^{(2)}`$.
According to the discussion in Section 4, the framings $`A^{(1)},A^{(2)}`$ are related by a gauge change
$$A^{(2)}=A^{(1)}g(s,v),$$
$`11`$
where $`v={}_{}{}^{t}(v^1,\mathrm{},v^4):M^+`$ is a solution to the integrable linear system
$$\text{d}v=\alpha _{\lambda }^{(1)}{}_{}{}^{}v,$$
$`12`$
for some $`\lambda `$, and $`s=\frac{r}{v^4}`$ with $`\text{d}rv^2\alpha _0^2v^3\alpha _0^3=0`$. Next, define the mapping $`g_\lambda ^{}(s,v):ML`$ by
$$g_\lambda ^{}(s,v)=[\left(\begin{array}{c}0\\ 0\\ 0\\ \frac{\lambda r}{v^1}\end{array}\right),\left(\begin{array}{cccc}\lambda ^1v^1& 0& 0& 0\\ v^2& 1& 0& 0\\ v^3& 0& 1& 0\\ \lambda v^4& \frac{\lambda v^2}{v^1}& \frac{\lambda v^3}{v^1}& \frac{\lambda }{v^1}\end{array}\right)].$$
Consider $`\overline{A}^{(1)}:=A^{(1)}g_\lambda ^{}(s,v)`$. The corresponding connection form $`\overline{\alpha }^{(1)}=(\theta ,\eta )`$ takes the form
$`{}_{}{}^{t}\theta =(0,\alpha _0^2+{\displaystyle \frac{\lambda r}{v^1}}\alpha _2^1,\alpha _0^3+{\displaystyle \frac{\lambda r}{v^1}}\alpha _3^1,0),`$
$`\eta =\left(\begin{array}{cccc}0& \frac{\lambda }{v^1}\alpha _2^1& \frac{\lambda }{v^1}\alpha _3^1& 0\\ v^1(\frac{1}{\lambda }1)\alpha _1^2& 0& \alpha _2^3+\frac{\lambda v^3}{v^1}\alpha _2^1\frac{\lambda v^2}{v^1}\alpha _3^1& \frac{\lambda }{v^1}\alpha _2^1\\ v^1(\frac{1}{\lambda }1)\alpha _1^3& \alpha _2^3+\frac{\lambda v^3}{v^1}\alpha _2^1\frac{\lambda v^2}{v^1}\alpha _3^1& 0& \frac{\lambda }{v^1}\alpha _3^1\\ 0& v^1(\frac{1}{\lambda }1)\alpha _1^2& v^1(\frac{1}{\lambda }1)\alpha _1^3& 0\end{array}\right).`$
We then have
###### Lemma
$`\overline{A}^{(1)}:=A^{(1)}g_\lambda ^{}(s,v)`$ is a curved flat framing such that $`[\overline{a}_{}^{(1)}{}_{0}{}^{},\overline{a}_{}^{(1)}{}_{4}{}^{}]=F^{(1)}`$.
Let $`F^{}=[\overline{a}_{}^{(1)}{}_{0}{}^{},\overline{a}_{}^{(1)}{}_{1}{}^{}]`$ be the first envelope of the congruence $`\overline{a}_{}^{(1)}{}_{0}{}^{}`$. We say that the $`L`$-isothermic map $`F^{}`$ is the superposition of $`F^{(1)}`$ and $`F^{(2)}`$ and write
$$F^{}=F^{(1)}_FF^{(2)}.$$
###### Theorem 5 (Permutability Theorem)
If an $`L`$-isothermic immersion $`F`$ has two Bianchi–Darboux transforms $`F^{(1)}`$ and $`F^{(2)}`$, then there is another $`L`$-isothermic immersion $`F^{}`$ which is a Bianchi–Darboux transform of $`F^{(1)}`$ and $`F^{(2)}`$ and is such that
$$F^{}=F^{(1)}_FF^{(2)}=F^{(2)}_FF^{(1)}.$$
###### Demonstration Proof
Write
$$A^{(1)}=A^{(2)}g(s,v)^1=A^{(2)}g(\widehat{s},\widehat{v}),$$
where
$$\widehat{s}=s,\widehat{v}={}_{}{}^{t}(v^1,\frac{v^2}{v^4},\frac{v^3}{v^4},\frac{1}{v^4}).$$
$`13`$
By a direct calculation it is verified that, if $`v:𝒰^+`$ is a solution to (12), then $`\widehat{v}`$ is a solution of
$$\text{d}\widehat{v}=\alpha _{\mu }^{(2)}{}_{}{}^{}\widehat{v},$$
where $`\mu `$ is given by
$$\mu =\lambda (\mu 1).$$
$`14`$
Thus, $`F^{(2)}_FF^{(1)}`$ is the $`L`$-isothermic immersion represented by the first envelope corresponding to $`\overline{A}^{(2)}=A^{(2)}g_\mu ^{}(\widehat{s},\widehat{v})`$. It is now easily seen that, if $`\widehat{v}`$, $`\widehat{s}`$, and $`\mu `$ are related to $`v`$, $`s`$, and $`\lambda `$ as in (13) and (14), then $`[\overline{a}_{}^{(2)}{}_{0}{}^{},\overline{a}_{}^{(2)}{}_{1}{}^{}]=[\overline{a}_{}^{(1)}{}_{0}{}^{},\overline{a}_{}^{(1)}{}_{1}{}^{}]`$. The situation is visualized in the following diagram:
$$\begin{array}{ccc}F& \stackrel{A^{(2)}}{}& F^{(2)}\\ A^{\left(1\right)}& & \overline{A}^{\left(2\right)}& & \\ F^{(1)}& \stackrel{\overline{A}^{(1)}}{}& F^{(1)}_FF^{(2)}\end{array}.$$ |
warning/0001/math0001058.html | ar5iv | text | # Any 3-manifold 1-dominates at most finitely many geometric 3-manifolds
support: both authors are supported by Outstanding Youth Grants of NNSFC and Grants of QSSTF.
§1 Introduction.
Maps between 3-manifolds has been studied by many people long times ago, and become an active subject again after Thurston’s revolution on 3-manifold theory. We refer to \[BW\], \[LWZ1\] for various results and references on the subject.
This paper addresses the following natural question which was raised around 1990, see also Kirby’s Problem List, \[K, 3.100\].
###### Question 1
Let $`M`$ be a closed orientable $`3`$-manifold. Are there at most finitely many closed, irreducible and orientable $`3`$-manifolds $`N`$ such that there exists a degree one map $`f:MN`$?
Remarks on the conditions in Question 1.
(i) If Poincare Conjecture fails, i.e., there is a homotopy 3-sphere $`N`$ which is not $`S^3`$, then one can get infinitely many reducible homotopy 3-spheres by doing connected sums on $`N`$. Since there always exists degree one map from a 3-manifold $`M`$ to a homotopy 3-sphere, the condition “irreducible” on the target $`N`$ is posed to avoid this unclear case.
(ii) The condition “closed” is posed on $`M`$ and $`N`$ just for simplicity. Indeed we can replace “closed” by “compact”, and meanwhile replace “degree one map” by “degree one proper map”. A map $`f:MN`$ between compact manifolds is proper if $`f^1(N)=M`$.
For simplicity, we adapt the following definition from \[BW\]. Let $`M`$ and $`N`$ be two compact orientable 3-manifolds. Say $`M`$ ($`1`$-)dominates $`N`$ if there is a proper map $`f:MN`$ of non-zero degree (degree $`1`$).
A closed orientable 3-manifold is called geometric if it admits one of the following geometries: $`H^3`$ (hyperbolic), $`\stackrel{~}{PSL_2()}`$, $`H^2\times E^1`$, Sol, Nil, $`E^3`$, $`S^2\times E^1`$, $`S^3`$ (spherical). Thurston’s geometrization conjecture claims that any closed, irreducible, and orientable 3-manifolds is either geometric or can be decomposed by the Jaco-Shalen-Johannson torus decomposition so that each piece is geometric. (For details see \[Th2\], \[Th3\] or \[Sc\].) All geometric 3-manifolds are precisely the Seifert manifolds except those carry hyperbolic geometry or Sol geometry, and all geometric manifolds have infinite fundamental group except those carry spherical geometry.
It is natural to study Question 1 when the targets are geometric first. There are many partial results of Question 1:
(i) The answer is affirmative if both the domain and the target are Seifert manifolds with infinite fundamental group \[Ro\], which is based on Waldhausen’s 3-manifold topology argument.
(ii) The answer is affirmative if the domain is non-Haken and the target is geometric \[RW\], which is based on Culler-Shalen’s character variety theory of 3-manifold groups. Since the domain is non-Haken, then the geometry of the target must be either hyperbolic or spherical.
Note also that there are additional conditions posed on the domains in (i) and (ii). Two substantial result to the Question 1 are obtained recently, where no additional conditions are posed on the domains.
(iii) The answer is affirmative if the targets are hyperbolic \[So\], which is based on the argument of Thurston’s original approach on the deformation of acylindrical manifolds.
(iv) The answer is affirmative if the target are spherical \[LWZ2\], which is based on the old knowledges of linking pair of 3-manifolds and of combinatorial groups.
In this paper we will prove the affirmative answer to Question 1 when the targets are all the remaining geometric 3-manifolds. The main result of this paper is the following.
###### Theorem 1
Any orientable closed $`3`$-manifold $`M`$ $`1`$-dominates at most finitely many closed orientable 3-manifolds which are either Seifert manifolds with infinite fundamental group or Sol manifolds.
Combining Theorem 1 with the results of \[So\] and \[LWZ2\] we obtain the following assertion.
###### Corollary 1
Any closed orientable $`3`$-manifold $`1`$-dominates at most finitely many geometric $`3`$-manifolds.
If an irreducible 3-manifold has non-trivial JSJ-decomposition, then each decomposition piece is either a hyperbolic 3-manifold or a Seifert manifolds with torus boundary. From the proof of Theorem 1, we have the following corollary, which should be useful in the discussion of non-trivial JSJ-torus decomposition case.
###### Corollary 2
Any compact orientable $`3`$-manifold $`M`$ dominates at most finitely many Seifert manifolds with non-empty boundary or zero Euler number.
In Section 2, we first explain that in the proof of the main result, one need only to deal with Seifert manifolds with orientable orbifold base and the torus bundle over the circle with Anosov monodromy. Then we present various known results about degree one map, Seifert manifolds, Thurston norm and volume of representations, including the brief descriptions of Thurston norm and of volume of representations. Those results will be used in the proof of the main results. The proof of the main result is given Section 3.
§2. Reductions and preliminary results.
Each Seifert manifold has an orbifold base which is either orientable or non-orientable.
###### Lemma 1
If there is a closed orientable 3-manifold 1-dominates infinitely many closed orientable Seifert manifolds with non-orientable orbifold base, then there is a closed orientable 3-manifold 1-dominates infinitely many closed orientable Seifert manifolds with orientable orbifold bases.
###### Demonstration Proof
Suppose $`f_j:MN_j`$ is degree one map for all $`j`$, $`p_j:N_jO_j`$ is the projection from the closed orientable Seifert manifolds onto the non-orientable orbifold base $`O_j`$, and $`N_i`$ and $`N_j`$ are not homeomorphic if $`ij`$. Let $`\stackrel{~}{q}_j:\stackrel{~}{O}_jO_j`$ be the unique orientable double cover of $`O_j`$ and $`q_j:\stackrel{~}{N}_jN_j`$ be the double covering which covers $`\stackrel{~}{q}_j`$. Then $`\stackrel{~}{N}_j`$ is a closed orientable Seifert manifold with orientable orbifold base $`\stackrel{~}{O}_j`$. Since $`f_j:MN_j`$ is of degree one, $`f_j:\pi _1(M)\pi _1(N_j)`$ is onto. This implies that the index of $`f_{}^1(\pi _1(\stackrel{~}{N}_j))`$ in $`\pi _1(M)`$ is two. Let $`\stackrel{~}{M}_j`$ be the double cover of $`M`$ corresponding to the subgroup $`f_{}^1(\pi _1(\stackrel{~}{N}_j))`$. Then $`f_j:MN_j`$ can be covered by a degree one map $`\stackrel{~}{f}_j:\stackrel{~}{M}_j\stackrel{~}{N}_j`$. Since any finitely presented group has only finitely many subgroup of given index, $`\pi _1(M)`$ has only finitely many subgroup of index 2. It follows that there are only finitely many homeomorphic types among $`\{\stackrel{~}{M}_j;j\}`$. By passing to a subsequence, we may assume all $`\stackrel{~}{M}_j=\stackrel{~}{M}`$ and we have degree one map $`f_j:\stackrel{~}{M}\stackrel{~}{N}_j,j`$. Since any double covering is a regular covering, each orientable Seifert manifold double covers at most finitely many Seifert manifolds by \[MS\]. It follows the homeomorphic types of $`\{\stackrel{~}{N}_j\}`$ are infinite. ∎
Each Sol manifold is either a tours bundle over the circle or a union of two twisted $`I`$-bundle over Klein bottle.
###### Lemma 2
If there is a closed orientable 3-manifold 1-dominates infinitely many Sol manifold which are unions of two twisted $`I`$-bundle over Klein bottle, then there is a closed orientable 3-manifolds 1-dominates infinitely many Sol manifolds which are torus bundle over $`S^1`$.
###### Demonstration Proof
Since each union of twisted $`I`$-bundle over Klein bottle is double covered by a torus bundle over the circle, the rest of the proof is the exactly same as that we did in the proof of Lemma 1. ∎
Let $`N`$ be an orientable Seifert fibered space with orientable orbifold base $`F_g`$ with $`n`$ exceptional fibers. Then $`N`$ has the standard form $`(g,b;a_1,b_1;a_2,b_2;\mathrm{}..;a_n,b_n)`$, $`a_i>b_i>0`$. There are two invariants associated with $`N`$: the Euler characteristic of the orbifold base
$$\chi _N=22g\underset{i=1}{\overset{n}{}}(1\frac{1}{a_i}),$$
and the Euler number of $`N`$
$$e(N)=b\underset{i=1}{\overset{n}{}}\frac{b_i}{a_i}.$$
We now give a brief description for the volume of representation (see \[Re\], \[Gr\], \[Th3\] for more details). Let $`G`$ be a semisimple Lie group and $`X=G/K`$, where $`K`$ is the maximum compact subgroup of $`G`$. For any orientable closed manifold $`M`$ and any representation $`\varphi :\pi _1(M)G`$, there is a flat $`X`$-bundle over $`M`$, $`M\times _\varphi X=\frac{\stackrel{~}{M}\times X}{\pi _1(M)}`$, with structure group $`G`$, where $`\stackrel{~}{M}`$ is the universal cover of $`M`$, $`\pi _1(M)`$ acting on the first factor $`\stackrel{~}{M}`$ by covering transformations, and by $`\varphi `$ on the second factor $`X`$. For simplicity, we assume that $`dimX=dimM=3`$ and $`X`$ is contractible. Let $`\omega ^{}`$ be the $`G`$-invariant volume form on $`X`$, which is a closed 3-form. Let $`q:\stackrel{~}{M}\times X`$ be the projection to the second factor. Then $`q^{}(\omega ^{})`$ is a $`\pi _1(M)`$-invariant closed 3-form on $`\stackrel{~}{M}\times X`$, and which induces a 3-form $`\omega `$ on $`M\times _\varphi X`$. Let $`s:MM\times _\varphi X`$ be a section. (Since $`X`$ is contractible, such a section exists and all such sections are homotopic.) We call $`_Ms^{}(\omega )`$ the volume of the representation $`\varphi `$, denoted by $`Vol(\varphi )`$, clearly it is independent of the choice of the section $`s`$. Define
$$Vol_G(M)=\text{max}\{|Vol(\varphi )|;\varphi :\pi _1(M)G\}.$$
Note if some $`\varphi :\pi _1(M)G`$ is discrete and faithful, then $`M`$ support the geometry of $`(G,X)`$ and $`Vol_G(M)=Vol(\varphi )`$. We get the famous Gromov norm in the case $`(G,X)=(PSL_2(),H^3)`$, and we are interested the case $`(G,X)=(PSL_2(),\stackrel{~}{PSL_2()})`$ in this paper. For short we use $`SV(M)`$ to denote $`Vol_{PSL_2()}(M)`$.
###### Lemma 3
Let $`M`$ and $`N`$ be closed orientable $`3`$-manifolds. If $`f:MN`$ is a degree one map, then
(1) $`TorH_1(M,)=ATorH_1(N,)`$, where $`TorH_1`$ is the torsion part of $`H_1`$.
If $`f:MN`$ is a map of degree $`d0`$, then
(2) $`SV(M)dSV(N)`$.
(3) $`[\pi _1(N):f_{}(\pi _1(M))]|d.`$
###### Demonstration Proof
For (1) see \[Br, 1.2.5 Theorem\]. For (2) see \[BG\] or \[Re\]. (3) is well-known and can be obtained directly by applying covering space argument. ∎
Let $`N`$ be a Seifert manifold with the standard form
$$(g,b;a_1,b_1;a_2,b_2;\mathrm{}\mathrm{};a_n,b_n),a_i>b_i>0.$$
###### Lemma 4
(1) $`N`$ supporting the geometry of either $`\stackrel{~}{PSL_2}`$, or Nil, or $`H^2\times E^1`$ is characterized by either $`e(N)0`$ and $`\chi _N<0`$, or $`e(N)0`$ and $`\chi _N=0`$, or $`e(N)=0`$ and $`\chi _N<0`$ respectively.
(2) If $`e(N)0`$, then the order of the torsion part of $`H_1(M,)`$,
$$|TorH_1(N,)|=\left|\left(\underset{i=1}{\overset{n}{}}a_i\right)\left(b+\underset{i=1}{\overset{n}{}}\frac{b_i}{a_i}\right)\right|=\left|e(N)\underset{i=1}{\overset{n}{}}a_i\right|.$$
(3) If $`N`$ supports the geometry of $`\stackrel{~}{PSL_2()}`$, Then
$$SV(N)=\left|\frac{\chi _N^2}{e(N)}\right|.$$
(4) The equation $`\chi _N=22g_{i=1}^n(1\frac{1}{a_i})=0`$ has only finitely many solutions $`(g,a_1,\mathrm{},a_n)`$.
(5) If $`\chi _N<0`$, then $`\chi _N\frac{1}{42}`$.
###### Demonstration Proof
For (1) see \[Sc\]. For (2) see \[LWZ1, 3.1\]. For (3) see \[BG\]. (4) and (5) are well-known and can be obtained by elementary algebra. ∎
Now we give a brief description on Thurston norm. In a closed oriented 3-manifold $`N`$, each element $`yH_2(N,)`$ can be represented by an embedded oriented surface $`F`$. Let $`\chi _{}(F)=\text{max}\{0,\chi (F)\}`$ if $`F`$ is connected, otherwise $`\chi _{}(F)=\chi _{}F_i`$, where $`F_i`$ are components of $`F`$. Then let
$$X(y)=min\{\chi _{}(F);F\text{ is an embedded surface representing }y\}.$$
Similarly we can define $`X_s(y)`$ if we replace “embedded surfaces” by “singular surfaces” in the definition of $`X`$ (see \[Th1\] for details). $`X`$ and $`X_s`$ can be extended to the second homology $`H_2`$ with real coefficent and are often called Thurston norm and Thurston singular norm respectively.
###### Lemma 5
(1) $`X`$ is a pseudonorm on $`H_2(M,)`$, in particular $`mX(y)nX(z)X(my+nz)mX(y)+nX(z)`$.
(2) $`X=X_s`$.
###### Demonstration Proof
For (1) see \[Th1\], and for (2) see \[Ga\].∎
Recall that there are only finitely many 3-manifolds support the geometries of $`S^2\times E^1`$ and $`E^3`$, and a torus bundle over the circle is a Sol manifold if and only if the gluing map is Anosov. With Lemma 1 and Lemma 2, (1) of Lemma 4 to prove Theorem 1, we need only to prove the following
###### Proposition 1
Any orientable closed $`3`$-manifold $`M`$ $`1`$-dominates at most finitely many closed orientable 3-manifolds $`N_j`$, where $`N_j`$ belongs to one of the following classes:
(a) Seifert manifolds with orientable orbifold bases with Euler number $`e=0`$ and Euler characteristic $`\chi <0`$.
(b) Seifert manifolds with orientable orbifold bases with Euler number $`e0`$ and Euler characteristic $`\chi 0`$.
(c) torus bundle over the circle with Anosov monodromy.
§3. Proof of the Theorems.
In this section we will prove Proposition 1. Suppose contrarily that there is an orientable closed $`3`$-manifold $`M`$ $`1`$-dominates infinitely many 3-manifolds $`N_j`$, where $`N_j`$ is subject to the conditions in Theorem 1. By passing to a subsequence we may assume that all $`N_i`$’s belong to one of the following classes:
(a) Seifert manifolds with orientable orbifold bases with Euler number $`e=0`$ and Euler characteristic $`\chi <0`$.
(b) Seifert manifolds with orientable orbifold bases with Euler number $`e0`$ and Euler characteristic $`\chi 0`$.
(c) torus bundle over the circle with Anosov monodromy.
We will show that none of those three cases can happen.
Let $`f_j:MN_j`$ be a degree one map defining 1-domination. By (1) of Lemma 3, we have
$$|TorH_1(M,)||TorH_1(N_j,)|.$$
$`1`$
In the first two case, we have the Seifert manifold
$$N_j=(g_j,b_j;a_{j1},b_{j1};\mathrm{}\mathrm{}a_{jn_j},b_{jn_j}),a_{ji}>b_{ji}>0.$$
Since $`f_i:\pi _1(M)\pi _1(N_j)`$ is surjective by (3) of Lemma 3, the rank of $`\pi _1(N_j)`$ is at most the rank of $`\pi _1(M)`$. The rank $`\pi _1(N_j)`$ is at least $`2g_j+n_j2`$ \[BZ\], so $`g_j`$’s and $`n_j`$’s are bounded. Passing to a subsequence we may assume that $`g_j=g`$ and $`n_j=n`$ and we have
$$N_j=(g,b_j;a_{j1},b_{j1};\mathrm{}\mathrm{}a_{jn},b_{jn}),a_{ji}>b_{ji}>0$$
$`2`$
Below we use $`e_j`$ to denote $`e(N_j)`$ and $`\chi _j`$ to denote $`\chi _{N_j}`$.
###### Demonstration Proof of Case (a)
Each homology class $`y`$ of $`H_2(N_j,)`$ can be presented by an incompressible surface. Since $`N_j`$ is irreducible Seifert manifold, each incompressible surface is either a vertical torus (foliated by Seifert circles), or a horizontal surfaces (transversal to all Seifert circles) \[p. 109, J\]. Since $`e_j=0`$, $`N_j`$ admits horizontal surfaces.
Let $`O_j`$ be the orbifold base of $`N_j`$ and $`C_j`$ be a regular fiber of $`N_j`$. Suppose also that $`O_j`$, $`C_j`$ and $`N_j`$ are compatible oriented.
Let $`F_j`$ be the horizontal surface of $`N_j`$, and $`p_j:F_jO_j`$ is the branched covering. Then we have $`\chi (F_j)=|d|\times \chi _j<0`$, where $`d=deg(p_j)`$ equals to $`F_jC_j`$, the algebraic intersection number of $`F_j`$ and a regular Seifert fiber of $`N_j`$. Note that $`|F_jC_j|`$, the absolute value of algebraic intersection number, is precisely the geometric intersection number.
Suppose further $`F_j`$ is a minimal genus horizontal surface of $`N_j`$, thus $`F_j`$ is characterized by that $`|F_jC_j|>0`$ is minimal.
Let $`X_j`$ be the Thurston norm on $`H_2(N_j,)`$. Let $`V_j=\{yH_2(N_j,);X_j(y)=0\}`$, which is generated by vertical tori. Then $`V_j`$ is a subgroup of $`H_2(N_j,)`$.
###### Lemma 6
$`H_2(N_j,)=[F_j]+V_j`$.
###### Demonstration Proof
Pick any homology class $`yH_2(N_j,)`$. If $`X_j(y)=0`$, then $`yV_j`$. Suppose $`X_j(y)0`$. Let $`F`$ be a oriented horizontal surface representing $`y`$ with $`\chi (F)=X_j(y)`$. We may assume that the degree of $`p_j:FO_j`$ is positive (otherwise replace $`y`$ by $`y`$). Then $`(l+1)X_j([F_j])>X_j(y)lX_j([F_j])`$ for some positive integer $`l`$, that is
$$F_jC_j>(FlF_j)C_j0.$$
Since $`F_jC_j`$ is minimal among all positive intersections, we have $`(FlF_j)C_j`$ is zero, and therefore the minimal genus incompressible surface which represents $`[FlF_j]`$ must be a union of tori. That is $`[FlF_j]V_j`$ and $`y=[lF_j]+[FlF_j]`$.∎
###### Lemma 7
If $`f:MN`$ is a map of degree $`d0`$, then $`f_{}:H_{}(M,)H_{}(N,)`$ is surjective.
###### Demonstration Proof
Recall that Poincare duality $`P:H^{nq}(N,)H_q(N,)`$ is given by $`z^{nq}z^{nq}[N]`$, where $`[N]`$ ($`[M]`$) is the fundamental class class of $`H_n(N,)`$ ($`H_n(M,)`$) and we also have the formula
$$f_{}(f^{}z^p[M])=z^pf_{}[M]$$
$`3`$
for any map $`f:MN`$.
Let $`z_qH_q(N,)`$. Let $`y_q=\frac{1}{d}f^{}P^1(z_g)[M]`$. Then by (3) we have
$$f_{}(y_q)=f_{}(\frac{1}{d}f^{}P_N^1(z_q)[M])$$
$$=\frac{1}{d}P^1(z_q)f_{}[M]=\frac{1}{d}P^1(z_q)d[N]=z_q.\mathit{}$$
Let $`X_M`$ be the Thurston norm on $`H_2(M,)`$, and $`X_{sj}`$ be the Thurston singular norm on $`H_2(N_j,)`$. Let $`z_1,\mathrm{},z_m`$ be a basis of $`H_1(M,)`$ and $`S_i`$ be a surface representing $`z_i`$ with $`\chi _{}(S_i)=X_M(z_i)`$, for $`i=1,\mathrm{},n`$.
Let $`y_i=[f_j(S_i)]=l_{ji}[F_j]+v_{ij}`$, where $`v_{ij}V_j`$. By (1) of Lemma 5, we have
$$X_j(y_i)=X_j(l_{ji}[F_j]+v_{ji})l_{ij}X_j([F_j])X_j(v_{ji})=l_{ij}X_j([F_j])$$
$`4`$
Then by (2) of Lemma 5 and the definition of Thurston singular norm, we have
$$X_j(y_i)=X_{sj}(y_i)\chi _{}(S_i)=X_M(z_i)$$
$`5`$
Combine (4) and (5), we have
$$X_M(z_i)X_j(y_i)l_{ji}X_j([F_j]).$$
$`6`$
Let $`L=\text{max}\{X_M(z_i);i=1,\mathrm{},m\}`$. If $`X_j([F_j])>L`$, then (6) implies that $`l_{ji}=0`$, and therefore $`y_i=[f_j(S_i)]=v_{ij}`$. It follows that $`f_j(H_2(M,))V_j`$. It contradicts Lemma 7 that $`f_j:H_2(M,)H_2(N_j,)`$ is surjective.
So $`LX_j(F_j)>0`$. By passing to a subsequence we may assume that all $`X_j(F_j)`$ are the same, therefore all $`F_j`$’s have the same homeomorphic types, denoted by $`S`$.
Cutting $`N_j`$ along the horizontal surface $`S`$, we obtained an $`I`$-bundle over $`S`$, and therefore $`N_j`$ can be presented as a surface bundle over $`S^1`$ with fiber $`S`$ and monodromy $`g_j:SS`$. Since $`N_j`$ is a Seifert manifold, $`g_j`$ must be a periodic map \[VI. 26., J\]. However it is well-known that there are only finitely many periodic maps on the given surface $`S`$ up to conjugacy. Since any two conjugated gluing map provide the homeomrphic 3-manifolds, there are only finitely many homeomorphic types among all $`N_j`$’s. We reach a contradiction. We have proved that Case (a) cannot happen.∎
###### Demonstration Proof of Corollary 2
First note that Lemma 7 is stated for any degree $`d0`$, and is still true for proper maps between manifolds with non-empty boundaries. Then note that for manifolds with boundary, Thurston norm was established and Lemma 5 is still valid (\[Th1\] and \[Ga\]). Finally note that if $`N_j`$ is a Seifert manifold with boundary, then $`N_j`$ always contains a horizontal embedded surface. With those three facts. The proof of Corollary 2 is the same as the proof of Case (a). ∎
###### Demonstration Proof for Case (b)
By (2) of Lemma 4, we have
$$|TorH_1(N_j,)|=\left|\left(a_{ji}\right)\left(b_j+\frac{b_{ji}}{a_{ji}}\right)\right|=\left|e_j\underset{i=1}{\overset{n}{}}a_{ji}\right|$$
$`7`$
If all $`a_{ji}`$’s are uniformly bounded, then all $`b_{ji}`$’s are uniformly bounded. Since we assume that all $`N_j`$’s are in different homeomorphic type, we must have that $`b_j`$ is unbounded.
Since $`a_{ji}>1`$ and $`|\frac{b_{ji}}{a_{ji}}|n`$, we have
$$\left|\left(a_{ji}\right)\left(b_j+\frac{b_{ji}}{a_{ji}}\right)\right||b_jn|$$
$`8`$
By (1) we have $`|TorH_1(N_j,)|`$ is bounded for all $`j`$, and by (7) and (8) we have $`|TorH_1(N_j,)|`$ is unbounded. We reach a contradiction.
By (4) of Lemma 4, we have ruled out the situation that $`\chi =0`$. Below we assume that $`\chi _j<0`$, i.e., all $`N_j`$ support the $`\stackrel{~}{PSL_2}`$ geometry.
Now we assume that some $`a_{ji}`$ tends to infinite as $`j`$ tends to infinite up to a subsequence. Then $`|a_{ji}|`$ tends to infinite as $`j`$ tends to infinite. To be not contradicted with (1) and (7), We must have
$$|e_j|=\left|b_j+\frac{b_{ji}}{a_{ji}}\right|0\text{as }j\mathrm{}$$
$`9`$
Since $`\chi _j<0`$, we have $`\chi _j\frac{1}{42}`$ by (5) of Lemma 4, and then $`|\chi _j|\frac{1}{42}`$. Then by (5) of Lemma 4, we have
$$SV(N_j)=\left|\frac{\chi _j^2}{e_j}\right|\left|\frac{1}{(42)^2e_j}\right|,$$
which is tends to infinite as $`j`$ tends to infinite. But by (2) of Lemma 4 we have
$$SV(M)SV(N_j)$$
$`11`$
i.e. $`SV(N_j)`$ is uniformly bounded for all $`j`$. We reach a contradiction again. We have proved that Case (b) cannot happen. ∎
###### Demonstration Proof of Case (c)
Now $`N_j`$ is a torus $`T`$ bundle over the circle with Anosov map $`g_j`$, denoted as $`(T,g_j)`$ sometimes. Fix a basis of $`H_1(T,)`$. Let $`SL_2()`$ be the group of 2 by 2 invertible integer matrices. Let $`A_jSL_2()`$ presents $`g_j`$ under the chosen basis of $`H_1(T,)`$. Then $`A_j`$ has two real eigenvalues of $`\lambda _j`$ and $`\lambda _j^1`$, with $`|\lambda _j|>1`$.
Using HHN extension one can calculate directly that
$$TorH_1(N_j,)=\frac{H_1(T,)}{IA_j},$$
where $`I`$ is the unit of $`SL_2()`$, and the first Betti number of $`N_j`$ is 1.
Then by linear algebra we have
$$|TorH_1(N_j,)|=|IA_j|=|(1\lambda _j)(1\lambda _j^1)|=|(2(\lambda _j+\lambda _j^1)|$$
$`12`$
By (1), $`|TorH_1(N_j,)|`$ is uniformly bounded, then by (12), the absolute value of the trace $`A_i`$, $`|\lambda _j+\lambda _j^1|`$, is uniformly bounded, say by some constant $`k>0`$.
Let $`(T,g)`$ and $`(T,g^{})`$ be two torus bundles over the circle with Anosov maps $`g`$ and $`g^{}`$. Let $`A`$ and $`A^{}`$ are matrices associated with $`g`$ and $`g^{}`$ under the given basis of $`H_1(T,)`$. If $`A=BA^{}B^1`$, for some $`BSL_2()`$, then $`g=hg^{}h^1`$, is induced by $`h:TT`$ is a homeomorphism realizing $`B`$. It follows easily that $`(T,g)`$ and $`(T,g^{})`$ are homeomorphic.
Now a contradiction in this case will follows by the following lemma.
###### Lemma 8
There are only finitely many conjugacy classes in $`SL_2()`$ representing Anosov maps with the absolute values of the traces are bounded by $`k>0`$.
###### Demonstration Proof
Let $`A=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right).`$ Since $`A`$ represents an Anosov map, $`bc0`$. Suppose $`|a|>k`$. Since $`|a+d|k`$, we have $`|d|<2|a|`$, and then $`|ad|<2a^2`$. The fact $`adbc=1`$ implies that $`|bc|<2a^2+1`$. In particular, either $`|b|`$ or $`|c|`$ is at most $`\sqrt{2}|a|`$. If $`|b|\sqrt{2}a`$, let $`C=\left(\begin{array}{cc}1& 0\\ \pm 1& 1\end{array}\right)`$, where we chose 1 if $`ab>0`$ and $`1`$ if $`ab<0`$. Then $`CAC^1=\left(\begin{array}{cc}ab& \\ & \end{array}\right)`$ and we can make $`|a\pm b||a|`$. If $`|c|\sqrt{2}a`$, let $`C=\left(\begin{array}{cc}1& \pm 1\\ 0& 1\end{array}\right)`$, $`CAC^1=\left(\begin{array}{cc}ac& \\ & \end{array}\right)`$ and we can make $`|a\pm c||a|`$. This concludes that if $`|a|>k`$, we always can get $`A_1=\left(\begin{array}{cc}a_1& b_1\\ c_1& d_1\end{array}\right)`$ which is conjugate to $`A`$ in $`SL_2()`$ and $`|a_1|<|a|`$.
Therefore to prove the Lemma, we may assume that $`|a|k`$. Then similarly since $`|a+d|k`$ we have $`|d|2k`$, and then $`|ad|2k^2`$. Since $`adbc=1`$, $`|bc|2k^2+1`$. In particular all entries are bounded by $`2k^2+1`$. Clearly there are only finitely many such elements in $`SL_2()`$.
We have proved that Case (c) cannot happen. ∎
We have completed the proof of Proposition 1, and therefore the proof of Theorem 1.
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Peking University, Beijing, China
and
East China Normal University, Shanghai, China |
warning/0001/nucl-th0001065.html | ar5iv | text | # Nuclear Saturation with QCD sum rules
## I Introduction
The research for the exact equilibrium properties of the nuclear matter is an old problem of nuclear physics and it continues to have a great interest nowadays. Quantum Hadrodynamics (QHD) is one of the models with larger success to describe the properties of the infinite nuclear matter as well as those of finite nuclei. QHD describes the nucleon-nucleon interaction through the mesons exchange ($`\pi `$, $`\sigma `$, $`\omega `$, $`\rho `$, etc.). Several calculations of nuclear structure using QHD and its extensions were made with success in the explanation of experimental data . For the purposes of this work only the simplest QHD model will be considered.
The basic QHD-model that explains the nuclear matter includes the nucleon ($`\psi `$) coupled with sigma ($`\sigma `$) and omega ($`\omega `$) mesons. In spite of the pion be the principal component of the nucleon-nucleon interaction, it is not included because the nuclear matter is an isotropic system with parity conservation. Analyzing the model we see that the real part of the scalar term ($`\sigma `$-meson) is typically of the order of several hundred MeV attractive while the real part of the time component of the vector term ($`\omega `$-meson) is typically of the order of several hundred MeV repulsive. However the energies involved in problems of nuclear structure are only of a few tens of MeV. That order of energy is obtained on the QHD model due to a large cancellation among the scalar and vector pieces and this process of cancelation can be controlled through an appropriate choice of the coupling constants. However, these so called QHD constants are quite different from those ones given in Bonn potential or empirical data . So it is interesting to discuss the validity of these coupling constants.
The success of the Walecka model, which is based on Dirac’s equation, is due to the mutual cancellation between the large scalar and vector potentials. However in some works it is discussed that the composed nature of the nucleon can suppress the scalar optical potential. However, as it is argued in , it can be shown that due to the nucleon being immersed in the nuclear media, such suppression does not exist. This result makes possible the use of the Dirac phenomenology for composed particles. Besides, recently the effective field theory (EFT) has verified that the QHD models are consistent with the symmetries of the Quantum Chromodynamics (QCD), the correct theory of the strong interaction. This last fact motivates the mixing between QHD and QCD accomplished in this work.
It is known that QCD is the correct theory of the strong interaction. However, there is not a perturbative treatment for QCD at energies involved in nuclear matter problems, so it is interesting to use the nonperturbative method given by QCD sum rules that was introduced by Shifman, Vainshtein, and Zakharov in the late 1970’s . The method consists in describing a correlation function in terms of hadron degree of freedom as well as quarks degree of freedom. This last one can be written in an operator product expansion (OPE) where the perturbative part (Wilson coefficients) is separate from the nonperturbative (condensates). The power of this technique is in the fact that the nonperturbative part is the same for all problems.
The objective of this work is to get the results obtained from QCD sum rules in the media and use them on QHD equations to analyze the necessary constants to obtain the saturation point of the infinite nuclear matter.
## II The QHD Model
The model used to describe the nuclear matter is from Serot and Walecka and it includes nucleons ($`\psi `$) interacting with $`\sigma `$ and $`\omega `$ mesons. So that the Lagrangian density is given by
$``$ $`=`$ $`\overline{\psi }(\gamma ^\mu _\mu +g_s\sigma g_v\gamma ^\mu \omega _\mu M)\psi +`$ (2)
$`+{\displaystyle \frac{1}{2}}(_\mu \sigma ^\mu \sigma m_s^2\sigma ^2){\displaystyle \frac{1}{2}}\left(G_{\mu \nu }G^{\mu \nu }m_v^2\omega ^\mu \omega _\mu \right),`$
where $`g_s`$ and $`g_v`$ are the mesons coupling constants, $`G_{\mu \nu }=_\mu \omega _\nu _\nu \omega _\mu `$ is the strength tensor for the vector field and $`M`$, $`m_s`$ and $`m_v`$ are the nucleon, $`\sigma `$-meson and $`\omega `$-meson masses, respectively. The Green function for the nucleon is given by
$`G(k)`$ $`=`$ $`(\gamma _\mu \stackrel{~}{k}^\mu M^{})\{{\displaystyle \frac{1}{\stackrel{~}{k}^2M^2+iϵ}}+`$ (4)
$`+{\displaystyle \frac{i\pi }{E^{}(k)}}\delta (k^0E(k))\theta (k_F|𝐤|)\},`$
where it was defined
$`\stackrel{~}{k}^\mu `$ $``$ $`k^\mu +\mathrm{\Sigma }_{(v)}^\mu ,`$ (5)
$`M^{}`$ $``$ $`M+\mathrm{\Sigma }_{(s)},`$ (6)
$`E^{}(k)`$ $``$ $`\sqrt{|\stackrel{~}{𝐤}|^2+M^2}.`$ (7)
Here $`\mathrm{\Sigma }_{(s)}`$ and $`\mathrm{\Sigma }_{(v)}^\mu `$ represent the scalar and vector self-energies, respectively.
In the more simple model only the contributions of the baryons from the Fermi sea will be considered. This is equivalent to disregard the contributions coming from antibarions (Dirac sea). So in agrement with the definition of energy-momentum tensor
$`\epsilon `$ $`=`$ $`\mathrm{\Psi }|\widehat{T}_{00}|\mathrm{\Psi }VEV`$ (8)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{m_s^2}{g_s^2}}\mathrm{\Sigma }_{(s)}^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{m_v^2}{g_v^2}}\mathrm{\Sigma }_{(v)}^2+{\displaystyle \frac{\gamma }{(2\pi )^3}}{\displaystyle _0^{k_F}}\text{d}^3k\sqrt{|𝐤|^2+M^2},`$ (9)
where $`VEV`$ is the vacuum expectation value of the $`\widehat{T}_{00}`$ and $`\gamma `$ is the spin-isospin degeneracy. The self-energies can be obtained through “tadpole” Feynman diagrams and are given by the following relationships
$`\mathrm{\Sigma }_{(s)}`$ $`=`$ $`M^{}M`$ (11)
$`=`$ $`{\displaystyle \frac{g_s^2}{m_s^2}}{\displaystyle \frac{\gamma M^{}}{4\pi ^2}}\left[k_FE_F^{}M^2\mathrm{ln}\left({\displaystyle \frac{k_F+E_F^{}}{M^{}}}\right)\right]`$ (12)
and
$$\mathrm{\Sigma }_{(v)}^\mu =\delta ^{\mu 0}\frac{g_v^2}{m_v^2}\frac{\gamma }{6\pi ^2}k_F^3,$$
(13)
The expression (12) should be solved in a self-consistent way. Therefore all the ingredients are available to calculate the energy density. This result is also known as Mean Field Theory (MFT), because it can be obtained with a mean field approach for the meson fields.
A simple analysis shows us that the self-energies $`\mathrm{\Sigma }_{(s)}`$ and $`\mathrm{\Sigma }_{(v)}`$ are proportional to the scalar and vectorial densities, respectively. Then the form presented by Eq. (LABEL:e-mft) represents an expansion on powers of these densities. On the other hand, in agreement with the EFT concepts, the energy density is a functional and has an expansion in powers of the densities (scalar, vectorial, etc.) organized through Georgi’s naive dimensional analysis (NDA). So, the energy densities have the same form that the expansions obtained by EFT. Therefore, since QHD has a foundation on the symmetries of QCD, the next step is to calculate the self-energies through the QCD methods. However, as the energy levels treated in the nuclear problems are in a nonperturbative regime for QCD, the sum rules is the appropriate method to be used.
## III QCD Sum Rules
The QCD sum rules method begins with the following time-ordered correlation function, defined by
$$\mathrm{\Pi }_{\alpha \beta }(q)i\text{d}^4x\text{ e}^{iqx}0|T\left[\eta _\alpha (x)\overline{\eta }_\beta (0)\right]|0,$$
(14)
where $`|0`$ is the physical nonperturbative vacuum state and $`\eta _\alpha (x)`$ is an interpolating field with the quantum numbers of a nucleon. As proposed by Ioffe , the proton field is given by
$$\eta (x)=ϵ_{abc}\left(u_a^TC\gamma _\mu u_b\right)\gamma _5\gamma ^\mu d_c,$$
(15)
where $`u`$ and $`d`$ are the quark fields, $`a,`$ $`b,`$ $`c`$ are the color indexes and $`C`$ is the charge-conjugation matrix. Now the correlation function can be written as an operator product expansion (OPE) whose nonperturbative part (condensates) can be separated from the perturbative one (Wilson coefficients). The OPE can be generated starting from the following expansion for the quark propagator
$`S_{ab}(x)`$ $`=`$ $`0|T\left[q_a(x)\overline{q}_b(0)\right]|0`$ (16)
$`=`$ $`i{\displaystyle \frac{\delta _{ab}}{2\pi ^2}}{\displaystyle \frac{\mathit{}}{x^4}}{\displaystyle \frac{\delta _{ab}}{12}}\overline{q}q_{vac}+\mathrm{},`$ (17)
where $`\overline{q}q_{vac}`$ is the quark condensate which can be determined from the Gell-Mann-Oakes-Renner relation,
$`2m_q\overline{q}q_{vac}=m_\pi ^2f_\pi ^2\left(1+𝒪(m_\pi ^2)\right).`$
In this relation, $`m_\pi =138`$ MeV is the pion mass, $`f_\pi =93`$ MeV is the pion decay constant and $`m_q=(m_u+m_d)/27\pm 2`$ MeV is the average of the up and down quark masses.
The same correlation function, Eq (14), can be written in a phenomenological way, which mean that the correlator has a hadronic description through the nucleon Green function, Eq (4). The match of the theoretical side (OPE) and the phenomenological one is the essence of the QCD sum rules. Using the fundamental state of the nuclear matter as vacuum and the interacting propagator in the phenomenological side, it is possible to obtain the following sum rules
$`M_N^{}`$ $`=`$ $`{\displaystyle \frac{8\pi ^2}{M^2}}\overline{q}q_{\rho _B},`$ (18)
$`\mathrm{\Sigma }_{(vQCD)}`$ $`=`$ $`{\displaystyle \frac{64\pi ^2}{3M^2}}q^{}q_{\rho _B}.`$ (19)
Here $`M^2`$ represent the borel mass with value near $`M^2=1`$GeV<sup>2</sup>, and
$`q^{}q_{\rho _B}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\rho _B,`$ (20)
$`\overline{q}q_{\rho _B}`$ $`=`$ $`\left(1{\displaystyle \frac{\sigma _B\rho _B}{m_\pi ^2f_\pi ^2}}+\mathrm{}\right)\overline{q}q_{vac}.`$ (21)
The expression (18) is a generalization of the Ioffe’s formula to finite density. The sigma term is estimate in Ref as $`\sigma _B45\pm 10`$ MeV. Taking the ratios of the expressions (18) and (19) in relation to the mass $`M_N`$ (Ioffe’s formula), they are obtained the following results:
$`\mathrm{\Sigma }_{(sQCD)}`$ $`=`$ $`{\displaystyle \frac{\sigma _BM_N}{m_\pi ^2f_\pi ^2}}\rho _B,`$ (22)
$`\mathrm{\Sigma }_{(vQCD)}`$ $`=`$ $`{\displaystyle \frac{8m_qM_N}{m_\pi ^2f_\pi ^2}}\rho _B.`$ (23)
Although a self-consistent relation is not obtained as in Eq (12), these expressions have a more fundamental nature based on quark and gluon degree of freedom. At this point, the idea is to use the expressions (22) and (23) to calculate the energy density of the nuclear matter.
## IV Numerical Results and Conclusion
Firstly it should be noted that only the ratios between the coupling constants and the mesons mass appear in the expressions where the energy density of the nuclear matter is calculated. Some values for these ratios are presented in the table I.
In the first line of table I the coupling constants were adjusted so that the MFT might reproduce the saturation properties of the nuclear matter. With the use of these coupling constants, the result obtained with MFT are the following (dotted curve on Fig. 1): saturation point on $`k_F=1.12`$ fm<sup>-1</sup> with a energy density by nucleon given by $`15.75`$ MeV. This result reproduces the exact equilibrium properties of the nuclear matter given by Refs , which is the expected result, once the coupling constants were chosen for such adjustment to occur. However, these coupling constants are very different of the empirical values .
The most accepted values for these constants in free space are found in the Bonn potential , for which the nucleon-nucleon interaction is adjusted to describe the scattering data. The Bonn coupling constants for the sigma and omega mesons are given on the second line of table I while the empirical values are presented on the third line of that table. It must be noted that, in MFT the value of the coupling constant for the scalar meson is higher than the respective value for the vector meson. But in the Bonn potential and in the empirical data there is an inversion in the magnitude of these constants. Using the Bonn’s values in Eq. (LABEL:e-mft), the result presented by the dashed curve on Fig. 1 is obtained. That curve shows us that the energy density for MFT does not present a good behavior. There is no saturation point, in other words, there is no formation of nuclear matter. The same result is obtained with the empirical constants. The problem is in the fact that the use of the Bonn or empirical coupling constants implicates in a small increase of the attraction but with a much larger increase in the repulsion. It can be argued that the difference between the QHD coupling constants and the Bonn values is due to the fact that the latter are obtained in free space while the QHD values are obtained in nuclear media, but this argument is not valid for empirical constants. However, using the constants given by Bonn potential, and the self-energies found through QCD sum rules, Eqs. (22) and (23), with $`\sigma _B=49`$ MeV and $`m_q=6.8`$ MeV, the result represented by the continuous line on Fig. 1 is obtained. That curve has a saturation point $`k_F=1.42`$ fm<sup>-1</sup> with $`15.75`$ MeV for the energy density by nucleon, which is the saturation point of the nuclear matter. Therefore the QCD sum rules allows us to reproduce the properties of nuclear matter with the simplest QHD model using more realistic values for the coupling constants.
Since the empiric constants are not so different from the Bonn values when they are compared to the QHD values, it can be that the effects of the density do not alter the values of these constants significantly. On the other hand, the exact determination of the parameters $`\sigma _B`$ and $`m_q`$ should allow a better evaluation of the coupling constants in the media for the various QHD models and consequently to determine with more precision their real variation in relation to the values of the constants in the vacuum. Furthermore this result is another indicative that the quark degree of freedom takes an important place in nuclear problems.
Finally, it is known that the QCD sum rules work with more accuracy when there is high transferred momentum. So I hope this mix between QCD sum rules and QHD model can be applied, with success, in systems in a regime of high density and high temperature as neutron stars. Besides, in these calculations the contributions of the gluon condensate and four quark condensate were not included. But it is known that the contribution of these terms for the sum rule is very small. However, if more precise answers are required, these terms have to be included as well as the most sophisticated versions of QHD. These and other questions are left for future works.
Acknowledgments
I would like to thank Prof. G. Krein for the discussions and suggestions and also FAPESP (Fundação de Amparo à Pesquisa do Estado de São Paulo) for the financial support. |
warning/0001/cond-mat0001307.html | ar5iv | text | # Pseudo gap in the density of states in cuprates
## Abstract
In the framework of the $`t`$-$`J`$ model for cuprates we analyze the development of a pseudo gap in the density of states (DOS), which at low doping starts to emerge for temperatures $`T<J`$ and persists up to the optimum doping. The analysis is based on numerical results for spectral functions obtained with the finite-temperature Lanczos method for finite two-dimensional clusters. We find that the pseudo gap scales with $`J`$ and is robust also in the presence of nearest neighbor repulsive interaction. Numerical results are additionally compared with the self consistent Born approximation (SCBA) results for hole-like (photoemission) and electron-like (inverse photoemission) spectra at $`T=0`$. The analysis is suggesting that the origin of the pseudo gap is in short-range antiferromagnetic (AFM) spin correlations and strong asymmetry between the hole and electron spectra in the underdoped regime.
In this paper we present the theoretical analysis of DOS in planar cuprates. As a prototype model we take the standard $`t`$-$`J`$ model, which incorporates strong electron correlations leading to AFM in undoped material and hindered motion of holes in doped system. The emphasis of the present study is on the pseudo gap found in recent angle-resolved photoemission (ARPES) experiments and also in some exact diagonalization studies . We add to the model also nearest neighbor repulsion $`V`$ term,
$`H=t{\displaystyle \underset{ijs}{}}(c_{js}^{}c_{is}+\mathrm{H}.\mathrm{c}.)+{\displaystyle \underset{ij}{}}[J𝐒_i𝐒_j+(V{\displaystyle \frac{J}{4}})n_in_j].`$
Here $`i,j`$ refer to planar sites on a square lattice and $`c_{is},c_{is}^{}`$ represent projected fermion operators forbidding double occupation of sites.
We study here the planar DOS, defined as $`𝒩(\omega )=2/N_𝐤A(𝐤,\omega \mu )`$, where $`A(𝐤,\omega )`$ is the electron spectral function , and $`\mu `$ denotes the chemical potential. First we calculate the DOS with the finite-temperature Lanczos method for clusters of $`N=18,20`$ sites doped with one hole, $`N_h=1`$. Note that $`𝒩^{}(\omega )`$ corresponds to adding a hole into the system and thus to the photoemission experiments, while $`𝒩^+(\omega )`$ represents the inverse photoemission (IPES) spectra.
In Fig. 1 we present $`𝒩(\omega )`$ for $`J/t=0.3,0.6`$ on a $`N=18`$ sites cluster. We note that the pseudo gap scales approximately as $`2J`$. The analysis at elevated temperatures shows that the gap slowly fills up and disappears at $`TJ`$. The gap remains robust also in the presence of the repulsive $`V`$ term, which on the other hand suppresses binding of hole pairs. Such an analysis thus suggests that the origin of the pseudo gap is in short-range AFM spin correlations rather than in the binding tendency of doped holes.
In Fig. 2(a) are shown spectra $`𝒩(\omega )`$ obtained on a $`N=\mathrm{\hspace{0.17em}20}`$ sites cluster. We compare these spectra with the DOS within the self-consistent Born approximation , obtained in the following manner. We assume that $`𝒩^{}(\omega )`$ can be approximated with the SCBA hole Green’s function for adding a hole to and antiferromagnetic reference system , Fig. 2(b). $`𝒩^+(\omega )`$ can be in SCBA correctly calculated and is presented in Fig. 2(c). The peaks in $`𝒩^+(\omega )`$ can well be explained with magnon structure of single hole ground state . As seen in Fig. 2(a) is the total DOS obtained with the SCBA (dashed lines) a reasonable approximation of numerical results.
We conclude stressing that the origin of the pseudo gap found in cuprates seems to be in the short range spin correlations of the reference AFM system, as well as in the strong asymmetry between the hole-like and electron-like spectra in underdoped systems. Namely, $`𝒩^+(\omega )`$ should scale linearly with doping $`c_h`$ but not changing substantially the width and form, while $`𝒩^{}(\omega )`$ away from chemical potential is less sensitive to $`c_h`$. Since $`\mu `$ lies in the pseudo gap, it is plausible that the pseudo gap observable in ARPES should also fill up with $`c_h`$, as found in experiments . |
warning/0001/hep-th0001134.html | ar5iv | text | # Untitled Document
IS QUANTUM SPACETIME INFINITE DIMENSIONAL ?
Carlos Castro
Center for Theoretical Studies of Physical Systems
Clark Atlanta University
Atlanta, GA. 30314
January 2000
Abstract
The Stringy Uncertainty relations, and corrections thereof, were explicitly derived recently from the New Relativity Principle that treats all dimensions and signatures on the same footing and which is based on the postulate that the Planck scale is the minimal length in Nature in the same vein that the speed of light was taken as the maximum velocity in Einstein’s theory of Special Relativity. A simple numerical argument is presented which suggests that Quantum Spacetime may very well be $`infinite`$ dimensional. A discussion of the repercusions of this new paradigm in Physics is given. A truly remarkably simple and plausible solution of the cosmological constant problem results from the New Relativity Principle : The cosmological constant is $`not`$ constant, in the same vein that Energy in Einstein’s Special Relativity is observer dependent. Finally, following El Naschie, we argue why the observed $`D=4`$ world might just be an $`average`$ dimension over the infinite possible values of the Quantum Spacetime and why the compactification mechanisms from higher to four dimensions in String theory may not be actually the right way to look at the world at Planck scales.
1. Preface
Before starting, we wish to say that readers already familiar with may skip the second section entirely. We deem it absolutely necessary to repeat the calculations that led us to the String Uncertainty Relations, corrections thereof, and the direct link between the Regge behaviour of string theory with the area quantization . In order to understand the main results of this work in section 3 one must follow closely section 2 . We apologize for having repeated the results of .
The New Relativity Principle that encompasses the ideas of Noncommutative C-spaces, , Polydimensional Covariance and Scale Relativity offers, in addition to the straightforward derivation of the String Uncertainty Relations, a truly remarkable and simple solution to the cosmological constant problem : the so called cosmological constant is not a constant. This is shown following the same arguments that Einstein gave when he showed that the energy was observer dependent : it is in the eye of the beholder. Finally, following El Naschie, we argue why the observed $`D=4`$ world might just be an $`average`$ dimension over the infinite possible values of the Quantum Spacetime and why the compactification mechanisms from higher to four dimensions in String theory may not be actually the right way to look at the world at Planck scales.
2. Introduction : String Uncertainty Relations from the New Relativity Principle
Recently we have proposed that a New Relativity principle may be operating in Nature which could reveal important clues to find the origins of $`M`$ theory . We were forced to introduce this new Relativity principle, where all dimensions and signatures of spacetime are on the same footing, to find a fully covariant formulation of the $`p`$-brane Quantum Mechanical Loop Wave equations. This New Relativity Principle, or the principle of Polydimensional Covariance as has been called by Pezzaglia, has also been crucial in the derivation of Papapetrou’s equations of motion of a spinning particle in curved spaces that was a long standing problem which lasted almost 50 years . A Clifford calculus was used where all the equations were written in terms of Clifford-valued multivector quantities; i.e one had to abandon the use of vectors and tensors and replace them by Clifford-algebra valued quantities, matrices, for example .
In this section we will explicitly derive the String Uncertainty Relations, and corrections thereof, directly from the Quantum Mechanical Wave equations on Noncommutative Clifford manifolds or C-spaces . There was a one-to-one correspondence between the nested hierarchy of point, loop, 2-loop, 3-loop,……p-loop histories encoded in terms of hypermatrices and wave equations written in terms of Clifford-algebra valued multivector quantities. This permits us to recast the QM wave equations associated with the hierarchy of nested p-loop histories, embedded in a target spacetime of $`D`$ dimensions , where the values of $`p`$ range from : $`p=0,1,2,3\mathrm{}\mathrm{}D1`$, as a $`single`$ QM line functional wave equation whose lines live in a Noncommutative Clifford manifold of $`2^D`$ dimensions. $`p=D1`$ is the the maximum value of $`p`$ that saturates the embedding spacetime dimension.
The line functional wave equation in the Clifford manifold, C-space is :
$$𝑑\mathrm{\Sigma }(\frac{\delta ^2}{\delta X(\mathrm{\Sigma })\delta X(\mathrm{\Sigma })}+^2)\mathrm{\Psi }[X(\mathrm{\Sigma })]=0.$$
$`(1)`$
where $`\mathrm{\Sigma }`$ is an invariant evolution parameter of $`l^D`$ dimensions generalizing the notion of the invariant proper time in Special Relativity linked to a massive point particle line ( path ) history :
$$(d\mathrm{\Sigma })^2=(d\mathrm{\Omega }_{p+1})^2+\mathrm{\Lambda }^{2p}(dx^\mu dx_\mu )+\mathrm{\Lambda }^{2(p1)}(d\sigma ^{\mu \nu }d\sigma _{\mu \nu })+\mathrm{\Lambda }^{2(p2)}(d\sigma ^{\mu \nu \rho }d\sigma _{\mu \nu \rho })+\mathrm{}\mathrm{}.$$
$`(2)`$
$`\mathrm{\Lambda }`$ is the Planck scale in $`D`$ dimensions. X$`(\mathrm{\Sigma })`$ is a Clifford-algebra valued ” line ” living in the Clifford manifold ( C-space) :
$$X=\mathrm{\Omega }_{p+1}+\mathrm{\Lambda }^px_\mu \gamma ^\mu +\mathrm{\Lambda }^{p1}\sigma _{\mu \nu }\gamma ^\mu \gamma ^\nu +\mathrm{\Lambda }^{p2}\sigma _{\mu \nu \rho }\gamma ^\mu \gamma ^\nu \gamma ^\rho +\mathrm{}\mathrm{}\mathrm{}$$
$`(3a)`$
The multivector X encodes in one single stroke the point history represented by the ordinary $`x_\mu `$ coordinates and the holographic projections of the nested family of 1-loop, 2-loop, 3-loop…p-loop histories onto the embedding coordinate spacetime planes given respectively by :
$$\sigma _{\mu \nu },\sigma _{\mu \nu \rho }\mathrm{}\mathrm{}\sigma _{\mu _1\mu _2\mathrm{}\mu _{p+1}}$$
$`(3b)`$
The scalar $`\mathrm{\Omega }_{p+1}`$ is the invariant proper $`p+1=D`$-volume associated with the motion of the ( maximal dimension ) p-loop across the $`D=p+1`$-dim target spacetime. There was a coincidence condition that required to equate the values of the center of mass coordinates $`x_\mu `$, for all the p -loops, with the values of the $`x^\mu `$ coordinates of the point particle path history. This was due to the fact that upon setting $`\mathrm{\Lambda }=0`$ all the p-loop histories collapse to a point history. The latter history is the baseline where one constructs the whole hierarchy. This also required a proportionality relationship :
$$\tau \frac{A}{\mathrm{\Lambda }}\frac{V}{\mathrm{\Lambda }^2}\mathrm{}\mathrm{}.\frac{\mathrm{\Omega }^{p+1}}{\mathrm{\Lambda }^p}.$$
$`(4)`$
$`\tau ,A,V\mathrm{}.\mathrm{\Omega }^{p+1}`$ represent the invariant proper time, proper area, proper volume,… proper $`p+1`$-dim volume swept by the point, loop, 2-loop, 3-loop,….. p-loop histories across their motion through the embedding spacetime, respectively. $`=T`$ is a quantity of dimension $`(mass)^{p+1}`$, the maximal $`p`$-brane tension ( $`p=D1`$) .
The wave functional $`\mathrm{\Psi }`$ is in general a Clifford-valued, hypercomplex number. In particular it could be a complex, quaternionic or octonionic valued quantity. At the moment we shall not dwell on the very subtle complications and battles associated with the quaternionic/octonionic extensions of Quantum Mechanics based on Division algebras and simply take the wave function to be a complex number. The line functional wave equation for lines living in the Clifford manifold ( C-spaces) are difficult to solve in general. To obtain the String Uncertainty Relations, and corrections thereof, one needs to simplify them. The most simple expression is to write the simplified wave equation in units $`\mathrm{}=c=1`$ :
$$[(\frac{^2}{x^\mu x_\mu }+\frac{\mathrm{\Lambda }^2}{2}\frac{^2}{\sigma ^{\mu \nu }\sigma _{\mu \nu }}+\frac{\mathrm{\Lambda }^4}{3!}\frac{^2}{\sigma ^{\mu \nu \rho }\sigma _{\mu \nu \rho }}+\mathrm{}\mathrm{})\mathrm{\Lambda }^{2p}^2]\mathrm{\Psi }[x^\mu ,\sigma ^{\mu \nu },\sigma ^{\mu \nu \rho },\mathrm{}..]=0$$
$`(5)`$
where we have dropped the first component of the Clifford multivector dependence, $`\mathrm{\Omega }^{p+1}`$, of the wave functional $`\mathrm{\Psi }`$ and we have replaced functional differential equations for ordinary differential equations. Had one kept the first component dependence $`\mathrm{\Omega }^{p+1}`$ on $`\mathrm{\Psi }`$ one would have had a cosmological constant contribution to the $``$ term as we will see below. Similar types of equations in a different context with only the first two terms of eq-(5), have also been written in .
The last equation contains the seeds of the String Uncertainty Relations and corrections thereof. Plane wave type solutions to eq-(5) are :
$$\mathrm{\Psi }=e^{i(k_\mu x^\mu +k_{\mu \nu }\sigma ^{\mu \nu }+k_{\mu \nu \rho }\sigma ^{\mu \nu \rho }+\mathrm{}\mathrm{}.)}.$$
$`(6)`$
where $`k_{\mu \nu },k_{\mu \nu \rho }\mathrm{}..`$ are the area-momentum, volume-momentum,….. $`p+1`$-volume-momentum conjugate variables to the holographic $`\sigma ^{\mu \nu },\sigma ^{\mu \nu \rho }\mathrm{}`$ coordinates respectively. These are the components of the Clifford-algebra valued $`multivector`$ K that admits an expansion into a family of antisymmetric tensors of arbitrary rank like the Clifford-algebra valued ”line” X did earlier in eq-(3a). The multivector K is nothing but the conjugate $`polymomentum`$ variable to X in C-space. Inserting the plane wave solution into the simplified wave equation yields the generalized dispersion relation, after reinserting the suitable powers of $`\mathrm{}`$ :
$$\mathrm{}^2(k^2+\frac{1}{2}\mathrm{\Lambda }^2(k_{\mu \nu })(k^{\mu \nu })+\frac{1}{3!}\mathrm{\Lambda }^4(k_{\mu \nu \rho })(k^{\mu \nu \rho })+\mathrm{}\mathrm{}..)\frac{\mathrm{\Lambda }^{2p}^2}{\mathrm{}^{2p}}=0.$$
$`(7)`$
this is just the generalization of the ordinary wave/particle dispersion relationship
$$p^2=\mathrm{}^2k^2.p^2m^2=0.$$
$`(8)`$
Had one included the $`\mathrm{\Omega }^{p+1}`$ dependence on $`\mathrm{\Psi }`$; i.e an extra piece $`exp[i\mathrm{\Omega }_{p+1}\lambda ]`$, where $`\lambda `$ is the cosmological constant of dimensions $`(mass)^{(p+1)}`$. The required $`\mathrm{\Lambda }^{2p}^2\mathrm{\Psi }/(\mathrm{\Omega }_{p+1})^2`$ term of the simplified wave equation (5) would have generated an extra term of the form $`\mathrm{\Lambda }^{2p}\lambda ^2`$. After reinserting the suitable powers of $`\mathrm{}`$, the cosmological constant term will precisely $`shift`$ the value of the $`\mathrm{\Lambda }^{2p}^2/\mathrm{}^{2p}`$ piece of eq-(7) to the value : $`(\frac{\mathrm{\Lambda }}{\mathrm{}})^{2p}(^2\lambda ^2)`$, which precisely has an overall dimension of $`m^2`$ as expected.
Hence, this will be then the ” vacuum ” contribution to $`maximal`$ $`p`$-brane tension ( $`p=D1`$) : $`=T_p`$ has overall units $`(mass)^{p+1}`$; i.e energy per $`p`$-dimensional volume. On dimensional grounds and due to the $`coincidence`$ condition referred above one has that :
$$(k_{\mu \nu })(k^{\mu \nu })=\beta _2(k^2)^2=\beta _2k^4.(k_{\mu \nu \rho })(k^{\mu \nu \rho })=\beta _3(k^3)^2=\beta _3k^6\mathrm{}\mathrm{}$$
$`(9)`$
where the proportionality factors in eq-(9) are the rank and dimension-dependent constants, $`\beta _2(D,r=2),\beta _3(D,r=3)\mathrm{}.`$ associated with the 2-vector, 3-vector,………components of the polymomentum K, respectively. $`\beta =1`$ for the first term in eq-(7), a rank one tensor : vector. The coincidence condition implies that upon setting $`\mathrm{\Lambda }=0`$ all the p-loop histories $`collapse`$ to a point history. In that case the areas, volumes, …hypervolumes collapse to $`zero`$ and the wave equation (5) reduces to the ordinary Klein-Gordon equation for a spin zero massive particle.
Factoring out the $`k^2`$ factor in (7), using the analog of the dispersion relation (8) and taking the square root, after performing the binomial/Taylor expansion of the square root, subject to the condition $`\mathrm{\Lambda }^2k^2<<1`$, one obtains an $`effective`$ energy dependent Planck ” constant ” that takes into account the Noncommutative nature of the Clifford manifold (C-space ) at Planck scales :
$$\mathrm{}_{eff}(k^2)=\mathrm{}(1+\frac{1}{2.2!}\beta _2\mathrm{\Lambda }^2k^2+\frac{1}{2.3!}\beta _3\mathrm{\Lambda }^4k^4+\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}.).$$
$`(10)`$
where we have included explicitly the $`D`$ and rank dependent coefficients $`\beta _1,\beta _2,\beta _3\mathrm{}`$ that arise in (9) due to the $`coincidence`$ condition and on dimensional analysis.
Arguments concerning an effective value of Planck’s ” constant ” related to higher derivative theories and the modified uncertainty relations have been given by . The advantage of this derivation based on the New Relativity Principle is that one automatically avoids the problems involving the ad hoc introduction of higher derivatives in Physics ( ghosts, …) .
The uncertainty relations for the coordinates-momenta follow from the Heisenberg-Weyl algebraic relation familiar in QM :
$$\mathrm{\Delta }x\mathrm{\Delta }p|<[\widehat{x},\widehat{p}]>|.[\widehat{x},\widehat{p}]=i\mathrm{}$$
$`(11)`$
Now we have that in C-spaces, $`x,p`$ must not, and should not, be interpreted as ordinary vectors of spacetime but as one of the many $`components`$ of the Clifford-algebra valued multivectors that ” coordinatize ” the Noncommutative Clifford Manifold, C-space. The Noncommutativity is $`encoded`$ in the $`effective`$ value of the Planck’s ” constant ” which $`modifies`$ the Heisenberg-Weyl $`x,p`$ algebraic commutation relations and, consequently, generates new uncertainty relations :
$$\mathrm{\Delta }x\mathrm{\Delta }p|<[\widehat{x},\widehat{p}]>|=<\mathrm{}_{eff}>=\mathrm{}(1+\frac{1}{2.2!}\beta _2\mathrm{\Lambda }^2<k^2>+\frac{1}{2.3!}\beta _3\mathrm{\Lambda }^4<k^4>+\mathrm{}\mathrm{}.)$$
$`(12)`$
Using the relations :
$$\mathrm{}k=p.<p^2>(\mathrm{\Delta }p)^2.<p^4>(\mathrm{\Delta }p)^4\mathrm{}..$$
$`(13)`$
one arrives at :
$$\mathrm{\Delta }x\mathrm{\Delta }p\mathrm{}+\frac{\beta _2\mathrm{\Lambda }^2}{4\mathrm{}}(\mathrm{\Delta }p)^2+\frac{\beta _3\mathrm{\Lambda }^4}{12\mathrm{}^3}(\mathrm{\Delta }p)^4+\mathrm{}\mathrm{}.$$
$`(14)`$
Finally, keeping the first two terms in the expansion in the r.h.s of eq- (14) one recovers the ordinary String Uncertainty Relation directly from the New Relativity Principle as promised :
$$\mathrm{\Delta }x\frac{\mathrm{}}{\mathrm{\Delta }p}+\frac{\beta _2\mathrm{\Lambda }^2}{4\mathrm{}}(\mathrm{\Delta }p).$$
$`(15)`$
which is just a reflection of the minimum distance condition in Nature and an inherent Noncommutative nature of the Clifford manifold ( C-space ). Eq-(15) yields a $`minimum`$ value of $`\mathrm{\Delta }x`$ of the order of the Planck length $`\mathrm{\Lambda }`$ that can be verified explicitly simply by minimizing eq-(15).
3. A Simple Argument Why Quantum Spacetime could be Infinite Dimensional
A Plausible Resolution of the Cosmological Constant Problem
So far the derivation of the String uncertainty relations from the New Relativity Principle has been straightforward. However, we wish to be more radical in our approach. An immediate question comes to mind :
$$Whydidwetruncatetheserieseqs(5,7)toafinitevalueoftheQuantumSpacetimedimension\mathrm{?}$$
If the New Relativity principle is true then we must include all dimensions for the Quantum Spacetime. It is all or nothing ! Taking this radical view will generate instead of the finite series of eq-(7) an $`infinite`$ series of the form :
$$\mathrm{}^2k^2\underset{r=1}{\overset{\mathrm{}}{}}\frac{\beta _r(r,D)}{r!}(k\mathrm{\Lambda })^{2(r1)}=lim_p\mathrm{}\frac{\mathrm{\Lambda }^{2p}(^2\lambda ^2)}{\mathrm{}^{2p}}.$$
$`(16)`$
where $`r=1,2,3,\mathrm{}\mathrm{}D`$ denotes the rank of the vector, 2-vector, 3-vector,…..associated with the Clifford-algebra valued polymomentum K conjugate to the Clifford-valued ” line ” in C-space : X$`(\mathrm{\Sigma })`$.
The sum of the infinite series depends on the infinite family of rank and dimension dependent coefficients $`\beta _r(D,r)`$ appearing in eq-(9,10,12). For simplicity purposes, and for the sake of the argument, we will take $`all`$ of the coefficients to have the simplest value of them all : 1. The infinite series yields :
$$\underset{r=1}{\overset{\mathrm{}}{}}\frac{(k\mathrm{\Lambda })^{2(r1)}}{r!}=\underset{r=1}{\overset{\mathrm{}}{}}\frac{z^{2(r1)}}{r!}=\underset{r^{}=0}{\overset{\mathrm{}}{}}\frac{z^{2r^{}}}{(r^{}+1)!}=\frac{e^{z^2}1}{z^2}.wherezk\mathrm{\Lambda }.$$
$`(17)`$
One recovers one of the confluent hypergeometric functions as the value of the sum. Therefore, after recasting the sum in terms of hyperbolic functions
$$\mathrm{}_{eff}^2=\mathrm{}^2\frac{e^{z^2}1}{z^2}\mathrm{}_{eff}^2=\mathrm{}^2e^{z^2/2}\frac{sinh(z^2/2)}{(z^2/2)}.$$
$`(18)`$
following the exact same steps as in the previous section one gets the full blown Uncertainty Relations for Quantum Spacetime due to the contributions of all extended objects : $`p=0,1,2,\mathrm{}..\mathrm{}`$ : :
$$\mathrm{\Delta }z\mathrm{\Lambda }\mathrm{\Delta }k\mathrm{\Delta }x\sqrt{2}\mathrm{\Lambda }\frac{e^{(\mathrm{\Delta }z)^2/4}}{(\mathrm{\Delta }z)^2}\sqrt{sinh[\frac{(\mathrm{\Delta }z)^2}{2}]}.$$
$`(19)`$
One can verify that the function :
$$x=x(z^2)\sqrt{2}\mathrm{\Lambda }\frac{e^{z^2/4}}{z^2}\sqrt{sinh(\frac{z^2}{2})}.$$
$`(20)`$
after differentiating it and equating it to zero, has a minimum/maximum for those values of $`z_o`$ such that satisfy :
$$tanh[z^2/2]=\frac{z^2}{4z^2}.$$
$`(21)`$
When $`z=0,\mathrm{}x=\mathrm{}`$ as expected in eq-(20), when the momentum is $`k=0,\mathrm{}`$. The minimum value of $`x`$ occurs when
$$1.2621<z_o<1.2626x_{min}1.2426\mathrm{\Lambda }.$$
$`(22)`$
Therefore, for momentum values precisely of the order of the Planck’s momentum : $`k_o(z_o/\mathrm{\Lambda })`$ that gives $`1.262k_P`$ , one reaches the minimum distance of $`1.2426\mathrm{\Lambda }`$ ! as it is required from the New Relativity Principle : Polydimensional Covariance and Scale Relativity .
Of course, one can always tune the infinite number of coefficients (16) in an infinite number of ways to reproduce the Planck scale as the minimum scale for arbitrary values of the momentum. A smaller subclass of infinite tuning possibilities appears when the Planck scale minimum occurs $`precisely`$ at Planck values of the momentum. It is remarkable that the simplicity arguments of setting all the values of the coefficients to 1 yields the desired results of attaining a minimum Planck scale for Planck values of the momentum. We believe this is not a numerical coincidence.
An immediate question arises : if there are many ways of selecting and tuning the coefficients $`\beta _r`$ in (16) to satisfy the requirements of attaining minimal Planck scale uncertainty at Planck scale momentum, i.e there is an infinite range of possible values for the $`[x,p]=i\mathrm{}_{eff}`$ commutation relations, is there a physical criteria to select a unique value of $`\mathrm{}`$ ?
We believe that the answer to this question may lie in Hopf algebraic structures at Planck scales . Recently there has been a lot of activity pertaining the Hopf algebraic structure underlying to perturbative QFT and the numerical ” miracles” of the Renormalization Group process. As Kreimer has pointed out, the iterated removal of nested divergences while maintining locality has to fulfill combinatorial properties summarized by Zimmermann’s forest formula. There is an underlying mathematical structure that is in no way accidental. For relations to low dimensional topology, number theory,…we refer to Kreimer et al . The authors have suggested that there is a Planck scale Hopf algebra as a particular example of a Noncommutative differential geometry at Planck scales where the Planck scale acts as a natural ultraviolet regulator . In fact, Majid found a $`[x,p]=i\mathrm{}_{eff}`$ commutation relation that bears a striking resemblance to ours. Thiemann has also argued that Planck scale should serve as natural regulator for matter QFT . Nottale has argued in numerous occasions that there is a deep link between the Renormalization Group process, that is based in scaling arguments, and Scale Relativity. In fact he has given a $`resolution`$ dependent effective Planck’s constant that Granik recently has shown to agree with eq-(10) up to the first terms . .
The confluent hypergeometric function that results after summing the infinite series (16) is no numerical accident. Gamma functions have long been known to be essential in the dimensional regularization procedures and in Veneziano’s formula that spawned String Theory. In addition, the $`crux`$ of including all dimensions in our calculation for the $`[x,p]`$ commutation relations is that one does not have to truncate/amputate the $`[x,p]`$ commutators to a $`finite`$ number of terms like it was done in . We have given in (19) the full blown Quantum Spacetime Uncertainty relations that are more general than the usual String Uncertainty relations; We are including the effects of all extended objects !
The main lesson from this numerical exercise is that Quantum Spacetime could be infinite dimensional if we invoke the New Relativity principle to the fullest potential within the context of Noncommutative Clifford manifolds, C-spaces and Quantum Groups ( Hopf algebras). This result that the Quantum Spacetime is infinite-dimensional has been advocated many times by within the context of Fractals, Scale Relativity and a Cantorian-Fractal spacetime : a transfinite infinite nested hierarchy of fractal Cantorian sets of infinite dimensionality. Quantum sets have been proposed long ago by Finkelstein in the formulation of Quantum Relativity . This straightforward numerical analysis is a strong indication that Quantum Spacetime could be infinite-dimensional and that it may indeed be fractal at its very core. Nature is Fractal. It is not a big surprise that Quantum Spacetime could be as well. Being fractal supports the view of Majid that Quantum Geometry is a Braided Categorical one. Since a Fractal Quantum Spacetime has fractal dimensions, it follows naturally that it should allow for fractional spins, charges, statistics,.. i.e The Quantum Geometric world has Braided Statistics.
Ordinary real numbers are no longer useful to describe the infinite dimensional Fractal Quantum Spacetime we are proposing. It is meaningless to assume that we can meausure a real number to infinite nonperiodic decimal places. It has been speculated for quite some time that due to the minimal Planck length, the geometry at Planck scales is Non-Archimedean. Therefore $`p`$-adic numbers are the natural ones to use at this scale. For a review of the mathematical applications of $`p`$-adic numbers in Physics and Fractals see . For the role of $`p`$-adics in the construction of TGD see .
If Fractal Quantum Spacetime is indeed $`infinite`$ dimensional we would have to drastically modify our naive perceptions that spacetime has a fixed dimension and reconsider the validity of the compactification arguments studied so far from higher to low dimensions. Dimensions are $`resolution`$ dependent . Instead we may be obliged to view $`D=4`$ only as an overall average dimension in the same way that the speed of the molecules inside a box at fixed temperature is distributed over a wide range of velocities and has an average one related to the temperature. El Naschie using a Gamma distribution for the ensemble of dimensions, fractal arguments and Astrophysical data results has obtained average dimensions close to $`D=4`$. The same ideas apply to the observed spacetime signature and to the resolution of the cosmological constant problem. The New Relativity principle treats all dimensions and signatures on the same footing.
To finalize we show why the cosmological ”constant” should not be treated as such : it is in the eye of the beholder . The key to a plausible and remarakaby simple solution to the cosmological ”constant” problem lies within eq-(16) which is nothing but an generalization of Einstein’s relation : $`E^2p^2=m^2`$. One simply shifts the cosmological constant term of (16) to the left hand side. Upon shifting it to the ”left” we have :
$$(\frac{\mathrm{\Lambda }}{\mathrm{}})^{2p}\lambda ^2+\mathrm{}^2k^2\underset{r=1}{\overset{\mathrm{}}{}}\frac{\beta _r(r,D)}{r!}(k\mathrm{\Lambda })^{2(r1)}=lim_p\mathrm{}\frac{\mathrm{\Lambda }^{2p}^2}{\mathrm{}^{2p}}.$$
$`(23)`$
The New Relativity principle which is based on the principles of Polydimensional Covariance which reshufle a string history for a $`5`$-brane history; a $`9`$-brane history for a $`5`$-brane history and so forth i.e the New Relativity principle is nothing but taking Chew’s bootstrap idea to the heart : each $`p`$-brane is made of all the others ! It is in this fashion why all dimensions ( and signatures) must be treated on the same footing. As Nottale and El Naschie have argued, dimensions are not absolute concepts in Quantum Spacetime, they are resolution dependent.
The New Relativity Principle ( Polydimensional Covariance ) states that the r.h.s (23) is truly an invariant ( like the proper time or proper rest mass of a particle) while the terms on the l.h.s are just the analogs of the squared-norm of a four-vector $`E^2p^2=m^2`$. Therefore, based on this simple analogy we propose that $`\lambda `$ is not a constant but instead is just one of the many observer dependent components of the polymomentum multivector K referred earlier in section 2. Hence we have that the combination :
$$(\frac{\mathrm{\Lambda }}{\mathrm{}})^{2p}\lambda ^2+\mathrm{}^2k^2\underset{r=1}{\overset{\mathrm{}}{}}\frac{\beta _r(r,D)}{r!}(k\mathrm{\Lambda })^{2(r1)}=(\frac{\mathrm{\Lambda }}{\mathrm{}})^{2p}\lambda ^2+\mathrm{}^2k^2\underset{r=1}{\overset{\mathrm{}}{}}\frac{\beta _r(r,D)}{r!}(k^{}\mathrm{\Lambda })^{2(r1)}=\mathrm{}.=\frac{\mathrm{\Lambda }^{2p}^2}{\mathrm{}^{2p}}.$$
$`(24)`$
is an invariant of this New Relativity Theory like
$$E^2p^2=E^2p^2=E^{\prime \prime 2}p^{\prime \prime 2}=\mathrm{}\mathrm{}.=m^2.$$
$`(25)`$
was in Special Relativity, where the squared of the maximal $`p`$-brane tension, associated with the spacetime filling $`p`$-brane, $`^2`$ plays identical role to the one played by $`m^2`$ in Einstein’s Relativity. Eq-(24) is remarkably simple. It relates the microscopic world quantities : Planck scale, cosmological ” constant ” on the left, with the total Quantum Spacetime Enery per Unit $`p`$-volume ( Elasticity of Spacetime quoting Zaharov ) associated with the Quantum Spacetime-filling maximal $`p`$-brane ( $`p+1=D`$ ) on the right. The essential terms required to match the left with the right are precisely provided by the infinite number of modes associated with the point-history, loop-history, 2-loop history, 3-loop history,…… p-loop history excitations OF the Quantum Spacetime given precisely by
$$\mathrm{}^2k^2\underset{r=1}{\overset{\mathrm{}}{}}\frac{\beta _r(r,D)}{r!}(k\mathrm{\Lambda })^{2(r1)}.$$
$`(26)`$
that led to the full blown Quantum Spacetime Uncertainty Relation (19). For interesting work on the cosmological constant we refer to . What remains is to find out what is the right Hopf Planck scale algebra that selects a unique value of the effective Planck ”constant ”. Perhaps there are several ?
Since in $`D=4`$ the Planck legth is given by :
$$\mathrm{\Lambda }=\sqrt{\frac{\mathrm{}G}{c^3}}.$$
$`(27)`$
the immediate question, similar to the one proposed long ago by Dirac, arises : isn’ it possible that $`\mathrm{},G,c`$ are not ”constants” in eq-(27) but they could vary in such a way as to leave the value of $`\mathrm{\Lambda }`$ invariant ??? In the past years there has been a lot of research activity in the Astrophysics community pondering if the speed of light varies in Cosmology . Nottale has given another explanation for the resolution of the cosmological constant problem based on Scale Relativity . His argument is essentially that it is meaningless to compare two values of energy densities at two different scales without including Scale-Relativistic effects. He explains why the $`10^{50},10^{60}`$ discrepancy is due to the Scale-Relativistic ”Lorentz” dilation factors.
It seems that one may be forced to demolish the old established ”idols” ( using Finkelstein’ terminoly ) of spacetime, dimension, cosmological constant,….in the same way that Relativity and Quantum Mechanics replaced the Cartesian-Newtonian paradigm. As we end the century, it is time perhaps to embrace a new paradigm in Physics that demolishes the concept of dimension as an idol. Number theory, Topology, Fractals, Cantor sets, $`p`$-adic analyis, QFT, Quantum Groups, Hopf algebras, Noncommutative Geometry…… seem all to be converging in disguised forms at the Planck scale upon looking at Quantum Spacetime with the magnifying glass of the New Relativity Theory based on Noncommuative C-spaces , the principle of Polydimensional Covariance and Scale Relativity : a magnifying glass lying deep inside the fuzzy crystal ball of imagination signaling what it may turn out to be a new paradigm in Physics.
Acknowledgements
We are indebted to E. Spallucci for a very constructive critical mail correspondence. To G. Chapline. A. Granik, L.Nottale , W. Pezzaglia, M. El Naschie and D. Finkelstein for discussions. Finally many thanks to C. Handy for his assistance and encouragement .
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warning/0001/cond-mat0001432.html | ar5iv | text | # Statistical mechanics of money
## I Introduction
The application of statistical physics methods to economics promises fresh insights into problems traditionally not associated with physics (see, for example, the recent review and book ). Both statistical mechanics and economics study big ensembles: collections of atoms or economic agents, respectively. The fundamental law of equilibrium statistical mechanics is the Boltzmann-Gibbs law, which states that the probability distribution of energy $`\epsilon `$ is $`P(\epsilon )=Ce^{\epsilon /T}`$, where $`T`$ is the temperature, and $`C`$ is a normalizing constant . The main ingredient that is essential for the textbook derivation of the Boltzmann-Gibbs law is the conservation of energy . Thus, one may generalize that any conserved quantity in a big statistical system should have an exponential probability distribution in equilibrium.
An example of such an unconventional Boltzmann-Gibbs law is the probability distribution of forces experienced by the beads in a cylinder pressed with an external force . Because the system is at rest, the total force along the cylinder axis experienced by each layer of granules is constant and is randomly distributed among the individual beads. Thus the conditions are satisfied for the applicability of the Boltzmann-Gibbs law to the force, rather than energy, and it was indeed found experimentally .
We claim that, in a closed economic system, the total amount of money is conserved. Thus the equilibrium probability distribution of money $`P(m)`$ should follow the Boltzmann-Gibbs law $`P(m)=Ce^{m/T}`$. Here $`m`$ is money, and $`T`$ is an effective temperature equal to the average amount of money per economic agent. The conservation law of money reflects their fundamental property that, unlike material wealth, money (more precisely the fiat, “paper” money) is not allowed to be manufactured by regular economic agents, but can only be transferred between agents. Our approach here is very similar to that of Ispolatov et al. . However, they considered only models with broken time-reversal symmetry, for which the Boltzmann-Gibbs law typically does not hold. The role of time-reversal symmetry and deviations from the Boltzmann-Gibbs law are discussed in detail in Sec. VII.
It is tempting to identify the money distribution $`P(m)`$ with the distribution of wealth . However, money is only one part of wealth, the other part being material wealth. Material products have no conservation law: They can be manufactured, destroyed, consumed, etc. Moreover, the monetary value of a material product (the price) is not constant. The same applies to stocks, which economics textbooks explicitly exclude from the definition of money . So, in general, we do not expect the Boltzmann-Gibbs law for the distribution of wealth. Some authors believe that wealth is distributed according to a power law (Pareto-Zipf), which originates from a multiplicative random process . Such a process may reflect, among other things, the fluctuations of prices needed to evaluate the monetary value of material wealth.
## II Boltzmann-Gibbs distribution
Let us consider a system of many economic agents $`N1`$, which may be individuals or corporations. In this paper, we only consider the case where their number is constant. Each agent $`i`$ has some money $`m_i`$ and may exchange it with other agents. It is implied that money is used for some economic activity, such as buying or selling material products; however, we are not interested in that aspect. As in Ref. , for us the only result of interaction between agents $`i`$ and $`j`$ is that some money $`\mathrm{\Delta }m`$ changes hands: $`[m_i,m_j][m_i^{},m_j^{}]=[m_i\mathrm{\Delta }m,m_j+\mathrm{\Delta }m]`$. Notice that the total amount of money is conserved in each transaction: $`m_i+m_j=m_i^{}+m_j^{}`$. This local conservation law of money is analogous to the conservation of energy in collisions between atoms. We assume that the economic system is closed, i. e. there is no external flux of money, thus the total amount of money $`M`$ in the system is conserved. Also, in the first part of the paper, we do not permit any debt, so each agent’s money must be non-negative: $`m_i0`$. A similar condition applies to the kinetic energy of atoms: $`\epsilon _i0`$.
Let us introduce the probability distribution function of money $`P(m)`$, which is defined so that the number of agents with money between $`m`$ and $`m+dm`$ is equal to $`NP(m)dm`$. We are interested in the stationary distribution $`P(m)`$ corresponding to the state of thermodynamic equilibrium. In this state, an individual agent’s money $`m_i`$ strongly fluctuates, but the overall probability distribution $`P(m)`$ does not change.
The equilibrium distribution function $`P(m)`$ can be derived in the same manner as the equilibrium distribution function of energy $`P(\epsilon )`$ in physics . Let us divide the system into two subsystems 1 and 2. Taking into account that money is conserved and additive: $`m=m_1+m_2`$, whereas the probability is multiplicative: $`P=P_1P_2`$, we conclude that $`P(m_1+m_2)=P(m_1)P(m_2)`$. The solution of this equation is $`P(m)=Ce^{m/T}`$; thus the equilibrium probability distribution of money has the Boltzmann-Gibbs form. From the normalization conditions $`_0^{\mathrm{}}P(m)𝑑m=1`$ and $`_0^{\mathrm{}}mP(m)𝑑m=M/N`$, we find that $`C=1/T`$ and $`T=M/N`$. Thus, the effective temperature $`T`$ is the average amount of money per agent.
The Boltzmann-Gibbs distribution can be also obtained by maximizing the entropy of money distribution $`S=_0^{\mathrm{}}𝑑mP(m)\mathrm{ln}P(m)`$ under the constraint of money conservation . Following original Boltzmann’s argument, let us divide the money axis $`0m\mathrm{}`$ into small bins of size $`dm`$ and number the bins consecutively with the index $`b=1,2,\mathrm{}`$ Let us denote the number of agents in a bin $`b`$ as $`N_b`$, the total number being $`N=_{b=1}^{\mathrm{}}N_b`$. The agents in the bin $`b`$ have money $`m_b`$, and the total money is $`M=_{b=1}^{\mathrm{}}m_bN_b`$. The probability of realization of a certain set of occupation numbers $`\{N_b\}`$ is proportional to the number of ways $`N`$ agents can be distributed among the bins preserving the set $`\{N_b\}`$. This number is $`N!/N_1!N_2!\mathrm{}`$ The logarithm of probability is entropy $`\mathrm{ln}N!_{b=1}^{\mathrm{}}\mathrm{ln}N_b!`$. When the numbers $`N_b`$ are big and Stirling’s formula $`\mathrm{ln}N!N\mathrm{ln}N`$ applies, the entropy per agent is $`S=(N\mathrm{ln}N_{b=1}^{\mathrm{}}N_b\mathrm{ln}N_b)/N=_{b=1}^{\mathrm{}}P_b\mathrm{ln}P_b`$, where $`P_b=N_b/N`$ is the probability that an agent has money $`m_b`$. Using the method of Lagrange multipliers to maximize the entropy $`S`$ with respect to the occupation numbers $`\{N_b\}`$ with the constraints on the total money $`M`$ and the total number of agents $`N`$ generates the Boltzmann-Gibbs distribution for $`P(m)`$ .
## III Computer simulations
To check that these general arguments indeed work, we performed several computer simulations. Initially, all agents are given the same amount of money: $`P(m)=\delta (mM/N)`$, which is shown in Fig. 1 as the double vertical line. One pair of agents at a time is chosen randomly, then one of the agents is randomly picked to be the “winner” (the other agent becomes the “loser”), and the amount $`\mathrm{\Delta }m0`$ is transferred from the loser to the winner. If the loser does not have enough money to pay ($`m_i<\mathrm{\Delta }m`$), then the transaction does not take place, and we proceed to another pair of agents. Thus, the agents are not permitted to have negative money. This boundary condition is crucial in establishing the stationary distribution. As the agents exchange money, the initial delta-function distribution first spread symmetrically. Then, the probability density starts to pile up at the impenetrable boundary $`m=0`$. The distribution becomes asymmetric (skewed) and ultimately reaches the stationary exponential shape shown in Fig. 1. We used several trading rules in the simulations: the exchange of a small constant amount $`\mathrm{\Delta }m=1`$, the exchange of a random fraction $`0\nu 1`$ of the average money of the pair: $`\mathrm{\Delta }m=\nu (m_i+m_j)/2`$, and the exchange of a random fraction $`\nu `$ of the average money in the system: $`\mathrm{\Delta }m=\nu M/N`$. Figures in the paper mostly show simulations for the third rule; however, the final stationary distribution was found to be the same for all rules.
In the process of evolution, the entropy $`S`$ increases in time and saturates at the maximal value for the Boltzmann-Gibbs distribution. This is illustrated by the top curve in Fig. 2 computed for the third rule of exchange. The bottom curve in Fig. 2 shows the time evolution of entropy for the first rule of exchange. The time scale for this curve is 500 times greater than for the top curve, so the bottom curve actually ends at the time $`10^6`$. The plot shows that, for the first rule of exchange, mixing is much slower than for the third one. Nevertheless, even for the first rule, the system also eventually reaches the Boltzmann-Gibbs state of maximal entropy, albeit over a time much longer than shown in Fig. 2.
One might argue that the pairwise exchange of money may correspond to a medieval market, but not to a modern economy. In order to make the model somewhat more realistic, we introduce firms. One agent at a time becomes a “firm”. The firm borrows capital $`K`$ from another agent and returns it with an interest $`rK`$, hires $`L`$ agents and pays them wages $`W`$, manufactures $`Q`$ items of a product and sell it to $`Q`$ agents at a price $`R`$. All of these agents are randomly selected. The firm receives the profit $`F=RQLWrK`$. The net result is a many-body exchange of money that still satisfies the conservation law.
Parameters of the model are selected following the procedure described in economics textbooks. The aggregate demand-supply curve for the product is taken to be $`R(Q)=V/Q^\eta `$, where $`Q`$ is the quantity people would buy at a price $`R`$, and $`\eta =0.5`$ and $`V=100`$ are constants. The production function of the firm has the conventional Cobb-Douglas form: $`Q(L,K)=L^\beta K^{1\beta }`$, where $`\beta =0.8`$ is a constant. In our simulation, we set $`W=10`$. By maximizing firm’s profit $`F`$ with respect to $`K`$ and $`L`$, we find the values of the other parameters: $`L=20`$, $`Q=10`$, $`R=32`$, and $`F=68`$.
However, the actual values of the parameters do not matter. Our computer simulations show that the stationary probability distribution of money in this model always has the universal Boltzmann-Gibbs form independent of the model parameters.
## IV Thermal machine
As explained in Introduction, the money distribution $`P(m)`$ should not be confused with the distribution of wealth. We believe that $`P(m)`$ should be interpreted as the instantaneous distribution of purchasing power in the system. Indeed, to make a purchase, one needs money. Material wealth normally is not used directly for a purchase. It needs to be sold first to be converted into money.
Let us consider an outside monopolistic vendor selling a product (say, cars) to the system of agents at a price $`p`$. Suppose that a certain small fraction $`f`$ of the agents needs to buy the product at a given time, and each agent who has enough money to afford the price will buy one item. The fraction $`f`$ is assumed to be sufficiently small, so that the purchase does not perturb the whole system significantly. At the same time, the absolute number of agents in this group is assumed to be big enough to make the group statistically representative and characterized by the Boltzmann-Gibbs distribution of money. The agents in this group continue to exchange money with the rest of the system, which acts as a thermal bath. The demand for the product is constantly renewed, because products purchased in the past expire after a certain time. In this situation, the vendor can sell the product persistently, thus creating a small steady leakage of money from the system to the vendor.
What price $`p`$ would maximize the vendor’s income? To answer this question, it is convenient to introduce the cumulative distribution of purchasing power $`𝒩(m)=N_m^{\mathrm{}}P(m^{})𝑑m^{}=Ne^{m/T}`$, which gives the number of agents whose money is greater than $`m`$. The vendor’s income is $`fp𝒩(p)`$. It is maximal when $`p=T`$, i. e. the optimal price is equal to the temperature of the system. This conclusion also follows from the simple dimensional argument that temperature is the only money scale in the problem. At the price $`p=T`$ that maximizes the vendor’s income, only the fraction $`𝒩(T)/N=e^1=0.37`$ of the agents can afford to buy the product.
Now let us consider two disconnected economic systems, one with the temperature $`T_1`$ and another with $`T_2`$: $`T_1>T_2`$. A vendor can buy a product in the latter system at its equilibrium price $`T_2`$ and sell it in the former system at the price $`T_1`$, thus extracting the speculative profit $`T_1T_2`$, as in a thermal machine. This example suggests that speculative profit is possible only when the system as a whole is out of equilibrium. As money is transferred from the high- to the low-temperature system, their temperatures become closer and eventually equal. After that, no speculative profit is possible, which would correspond to the “thermal death” of the economy. This example brings to mind economic relations between developed and developing countries, with manufacturing in the poor (low-temperature) countries for export to the rich (high-temperature) ones.
## V Models with debt
Now let us discuss what happens if the agents are permitted to go into debt. Debt can be viewed as negative money. Now when a loser does not have enough money to pay, he can borrow the required amount from a reservoir, and his balance becomes negative. The conservation law is not violated: The sum of the winner’s positive money and loser’s negative money remains constant. When an agent with a negative balance receives money as a winner, she uses this money to repay the debt until her balance becomes positive. We assume for simplicity that the reservoir charges no interest for the lent money. However, because it is not sensible to permit unlimited debt, we put a limit $`m_d`$ on the maximal debt of an agent: $`m_i>m_d`$. This new boundary condition $`P(m<m_d)=0`$ replaces the old boundary condition $`P(m<0)=0`$. The result of a computer simulation with $`m_d=800`$ is shown in Fig. 3 together with the curve for $`m_d=0`$. $`P(m)`$ is again given by the Boltzmann-Gibbs law, but now with the higher temperature $`T=M/N+m_d`$, because the normalization conditions need to be maintained including the population with negative money: $`_{m_d}^{\mathrm{}}P(m)𝑑m=1`$ and $`_{m_d}^{\mathrm{}}mP(m)𝑑m=M/N`$. The higher temperature makes the money distribution broader, which means that debt increases inequality between agents .
Imposing a sharp cutoff at $`m_d`$ may be not quite realistic. In practice, the cutoff may be extended over some range depending on the exact bankruptcy rules. Over this range, the Boltzmann-Gibbs distribution would be smeared out. So we expect to see the Boltzmann-Gibbs law only sufficiently far from the cutoff region. Similarly, in experiment , some deviations from the exponential law were observed near the lower boundary of the distribution. Also, at the high end of the distributions, the number of events becomes small and statistics poor, so the Boltzmann-Gibbs law loses applicability. Thus, we expect the Boltzmann-Gibbs law to hold only for the intermediate range of money not too close either to the lower boundary or to the very high end. However, this range is the most relevant, because it covers the great majority of population.
Lending creates equal amounts of positive (asset) and negative (liability) money . When economics textbooks describe how “banks create money” or “debt creates money” , they do not count the negative liabilities as money, and thus their money is not conserved. In our operational definition of money, we include all financial instruments with fixed denomination, such as currency, IOUs, and bonds, but not material wealth or stocks, and we count both assets and liabilities. With this definition, money is conserved, and we expect to see the Boltzmann-Gibbs distribution in equilibrium. Unfortunately, because this definition differs from economists’ definitions of money (M1, M2, M3, etc. ), it is not easy to find the appropriate statistics. Of course, money can be also emitted by a central bank or government. This is analogous to an external influx of energy into a physical system. However, if this process is sufficiently slow, the economic system may be able to maintain quasi-equilibrium, characterized by a slowly changing temperature.
We performed a simulation of a model with one bank and many agents. The agents keep their money in accounts on which the bank pays interest. The agents may borrow money from the bank, for which they must pay interest in monthly installments. If they cannot make the required payments, they may be declared bankrupt, which relieves them from the debt, but the liability is transferred to the bank. In this way, the conservation of money is maintained. The model is too elaborate to describe it in full detail here. We found that, depending on the parameters of the model, either the agents constantly lose money to the bank, which steadily reduces the agents’ temperature, or the bank constantly loses money, which drives down its own negative balance and steadily increases the agents’ temperature.
## VI Boltzmann equation
The Boltzmann-Gibbs distribution can be also derived from the Boltzmann equation , which describes the time evolution of the distribution function $`P(m)`$ due to pairwise interactions:
$`{\displaystyle \frac{dP(m)}{dt}}={\displaystyle }{\displaystyle }\{w_{[m,m^{}][m\mathrm{\Delta },m^{}+\mathrm{\Delta }]}P(m)P(m^{})`$ (1)
$`+w_{[m\mathrm{\Delta },m^{}+\mathrm{\Delta }][m,m^{}]}P(m\mathrm{\Delta })P(m^{}+\mathrm{\Delta })\}dm^{}d\mathrm{\Delta }.`$ (2)
Here $`w_{[m,m^{}][m\mathrm{\Delta },m^{}+\mathrm{\Delta }]}`$ is the rate of transferring money $`\mathrm{\Delta }`$ from an agent with money $`m`$ to an agent with money $`m^{}`$. If a model has time-reversal symmetry, then the transition rate of a direct process is the same as the transition rate of the reversed process, thus the $`w`$-factors in the first and second lines of Eq. (1) are equal. In this case, the Boltzmann-Gibbs distribution $`P(m)=C\mathrm{exp}(m/T)`$ nullifies the right-hand side of Eq. (1); thus this distribution is stationary: $`dP(m)/dt=0`$ .
## VII Non-Boltzmann-Gibbs distributions
However, if time-reversal symmetry is broken, the two transition rates $`w`$ in Eq. (1) may be different, and the system may have a non-Boltzmann-Gibbs stationary distribution or no stationary distribution at all. Examples of such kind were studied in Ref. . One model was called the multiplicative random exchange. In this model, a randomly selected loser $`i`$ loses a fixed fraction $`\alpha `$ of his money $`m_i`$ to a randomly selected winner $`j`$: $`[m_i,m_j][(1\alpha )m_i,m_j+\alpha m_i]`$. If we try to reverse this process and appoint the winner $`j`$ to become a loser, the system does not return to the original configuration $`[m_i,m_j]`$: $`[(1\alpha )m_i,m_j+\alpha m_i][(1\alpha )m_i+\alpha (m_j+\alpha m_i),(1\alpha )(m_j+\alpha m_i)]`$. Except for $`\alpha =1/2`$, the exponential distribution function is not a stationary solution of the Boltzmann equation derived for this model in Ref. . Instead, the stationary distribution has the shape shown in Fig. 4 for $`\alpha =1/3`$, which we reproduced in our numerical simulations. It still has an exponential tail end at the high end, but drops to zero at the low end for $`\alpha <1/2`$. Another example of similar kind was studied in Ref. , which appeared after the first version of our paper was posted as cond-mat/0001432 on January 30, 2000. In that model, the agents save a fraction $`\lambda `$ of their money and exchange a random fraction $`ϵ`$ of their total remaining money: $`[m_i,m_j][\lambda m_i+ϵ(1\lambda )(m_i+m_j),\lambda m_j+(1ϵ)(1\lambda )(m_i+m_j)]`$. This exchange also does not return to the original configuration after being reversed. The stationary probability distribution was found in Ref. to be nonexponential for $`\lambda 0`$ with a shape qualitatively similar to the one shown in Fig. 4.
Another interesting example of a non-Boltzmann-Gibbs distribution occurs in a model with taxes and subsidies. Suppose a special agent (“government”) collects a fraction (“tax”) of every transaction in the system. The collected money is then equally divided between all agents of the system, so that each agent receives the subsidy $`\delta m`$ with the frequency $`1/\tau _s`$. Assuming that $`\delta m`$ is small and approximating the collision integral with a relaxation time $`\tau _r`$ , we obtain the following Boltzmann equation
$$\frac{P(m)}{t}+\frac{\delta m}{\tau _s}\frac{P(m)}{m}=\frac{P(m)\stackrel{~}{P}(m)}{\tau _r},$$
(3)
where $`\stackrel{~}{P}(m)`$ is the equilibrium Boltzmann-Gibbs function. The second term in the left-hand side of Eq. (3) is analogous to the force applied to electrons in a metal by an external electric field . The approximate stationary solution of Eq. (3) is the displaced Boltzmann-Gibbs one: $`P(m)=\stackrel{~}{P}(m(\tau _r/\tau _s)\delta m)`$. The displacement of the equilibrium distribution $`\stackrel{~}{P}(m)`$ by $`(\tau _r/\tau _s)\delta m`$ would leave an empty gap near $`m=0`$. This gap is filled by interpolating between zero population at $`m=0`$ and the displaced distribution. The curve obtained in a computer simulation of this model (Fig. 5) qualitatively agrees with this expectation. The low-money population is suppressed, because the government, acting as an external force, “pumps out” that population and pushes the system out of thermodynamic equilibrium. We found that the entropy of the stationary state in the model with taxes and subsidies is few percents lower than without.
These examples show that the Boltzmann-Gibbs distribution is not fully universal, meaning that it does not hold for just any model of exchange that conserves money. Nevertheless, it is universal in a limited sense: For a broad class of models that have time-reversal symmetry, the stationary distribution is exponential and does not depend on the details of a model. Conversely, when time-reversal symmetry is broken, the distribution may depend on model details. The difference between these two classes of models may be rather subtle. For example, let us change the multiplicative random exchange from a fixed fraction of loser’s money to a fixed fraction of the total money of winner and loser. This modification retains the multiplicative idea that the amount exchanged is proportional to the amount involved, but restores time-reversal symmetry and the Boltzmann-Gibbs distribution. In the model with $`\mathrm{\Delta }m=1`$ discussed in the next Section, the difference between time-reversible and time-irreversible formulations amounts to the difference between impenetrable and absorbing boundary conditions at $`m=0`$. Unlike in physics, in economy there is no fundamental requirement that interactions have time-reversal symmetry. However, in the absence of detailed knowledge of real microscopic dynamics of economic exchange, the semiuniversal Boltzmann-Gibbs distribution appears to be a natural starting point.
Moreover, deviations from the Boltzmann-Gibbs law may occur only if the transition rates $`w`$ in Eq. (1) explicitly depend on the agents money $`m`$ or $`m^{}`$ in an asymmetric manner. In another simulation, we randomly preselected winners and losers for every pair of agents $`(i,j)`$. In this case, money flows along directed links between the agents: $`ijk`$, and time-reversal symmetry is strongly broken. This model is closer to the real economy, in which, for example, one typically receives money from an employer and pays it to a grocer, but rarely the reverse. Nevertheless, we still found the Boltzmann-Gibbs distribution of money in this model, because the transition rates $`w`$ do not explicitly depend on $`m`$ and $`m^{}`$.
## VIII Nonlinear Boltzmann equation vs. linear master equation
For the model where agents randomly exchange the constant amount $`\mathrm{\Delta }m=1`$, the Boltzmann equation is:
$`{\displaystyle \frac{dP_m}{dt}}`$ $`=`$ $`P_{m+1}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}P_n+P_{m1}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}P_n`$ (5)
$`P_m{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}P_nP_m{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}P_n`$
$`=`$ $`(P_{m+1}+P_{m1}2P_m)+P_0(P_mP_{m1}),`$ (6)
where $`P_mP(m)`$ and we have used $`_{m=0}^{\mathrm{}}P_m=1`$. The first, diffusion term in Eq. (6) is responsible for broadening of the initial delta-function distribution. The second term, proportional to $`P_0`$, is essential for the Boltzmann-Gibbs distribution $`P_m=e^{m/T}(1e^{1/T})`$ to be a stationary solution of Eq. (6). In a similar model studied in Ref. , the second term was omitted on the assumption that agents who lost all money are eliminated: $`P_0=0`$. In that case, the total number of agents is not conserved, and the system never reaches any stationary distribution. Time-reversal symmetry is violated, since transitions into the state $`m=0`$ are permitted, but not out of this state.
If we treat $`P_0`$ as a constant, Eq. (6) looks like a linear Fokker-Planck equation for $`P_m`$, with the first term describing diffusion and the second term an external force proportional to $`P_0`$. Similar equations were studied in Ref. . Eq. (6) can be also rewritten as
$$\frac{dP_m}{dt}=P_{m+1}(2P_0)P_m+(1P_0)P_{m1}.$$
(7)
The coefficient $`(1P_0)`$ in front of $`P_{m1}`$ represents the rate of increasing money by $`\mathrm{\Delta }m=1`$, and the coefficient 1 in front of $`P_{m+1}`$ represents the rate of decreasing money by $`\mathrm{\Delta }m=1`$. Since $`P_0>0`$, the former is smaller than the latter, which results in the stationary Boltzmann-Gibbs distributions $`P_m=(1P_0)^m`$. An equation similar to Eq. (7) describes a Markov chain studied for strategic market games in Ref. . Naturally, the stationary probability distribution of wealth in that model was found to be exponential .
Even though Eqs. (6) and (7) look like linear equations, nevertheless the Boltzmann equation (1) and (5) is a profoundly nonlinear equation. It contains the product of two probability distribution functions $`P`$ in the right-hand side, because two agents are involved in money exchange. Most studies of wealth distribution have the fundamental flaw that they use a single-particle approach. They assume that the wealth of an agent may change just by itself and write a linear master equation for the probability distribution. Because only one particle is considered, this approach cannot adequately incorporate conservation of money. In reality, an agent can change money only by interacting with another agent, thus the problem requires a two-particle probability distribution function. Using Boltzmann’s molecular chaos hypothesis, the two-particle function is factorized into a product of two single-particle distributions functions, which results in the nonlinear Boltzmann equation. Conservation of money is adequately incorporated in this two-particle approach, and the universality of the exponential Boltzmann-Gibbs distribution is transparent.
## IX Conclusions
Everywhere in the paper we assumed some randomness in the exchange of money. The results of our paper would apply the best to the probability distribution of money in a closed community of gamblers. In more traditional economic studies, the agents exchange money not randomly, but following deterministic strategies, such as maximization of utility functions . The concept of equilibrium in these studies is similar to mechanical equilibrium in physics, which is achieved by minimizing energy or maximizing utility. However, for big ensembles, statistical equilibrium is a more relevant concept. When many heterogeneous agents deterministically interact and spend various amounts of money from very little to very big, the money exchange is effectively random. In the future, we would like to uncover the Boltzmann-Gibbs distribution of money in a simulation of a big ensemble of economic agents following realistic deterministic strategies with money conservation taken into account. That would be the economics analog of molecular dynamics simulations in physics. While atoms collide following fully deterministic equations of motion, their energy exchange is effectively random due to complexity of the system and results in the Boltzmann-Gibbs law.
We do not claim that the real economy is in equilibrium. (Most of the physical world around us is not in true equilibrium either.) Nevertheless, the concept of statistical equilibrium is a very useful reference point for studying nonequilibrium phenomena.
The authors are grateful to M. Gubrud for helpful discussion and proofreading of an earlier version of the manuscript.
Note added: After the paper had been submitted for publication, we have learned about the book by Aoki , who applied many ideas of statistical physics to economics, albeit not specifically to the distribution of money. |
warning/0001/nlin0001042.html | ar5iv | text | # Measuring Information Transfer
## Abstract
An information theoretic measure is derived that quantifies the statistical coherence between systems evolving in time. The standard time delayed mutual information fails to distinguish information that is actually exchanged from shared information due to common history and input signals. In our new approach, these influences are excluded by appropriate conditioning of transition probabilities. The resulting transfer entropy is able to distinguish driving and responding elements and to detect asymmetry in the coupling of subsystems.
The time evolution of a system may be called irregular if it generates information at a non-zero rate. For stochastic or deterministically chaotic systems, this is quantified by the entropy. For a system consisting of more than one component, important information on its structure can be obtained by measuring to which extent the individual components contribute to information production and at what rate they exchange information among each other. This paper proposes a method to answer the latter question on the basis of time series observations.
Many authors have used mutual information to quantify the overlap of the information content of two (sub-) systems. Unfortunately, mutual information neither contains dynamical nor directional information. Introducing a time delay in one of the observations is an important, if somewhat arbitrary, improvement in this respect, but still does not explicitly distinguish information that is actually exchanged from that due to the response to a common input signal or history.
The purpose of this paper is to motivate and derive an alternative information theoretic measure, to be called transfer entropy, that shares some of the desired properties of mutual information but takes the dynamics of information transport into account. With minimal assumptions about the dynamics of the system and the nature of their coupling one will be able to quantify the exchange of information between two systems, separately for both directions, and, if desired, conditional to common input signals.
This work may be seen in the context of a considerable number of recently proposed measures for the nonlinear coherence of signals, used to study generalized synchronization phenomena in many contects, most notably in physiological systems. While these measures are often very powerful for a specific set of applications, it is also important to aim at an understanding of the underlying theoretical concepts. In the generic case that neither of the systems, nor their coupling may be assumed to be deterministic, information theory seems to be an appropriate starting point.
Let us briefly recall the most basic concepts of information theory. The average number of bits needed to optimally encode independent draws of the discrete variable $`I`$ following a probability distribution $`p(i)`$ is given by the Shannon entropy
$`H_I={\displaystyle \underset{i}{}}p(i)\mathrm{log}_2p(i)`$
where the sum extends over all states $`i`$ the process can assume. The base of the logarithm only determines the units used for measuring information and will be dropped henceforth.
In order to construct an optimal encoding that uses just as many bits as given by the entropy, it is necessary to know the probability distribution $`p(i)`$. The excess number of bits that will be coded if a different distribution $`q(i)`$ is used instead of $`p(i)`$ is given by the Kullback entropy $`K_I=_ip(i)\mathrm{log}p(i)/q(i)`$. We will later also need the Kullback entropy for conditional probabilities $`p(i|j)`$. For a single state $`j`$ we have $`K_j=_ip(i|j)\mathrm{log}p(i|j)/q(i|j)`$. Summation over $`j`$ with respect to $`p(j)`$ yields
$$K_{I|J}=\underset{i,j}{}p(i,j)\mathrm{log}\frac{p(i|j)}{q(i|j)}.$$
(1)
The mutual information of two processes $`I`$ and $`J`$ with joint probability $`p_{IJ}(i,j)`$ can be seen as the excess amount of code produced by erroneously assuming that the two systems are independent, i.e. assuming $`q_{IJ}(i,j)=p_I(i)p_J(j)`$ instead of $`p_{IJ}(i,j)`$. The corresponding Kullback entropy is
$$M_{IJ}=p(i,j)\mathrm{log}\frac{p(i,j)}{p(i)p(j)}$$
(2)
which is the well known formula for the mutual information. Here and in the following, we omitted the summation index and the subscript of the probabilities specifying the process. This derivation shows that mutual information is a natural way to quantify the deviation from independence of two processes. We have $`M_{IJ}=H_I+H_JH_{IJ}0`$. Note that $`M_{IJ}`$ is symmetric under the exchange of $`I`$ and $`J`$ and therefore does not contain any directional sense.
A related, non-symmetric quantity is the conditional entropy $`H_{I|J}=p(i,j)\mathrm{log}p(i|j)=H_{IJ}H_J`$. However, since $`H_{I|J}H_{J|I}=H_IH_J`$, it is non-symmetric only due to the different individual entropies and not due to information flow. Mutual information can be given a directional sense in a somwhat ad-hoc way by introducing a time lag in either one of the variables and compute e.g.
$`M_{IJ}(\tau )={\displaystyle p(i_n,j_{n\tau })\mathrm{log}\frac{p(i_n,j_{n\tau })}{p(i)p(j)}}.`$
As we will see below, considering the two systems at different times occurs naturally as soon as transition probabilities are introduced. This will yield a more justified approach to measuring information transfer that explicitly incorporates directional, dynamical structure.
One can incorporate dynamical structure by studying transition probabilities rather than static probabilities. Consider a system that may be approximated by a stationary Markov process of order $`k`$, that is, the conditional probability to find $`I`$ in state $`i_{n+1}`$ at time $`n+1`$ observes $`p(i_{n+1}|i_n,\mathrm{},i_{nk+1})=p(i_{n+1}|i_n,\mathrm{},i_{nk})`$. Henceforth we will use the shorthand notation $`i_n^{(k)}=(i_n,\mathrm{},i_{nk+1})`$ for words of length $`k`$, or $`k`$ dimensional delay embedding vectors.
The average number of bits needed to encode one additional state of the system if all previous states are known is given by the entropy rate
$$h_I=p(i_{n+1},i_n^{(k)})\mathrm{log}p(i_{n+1}|i_n^{(k)}).$$
(3)
Since $`p(i_{n+1}|i_n^{(k)})=p(i_{n+1}^{(k+1)})/p(i_n^{(k)})`$, this is just the difference between the Shannon entropies of the processes given by $`k+1`$ and $`k`$ dimensional delay vectors constructed from $`I`$: $`h_I=H_{I^{(k+1)}}H_{I^{(k)}}`$.
If $`I`$ is obtained by coarse graining a continuous system $`X`$ at resolution $`ϵ`$, the entropy $`H_X(ϵ)`$ and entropy rate $`h_X(ϵ)`$ will depend on the partitioning and in general diverge like $`\mathrm{log}ϵ`$ when $`ϵ0`$. However, for the special case of a deterministic dynamical system, $`lim_{ϵ0}h_X(ϵ)=h_{KS}`$ may exist and is then called the Kolmogorov–Sinai entropy. (For non-Markov systems, also the limit $`k\mathrm{}`$ needs to be taken.) Confusingly, the opposite is true for the mutual information. For generic noisy interdependence, $`lim_{ϵ0}M_{XY}(ϵ)`$ is finite and independent of the partition, but for deterministically coupled processes, $`M_{XY}(ϵ)`$ will diverge as $`ϵ0`$.
The Shannon entropy and its generalization, the mutual information, are properties of the static probability distributions while the dynamics of the processes is contained in the transition probabilities. For the study of the dynamics of shared information between processes it is therefore desirable to generalize the entropy rate, rather than Shannon entropy, to more than one system. In the next section I will propose such generalizations, in particular one that is non-symmetric under the exchange of the two processes.
The most straightforward way to construct a mutual information rate by generalizing $`h_I`$ to two processes $`(I,J)`$ is again by measuring the deviation from independence. The corresponding Kullback entropy is sometimes called transinformation and is still symmetric under the exchange of $`I`$ and $`J`$. It is therefore preferable to measuring the deviation from the generalized Markov property
$`p(i_{n+1}|i_n^{(k)})=p(i_{n+1}|i_n^{(k)},j_n^{(l)}).`$
In the absence of information flow from $`J`$ to $`I`$, the state of $`J`$ has no influence on the transition probabilities on system $`I`$. The incorrectness of this assumption can again be quantified by a Kullback entropy (1) by which we define the transfer entropy:
$$T_{JI}=p(i_{n+1},i_n^{(k)},j_n^{(l)})\mathrm{log}\frac{p(i_{n+1}|i_n^{(k)},j_n^{(l)})}{p(i_{n+1}|i_n^{(k)})}.$$
(4)
This is the central concept of this paper. The most natural choices for $`l`$ are $`l=k`$ or $`l=1`$. Usually, the latter is preferable for computational reasons. $`T_{JI}`$ is now explicitly non-symmetric since it measures the degree of dependence of $`I`$ on $`J`$ and not vice versa.
For coarse grained states $`(I,J)`$ of continuous systems $`(X,Y)`$, the limit $`lim_{ϵ0}T_{YX}(ϵ)`$ is finite and independent of the partition, except for the case of deterministic coupling, when $`T_{YX}(ϵ)`$ diverges as $`ϵ0`$. In this respect, transfer entropy behaves like mutual information. If computationally feasible, the influence of a known common driving force $`Z`$ may be excluded by conditioning the probabilities under the logarithm to $`z_n`$ as well.
For numerical and practical applications, the limit $`ϵ0`$ is not obtainable and has to be replaced appropriately. Either one can study transfer entropy as a function of the resolution, or one can fix a resolution for the scope of a study. Furthermore, there are several methods of coarse graining and a partition consisting of a fixed mesh of boxes is not always the best choice. Fixed boxes are only suitable in cases where data can be produced with little effort and small statistical errors at reasonable speed of computation are desired.
For time series applications, an alternative implementation using generalized correlation integrals is preferable. Mutual information and redundancies have been generalized for their estimation by order $`q`$ correlation integrals . It is possible to follow the same arguments in generalizing transfer entropy. However, for the computationally most attractive case $`q=2`$, we would have to give up positivity of $`T_{IJ}`$. Instead, we propose an implementation of the definition (4) where the probability measure $`p(i_{n+1},i_n^{(k)},j_n^{(l)})`$ is realized by a sum over all available realizations of $`(x_{n+1},x_n^{(k)},y_n^{(l)})`$ in a time series. The transition probabilities are expressed by joint probabilities and then obtained by kernel estimation, e.g.
$`\widehat{p}_r(x_{n+1},x_n,y_n)={\displaystyle \frac{1}{N}}{\displaystyle \underset{n^{}}{}}\mathrm{\Theta }\left(\left|(\begin{array}{c}x_{n+1}x_{n^{}+1}\\ x_nx_n^{}\\ y_ny_n^{}\end{array})\right|r\right).`$
We use the step kernel $`\mathrm{\Theta }(x>0)=1`$; $`\mathrm{\Theta }(x0)=0`$. The norm $`||`$ can be simply the maximum distance but other norms and kernels can be considered. In particular, different overall scales of $`X`$ and $`Y`$ can be accounted for by using appropriate weights. Similarly to standard dimension and entropy calculations, fast neighbour search strategies are advisable for all but the smallest data sets. Dynamically correlated pairs should be excluded as usual. Since these technical issues are the same as in many nonlinear time series methods, the reader is referred to the discussion in the literature .
In order to demonstrate the use of transfer entropy, let us study three examples, two spatio-temporal systems and a bi-variate physiological time series. In a one dimensional lattice of unidirectionally coupled maps
$$x_{n+1}^m=f(ϵx_n^{m1}+(1ϵ)x_n^m),$$
(5)
information can be transported only in the direction of increasing $`m`$. One of the simplest cases is given by the tent map, $`f(x<0.5)=2x`$; $`f(x0.5)=22x`$. Let us study coarse grained states $`I^m`$ with $`i_n^m`$ defined by a partition at $`x_0=0.5`$. At zero coupling, all static and transfer probabilities are equal to 1/2, $`M(\tau )=0`$ for all values of $`\tau `$, and also $`T_{I^{m1}I^m}=T_{I^mI^{m1}}=0`$. For nonzero coupling, we still have $`T_{I^mI^{m1}}=0`$, but $`T_{I^{m1}I^m}`$ becomes positive. For small coupling, it can be assumed that the invariant density at a single site is essentially unchanged whence the transition probabilities $`p(I_{n+1}^m|I_n^m,I_n^{m1})`$ are changed by an amount proportional to $`ϵ`$. In particular, $`p(0|0,0)`$, $`p(0|1,1)`$, $`p(1|0,1)`$, and $`p(1|1,0)`$ are increased by a factor $`1+\alpha ϵ`$ with $`\alpha =O(1)`$. All others are decreased by that amount. Evaluating (4) in lowest order of $`ϵ`$ with $`k=l=1`$, we obtain $`T_{I^{m1}I^m}=\alpha ^2ϵ^2/\mathrm{ln}(2)+O(ϵ^4)`$. For this particular case, the changes in $`p(i_{n+1}^m,i_n^{m1})`$ exactly cancel out and the mutual information is zero. Figure 1 shows a numerical verification of these results for a spatially periodic lattice of 100 maps. Averages of 10 runs of $`10^5`$ iterates after $`10^5`$ transients are shown. The transfer entropy $`T_{I^mI^{m1}}`$ and both directions of $`M(\tau =1)`$ were found consistent with zero and are therefore not shown.
The situation is more complicated for the Ulam map $`f(x)=2x^2`$ and non-small coupling. For each coupling, a bi-variate time series was generated using a lattice of $`100`$ points (random initial conditions) and recording 10000 iterates of $`x_n^1`$ and $`x_n^2`$ after $`10^5`$ steps of transients. Correlation sums at $`r=0.2`$ were used to compute mutual information in both directions, $`M_{X^1,X^2}(\tau =1)`$ and $`M_{X^2,X^1}(\tau =1)`$, as well as transfer entropies $`T_{X^1X^2}`$ and $`T_{X^2X^1}`$ with $`k=l=1`$. Neighbors closer in time than 100 iterates were excluded from the kernel estimation.
Figure 2 shows $`M`$ and $`T`$ as functions of the coupling strength. Both $`M`$ and $`T`$ are able to detect the anisotropy since the information is consistently larger in the positive direction. The lattice undergoes a number of bifurcations when the coupling is changed. Around $`ϵ=0.18`$, the asymptotic state is of temporal and spatial period two. For this case, the mutual information is found to be 1 bit. This is correct although information is neither produced nor exchanged and reflects the static correlation between the sites. The transfer entropy finds a zero rate of information transport, as desired. Around this pariodic window, the mutual information is non-zero in both directions and the signature of the unidirectional coupling is less pronounced. Around $`ϵ=0.82`$, the lattice settles to a (spatially inhomogenious) fixed point state. Here both measures correctly show zero information transfer. The most important finding, however, is that the transfer entropy for the negative direction remains consistent with zero for all couplings, reflecting the causality in the system.
As a last example, take a bi-variate time series (see Fig. 3) of the breath rate and instantaneous heart rate of a sleeping human suffering from sleep apnea (part of data set B of the Santa Fe Institute time series contest held in 1991 ). Figure 4 shows that while time delayed mutual information is almost symmetric between both series, the transfer entropy indicates a stronger flow of information from the heart rate to the breath rate than vice versa over a significant range of length scales $`r`$. Note that for small $`r`$, the curves deflect down to zero due to the finite sample size. This result is consistent with the observation that the patient breathes in bursts which seem to occur whenever the heart rate crosses some threshold. Certainly, both signals could instead be responding to a common external trigger.
In conclusion, the new transfer entropy is able to detect the directed exchange of information between two systems. Unlike mutual information, it is designed to ignore static correlations due to the common history or common input signals. Most prominent applications include multivariate analysis of time series and the study of spatially extended systems.
Several authors have proposed to use time delayed mutual information $`M(\mathrm{\Delta }l,\tau )`$ as a function of spatial distance $`\mathrm{\Delta }l`$ and temporal delay $`\tau `$ to define a velocity of information transport in spatio-temporal systems. Often, one finds that $`M(\mathrm{\Delta }l,\tau )`$ for fixed $`\mathrm{\Delta }i`$ reaches a local maximum at some lag $`\tau ^{}`$. Hence a velocity can be defined by the ratio $`\mathrm{\Delta }i/\tau ^{}`$, in particular if that ratio is fairly constant over the resolvable range of values for $`\mathrm{\Delta }i`$. This reasoning has been challenged by giving an example where the above interpretation implies super-luminar communication. In fact, much of the common information is due to the common history that allows the lattice to partially synchronize. Preliminary results indicate that appropriate conditioning for the common history by replacing time delayed mutual information by a variant of Eq.(4) resolves this aparent paradox. However, conditioning with respect to a large number of variables poses immense numerical problems whence this study will be concluded at a later time.
Part of this work has been supported by the SFB 237 of the Deutsche Forschungsgemeinschaft. |
warning/0001/math0001056.html | ar5iv | text | # Failure of Brown representability in derived categories
## Introduction
The introduction is written for the reader who knows about derived categories, but is not necessarily familiar with previous articles by the authors and their close friends. We begin with a sketch of the work done in the last 10 years, generalising results from homotopy theory to derived categories. The experts may want to skip this, and go directly to Notation 0.4 on page 0.4. After the very general survey, will come a much more focused one. We will give, in some detail, the history of the results on generalising the theorem of Brown and Adams to derived categories. Then we will explain the two open problems, which we settle in this article. Finally, we will give the nature of our counterexamples.
Let $`𝒯`$ be a triangulated category. The representable functors $`𝒯(,X)`$ are all homological; that is, they take triangles to long exact sequences. Given a triangulated subcategory $`𝒮𝒯`$, we can restrict a representable functor on $`𝒯`$ to a functor on $`𝒮`$. We denote the restriction by $`𝒯(,X)|_𝒮`$. All such functors are clearly homological.
The most interesting version of this, is where $`𝒯`$ is a triangulated category with coproducts, and $`𝒮`$ is the full subcategory $`𝒯^c`$ of all compact objects in $`𝒯`$.
###### Definition 0.1.
An object $`c𝒯`$ is called compact, if the functor $`𝒯(c,)`$ commutes with coproducts.
We remind the reader that for $`𝒯`$ the homotopy category of spectra, $`𝒯^c𝒯`$ is the subcategory of finite spectra. For $`𝒯=D(R)`$, the unbounded derived category of right $`R`$–modules, $`𝒯^c`$ turns out to be the subcategory of perfect complexes, that is, complexes isomorphic to finite complexes of finitely generated projective $`R`$–modules. For a more detailed discussion of examples, where $`𝒯`$ is the unbounded derived category of coherent sheaves on a scheme, see Sections 1 and 2 in .
Since the functor $`𝒯(,X)|_{𝒯^c}`$ plays a major role in what follows, we adopt a shorthand for it. We will write
$$𝐲X=𝒯(,X)|_{𝒯^c}.$$
The subject we will be studying began with a theorem of Adams .
###### Theorem 0.2.
(Adams, 1971) Let $`𝒯`$ be the homotopy category of spectra, and $`𝒯^c`$ the subcategory of finite spectra. Then any homological functor $`\{𝒯^c\}^{op}𝒜b`$ is isomorphic to $`𝐲X`$, for some object $`X𝒯`$. Furthermore, any natural transformation of functors
$$\begin{array}{ccc}𝐲X& & 𝐲Y\end{array}$$
is induced by some (non–unique) map $`XY`$.
###### Remark 0.3.
This theorem is usually referred to as “Brown representability”. The reason for this is that, 10 years earlier, Brown proved a special case. In Brown’s theorem, there was a countability hypothesis on the functor.
Calling this theorem “Brown representability” is somewhat confusing, since in the same paper, Brown proved another result, somewhat related. He showed that, if $`𝒯`$ is the homotopy category of spectra, and $`H:𝒯^{op}𝒜b`$ is a homological functor taking coproducts to products, then $`H`$ is representable. There are two theorems here, one about homological functors on $`𝒯^{op}`$, and another about homological functors on the subcategory $`\{𝒯^c\}^{op}`$. And both theorems usually go under the name Brown representability. Neither theorem is a special case of the other. In the literature, one sometimes distinguishes them by calling the theorem about functors on $`\{𝒯^c\}^{op}`$ “Brown representability for homology”, while the theorem about functors on $`𝒯^{op}`$ goes by the name “Brown representability for cohomology”.
The reason for this strange terminology is the following. In many interesting cases, the category $`𝒯^c`$ is self dual. Thus functors $`\{𝒯^c\}^{op}𝒜b`$ are the same as functors $`𝒯^c𝒜b`$, and these correspond to functors $`𝒯𝒜b`$ respecting coproducts. Thus “Brown representability for homology” may be viewed as a theorem about covariant homological functors $`𝒯𝒜b`$, respecting coproducts.
In hindsight, it seems natural to ask how these statements generalise to other triangulated categories, in particular, the derived category of a ring. Surprisingly, questions of this sort were not asked until the 1980’s.
Even then, the first questions to be asked were: To what extent can results about rings be generalised to homotopy theory. The first to suggest that this might be a fruitful pursuit was probably Waldhausen. Waldhausen proposed that techniques from homological algebra—Hochschild homology and cohomology, trace maps, and cyclic versions of these—should all be done in the context of $`E^{\mathrm{}}`$ ring spectra. The work that followed, by Goodwillie, Bökstedt, Hsiang, Madsen and many others since, showed how good the idea was.
The idea that translating results from homotopy theory to derived categories could be worthwhile came later. The first paper we are aware of is Hopkins’ ; in it, one has a derived category version of the nilpotence theorem. But it was really only in Bökstedt–Neeman’s that the first attempt was made, to use homotopy theoretic techniques to solve standard problems on derived categories. In the 1990’s, we have seen explosive growth in the subject. In and , Neeman applied techniques coming from homotopy theory to the study of, respectively, the localisation theorem in $`K`$–theory and to Grothendieck duality. The articles by Rickard , Benson, Carlson and Rickard , , , Benson and Krause , Krause , and Benson and Gnacadja , give beautiful applications to group cohomology. Keller , , applies the techniques to the study of cyclic homology. And Voevodsky , and , Suslin–Voevodsky and Morel and , have produced a string of results, which apply homotopy theory to the study of motives.
Along with the applications, came the study of the degree to which the theorems extend. Homotopy theorists, over a period of 30 years, developed certain tools to handle the category of spectra. It became interesting to know which parts of these tools work, in the new and greater generality. This has also led to a series of papers. Hovey, Palmieri and Strickland set up a convenient axiomatic formalism. Without going into detail, we remind the reader of the work of Beligiannis , Christensen , Christensen–Strickland , Franke , Keller , Krause , , , Krause and Reichenbach , and Neeman , and .
This skimpy historical survey was intended to explain why people have studied whether Brown representability generalises to derived categories. As we mentioned in Remark 0.3, the term Brown representability is used to cover two theorems. Brown representability for cohomology is a characterisation of representable functors $`𝒯^{op}𝒜b`$, while Brown representability for homology is a more complicated statement about functors $`\{𝒯^c\}^{op}𝒜b`$. Of these, the generalisation of Brown representability for cohomology is very well understood. The best and most recent results were obtained independently by Franke and Neeman , and one of the remarkable aspects of their theorems, is that they prove new results even in homotopy theory. The theorems tell us, that Brown representability for cohomology generalises to the categories of $`E`$–acyclic spectra and $`E`$–local spectra, for any homology theory $`E`$.
This paper addresses the less well understood problem, of Brown representability for homology. In the remainder of the Introduction, we will do two things. First, we will go through the history of this problem in detail, explaining what was already known. Then, we will outline the counterexamples and results obtained in this article. But before we start, we need to establish some notation.
###### Notation 0.4.
All rings will be associative, with unit. All $`R`$–modules will be right, unitary modules. The ring $`R`$ is called hereditary if its global dimension is at most $`1`$. The triangulated category $`𝒯=D(R)`$ will be the unbounded derived category of right $`R`$–modules. The category $`𝒯^c`$ is, as above, the full subcategory of compact objects in $`𝒯`$.
We will denote the category of right $`R`$–modules by the symbol $`\text{Mod-}R`$. The subcategory of finitely presented $`R`$–modules will be denoted $`\text{mod-}R`$. The category of all additive functors $`\{𝒯^c\}^{op}𝒜b`$ will be denoted $`\text{Mod-}𝒯^c`$, while the category of all additive functors $`\{\text{mod-}R\}^{op}𝒜b`$ will bear the name $`\text{Mod}(\text{mod-}R)`$.
When speaking of objects of the category $`\text{Mod-}𝒯^c`$, that is, of functors $`\{𝒯^c\}^{op}𝒜b`$, we frequently wish to single out the ones that are homological, that is, take triangles to long exact sequences. We will feel free to interchangeably use the adjectives “homological”, “exact” or “flat”. We remind the reader that an object of $`\text{Mod-}𝒯^c`$ is exact if and only if it is a filtered colimit of representable functors. Furthermore, the representable functors are projective. (We use the term “representable” to mean functors of the form $`𝐲C`$, with $`C`$ compact. In the literature, people sometimes call all functors $`𝐲X`$ representable.)
We also need to recall the notion of purity for $`R`$–modules. A short exact sequence of $`R`$–modules
$$\begin{array}{ccccccccc}0& & A& & B& & C& & 0\end{array}$$
is called pure exact, if it remains exact when tensored with an arbitrary left $`R`$–module. Equivalently, it is a pure exact sequence if, for every finitely presented module $`P`$, the functor $`\text{Hom}(P,)`$ takes it to an exact sequence
$$\begin{array}{ccccccccc}0& & \text{Hom}(P,A)& & \text{Hom}(P,B)& & \text{Hom}(P,C)& & 0.\end{array}$$
An $`R`$–module $`P`$ is called pure projective, if the functor $`Hom(P,)`$ takes pure exact sequences to exact sequences. A module $`P`$ is pure projective if and only if it is a summand of a coproduct of finitely presented modules. The pure projective dimension of an $`R`$–module $`M`$ is defined to be the length of its shortest pure resolution by pure projectives.
A module $`I`$ is called pure injective, if the functor $`Hom(,I)`$ takes pure exact sequences to exact sequences. The pure injective dimension of a module $`I`$ is the length of the shortest pure resolution by pure injectives. The pure global dimension of $`R`$, denoted $`pgldimR`$, is the supremum over all $`M`$, of the pure projective dimension of $`M`$. This equals the supremum of the pure injective dimensions. We refer the reader to for a more thorough discussion, with proofs.
Finally, recall our shorthand: for $`X𝒯`$, we write $`𝐲X`$ for the exact=homological=flat functor $`𝒯(,X)|_{𝒯^c}`$. It is also convenient to make a definition which is not so standard:
###### Definition 0.5.
(Beligiannis ) The pure global dimension of $`𝒯`$, denoted $`pgldim𝒯`$, is defined to be the supremum, over all $`X𝒯`$, of the projective dimension in $`\text{Mod-}𝒯^c`$ of the object $`𝐲X`$.
The following proposition will be useful.
###### Proposition 0.6.
(Beligiannis \[5, Prop. 11.2\]. The proof is based on an idea by Jensen, which appeared in a paper by Simson \[42, Thm. 2.7\].)
The pure global dimension of $`𝒯`$ is also the supremum over all homological=exact functors $`F`$, of the projective dimension of $`F`$. Note that, as we will discover in this article, there can be more $`F`$’s than $`𝐲X`$’s.
Let $`𝒯`$ be a triangulated category with coproducts, and $`𝒯^c𝒯`$ the full subcategory of compact objects. We adopt the following notation:
The category $`𝒯`$ satisfies \[BRO\] if every exact functor $`\{𝒯^c\}^{op}𝒜b`$ is of the form $`𝐲X`$, for some $`X𝒯`$.
The category $`𝒯`$ satisfies \[BRM\] if every natural transformation $`𝐲X𝐲X^{}`$ is induced by a map $`XX^{}`$.
The theorem of Adams (see 0.2) says, that if $`𝒯`$ is the homotopy category of spectra, then both \[BRO\] and \[BRM\] hold in $`𝒯`$. In , Neeman found a necessary and sufficient condition for this to generalise, to arbitrary compactly generated $`𝒯`$’s. For this article, in the statements that follow, assume $`𝒯=D(R)`$ is the derived category of a ring $`R`$.
###### Theorem 0.7.
The following are equivalent:
1. Both \[BRM\] and \[BRO\] hold in $`𝒯`$
2. $`pgldim𝒯1`$.
The direction (i)$``$(ii) was also observed in . Beligiannis, using his Proposition 0.6 above, recently showed:
###### Theorem 0.8.
\[5, Theorem 11.8\] \[BRM\]$``$\[BRO\].
Neeman also showed that when $`R`$ is countable, \[BRM\] (and therefore also \[BRO\]) holds.
Keller produced the first example, where \[BRM\] fails. It may be found in Neeman’s . The example hinges on the following observation. If \[BRM\] holds, then by Theorem 0.7, we have $`pgldim𝒯1`$. That is, for any object $`X𝒯`$, $`𝐲X`$ has projective dimension at most $`1`$. If $`R`$ is a noetherian ring, this means that the cohomology modules $`H^iY`$ have pure projective dimension at most $`1`$. For a counterexample, one needs only produce an object $`Y𝒯=D(R)`$, so that its cohomology is of pure projective dimension greater than $`1`$.
The most recent progress preceding this article is a theorem of Beligiannis:
###### Theorem 0.9.
\[5, Remark 11.12\] \[BRO\] holds, whenever $`pgldim𝒯2`$.
This leaves several obvious questions:
What is the precise relation between the pure global dimension of $`R`$, denoted $`pgldimR`$, and the pure global dimension of $`𝒯`$, denoted $`pgldim𝒯`$?
Just how closely are the two related to \[BRM\] and \[BRO\]?
Does \[BRO\] hold in general?
In this article, we make progress on these questions. Regarding Q1, we prove that for many rings $`pgldimRpgldim𝒯`$, and that for hereditary rings this is an equality. Then we give examples to show that in general the inequality can be strict.
Regarding Q2, we give a precise relationship between pure global dimension, \[BRO\] and \[BRM\] for hereditary rings. Then we give examples to show that in general no such simple relationship holds. At the same time we show that \[BRO\] can fail, answering Q3. For example, it fails for $`R=k[x,y]`$ when $`k`$ has cardinality at least $`\mathrm{}_3`$ (Example 2.12).
Here is a more detailed overview of these results. We begin with an easy proposition giving our positive results about Q1. It is followed by a description of our counterexamples. We end with our positive results about Q2.
Proposition 1.4
1. Suppose that $`R`$ is a coherent ring, and that all finitely presented $`R`$–modules are of finite projective dimension. (This hypothesis holds when $`R`$ is noetherian of finite global dimension.) Then we have
$$pgldimRpgldimD(R).$$
2. Suppose that $`R`$ is hereditary. Then we have
$$pgldimR=pgldimD(R).$$
Weaker versions of this proposition were known before, and the inequality was after all at the basis of Keller’s counterexample to \[BRM\]. The really new result we show in this article is that, for some $`R`$, the inequality can be strict; Example 1.5 gives such an $`R`$. The idea of the counterexample is to produce two rings $`R`$ and $`S`$, of different pure global dimensions, but with $`D(R)D(S)`$. Then $`pgldimD(R)=pgldimD(S)`$ must be at least the maximum, and strictly bigger than the minimum, of $`pgldimR`$ and $`pgldimS`$. These rings, due to Thomas Bruestle, are finite-dimensional non-commutative $`k`$-algebras described by means of quivers.
Even more, we show that in general the answer to Q3 is negative: \[BRO\] can fail. It fails for the rings $`R`$ and $`S`$ mentioned above when the cardinality of $`k`$ is at least $`\mathrm{}_2`$, for the ring $`k[x,y]`$ when $`|k|\mathrm{}_3`$, and also for the ring $`T=kX,Y`$ of polynomials in two non-commuting variables when $`|k|\mathrm{}_2`$. (In particular, since it is consistent with ZFC that $`||=\mathrm{}_3`$, it is impossible to prove \[BRO\] using ZFC when $`R=[x,y]`$.) The proof that these are counterexamples is presented in Section 2. Our method is to find an exact sequence
$$\begin{array}{ccccccccc}0& & 𝐲A& & F& & 𝐲B& & 0\end{array}$$
in $`\text{Mod-}𝒯^c`$, and show that $`F`$ is not isomorphic to $`𝐲Y`$ for any $`Y`$. The idea is to study the extension group $`Ext^1(𝐲B,𝐲A)`$. We get a handle on this group using several spectral sequences. The precise statement of our theorem is:
Theorem 2.11. Let $`R`$ be an associative ring. Assume that $`R`$ is coherent, and that every finitely presented $`R`$–module has a finite projective resolution. Suppose there exists an $`R`$–module $`N`$ so that
$$\text{pure inj dim}(N)\text{inj dim}(N)2.$$
Then \[BRO\] fails for in $`D(R)`$. This means that there exists a homological functor $`F:\{𝒯^c\}^{op}𝒜b`$, which is not the restriction of any representable. That is, there exists no $`Y`$ with $`𝐲Y=F`$.
What is mysterious here, is that given a homological $`F`$, we cannot directly tell whether it is of the form $`𝐲X`$. We have no criterion to distinguish $`𝐲X`$’s from other homological functors. In fact, Beligiannis’ Proposition 0.6 tells us, that given any homological $`F`$, there exists a $`𝐲X`$ of projective dimension greater than or equal to that of $`F`$; projective dimension will not distinguish $`𝐲X`$’s from other homological functors. What we do amounts to finding a trick, to get around this problem.
For general rings, this is all we can say. We can give a refinement of the results for hereditary rings $`R`$; recall that $`R`$ is hereditary if its global dimension is $`1`$. Examples of hereditary rings are commutative principal ideal domains, and non-commutative polynomial rings.
Theorem 2.13. Let $`R`$ be a hereditary ring. Then
1. \[BRM\] holds in $`𝒯`$ if and only if the pure global dimension of $`R`$ is at most $`1`$; and
2. \[BRO\] holds in $`𝒯`$ if and only if the pure global dimension of $`R`$ is at most $`2`$.
We conclude the paper with the observation (Lemma 2.14) that any counterexample to \[BRO\] must take values in infinite-dimensional vector spaces.
*Acknowledgements.* The authors would like to thank Apostolos Beligiannis, Thomas Brüstle, Henning Krause and Michel Van den Bergh for helpful conversations. The first and second authors thank the third author, and the Centre for Mathematics and its Applications at the Australian National University, for providing a friendly and productive setting while this work was carried out.
## 1. Pure global dimension: module categories versus derived categories
Let $`R`$ be an associative ring. We denote by $`𝒯`$ the unbounded derived category $`D(R)`$ of the category of (right) $`R`$–modules, and by $`𝒯^c`$ the full subcategory of compact objects. Recall that a complex is a compact object of $`𝒯`$ iff it is quasi-isomorphic to a bounded complex of finitely generated projective $`R`$–modules. Here and elsewhere, we identify the category $`\text{Mod-}R`$ of $`R`$–modules with the subcategory of $`𝒯`$ consisting of complexes concentrated in degree $`0`$.
###### Lemma 1.1.
The following are equivalent:
1. $`R`$ is coherent and each finitely presented $`R`$–module is of finite projective dimension.
2. Each finitely presented $`R`$–module is compact when viewed as an object of $`D(R)`$.
3. A complex $`X`$ is compact iff each $`H^nX`$ is finitely presented and $`H^nX0`$ for all but finitely many $`n`$.
###### Remark 1.2.
In particular, the conditions of the lemma are satisfied if $`R`$ is noetherian and of finite global dimension. They are also satisfied by any hereditary ring, that is, any ring of global dimension at most 1.
###### Proof.
We will prove (i)$``$(ii), and then that (i)+(ii)$``$(iii). But first, we remind the reader that a ring is coherent iff the kernel of every map between finitely generated projective modules is finitely presented. We will also use the easy fact that a module is a compact object of $`D(R)`$ iff it admits a finite resolution by finitely generated projective objects.
Assume (i) holds. Let $`M`$ be a finitely presented module. Since $`R`$ is coherent, $`M`$ admits a resolution by finitely generated projective modules. Since $`M`$ is of finite projective dimension, this resolution may be chosen to be finite. So $`M`$ is compact in $`D(R)`$. That is, (ii) follows.
Suppose that (ii) holds. Then each finitely presented module admits a finite resolution by finitely generated projectives, and so in particular has finite projective dimension. Now let $`K`$ be the kernel of a map $`f:P_1P_0`$ between finitely generated projectives. Let $`C`$ be the cokernel of $`f`$. In $`D(R)`$, we have the canonical triangle
$$\mathrm{\Sigma }KPC\mathrm{\Sigma }^2K,$$
where $`P`$ is the complex $`P_1P_0`$. By assumption, $`P`$ and $`C`$ are compact. Hence $`K`$ is compact. So it admits a finite resolution by finitely generated projective objects. In particular, it is finitely presented. Thus $`R`$ is coherent; (i) holds.
Thus far, we have proved (i)$``$(ii). Assume these equivalent conditions hold; we wish to prove (iii). Let $`X`$ be a compact object in $`D(R)`$. It is isomorphic to a finite complex of finitely generated projective modules. By (i), $`R`$ is coherent; hence $`H^nX`$ is finitely presented for all $`n`$. And since the complex $`X`$ is finite, $`H^nX0`$ for all but finitely many $`n`$.
Suppose now that $`H^nX`$ is finitely presented for all $`n`$, and that $`H^nX0`$ for all but finitely many $`n`$. The $`t`$–structure on $`D(R)`$ gives us triangles
$$\begin{array}{ccccccc}X^n& & X& & X^{>n}& & \mathrm{\Sigma }X^n\end{array}$$
and these allow us to assemble $`X`$ from its homology. Now $`H^nX`$ is finitely presented for all $`n`$, and by (ii) it is compact. This forces $`X`$, an iterated extension of compact objects, to also be compact. We conclude that (iii) holds.
Finally, (iii)$``$(ii) is immediate. ∎
Recall that the functor $`𝐲:𝒯\text{Mod-}𝒯^c`$ sends an object $`X𝒯`$ to the functor
$$𝐲X=𝒯(,X)|_{𝒯^c}.$$
For $`i`$ and $`F\text{Mod-}𝒯^c`$, we define the $`i^{th}`$ homology of $`F`$ by
$$H^iF=F(\mathrm{\Sigma }^iR).$$
The functor $`H^i:\text{Mod-}𝒯^c\text{Mod-}R`$ extends the homology functor on $`𝒯`$ in the sense that we have a canonical isomorphism $`H^i𝐲=H^i`$.
An object $`G`$ in the category $`\text{Mod-}𝒯^c`$ is called finitely presented, if there exists an exact sequence
$$𝐲X𝐲YG0.$$
The full subcategory of all finitely presented objects in $`\text{Mod-}𝒯^c`$ is known to be an abelian category; see for example Freyd’s . As in the case of a module category, a sequence
$$0F_1F_2F_30$$
of $`\text{Mod-}𝒯^c`$ is called pure exact if the sequence
$$0Hom(G,F_1)Hom(G,F_2)Hom(G,F_3)0$$
is exact for each finitely presented functor $`G`$. (In particular, the sequence is then exact.)
###### Lemma 1.3.
Suppose that the conditions of Lemma 1.1 hold.
1. The functor $`𝐲:\text{Mod-}R\text{Mod-}𝒯^c`$ commutes with filtered colimits. It takes pure projective $`R`$–modules to projective objects of $`\text{Mod-}𝒯^c`$. It transforms pure exact sequences of $`R`$–modules into pure exact sequences in $`\text{Mod-}𝒯^c`$.
2. For each $`i`$, the functor $`H^i`$ commutes with filtered colimits. It takes projective objects of $`\text{Mod-}𝒯^c`$ to pure projective $`R`$–modules. It transforms pure exact sequences of $`\text{Mod-}𝒯^c`$ into pure exact sequences of $`R`$–modules.
###### Proof.
(i) Let $`M_\lambda `$ be a filtered system of $`R`$–modules. Clearly, if $`P=\mathrm{\Sigma }^iR`$ for some $`i`$, the canonical map
$$\underset{}{colim}𝒯(P,M_\lambda )𝒯(P,\underset{}{colim}M_\lambda )$$
is bijective. Since both sides are cohomological functors of $`P`$, this map is still bijective if $`P`$ is any compact object of $`𝒯`$, since $`𝒯^c`$ is the thick subcategory generated by $`R`$. This means that $`𝐲`$ takes $`\underset{}{colim}M_\lambda `$ to $`\underset{}{colim}𝐲M_\lambda `$.
Each pure projective $`R`$–module is a direct factor of a coproduct of finitely presented modules. Since the functor $`𝐲`$ commutes with coproducts, it is enough to show that $`𝐲M`$ is projective if $`M`$ is finitely presented. But in this case, $`M`$ is compact in $`𝒯`$, by our assumption on the ring $`R`$. So $`𝐲M`$ is projective since it is even representable.
Now let
$$0LMN0$$
be a pure exact sequence of $`R`$–modules. Clearly, if $`N`$ is finitely presented, the sequence splits. An arbitrary module $`N`$ is a filtered colimit of finitely presented modules. Thus the sequence is a filtered colimit of split sequences. Since the functor $`𝐲`$ commutes with filtered colimits, the image of the sequence is also a filtered colimit of split sequences. Thus it is pure.
(ii) By definition, the functor $`H^i`$ is evaluation at $`\mathrm{\Sigma }^iR`$. Thus it commutes with colimits. The projective objects of $`\text{Mod-}𝒯^c`$ are direct factors of coproducts of representable functors, and the functor $`H^i`$ commutes with coproducts. So it is enough to show that $`H^i𝐲P=H^iP`$ is pure projective for $`P𝒯^c`$. This is clear since $`H^iP`$ is finitely presented, by our assumption on the ring $`R`$.
Let
$$0F_1F_2F_30$$
be a pure exact sequence of $`\text{Mod-}𝒯^c`$. Clearly if $`F_3`$ is finitely presented, the sequence splits. In the general case, $`F_3`$ is a filtered colimit of a system of finitely presented functors. So the sequence is a filtered colimit of split sequences. Since the functor $`H^i`$ commutes with filtered colimits, this implies the last assertion. ∎
The pure global dimension of the derived category $`D(R)=𝒯`$ is by definition the supremum of the projective dimensions of the functors $`𝐲X`$, $`X𝒯`$. We write $`pgldim`$ for ‘pure global dimension’. Part (ii) of the following lemma is due to Beligiannis \[5, Prop. 12.8\].
###### Proposition 1.4.
Suppose that the conditions of Lemma 1.1 hold.
1. Let $`M`$ be an $`R`$–module. Then the projective dimension of $`𝐲M`$ equals the pure projective dimension of $`M`$. Hence we have
$$pgldimRpgldimD(R).$$
2. Suppose that $`R`$ is hereditary. Then we have
$$pgldimR=pgldimD(R).$$
###### Proof.
(i) The first part of the preceding lemma shows that the functor $`𝐲`$ takes pure projective resolutions of a module $`M`$ to projective resolutions of $`𝐲M`$. Hence the projective dimension of $`𝐲M`$ is no more than the pure projective dimension of $`M`$. Conversely, let
$$\mathrm{}Q_1Q_0𝐲M0$$
be a projective resolution of $`𝐲M`$. If $`M`$ is finitely presented, then $`𝐲M`$ is projective, so the resolution is nullhomotopic. An arbitrary $`M`$ is still a filtered colimit of finitely presented modules. So for arbitrary $`M`$ the resolution is a filtered colimit of nullhomotopic complexes. Thus it is a pure exact sequence. By the second part of the above lemma, its image under $`H^0`$ is a pure projective resolution of $`H^0𝐲M=M`$. Thus the pure projective dimension of $`M`$ is no more than the projective dimension of $`𝐲M`$.
(ii) By part (i), it suffices to prove that $`pgldimRpgldimD(R)`$. Let $`XD(R)`$. Since $`R`$ is hereditary, the object $`X`$ is isomorphic in $`D(R)`$ to the coproduct of the $`\mathrm{\Sigma }^iH^iX`$, $`i`$; see Lemma 6.7, on page 153 of . Hence the projective dimension of $`𝐲X`$ is no greater than the supremum of the projective dimensions of the $`𝐲H^iX`$. These are bounded by $`pgldimR`$ thanks to part (i). ∎
###### Example 1.5.
Let $`k`$ be a field and let $`t`$ be the cardinal such that $`\mathrm{}_t=\mathrm{max}(|k|,\mathrm{}_0)`$. So $`t`$ is $`0`$ if $`k`$ is finite or countable, $`1`$ if $`k`$ has the smallest uncountable cardinality, etc. Building on an example due to Th. Bruestle we will exhibit a $`k`$-algebra $`R`$ such that the inequality
$$pgldimRpgldimD(R)$$
is strict. Our example is based on the observation that there are algebras with equivalent derived categories but widely differing pure global dimensions. More precisely, we will exhibit a finite-dimensional $`k`$-algebra $`R`$ with $`pgldimR=0`$ such that $`D(R)`$ is triangle equivalent to $`D(S)`$ for a finite-dimensional hereditary $`k`$-algebra $`S`$ whose pure global dimension is $`t+1`$ ($`\mathrm{}`$ if $`t`$ is infinite). Thus we have
$$pgldimR<pgldimS=pgldimD(S)=pgldimD(R),$$
where we have used part (ii) of the above proposition for the first equality.
Thus Theorem 2.13 implies that \[BRM\] fails for $`D(R)`$ when $`t1`$ and that \[BRO\] fails for $`D(R)`$ when $`t2`$, even though $`R`$ has pure global dimension $`0`$.
The algebras $`R`$ and $`S`$ are due to Th. Bruestle. We will define them using the language of quivers with relations (cf. , , ). Here is all we need: A quiver is an oriented graph. It is thus given by a set $`Q_0`$ of points, a set $`Q_1`$ of arrows, and two maps $`s,t:Q_1Q_0`$ associating with each arrow its source and its target. A simple example is the quiver
$$\stackrel{}{A}_{10}:1\stackrel{\alpha _1}{}2\stackrel{\alpha _2}{}3\mathrm{}8\stackrel{\alpha _8}{}9\stackrel{\alpha _9}{}10.$$
A path in a quiver $`Q`$ is a sequence $`(y|\beta _r|\beta _{r1}|\mathrm{}|\beta _1|x)`$ of composable arrows $`\beta _i`$ with $`s(\beta _1)=x`$, $`s(\beta _i)=t(\beta _{i1})`$, $`2ir`$, $`t(\beta _r)=y`$. In particular, for each point $`xQ_0`$, we have the lazy path $`(x|x)`$. It is neutral for the obvious composition of paths. The quiver algebra $`kQ`$ has as its basis all paths of $`Q`$. The product of two basis elements equals the composition of the two paths if they are composable and $`0`$ otherwise. For example, the quiver algebra of $`Q=\stackrel{}{A}_{10}`$ is isomorphic to the algebra of lower triangular $`10\times 10`$ matrices.
The construction of the quiver algebra $`kQ`$ is motivated by the (easy) fact that the category of left $`kQ`$-modules is equivalent to the category of all diagrams of vector spaces of the shape given by $`Q`$. It is not hard to show that each quiver algebra is hereditary. It is finite-dimensional over $`k`$ iff the quiver has no oriented cycles.
Gabriel showed that the quiver algebra of a finite quiver has only a finite number of $`k`$–finite-dimensional indecomposable modules (up to isomorphism) iff the underlying graph of the quiver is a disjoint union of Dynkin diagrams of type $`A`$, $`D`$, $`E`$.
The above example has underlying graph of Dynkin type $`A_{10}`$ and thus its quiver algebra has only a finite number of finite-dimensional indecomposable modules.
An ideal $`I`$ of a finite quiver $`Q`$ is admissible if for some $`N`$ we have
$$(kQ_1)^NI(kQ_1)^2,$$
where $`(kQ_1)`$ is the two-sided ideal generated by all paths of length $`1`$. A *quiver $`Q`$ with relations $`R`$* is a quiver $`Q`$ with a set $`R`$ of generators for an admissible ideal $`I`$ of $`kQ`$. The algebra $`kQ/I`$ is then the *algebra associated with $`(Q,R)`$*. Its category of left modules is equivalent to the category of diagrams of vector spaces of shape $`Q`$ obeying the relations in $`R`$. The algebra $`kQ/I`$ is finite-dimensional (since $`I`$ contains all paths of length at least $`N`$), hence artinian and noetherian. By induction on the number of points one can show that if the quiver $`Q`$ contains no oriented cycle, then the algebra $`kQ/I`$ is of finite global dimension.
One can show that every finite-dimensional algebra over an algebraically closed field is Morita equivalent to the algebra associated with a quiver with relations and that the quiver is unique (up to isomorphism).
Now we let $`R`$ be the finite-dimensional $`k`$-algebra associated with the above quiver $`\stackrel{}{A}_{10}`$ and the relation $`\alpha _8\alpha _7\mathrm{}\alpha _1`$ (no $`\alpha _9`$!). The algebra $`R`$ is a quotient of $`k\stackrel{}{A}_{10}`$ and thus it admits only a finite number of indecomposable finite-dimensional modules. By a result of Auslander and Tachikawa , this is equivalent to $`pgldimR=0`$.
Let $`S`$ be the quiver algebra of the quiver
$$E:\begin{array}{ccccccccc}2\hfill & 3\hfill & 4\hfill & 5\hfill & 6\hfill & 7\hfill & 8\hfill & 9\hfill & 10\hfill \\ & & & & & & \hfill & & \\ & & & & & & 1.\hfill & & \end{array}$$
Thus $`S`$ is finite-dimensional over $`k`$ and hereditary. By Theorem 4.1 of Baer-Lenzing’s , we have $`pgldimS=t+1`$ ($`\mathrm{}`$ if $`t`$ is infinite).
Finally, we need to show that $`R`$ and $`S`$ have equivalent derived categories. Indeed, the algebra $`R`$ admits a tilting complex with endomorphism ring $`S`$ so that the equivalence follows from Rickard’s Morita theorem for derived categories . To describe the tilting complex, let $`P_i=e_iR`$ be the projective $`R`$–module associated with the idempotent $`e_i=(i|i)`$ (the lazy path). It is easy to compute the morphism spaces between these modules: Indeed, we have $`Hom(e_iR,e_jR)=e_jRe_i`$ and this space identifies with the vector space on the set of paths from $`i`$ to $`j`$ divided by the subspace of linear combinations of paths lying in the ideal of relations. For example, for $`ij`$, the path from $`i`$ to $`j`$ yields a canonical morphism $`P_iP_j`$, which vanishes iff $`(i,j)=(1,9)`$ or $`(i,j)=(1,10)`$. The tilting complex $`T`$ is now the sum of the complexes
$`T_2=(P_1P_2),T_3=(P_1P_3),\mathrm{},T_8=(P_1P_8),`$
$`T_1=(P_10),T_9=(0P_9),T_{10}=(0P_{10}),`$
where the first term of each complex is in degree $`0`$. Using the description of the morphism spaces between the $`P_i`$ it is not hard to check that, in the homotopy category of right $`R`$-modules, we do have $`Hom(T_i,T_j[l])=0`$ for all $`i,j`$ and all $`l0`$, and that the endomorphism ring of $`T`$ is indeed isomorphic to $`S`$. For example, the canonical idempotent $`(i|i)`$ of the quiver $`E`$ corresponds to the idempotent of $`End(T)`$ arising from the identity of $`T_i`$ and the arrow $`89`$ of $`E`$ corresponds to the obvious morphism of complexes
$$\begin{array}{ccc}\hfill P_1& & P_8\hfill \\ \hfill & & \hfill \\ \hfill 0& & P_9\hfill \end{array}$$
which is well-defined thanks to the relation $`\alpha _8\alpha _7\mathrm{}\alpha _1`$ that we imposed.
## 2. Failure of Brown representability
In this section, $`R`$ will be a ring satisfying the equivalent conditions of Lemma 1.1. In particular, all the theorems hold if $`R`$ is a noetherian ring of finite global dimension, or if $`R`$ is hereditary. We begin by reminding ourselves of a standard spectral sequence.
###### Lemma 2.1.
Let $`𝒜`$ be an abelian category satisfying AB5, and with enough projectives. Suppose that $`X`$ and $`Y`$ are objects of $`𝒜`$ and that $`X=\underset{}{colim}X_\lambda `$ expresses $`X`$ as a filtered colimit of objects $`X_\lambda 𝒜`$. Then there is a spectral sequence, converging to $`Ext^{i+j}(X,Y)`$, whose $`E^2`$ term is
$$\underset{}{lim}^i\stackrel{j}{Ext}(X_\lambda ,Y).$$
###### Proof.
There is a standard chain complex which computes the derived functors of $`\underset{}{colim}`$. Since the abelian category $`𝒜`$ satisfies AB5, the derived functors of filtered colimits vanish, and we deduce an exact sequence in $`𝒜`$
$$\begin{array}{ccccccccc}\mathrm{}& & \underset{\lambda \mu }{}X_\lambda & & \underset{\lambda }{}X_\lambda & & X& & 0.\end{array}$$
This gives us a resolution of $`X`$ in $`𝒜`$, and the spectral sequence is just the spectral sequence of the functor $`Ext^{}(,Y)`$ applied to this resolution. ∎
In the following, we write $`\text{mod-}R`$ for the category of finitely presented $`R`$–modules and $`\text{Mod}(\text{mod-}R)`$ for the category of contravariant additive functors from $`\text{mod-}R`$ to $`𝒜b`$. The object
$$\text{Mod-}R(,M)|_{\text{mod-}R}$$
of $`\text{Mod}(\text{mod-}R)`$ will be denoted $`𝐳M`$.
###### Lemma 2.2.
Let $`R`$ be a ring, and let $`M_\lambda `$ be a filtered diagram of $`R`$–modules with colimit $`M`$. Then
1. $`𝐲M=\underset{}{colim}𝐲M_\lambda `$ in $`\text{Mod-}𝒯^c`$.
2. $`𝐳M=\underset{}{colim}𝐳M_\lambda `$ in $`\text{Mod}(\text{mod-}R)`$.
###### Proof.
(i) was proved in Lemma 1.3 (i). The second statement is more familiar in the equivalent form, which states that $`\text{Mod-}R(K,M)=\underset{}{colim}\text{Mod-}R(K,M_\lambda )`$ for any finitely presented $`K`$. This is not hard to prove. ∎
###### Reminder 2.3.
Let $`R`$ be a ring and let $`M`$ be an $`R`$–module. Consider the filtered diagram of finitely presented modules $`M_\lambda `$ equipped with a map to $`M`$. Then $`M`$ is the colimit of this diagram; we already used this in the proof of Proposition 1.4(i). This is the setting in which we will apply Lemma 2.2.
The following lemma is well known; the proof may be found, for example, in Theorem 2.8 of Simson’s . We include a sketch of the proof for the reader’s convenience.
###### Lemma 2.4.
Let $`R`$ be a ring satisfying the conditions of Lemma 1.1, and let $`M`$ be an $`R`$–module. As mentioned in Remark 2.3, $`M`$ is the filtered colimit of all finitely presented modules $`M_\lambda `$ mapping to $`M`$.
1. Let $`F`$ be an object of $`\text{Mod-}𝒯^c`$. That is, $`F`$ is a functor $`\{𝒯^c\}^{op}𝒜b`$. Then the group $`Ext^i(𝐲M,F)`$ of extensions in $`\text{Mod-}𝒯^c`$ is isomorphic to $`\underset{}{lim}^iF(M_\lambda )`$.
2. Let $`F`$ be an object of $`\text{Mod}(\text{mod-}R)`$. That is, $`F`$ is a functor $`\{\text{mod-}R\}^{op}𝒜b`$. Then the group $`Ext^i(𝐳M,F)`$ of extensions in $`\text{Mod}(\text{mod-}R)`$ is isomorphic to $`\underset{}{lim}^iF(M_\lambda )`$.
###### Proof.
(i): By Lemma 2.2, $`𝐲M`$ is the colimit of $`𝐲M_\lambda `$ in $`\text{Mod-}𝒯^c`$. Lemma 2.1 then tells us that we get a spectral sequence with $`E^2`$ term
$$\underset{}{lim}^i\stackrel{j}{Ext}(𝐲M_\lambda ,F)$$
converging to the group $`Ext^{i+j}(𝐲M,F)`$ of extensions in $`\text{Mod-}𝒯^c`$. The functor $`𝐲M_\lambda `$ is representable, since by our hypothesis on $`R`$ the module $`M_\lambda `$ is compact. Thus $`𝐲M_\lambda `$ is projective, the $`Ext^j`$ terms vanish unless $`j=0`$, the spectral sequence collapses, and the desired isomorphism follows.
The proof of (ii) is similar. ∎
###### Remark 2.5.
In part (i) of Lemma 2.4, we computed the extensions of $`𝐲M`$ by $`F`$. This interests us most in the case where $`F=𝐲\mathrm{\Sigma }^jN`$, with $`N`$ an $`R`$–module. In this case, the computation tells us that we have isomorphisms
$$\stackrel{i}{Ext}(𝐲M,𝐲\mathrm{\Sigma }^jN)=\underset{}{lim}^i𝒯(M_\lambda ,\mathrm{\Sigma }^jN)=\underset{}{lim}^i\underset{R}{\overset{j}{Ext}}(M_\lambda ,N).$$
In part (ii) of Lemma 2.4, we computed the extensions of $`𝐳M`$ by $`F`$. This interests us most in the case where $`F=𝐳N`$, with $`N`$ an $`R`$–module. In this case, the computation tells us that we have an isomorphism
$$\stackrel{i}{Ext}(𝐳M,𝐳N)=\underset{}{lim}^i\underset{R}{Hom}(M_\lambda ,N).$$
Moreover the group $`Ext^i(𝐳M,𝐳N)`$ of extensions in $`\text{Mod}(\text{mod-}R)`$ can be identified with the group $`PExt^i(M,N)`$; see . We deduce that
$$\stackrel{i}{PExt}(M,N)=\underset{}{lim}^i\underset{R}{Hom}(M_\lambda ,N).$$
###### Corollary 2.6.
If $`M`$ and $`N`$ are $`R`$–modules and $`j>0`$, then every map $`𝐲\mathrm{\Sigma }^jM𝐲N`$ vanishes. Moreover, maps $`𝐲M𝐲N`$ are in one-to-one correspondence with maps of $`R`$–modules $`MN`$.
###### Proof.
For $`j>0`$, we must show that any map $`𝐲M𝐲\mathrm{\Sigma }^jN`$ vanishes. But by Remark 2.5, the group of such maps is
$$\underset{}{lim}^0\underset{R}{\overset{j}{Ext}}(M_\lambda ,N),$$
which vanishes because there are no extensions of negative degree.
The group of maps $`𝐲M𝐲N`$ is exactly
$$\underset{}{lim}^0\underset{R}{\overset{0}{Ext}}(M_\lambda ,N),$$
which is $`Hom_R(M,N)`$. ∎
###### Lemma 2.7.
Let $`F`$ be an object in $`\text{Mod-}𝒯^c`$, that is, a contravariant additive functor from $`𝒯^c`$ to $`𝒜b`$. Suppose there exists an integer $`j>0`$, $`R`$–modules $`M`$ and $`N`$, and a short exact sequence in $`\text{Mod-}𝒯^c`$
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^jN& \stackrel{\alpha }{}& F& \stackrel{\beta }{}& 𝐲M& & 0.\end{array}$$
Then this sequence is unique up to isomorphism.
###### Proof.
The integer $`j`$ and the modules $`M`$ and $`N`$ are clearly determined by the homology of $`F`$. In Corollary 2.6 we saw that any map $`𝐲\mathrm{\Sigma }^jN𝐲M`$ vanishes. Therefore, given any map $`\gamma :𝐲\mathrm{\Sigma }^jNF`$, the composite
$$\begin{array}{ccccc}𝐲\mathrm{\Sigma }^jN& \stackrel{\gamma }{}& F& \stackrel{\beta }{}& 𝐲M\end{array}$$
vanishes, and hence $`\gamma `$ must factor through $`\alpha `$. Dually, any map $`F𝐲M`$ must factor through $`\beta `$. This shows that the given exact sequence is unique. ∎
###### Lemma 2.8.
Let $`F`$ be an object of $`\text{Mod-}𝒯^c`$, and suppose there exists an integer $`j>0`$, $`R`$–modules $`M`$ and $`N`$, and a short exact sequence in $`\text{Mod-}𝒯^c`$
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^jN& \stackrel{\alpha }{}& F& \stackrel{\beta }{}& 𝐲M& & 0.\end{array}$$
The functor $`F`$ will be of the form $`𝐲Y`$ if and only if the short exact sequence comes from a triangle. That is, if and only if there exists a triangle in $`𝒯`$
$$\begin{array}{ccccccc}\mathrm{\Sigma }^jN& & Y& & M& \stackrel{}{}& \mathrm{\Sigma }^{j+1}N\end{array}$$
with $``$ a phantom map, so that the sequence
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^jN& \stackrel{\alpha }{}& F& \stackrel{\beta }{}& 𝐲M& & 0\end{array}$$
is obtained by restricting the representable functors to $`𝒯^c`$.
We remind the reader that a map $`WX`$ in $`𝒯`$ is called *phantom* if the composite $`CWX`$ is zero for each compact object $`C`$ and each map $`CW`$.
###### Proof.
The implication $``$ is trivial. If the triangle exists and is isomorphic to the short exact sequence of functors on $`𝒯^c`$, then $`F`$ is the restriction of a representable functor on $`𝒯`$. We wish to prove $``$. We suppose therefore that the short exact sequence of functors is given, and that $`F`$ is the restriction of a representable. We want to produce a triangle.
The short exact sequence
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^jN& \stackrel{\alpha }{}& F& \stackrel{\beta }{}& 𝐲M& & 0\end{array}$$
permits us easily to compute $`F(\mathrm{\Sigma }^nR)`$, for all $`n`$. We have
$$F(\mathrm{\Sigma }^nR)=\{\begin{array}{ccc}M\hfill & & \text{if }n=0\hfill \\ N\hfill & & \text{if }n=j\hfill \\ 0\hfill & & \text{otherwise. }\hfill \end{array}$$
But if $`F=𝐲Y`$, then $`F(\mathrm{\Sigma }^nR)=H^n(Y)`$. The above computes for us the cohomology of $`Y`$, as an object in $`D(R)=𝒯`$.
There is a $`t`$-structure truncation on $`D(R)`$, giving a triangle
$$\begin{array}{ccccccc}Y^1& & Y& & Y^0& \stackrel{}{}& \mathrm{\Sigma }Y^1,\end{array}$$
and our homology computation shows that $`Y^1`$ and $`Y^0`$ each have only one non-zero cohomology group. The triangle is therefore of the form
$$\begin{array}{ccccccc}\mathrm{\Sigma }^jN& & Y& & M& \stackrel{}{}& \mathrm{\Sigma }^{j+1}N.\end{array}$$
We deduce an exact sequence
$$\begin{array}{ccccc}𝐲\mathrm{\Sigma }^jN& & 𝐲Y& & 𝐲M.\end{array}$$
Now recall that $`𝐲Y=F`$, and that by the proof of Lemma 2.7, any map $`𝐲\mathrm{\Sigma }^jNF`$ factors through $`\alpha `$, and any map $`F𝐲M`$ factors through $`\beta `$. The exact sequence coming from the triangle therefore factors through
$$\begin{array}{ccccccccc}& & 𝐲\mathrm{\Sigma }^jN& & & & & & \\ & & f& & & & & & & & \\ 0& & 𝐲\mathrm{\Sigma }^jN& \stackrel{\alpha }{}& F& \stackrel{\beta }{}& 𝐲M& & 0.\\ & & & & & & g& & & & \\ & & & & & & 𝐲M& & \end{array}$$
By Corollary 2.6, the morphisms $`f`$ and $`g`$ in the diagram above come from maps of modules $`NN`$ and $`MM`$. Evaluating the functors at $`R`$ and $`\mathrm{\Sigma }^jR`$, we compute that both $`f`$ and $`g`$ are isomorphisms. Hence the triangle gives rise to the short exact sequence of functors, and $``$ must be a phantom map. ∎
Next comes a spectral sequence argument. To help the reader, we will first do the easy, baby case.
###### Proposition 2.9.
Let $`R`$ be a ring satisfying the conditions of Lemma 1.1. Let $`N`$ be an $`R`$–module with injective dimension at most $`1`$ and pure injective dimension at least $`3`$. Then in $`\text{Mod-}𝒯^c`$ there exists a homological functor $`F:\{𝒯^c\}^{op}𝒜b`$ which is not the restriction of any representable. That is, there exists no $`Y`$ with $`𝐲Y=F`$.
###### Example 2.10.
Let $`k`$ be a field and $`R`$ the algebra of the quiver $`E`$ of Example 1.5 (we called it $`S`$ there). Then $`R`$ is finite-dimensional over $`k`$ and hereditary, since it is the quiver algebra of a finite quiver. So all $`R`$–modules are of injective dimension at most $`1`$. Assume that $`k`$ is infinite of cardinality $`\mathrm{}_t`$. Then by , the pure global dimension of $`R`$ equals $`t+1`$ ($`\mathrm{}`$ if $`t`$ is infinite). Thus when $`t2`$ there does exist an $`R`$–module satisfying the assumptions of the proposition.
Similarly, the ring $`kX,Y`$ of polynomials in two non-commuting variables is an example when $`t2`$.
To obtain examples where $`R`$ is commutative, we will need to use Theorem 2.11, which is a refined version of the above proposition.
###### Proof.
Because $`N`$ is of pure injective dimension at least $`3`$, there exists a module $`M`$ and integer $`n3`$, so that $`PExt^n(M,N)0`$. If $`n>3`$, choose a pure exact sequence
$$\begin{array}{ccccccccc}0& & M^{}& & P& & M& & 0,\end{array}$$
with $`P`$ pure projective. Then $`PExt^n(M,N)=PExt^{n1}(M^{},N)`$. By a sequence of such dimension shifts, we may find an $`M`$ so that
$$\stackrel{3}{PExt}(M,N)0.$$
By Remark 2.3, we may express $`M`$ as a filtered colimit of finitely presented modules $`M_\lambda `$. By Lemma 2.1, applied this time to the category of $`R`$–modules, there is a spectral sequence with $`E^2`$ term $`\underset{}{lim}^iExt_R^j(M_\lambda ,N)`$ converging to $`Ext_R^{i+j}(M,N)`$. We will now compute in this spectral sequence.
In Remark 2.5, we computed that
$$\underset{}{lim}^3\underset{R}{\overset{0}{Ext}}(M_\lambda ,N)=\stackrel{3}{PExt}(M,N),$$
and by the above, this does not vanish. On the other hand, we know that $`Ext_R^3(M,N)=0`$, since by hypothesis $`N`$ is of injective dimension at most $`1`$. It follows that one of the differentials in the spectral sequence into the term
$$\underset{}{lim}^3\underset{R}{\overset{0}{Ext}}(M_\lambda ,N)$$
must be non-zero.
But there are only two differentials into this term, one from $`\underset{}{lim}^1Ext^1`$ and one from $`\underset{}{lim}^0Ext^2`$. The latter vanishes, since by hypothesis $`N`$ is of injective dimension at most $`1`$. It follows that
$$\underset{}{lim}^1\underset{R}{\overset{1}{Ext}}(M_\lambda ,N)0.$$
But in Lemma 2.4 we showed that this is the group of extensions, in $`\text{Mod-}𝒯^c`$,
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }N& & F& & 𝐲M& & 0.\end{array}$$
The group does not vanish so we may choose a non-trivial extension. Since $`F`$ is the extension of two homological functors, $`F`$ must be homological. Now we will show that $`F`$ cannot be isomorphic to a functor $`𝐲Y`$.
Lemma 2.8 tells us that if $`F`$ is isomorphic to $`𝐲Y`$, then there is a triangle in $`𝒯`$
$$\begin{array}{ccccccc}\mathrm{\Sigma }N& & Y& & M& \stackrel{}{}& \mathrm{\Sigma }^2N\end{array}$$
so that the exact sequence of functors above is isomorphic to the one obtained from the triangle. But the map $`:M\mathrm{\Sigma }^2N`$ is an element of
$$\stackrel{2}{Ext}(M,N)=0,$$
and therefore the triangle splits. The exact sequence of functors is not split, and we conclude that $`F`$ cannot be isomorphic to any $`𝐲Y`$. ∎
The next Theorem is the more macho computation with the same spectral sequence.
###### Theorem 2.11.
Let $`R`$ be a ring satisfying the conditions of Lemma 1.1. Suppose there exists an $`R`$–module $`N`$ so that
$$\text{pure inj dim}(N)\text{inj dim}(N)2.$$
Then \[BRO\] fails for in $`D(R)`$. This means that there exists a homological functor $`F:\{𝒯^c\}^{op}𝒜b`$ which is not the restriction of any representable. That is, there exists no $`Y`$ with $`𝐲Y=F`$.
###### Proof.
Let $`N`$ be a module satisfying the hypotheses. Let $`n=\text{inj dim}(N)`$. Then $`\text{pure inj dim}(N)n+2`$. As in the proof of Proposition 2.9, we may choose a module $`M`$ with $`PExt^{n+2}(M,N)0.`$ We may also express $`M`$ as a filtered colimit of finitely presented modules $`M_\lambda `$.
Lemma 2.1 gives us a spectral sequence, whose $`E^2`$ term is
$$\underset{}{lim}^i\underset{R}{\overset{j}{Ext}}(M_\lambda ,N),$$
which converges to $`Ext_R^{i+j}(M,N)`$. Once again, we have that
$$\underset{}{lim}^{n+2}\underset{R}{\overset{0}{Ext}}(M_\lambda ,N)=\stackrel{n+2}{PExt}(M,N),$$
and this does not vanish, by the choice of $`M`$. But $`Ext_R^{n+2}(M,N)=0`$, since $`N`$ is of injective dimension at most $`n`$, so there must be a non-zero differential into the term
$$\underset{}{lim}^{n+2}\underset{R}{\overset{0}{Ext}}(M_\lambda ,N).$$
Now observe that
$$\underset{}{lim}^0\underset{R}{\overset{n+1}{Ext}}(M_\lambda ,N)=0,$$
since $`N`$ is of injective dimension at most $`n`$. It follows that for some $`i`$ with $`1in`$, there is a non-zero differential in the spectral sequence, from
$$\underset{}{lim}^i\underset{R}{\overset{n+1i}{Ext}}(M_\lambda ,N)$$
to the term $`\underset{}{lim}^{n+2}Ext_R^0(M_\lambda ,N)0`$.
Now recall the construction of our spectral sequence, from Lemma 2.1. Since $`M`$ is the filtered colimit of $`M_\lambda `$, there is an exact resolution of $`M`$
$$\begin{array}{ccccccccc}\mathrm{}& & \underset{\lambda \mu }{}M_\lambda & & \underset{\lambda }{}M_\lambda & & M& & 0.\end{array}$$
This resolution is a pure exact resolution by pure projectives. (It is pure exact because it remains exact in the category $`\text{Mod}`$($`\text{mod-}R`$). And direct sums of finitely presented modules $`M_\lambda `$ are pure projective.) By Lemma 1.3, it becomes an exact resolution by projectives in the category $`\text{Mod-}𝒯^c`$.
To simplify the notation, we will write the above resolution as
$$\begin{array}{ccccccccccc}\mathrm{}& & P_2& & P_1& & P_0& & M& & 0.\end{array}$$
Let $`K_i`$ stand for the image of the map $`P_iP_{i1}`$. In Lemma 2.4 we showed that
$$\underset{}{lim}^i\underset{R}{\overset{n+1i}{Ext}}(M_\lambda ,N)$$
is the group of extensions
$$\stackrel{i}{Ext}(𝐲M,𝐲\mathrm{\Sigma }^{n+1i}N).$$
But since the pure exact sequence
$$\begin{array}{ccccccccccccc}0& & K_{i1}& & P_{i2}& & \mathrm{}& & P_0& & M& & 0\end{array}$$
remains exact in $`\text{Mod-}𝒯^c`$, and the middle modules map to projectives in $`\text{Mod-}𝒯^c`$, we deduce that the above extension group is isomorphic to
$$\stackrel{1}{Ext}(𝐲K_{i1},𝐲\mathrm{\Sigma }^{n+1i}N).$$
In other words, an element of the group
$$\underset{}{lim}^i\underset{R}{\overset{n+1i}{Ext}}(M_\lambda ,N)$$
may be thought of as a short exact sequence in $`\text{Mod-}𝒯^c`$
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^{n+1i}N& & F& & 𝐲K_{i1}& & 0.\end{array}$$
We know that in the spectral sequence, for some $`1in`$, there is a non-zero differential
$$\begin{array}{ccc}\underset{}{lim}^iExt_R^{n+1i}(M_\lambda ,N)E& \stackrel{\gamma }{}& \underset{}{lim}^{n+2}Ext_R^0(M_\lambda ,N)\end{array},$$
for a subgroup $`E\underset{}{lim}^iExt_R^{n+1i}(M_\lambda ,N)`$. What we will now show is that, if $`\gamma (x)0`$, then $`x`$ corresponds to an exact sequence
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^{n+1i}N& & F& & 𝐲K_{i1}& & 0\end{array}$$
where $`F`$ is not isomorphic to any $`𝐲Y`$. Expressing the same thing slightly differently, we will show that if $`x\underset{}{lim}^iExt_R^{n+1i}(M_\lambda ,N)`$ comes from an exact sequence of functors with $`F=𝐲Y`$, then $`\gamma (x)=0`$.
Suppose therefore that we are given a short exact sequence in $`\text{Mod-}𝒯^c`$
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^{n+1i}N& & 𝐲Y& & 𝐲K_{i1}& & 0.\end{array}$$
We need to show that in the spectral sequence, the differential $`\gamma `$ annihilates $`x`$. By Lemma 2.8, the exact sequence of functors comes from a triangle
$$\begin{array}{ccccccc}\mathrm{\Sigma }^{n+1i}N& & Y& & K_{i1}& \stackrel{}{}& \mathrm{\Sigma }^{n+2i}N\end{array}$$
with $``$ a phantom map. From the definition of the modules $`K_i`$, we have a pure exact sequence of $`R`$–modules
$$\begin{array}{ccccccccc}0& & K_i& & P_{i1}& & K_{i1}& & 0.\end{array}$$
This exact sequence gives a triangle in $`𝒯=D(R)`$. The fact that $`:K_{i1}\mathrm{\Sigma }^{n+2i}N`$ is phantom tells us that the composite
$$\begin{array}{ccccc}P_{i1}& & K_{i1}& \stackrel{}{}& \mathrm{\Sigma }^{n+2i}N\end{array}$$
must vanish, since $`P_{i1}`$ is a coproduct of compact objects. But then the map $``$ must factor as
$$\begin{array}{ccccc}K_{i1}& & \mathrm{\Sigma }K_i& & \mathrm{\Sigma }^{n+2i}N.\end{array}$$
Thus if an element $`x\underset{}{lim}^iExt_R^{n+1i}(M_\lambda ,N)`$ comes from a short exact sequence
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^{n+1i}N& & F& & 𝐲K_{i1}& & 0\end{array}$$
with $`F𝐲Y`$, then $`Y`$ is determined by a class
$$y\underset{R}{\overset{n+1i}{Ext}}(K_i,N).$$
In conclusion, we deduce the following. Let us define $`K_0=0`$. We have a map of chain complexes
$$\begin{array}{ccccccccccccc}\mathrm{}& & P_i& & P_{i1}& & \mathrm{}& & P_1& & P_0& & 0\\ & & & & & & & & & & & & & & \\ \mathrm{}& \stackrel{0}{}& K_i& \stackrel{0}{}& K_{i1}& \stackrel{0}{}& \mathrm{}& \stackrel{0}{}& K_1& \stackrel{0}{}& K_0& & 0\text{ .}\end{array}$$
Hence there is a map of spectral sequences in hypercohomology. On the $`E^2`$ term, it is
$$\begin{array}{ccc}Ext_R^j(K_i,N)& & \underset{}{lim}^iExt^j(M_\lambda ,N).\end{array}$$
The whole point is that the spectral sequence on the left degenerates at $`E^1`$, since it comes from a complex with zero differentials. We have shown that if $`x\underset{}{lim}^iExt^{n+1i}(M_\lambda ,N)`$ corresponds to an extension
$$\begin{array}{ccccccccc}0& & 𝐲\mathrm{\Sigma }^{n+1i}N& & F& & 𝐲K_{i1}& & 0\end{array}$$
with $`F𝐲Y`$, then $`x`$ is the image of some $`y`$ from the trivial spectral sequence. Therefore, all differentials out of $`x`$ vanish. ∎
###### Example 2.12.
Let $`k`$ be an infinite field of cardinality $`\mathrm{}_t`$. Then by , the polynomial ring $`R=k[x,y]`$ is of pure global dimension $`t+1`$ ($`\mathrm{}`$ if $`t`$ is infinite). On the other hand, it is of global dimension $`2`$. Hence there do exist modules $`N`$ over $`R=k[x,y]`$, satisfying the assumptions of the theorem when $`t`$ is at least 3.
We can give a refinement of our results for when the ring $`R`$ is hereditary; recall that $`R`$ is hereditary if its global dimension is $`1`$. Examples of hereditary rings are commutative principal ideal domains, and non-commutative polynomial rings.
###### Theorem 2.13.
Let $`R`$ be a hereditary ring. Then
1. \[BRM\] holds in $`𝒯`$ if and only if the pure global dimension of $`R`$ is at most $`1`$; and
2. \[BRO\] holds in $`𝒯`$ if and only if the pure global dimension of $`R`$ is at most $`2`$.
###### Proof.
(i) holds by Neeman’s theorem 0.7, combined with the equality we prove in Proposition 1.4: for hereditary rings
$$pgldimR=pgldimD(R).$$
For (ii), note that Beligiannis’ result (Theorem 0.9) tells us, that \[BRO\] holds if $`pgldimD(R)2`$. The converse comes from Proposition 2.9 which says that if $`N`$ is an $`R`$–module of injective dimension $`1`$ and $`PExt^3(M,N)0`$, then \[BRO\] fails for $`𝒯=D(R)`$. Thus if $`R`$ is hereditary but of pure global dimension $`3`$, \[BRO\] must fail. (Here we have used the easy fact that every hereditary ring is coherent.) ∎
Let $`k`$ be a field. In our counterexamples, we always consider $`k`$-linear triangulated categories $`𝒯`$. When $`𝒯`$ is $`k`$-linear, an additive functor $`\{𝒯^c\}^{op}𝒜b`$ always extends uniquely to a $`k`$-linear functor $`\{𝒯^c\}^{op}\text{Mod-}k`$, so we can restrict attention to such $`k`$-linear functors. The following lemma shows that our counterexamples must take values in infinite-dimensional vector spaces. The idea of the double dual used in the proof is due to M. Van den Bergh.
###### Lemma 2.14.
Let $`k`$ be a field and
$$F:\{𝒯^c\}^{op}\text{mod-}k$$
an exact functor which takes its values in the category $`\text{mod-}k`$ of finite-dimensional vector spaces. Then $`F`$ is of the form $`𝐲X`$ for some $`X𝒯`$.
###### Proof.
Denote by $`D`$ the functor which takes a vector space to its dual. Then the functor $`G=DF`$ is exact and covariant. Let
$$\stackrel{~}{G}:𝒯\text{Mod-}k$$
be the Kan extension of $`G`$ to $`𝒯`$. Thus, for $`Y𝒯`$, we have
$$\stackrel{~}{G}(Y)=\underset{}{colim}GC,$$
where the colimit is taken over the category of arrows $`CY`$ from a compact $`C`$ to $`Y`$. A moment’s thought will convince the reader that $`\stackrel{~}{G}`$ is exact and commutes with coproducts (cf. Prop. 2.3 of ). Hence $`D\stackrel{~}{G}`$ is exact and takes coproducts to products. By Brown’s theorem, it is representable: We have
$$D\stackrel{~}{G}=𝒯(,X)$$
for some $`X𝒯`$. We claim that $`𝐲X=F`$. Indeed, the restriction of $`D\stackrel{~}{G}`$ to $`𝒯^c`$ is isomorphic to $`DDF`$, and this functor is isomorphic to $`F`$ because $`FC`$ is finite-dimensional for all $`C𝒯^c`$. ∎ |
warning/0001/hep-th0001182.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It is now commonly admitted that our current concepts about space and time have to be changed when exploring space-time at a very small scale. Indeed, one can show that it is impossible to locate a particle with an arbitrarily small uncertainty when taking both into account the principles of quantum mechanics and general relativity . Roughly speaking, one can say that measurements of coordinates on space-time are subject to uncertainty relations, thus ruining all geometrical concepts that have proved to be a guidance principle in elaborating many physical theories.
Following the previously alluded analogy with quantum mechanics, one can try to solve this puzzle by assuming that the coordinates themselves are noncommuting objects. Thus, the natural extension of geometrical ideas to this new type of coordinates has been called “noncommutative geometry”. Following even closer the ideas and methods of quantum mechanics, we are led to assume that these noncommutative coordinates are represented as a subalgebra of the algebra of operators acting on a Hilbert space. This is the framework of the theory pioneered by A. Connes , which allows us to make use of the powerful tools of functional analysis. Within this framework, an analogue of gauge theory has been developed, even with non trivial topological properties, and it has already proved to be useful in various areas of physics, ranging from the classical description of the Higgs sector of the standard model (see for a review) to recent ideas in string theory (see and references therein).
This last example involves what we will call NonCommutative Yang-Mills (NCYM) theories in the sequel and can be thought of as a generalization of non-abelian gauge theories, whose gauge symmetry and interactions involves the noncommutative nature of the coordinates. A first example of such a theory appeared almost ten years ago, when Connes and Rieffel developed classical two dimensional Yang-Mills theory on the noncommutative torus . This idea has also been generalized to higher dimensions .
This naturally raises the question of the quantization of such theories, which has been tackled, at the one loop order, on the tori in and on noncommutative $`\mathrm{}^\mathrm{𝔻}`$ in and . Although these theories turned out to be non local, i.e. their interacting vertices involve trigonometric functions of the incoming momenta, it turns that the one-loop behavior is quite similar to the standard non-abelian case. This relies on an older work of Filk , who proved that the trigonometric factor of any planar graph (in a sense to be defined below) does not involve trigonometric functions of the internal momenta flowing into the loops, thus exhibiting the same divergence as the standard theory.
However, this is not true for non planar diagrams whose trigonometric factor does involve a phase depending on the internal momenta. Obviously such a phase softens the ultraviolet behavior of the corresponding diagram and it has been conjectured by several authors that such a diagram is in fact finite .
Nevertheless, it has been pointed out that this is not always the case . Indeed, if the non planar diagram contains some special kind of non planar subdiagrams whose standard degree of divergence is strictly greater that zero, which is the case in a scalar field theory, the small momentum behavior of these diagrams yields a new kind of infrared divergence intimately tied up with the non locality.
This short paper is devoted to a survey of this problem in the simple case of multiloop corrections to the ghost propagator involving only nested and disjoint subdivergent one loop corrections to the ghost propagator. In the following section we shall briefly review the problem raised in and then we shall present an explicit computation of the corresponding diagrams, postponing a complete probe into the renormalization of NCYM theory to a future publication.
## 2 Small momentum singularities induced by non planar diagrams
Before entering into the details of NCYM theory, let us recall that the noncommutative $`\mathrm{}^\mathrm{𝔻}`$ is the algebra generated by $`D`$ hermitean elements $`x_\mu `$ with commutator $`[x_\mu ,x_\nu ]=2\mathrm{i}\theta _{\mu \nu }`$, where $`\theta _{\mu \nu }`$ denotes a real antisymmetric matrix which we will assume to be of maximal rank for convenience. Furthermore, one introduces Fourier modes $`U(k)=\mathrm{e}^{\mathrm{i}kx}`$, with $`kx=k^\mu x_\mu `$. We will always think of a smooth and at infinity rapidly decreasing function as a Fourier transform
$$f=d^Dkf(k)U(k),$$
where $`kf(k)`$ is itself a smooth and rapidly decreasing function on standard $`\mathrm{}^\mathrm{𝔻}`$. The commutation relations of the coordinates endow the algebra with the star product
$$f_\theta g:=d^Dkd^Dlf(k)g(l)U(k)U(l),$$
(1)
which yields
$$(f_\theta g)(k)=d^Dlf(kl)g(l)\mathrm{e}^{\mathrm{i}\theta (k,l)}$$
with $`\theta (k,l)=\theta _{\mu \nu }k^\mu l^\nu `$. Finally, this algebra is equipped with the analogue of an integral defined as
$$f:=f(0)$$
(2)
and partial derivatives
$$_\mu f:=d^Dk\mathrm{i}k_\mu f(k)U(k),$$
(3)
which satisfy most of the properties of their commutative counterparts: positivity and definiteness of the integral, Leibniz rule, commutativity of partial derivatives and integration by part, together with the tracial property of the integral
$$f_\theta g=g_\theta f,$$
which proves to be fundamental in the construction of gauge invariant theories.
At that point, two additional remarks are in order. First of all, let us notice that the matrix $`\theta _{\mu \nu }`$ explicitly breaks Lorentz invariance (or its euclidian counterpart), which is reduced to the transformations commuting with $`\theta `$, whereas translational invariance is preserved. Furthermore, as $`\theta _{\mu \nu }`$ is dimensionful, it also breaks scale invariance already at the classical level, for instance, in the case of four-dimensional NCYM theory or for a two-dimensional scalar field theory.
From now on, one easily constructs scalar field theories, like $`\varphi ^4`$, whose euclidian action is
$$S[\varphi ]=\left(\frac{1}{2}_\mu \varphi _\theta _\mu \varphi +\frac{m^2}{2}\varphi _\theta \varphi +\frac{g}{4!}\varphi _\theta \varphi _\theta \varphi _\theta \varphi \right),$$
(4)
or the NCYM action
$`S[A_\mu ]`$ $`={\displaystyle \frac{1}{4}}{\displaystyle F_{\mu \nu }_\theta F^{\mu \nu }},\text{with}`$ (5)
$`F_{\mu \nu }`$ $`=_\mu A_\nu _\nu A_\mu +g(A_\mu _\theta A_\nu A_\nu _\theta A_\mu ).`$
Because of the tracial properties of the integration, the latter enjoys invariance under noncommutative gauge transformations
$$\delta _\lambda A_\mu =g(\lambda _\theta A_\mu A_\mu _\theta \lambda )_\mu \lambda .$$
Perturbative quantization of these theories is easily performed within a formal functional integral point of view: The quadratic parts of the actions are equal to their commutative counterparts, whereas the interactions are non local and exhibit trigonometric functions of the incoming momenta of the interaction vertices. Thus, the total contribution of any Feynman diagram can be written as the product of a rational function by a trigonometric function. Because trigonometric functions are bounded, the standard rules of powercounting are unchanged so that Weinberg’s convergence theorem remains valid.
It has been shown that for planar diagrams the trigonometric function is independent of the internal momenta so that the Feynman integral reduces to the one encountered in a commutative field theory. For non planar diagrams the situation is more involved and we mainly have to distinguish two cases: whether the non planarity results from crossing of internal lines (i.e. from the non planar character of the amputed diagram), which we call type I diagrams, or whether it solely comes from crossing of internal lines whith external lines (type E). The latter case is much more tricky because the phase vanishes when the corresponding external momenta satisfy some particular relation.
The corresponding Feynman integral has been evaluated in within Schwinger’s regularization scheme. It turns out that a type I non planar diagram whose powercounting subdivergent non planar diagrams are all of type I will be convergent, the corresponding singularity when the Schwinger parameters goes to zero being removed. If the diagram is of type E, but does not contain any subdivergence of type E, it is non singular, except for some exceptional values of its external momenta.
Most of the trouble comes from the insertion of type E non planar subdivergences. Indeed, the latter correspond to Feynman integrals of the type
$$d^{nD}kR(k_1,\mathrm{},k_n,p_1,\mathrm{},p_N)\mathrm{e}^{\mathrm{i}(\theta (k_1,P_1)+\mathrm{}+\theta (k_n,P_n))},$$
(6)
where $`k_1,\mathrm{},k_n`$ are the independent internal momenta, $`p_1,\mathrm{},p_N`$ are the external momenta and $`P_1,\mathrm{},P_n`$ are linear combinations of the external momenta. The rational function $`R(k_1,\mathrm{},k_n,p_1,\mathrm{},p_N)`$ is responsible for the subdivergence.
When all the $`P_i`$’s which couple to internal momenta belonging to divergent integrals do not vanish, the corresponding integral can be considered as finite, being regularized by the oscillatory factor. Indeed, a computation with a cut-off introduced within Schwinger’s parametric formula yields such a finiteness. However, whenever one of the $`P_i`$’s coupled to a divergent loop integral vanishes, then the corresponding Feynman integral is just a usual divergent loop integral and yields a singularity. When inserted into a larger diagram, this could create some trouble when integrating over momenta approaching the subspace $`P_i=0`$. The question is how fast the divergence appears compared with the smoothening property of the integration measure. A quadratic divergence seems to destroy the renormalizability whereas a logarithmic divergence could be harmless. This conjecture is supported by our simple example.
In a mathematically more satisfying manner, one can also consider such an integral as a well defined distribution which is nothing but the Fourier transform of the rational fraction $`R`$. However, we still have to face a problem when inserting the distribution into a larger Feynman diagram. In particular, the example described in corresponds to a distribution which is nothing but the Feynman propagator and the infrared troubles when $`p0`$ are quite similar to the usual small $`x`$ singularities encountered in QFT.
Finally, let us point out that we did not encounter such a problem when quantizing NCYM theory on a torus . The construction of the latter is similar to that of NCYM on $`\mathrm{}^{\mathrm{}}`$ except for the quantization of the momenta appearing in the Fourier transform. As a consequence, there is no singularity when $`p0`$ since such a limit cannot be taken. However, we noticed that a one loop non planar diagram with external momenta $`p`$ exhibits an extra UV singularity of the type $`\delta (p)/ϵ`$. Quite miraculously, all these singularities turn out to cancel, leaving us with finite one loop renormalized correlators.
In the next section, we shall evaluate the non planar contribution to some of the one loop corrections to the ghost self energy using Bessel functions, showing that they do not lead to any singularities when inserted into larger diagrams. We shall use the Feynman rules for NCYM theory without deriving them, the latter being obtained by replacing the structure constant of non-abelian gauge theory $`f_{abc}`$ by $`2\mathrm{i}\mathrm{sin}\theta (p,q)`$ .
## 3 A simple example. One-loop calculation
Our goal is to compute the following 1-loop correction to the ghost propagator in NCYM theory:
(7)
Wavy lines represent gluons and straight lines ghosts. The Feynman rules derived in and (for the noncommutative torus) lead to the integral
$`I_1`$ $`={\displaystyle d^4k\mathrm{\hspace{0.33em}4}g^2\mathrm{}(pk)^\mu \frac{(1)}{k^2}\left(\delta _{\mu \nu }(1\alpha )\frac{k_\mu k_\nu }{k^2}\right)\frac{(1)}{(k+p)^2}(p)^\nu \mathrm{sin}\theta (k,p)\mathrm{sin}\theta (k,k+p)}`$
$`=4g^2\mathrm{}{\displaystyle d^4k\mathrm{sin}^2\theta (k,p)\left(\frac{p^2+\alpha pk}{k^2(k+p)^2}(1\alpha )\frac{(pk)^2}{k^2k^2(k+p)^2}\right)}.`$ (8)
We work in a $`D=4`$ dimensional euclidian momentum space with metric $`g_{\mu \nu }=\delta _{\mu \nu }`$ and use obvious abbreviations such as $`pk=g_{\mu \nu }p^\mu k^\nu `$. Using Feynman parameters
$$\frac{1}{A^rB^s}=\frac{\mathrm{\Gamma }(r+s)}{\mathrm{\Gamma }(r)\mathrm{\Gamma }(s)}_0^1\frac{x^{r1}(1x)^{s1}dx}{(Ax+B(1x))^{r+s}}$$
(9)
we obtain
$$I_1=4g^2\mathrm{}d^4k\mathrm{sin}^2\theta (k,p)_0^1𝑑x\left(\frac{(p^2+\alpha pk)}{(k^2+2pkx+p^2x)^2}\frac{2(1\alpha )(1x)(pk)^2}{(k^2+2pkx+p^2x)^3}\right).$$
In the denominator we write $`qk`$ instead of $`pk`$ so that we reproduce the $`k`$’s in the numerator by differentiation with respect to $`q`$:
$`I_1=4g^2\mathrm{}{\displaystyle }d^4k\mathrm{sin}^2\theta (k,p)(p^2{\displaystyle _0^1}`$ $`{\displaystyle \frac{dx}{(k^2+2qkx+p^2x)^2}}p^\mu {\displaystyle \frac{}{q^\mu }}{\displaystyle _0^1}{\displaystyle \frac{\alpha dx}{2x(k^2+2qkx+p^2x)}}`$
$`p^\mu p^\nu {\displaystyle \frac{}{q^\mu }}{\displaystyle \frac{}{q^\nu }}{\displaystyle _0^1}{\displaystyle \frac{(1\alpha )(1x)dx}{4x^2(k^2+2qkx+p^2x)}}\left)\right|_{q=p}.`$
Using $`\frac{1}{A^n}=\frac{1}{\mathrm{\Gamma }(n)}_0^{\mathrm{}}𝑑tt^{n1}\mathrm{e}^{tA}`$ we rewrite the integral into
$`I_1`$ $`=4g^2\mathrm{}{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t\left(p^2t+{\displaystyle \frac{\alpha }{2x}}p^\mu {\displaystyle \frac{}{q^\mu }}+{\displaystyle \frac{(1\alpha )(1x)}{4x^2}}p^\mu p^\nu {\displaystyle \frac{}{q^\mu }}{\displaystyle \frac{}{q^\nu }}\right)K[t,p,q,x]|_{q=p},`$
$`K[t,p,q,x]:={\displaystyle d^4k\mathrm{e}^{t(k^2+2qkx+p^2x)}\mathrm{sin}^2\theta (k,p)}.`$
Developing the sine into a Fourier series we obtain for the kernel
$`K[t,p,q,x]={\displaystyle }d^4k(\frac{1}{2}\mathrm{e}^{t(k+qx)^2t(p^2xq^2x^2)}`$ $`\frac{1}{4}\mathrm{e}^{t(k+qx+\mathrm{i}\theta (p)/t)^2t(p^2xq^2x^2)2\mathrm{i}x\theta (p,q)pp/t}`$
$`\frac{1}{4}\mathrm{e}^{t(k+qx\mathrm{i}\theta (p)/t)^2t(p^2xq^2x^2)+2\mathrm{i}x\theta (p,q)pp/t}),`$
where $`\theta (p)^\mu :=\theta ^{\mu \alpha }p_\alpha `$ and $`pp:=g_{\mu \nu }\theta ^{\mu \alpha }p_\alpha \theta ^{\nu \beta }p_\beta (\theta (p))^2`$. Then it is easy to perform the Gaussian integration:
$$K[t,p,q,x]=\frac{\pi ^2}{2t^2}\left(\mathrm{e}^{t(p^2xq^2x^2)}\mathrm{e}^{t(p^2xq^2x^2)pp/t}\mathrm{cos}2x\theta (p,q)\right).$$
(10)
We can now perform the differentiations with the following result:
$`I_1`$ $`=\pi ^2g^2\mathrm{}{\displaystyle _0^1}dx{\displaystyle _0^{\mathrm{}}}dt({\displaystyle \frac{1}{t}}p^2(3\alpha x\alpha 1x)+2(p^2)^2(1\alpha )x^2(1x))\times `$
$`\times \left(\mathrm{e}^{tp^2x(1x)}\mathrm{e}^{tp^2x(1x)pp/t}\right)`$ (11)
$`=p^2\pi ^2g^2\mathrm{}{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t\left({\displaystyle \frac{1}{t}}(3\alpha x\alpha 1x)+2(1\alpha )x\right)\left(\mathrm{e}^t\mathrm{e}^{tx(1x)p^2pp/t}\right).`$
The integration over $`t`$ diverges logarithmically. We define the projection $`t_p^2`$ onto the divergent part by
$`t_p^2(I_1)`$ $`:=p^2\pi ^2g^2\mathrm{}{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}(3\alpha x\alpha 1x)\mathrm{e}^t=p^2\pi ^2g^2\mathrm{}{\displaystyle \frac{3\alpha }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\mathrm{e}^tdt}{t}}.`$ (12)
The projection to the convergent part is
$`R_1`$ $`:=(1t_p^2)(I_1)`$ (13)
$`=p^2\pi ^2g^2\mathrm{}((1\alpha ){\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t\left({\displaystyle \frac{3\alpha x\alpha 1x}{t}}+2(1\alpha )x\right)\mathrm{e}^{tx(1x)p^2pp/t}).`$
Instead of introducing a cutoff as in we prefer such a momentum subtraction before performing the divergent integral, in analogy to the BPHZ scheme, as the notation $`t_p^2`$ for the projection indicates. We believe this is advantageous for the renormalizability proof to all orders based on Zimmermann’s forest formula. Please notice the difference between the power counting degree 2 entering the forest formula and the actually only logarithmic divergence in (12). This is crucial for the insertion as subdivergences and gives the reason why we will obtain local counterterms to all orders whereas there are true quadratic divergences in the scalar theory in . It should be not difficult to extend such a subtraction scheme to the entire NCYM theory. It is however important to define $`t_p^d`$ as the projector onto the strictly divergent part of an integral in order to produce local counterterms. This assumes one can prove that all divergent integrations give rise to such local terms, as we will do in this paper for repeated insertions of the ghost propagator.
The $`x`$-integration in (13) yields $`_0^1𝑑xx\mathrm{e}^{tax(1x)/t}=\frac{1}{2}_0^1𝑑x\mathrm{e}^{tax(1x)/t}`$ for any $`a`$ so that
$`R_1`$ $`=p^2\pi ^2g^2\mathrm{}\left((1\alpha )+{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑t\left({\displaystyle \frac{3\alpha }{2t}}(1\alpha )\right)\mathrm{e}^{tx(1x)p^2pp/t}\right).`$
The $`t`$-integration leads to Bessel functions:
$`R_1`$ $`=p^2\pi ^2g^2\mathrm{}(1\alpha )\left(1{\displaystyle _0^1}𝑑x\mathrm{\hspace{0.17em}2}\sqrt{x(1x)p^2pp}K_1[2\sqrt{x(1x)p^2pp}]\right)`$
$`+p^2\pi ^2g^2\mathrm{}(3\alpha ){\displaystyle _0^1}𝑑xK_0[2\sqrt{x(1x)p^2pp}].`$ (14)
The Bessel function $`K_0[y]`$ diverges logarithmically to $`+\mathrm{}`$ for $`y0`$ and converges exponentially to $`0`$ for $`y+\mathrm{}`$. Hence, for any exponent $`r>0`$ there exists a number $`c_r^0>0`$ such that
$$K_0[y]c_r^0/y^r0<y<\mathrm{}.$$
(15)
This will be proven algebraically in the Appendix. Graphically the situation is sketched in Figure 1.
The function $`yK_1[y]`$ approaches 1 for $`y0`$ and converges exponentially to 0 for $`y+\mathrm{}`$. It is nevertheless convenient to regard it as $`K_0[y]`$ before: For any exponent $`r>0`$ there exists a number $`c_r^1>0`$ such that
$$yK_1[y]c_r^1/y^r0<y<\mathrm{}.$$
(16)
The corresponding graphic is shown in Figure 2.
Now the $`x`$-integration is easy to perform. It converges for $`0<r<1`$, which means that infrared divergences are absent (for non exceptional momenta).
We restrict ourselves to the case where the rank of the tensor $`g_{\mu \nu }\theta ^{\mu \alpha }\theta ^{\nu \beta }`$ equals the space-time dimension (maximal noncommutativity). Then there exists some parameter $`m_P`$ of dimension of a mass (the ‘Planck mass’) such that
$$pp\frac{p^2}{m_P^4}\sqrt{p^2pp}\frac{p^2}{m_P^2}.$$
(17)
This yields the estimation
$$R_1=p^2\pi ^2g^2\mathrm{}\left((1\alpha )+O\left(P_r^1(\alpha )\left(\frac{m_P^2}{p^2}\right)^r\right)\right),$$
(18)
where $`P_r^n(\alpha )`$ is a polynomial of homogeneous degree $`n`$ in $`(1\alpha )`$ and $`(3\alpha )`$ with coefficients of order 1.
## 4 Higher loop order calculation
Now we insert $`n`$ of these 1-loop propagator corrections into a propagator correction, giving an $`n+1`$ loop diagram:
(19)
Using (8) and the Feynman rule for the ghost propagator, the corresponding integral is of order
$`I_{n+1}=4g^2\mathrm{}(\pi ^2g^2\mathrm{})^n{\displaystyle d^4k}`$ $`\mathrm{sin}^2\theta (k,p)({\displaystyle \frac{p^2+\alpha pk}{k^2(k+p)^2}}(1\alpha ){\displaystyle \frac{(pk)^2}{k^2k^2(k+p)^2}})\times `$
$`\times {\displaystyle \underset{j=0}{\overset{n}{}}}P_r^j(\alpha )(1\alpha )^{nj}\left({\displaystyle \frac{m_P^2}{(k+p)^2}}\right)^{jr},`$ (20)
with $`P_r^0(\alpha )=1`$. Introduction of Feynman and Schwinger parameters as before and use of $`x\mathrm{\Gamma }(x)=\mathrm{\Gamma }(x+1)`$ leads to
$`I_{n+1}`$ $`=4g^2\mathrm{}(\pi ^2g^2\mathrm{})^n{\displaystyle \underset{j=0}{\overset{n}{}}}(1+jr)P_r^j(\alpha )(1\alpha )^{nj}m_P^{2jr}{\displaystyle }d^4k\mathrm{sin}^2\theta (k,p)\times `$
$`\times {\displaystyle _0^1}dx({\displaystyle \frac{x^{jr}(p^2+\alpha pk)}{(k^2+2pkx+p^2x)^{2+jr}}}{\displaystyle \frac{(2+jr)(1\alpha )(1x)x^{jr}(pk)^2}{(k^2+2pkx+p^2x)^{3+jr}}})`$
$`=4g^2\mathrm{}(\pi ^2g^2\mathrm{})^n{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{1}{\mathrm{\Gamma }(1+jr)}}P_r^j(\alpha )(1\alpha )^{nj}m_P^{2jr}\times `$
$`\times {\displaystyle _0^1}dx{\displaystyle _0^{\mathrm{}}}dt(x^{jr}t^{1+jr}p^2+\frac{1}{2}\alpha x^{jr1}t^{jr}p^\mu {\displaystyle \frac{}{q^\mu }}`$
$`+\frac{1}{4}(1\alpha )(1x)x^{jr2}t^{jr}p^\mu p^\nu {\displaystyle \frac{}{q^\mu }}{\displaystyle \frac{}{q^\nu }})K[t,p,q,x]|_{q=p}`$
$`=(\pi ^2g^2\mathrm{})^{n+1}{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{1}{\mathrm{\Gamma }(1+jr)}}P_r^j(\alpha )(1\alpha )^{nj}m_P^{2jr}\times `$
$`\times {\displaystyle _0^1}dx{\displaystyle _0^{\mathrm{}}}dt((3\alpha xx\alpha 1)x^{jr}t^{jr1}p^2+2(1\alpha )(1x)x^{jr+2}t^{jr}(p^2)^2)\times `$
$`\times \left(\mathrm{e}^{tp^2x(1x)}\mathrm{e}^{tp^2x(1x)pp/t}\right)`$
$`=(\pi ^2g^2\mathrm{})^{n+1}p^2{\displaystyle \underset{j=0}{\overset{n}{}}}{\displaystyle \frac{1}{\mathrm{\Gamma }(1+jr)}}P_r^j(\alpha )(1\alpha )^{nj}\left({\displaystyle \frac{m_P^2}{p^2}}\right)^{jr}{\displaystyle _0^1}{\displaystyle \frac{dx}{(1x)^{jr}}}\times `$
$`\times {\displaystyle _0^{\mathrm{}}}dt((3\alpha xx\alpha 1)t^{jr1}+2(1\alpha )xt^{jr})(\mathrm{e}^t\mathrm{e}^{tx(1x)p^2pp/t}).`$ (21)
The only divergent integral is for $`j=0`$ the projection
$`t_p^2(I_{n+1})`$ $`:=(\pi ^2g^2\mathrm{})^{n+1}p^2(1\alpha )^n{\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}(3\alpha xx\alpha 1)\mathrm{e}^t`$
$`=`$ $`(\pi ^2g^2\mathrm{})^{n+1}p^2(1\alpha )^n{\displaystyle \frac{3\alpha }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dt}{t}}\mathrm{e}^t.`$ (22)
The convergent part can be evaluated to
$`R_{n+1}`$ $`=(1t_p^2)(I_{n+1})`$
$`=p^2(\pi ^2g^2\mathrm{})^{n+1}(1\alpha )^{n+1}\left(1{\displaystyle _0^1}𝑑x\mathrm{\hspace{0.17em}2}\sqrt{x(1x)p^2pp}K_1[2\sqrt{x(1x)p^2pp}]\right)`$
$`p^2(\pi ^2g^2\mathrm{})^{n+1}(1\alpha )^n(3\alpha ){\displaystyle _0^1}𝑑xK_0[2\sqrt{x(1x)p^2pp}]`$
$`+p^2(\pi ^2g^2\mathrm{})^{n+1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{P_r^j(\alpha )}{jr(2jr)}}(1\alpha )^{nj}(3\alpha )\left({\displaystyle \frac{m_P^2}{p^2}}\right)^{jr}`$ (23)
$`+p^2(\pi ^2g^2\mathrm{})^{n+1}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{1}{\mathrm{\Gamma }(1+jr)}}P_r^j(\alpha )(1\alpha )^{nj}\left({\displaystyle \frac{m_P^2}{p^2}}\right)^{jr}\times `$
$`\times {\displaystyle _0^1}{\displaystyle \frac{dx}{(1x)^{jr}}}{\displaystyle _0^{\mathrm{}}}dt((3\alpha xx\alpha 1)t^{jr1}+2(1\alpha )xt^{jr})\mathrm{e}^{tx(1x)p^2pp/t}.`$
The last two lines range from zero (for $`p=\mathrm{}`$) to minus the value of the third last line for $`p=0`$. Thus, in our estimation we have to neglect the last two lines. The remaining integrals over $`K_0`$ and $`K_1`$ are familiar to us, see Figures 1 and 2, and we choose the essential exponents in (15) and (16) to be $`(n+1)r`$ instead of $`r`$. Then we arrive at
$$R_{n+1}=p^2(\pi ^2g^2\mathrm{})^{n+1}\left((1\alpha )^{n+1}+O\left(\underset{j=1}{\overset{n+1}{}}P_{r,j}^{n+1}(\alpha )\left(\frac{m_P^2}{p^2}\right)^{jr}\right)\right).$$
(24)
But this was precisely our starting point we inserted into (8) to obtain (20). Hence, (24) provides the structure of any renormalized $`n+1`$ loop graph made of ghost propagator corrections. The counterterm of such an $`n`$-loop graph is given by (22), and we see explicitly that the Feynman graphs made of nested 1-loop ghost propagator corrections
(25)
are renormalized by local counterterms for any order in $`\mathrm{}`$. Here, locality means that the momentum dependence of the counterterm and the kinetic part of the ghost action are identical. In order to renormalize an $`n`$-loop graph (to avoid infrared divergences) the critical exponent has to be chosen $`0<r<1/n`$.
The essential step in this proof was the observation that the noncommutative Feynman graphs under consideration evaluate to Bessel functions, which can be estimated by a power law. It seems plausible that any Feynman graph of noncommutative Yang-Mills theory evaluates to Bessel functions, and applying the same techniques it should be possible to show that local counterterms suffice to renormalize this model.
## Appendix: Proof of Eq. (15)
We prove that for each $`r>0`$ there is a number $`c_r>0`$ such that
$$C_0(x):=\frac{c_r}{x^r}K_0(x)0<x<\mathrm{}.$$
The Bessel function $`K_0(x)`$ is one of the two solutions of the differential equation
$$xK_0^{\prime \prime }+K_0^{}xK_0=0,0<x<\mathrm{}.$$
(26)
It is however more convenient to consider the function
$$K(x):=\sqrt{x}K_0(x),K^{\prime \prime }+\frac{14x^2}{4x^2}K=0,$$
and compare it with
$$C(x):=\sqrt{x}C_0(x),C^{\prime \prime }+\frac{14r^2}{4x^2}C=0.$$
The derivative of the Wronskian $`W(K,C):=K^{}CC^{}K`$ is
$$W^{}=\frac{x^2r^2}{x^2}KC,\begin{array}{cc}W^{}>0\hfill & \text{ for }x>r\hfill \\ W^{}<0\hfill & \text{ for }x<r\hfill \end{array}$$
The asymptotic development shows that for the solution $`K_0`$ of (26) one has $`W(x)<0`$ for $`x\mathrm{}`$ and $`W(x)>0`$ for $`x0`$. Therefore, there is only one zero of the Wronskian, at $`x=x_r`$, as illustrated in Figure 3.
We choose the normalization
$$K_0(x_r)=c_rx_r^r,K_0^{}(x_r)=rc_rx_r^{r1}$$
so that $`W(x_r)=0`$ and $`K(x_r)=C(x_r)`$. Then we can integrate
$`x>x_rW(x)<0:`$ $`{\displaystyle _{x_r}^x}{\displaystyle \frac{K^{}(x)}{K(x)}}𝑑x<{\displaystyle _{x_r}^x}{\displaystyle \frac{C^{}(x)}{C(x)}}𝑑x`$ $`K(x)<C(x)`$
$`x<x_rW(x)>0:`$ $`{\displaystyle _x^{x_r}}{\displaystyle \frac{K^{}(x)}{K(x)}}𝑑x>{\displaystyle _x^{x_r}}{\displaystyle \frac{C^{}(x)}{C(x)}}𝑑x`$ $`K(x)<C(x)`$
This finishes the proof. |
warning/0001/quant-ph0001007.html | ar5iv | text | # Universal quantum limits on single channel information, entropy and heat flow
## I Introduction
In a recent experiment , Schwab et al. succeeded in measuring for the first time the thermal conductance quantum for a suspended, dielectric wire of submicron cross section. In accordance with predictions , the thermal conductance was found to approach the limiting value $`4\times \pi k_B^2T/6\mathrm{}4\times 10^{12}T\text{W K}^1`$ as the wire thermal reservoirs were cooled such that the dominant phonon wavelength became comparable to the wire cross section. The factor of 4 is just the number of independent vibrational mode branches of the wire satisfying $`\omega (k)0`$ as $`k0`$ (see, e.g., Ref. ). Only such modes can have non-negligible phonon occupation numbers as $`T0`$, giving four available channels for heat transport. The single channel thermal conductance can never exceed the thermal conductance quantum $`\pi k_B^2T/6\mathrm{}`$. The conductance quantum can only be attained for ballistic transport (i.e., no scattering) as was achieved in the experiment.
In common with the quantum limits for other single channel, linear transport coefficients, such as the electronic conductance quantum $`e^2/h`$ , the thermal conductance quantum does not depend on the form of the $`\omega (k)`$ dispersion relations, a consequence of the cancellation of the group velocity and density of states factors in the formula for the one dimensional heat current. Wires made from different insulating materials and with different cross section geometries will therefore all have the same limiting single channel thermal conductance value for ballistic transport at low temperatures. For this reason, the conductance quantum is often termed ‘universal’.
The thermal conductance is in fact universal in a much wider sense. For a single channel connecting two heat reservoirs with (quasi)particles obeying fractional statistics according to Haldane’s definition (which generalizes Bose and Fermi statistics) , it was recently found that the maximum, limiting thermal conductance quantum is independent of the particle statistics as well . For example, in the case of an ideal electron gas, the limiting single channel thermal conductance coincides with the above thermal conductance quantum for phonons. While dimensional analysis would lead us to expect the same factor $`k_B^2T/\mathrm{}`$ to occur independently of the statistics, there is no a priori reason to expect the same numerical factor $`\pi /6`$ as well, given that the latter results from integrating with respect to the energy the expansion to first order in small temperature differences of the thermal reservoir distributions, which have qualitatively different forms for particles obeying different statistics. This remarkable property is unique to the thermal conductance: all other single channel transport coefficients depend on the particle statistics.
In an earlier and unrelated investigation concerning the quantum limits on single channel information and entropy flow , Pendry showed that the bound $`\dot{S}^2/\dot{E}\pi k_B^2/3\mathrm{}`$, involving the averaged single channel entropy and energy currents, is obeyed for both bosons and fermions. The striking resemblance between this bound and that for the single channel thermal conductance suggests the possible existence of an universal and more general, attainable bound relating the entropy and energy currents, from which the thermal conductance bound would follow as a special case. In particular, there is the possibility of a bound which would be independent of the channel materials properties and particle statistics and which would apply even far from equilibrium where the temperatures (and perhaps also the chemical potentials) of the two heat reservoirs connecting the ends of the channel are significantly different.
Given that entropy and information are closely related, the existence of such a universal upper bound on the entropy flow rate would in turn suggest the existence of an optimum capacity for single channel classical information transmission, which is also universal in the wider sense (i.e., independent of channel materials properties and particle statistics). This is in fact the main subject of Ref. . However, there the analysis is restricted to situations in which the channel is noiseless, with the information encoded and decoded in terms of the boson/fermion number eigenstates. A proper determination of the optimum capacity would consider all possible input quantum states for encoding letters and all possible detection schemes at the output. The crucial result which allows one to generalize the analysis of Ref. is Holevo’s theorem , which bounds the mutual information between channel output and input with a quantity involving the quantum entropies of the input states. Caves and Drummond have carried out the more general analysis for particles obeying Bose statistics only and confirm Pendry’s upper bound as the optimum channel capacity.
Thus, there is the possibility of an optimum, universal limiting capacity which bounds all possible methods of encoding and detection, and which is independent of the physical properties of the channel. We emphasize that the existence of such a single channel optimum capacity is suggested by the established existence of the universal thermal conductance quantum.
The fact that the thermal conductance by its very definition requires that the channel be connected at each end to a heat reservoir which can act both as an emitter and absorber of quanta, suggests that in order to attain the above-conjectured optimum capacity, a generalization of the mutual information will be required in which there is a sender/receiver pair at both ends of the channel. The pairs thus share the same channel and information can now flow in opposite directions. In fact, as we shall see, with the exception of bosons the attainability of Pendry’s upper bound for particles obeying arbitrary statistics necessarily requires that the chemical potentials of the two reservoirs coincide and be nonzero, so that both reservoirs are sources of particles.
In the following sections, we provide evidence for the validity of the two conjectures outlined above: namely (1) the existence of an universal bound relating the entropy and energy flow rates of a single quantum channel, with the universal thermal conductance bound following as a special case, and (2) the existence of an universal bound on the optimum capacity for single-channel information transmission, subject to certain constraints on the channel and input states.
In Sec. II, we first show that Pendry’s bound holds for particles obeying arbitrary statistics under quite general conditions for the reservoir temperatures and chemical potentials. We then introduce a less general but tighter bound, which requires that the chemical potentials of the two reservoirs coincide, and replaces the energy current with the heat current $`\dot{E}\mu \dot{N}`$. We show that the thermal conductance bound follows as a special case from this latter bound when the temperature difference between the two reservoirs approaches zero.
In Sec. III, we first repeat the analysis of Caves and Drummond in the conventional case for unidirectional information transmission, but generalizing to particles obeying fractional statistics, and obtain the limiting statistics-dependent optimum capacity. We then generalize the definition of mutual information and Holevo’s theorem in a simple way to allow for non-interfering, two-way information flow. Using this generalized definition we show, subject to certain constraints on the channel and input states, that the limiting optimum capacity is now independent of the particle statistics and coincides with that of the bosonic case. In the final part of the section, we introduce a further generalization of the mutual information which allows for the possibility of interference between the ‘left-moving’ and ‘right-moving’ information flows.
In Sec. IV, we conclude and also briefly outline various open problems which have a bearing on the two conjectures.
## II entropy bounds
As our generic model structure, we consider some confining ‘wire’ which supports particles obeying a given statistics and which is connected adiabatically at each end to two particle reservoirs characterized by temperatures $`T_L`$ and $`T_R`$ and chemical potentials $`\mu _L`$ and $`\mu _R`$, where the subscripts $`L`$ and $`R`$ denote the left and right reservoirs, respectively. The device of Ref. is one possible realisation of the model structure in the case of phonons. Typically, a wire will provide several available parallel channels for given reservoir chemical potential and temperature values. However, in the case of ballistic transport the channel currents don’t interfere with each other and thus can be treated independently. We will restrict ourselves to ballistic transport in the present investigation.
The distribution function for particles obeying fractional statistics is
$$f_g(E)=\left[w\left(\frac{E\mu }{k_BT}\right)+g\right]^1,$$
(1)
where the function $`w(x)`$ satisfies
$$w(x)^g[1+w(x)]^{1g}=e^x.$$
(2)
The parameter $`g`$, assumed to be a rational number, determines the statistics. From these equations, we can see immediately that $`g=0`$ describes bosons and $`g=1`$ fermions. The left(right) components of the single channel energy and entropy currents are
$$\dot{E}_{L(R)}=\frac{(k_BT_{L(R)})^2}{2\pi \mathrm{}}_{x_{L(R)}^0}^{\mathrm{}}𝑑x\left(x+\mu _{L(R)}/k_BT_{L(R)}\right)f_g(x)$$
(3)
and
$`\dot{S}_{L(R)}=`$ $``$ $`{\displaystyle \frac{k_B^2T_{L(R)}}{2\pi \mathrm{}}}{\displaystyle _{x_{L(R)}^0}^{\mathrm{}}}dx\{f_g\mathrm{ln}f_g+(1gf_g)\mathrm{ln}(1gf_g)`$ (4)
$``$ $`[1+(1g)f_g]\mathrm{ln}[1+(1g)f_g]\},`$ (5)
where $`x_{L(R)}^0=\mu _{L(R)}/k_BT_{L(R)}`$, and where we define the energy origin such that the minimum energy of a channel particle is zero, i.e., the energy is given by the longitudinal kinetic component. The total energy and entropy channel currents are then just $`\dot{E}=\dot{E}_L\dot{E}_R`$ and $`\dot{S}=\dot{S}_L\dot{S}_R`$, respectively.
The first conjectured bound involving these single channel entropy and energy currents is
$$\dot{S}^2\frac{\pi k_B^2}{3\mathrm{}}\dot{E},$$
(6)
provided $`T_L>T_R`$ and $`\mu _L\mu _R`$. We have numerically tested this bound extensively in $`\mu `$ and $`T`$ parameter space, for several rational values of the statistical parameter $`g`$ ranging between between zero and one. Fig. 1 gives an initial idea of the bound by showing the dependence of the ratio $`3\mathrm{}\dot{S}^2/\pi k_B^2\dot{E}`$ on a selected parameter range.
In the case of bosons with constant $`\mu _L=\mu _R=0`$ (e.g., photons or phonons), evaluating Eqs. (3) and (5) gives $`\dot{E}_{L(R)}=\pi (k_BT_{L(R)})^2/12\mathrm{}`$ and $`\dot{S}_{L(R)}=\pi k_B^2T_{L(R)}/6\mathrm{}`$, respectively. Setting $`T_R=0`$ and eliminating $`T_L`$ by solving for $`\dot{S}(=\dot{S}_L)`$ in terms of $`\dot{E}(=\dot{E}_L)`$, we obtain equality in bound (6). For all other physically achievable parameter choices, we have strict inequality in (6). The key point, however, is that the bound can be approached arbitrarily closely, no matter the particle statistics. For example, in the case of bosons with non-constant reservoir chemical potentials, the bound is approached asymptotically in the degenerate limit $`x_L^0=\mu _L/k_BT_L0^{}`$, with $`\mu _R=0`$ and $`T_R=0`$. For particles with $`g>0`$, the bound is approached asymptotically in the degenerate limit $`x_L^0=\mu _L/k_BT_L+\mathrm{}`$, with $`\mu _L=\mu _R`$ and $`T_R=0`$. That these are the correct conditions for approaching the upper bound can be seen more clearly after transforming the integrals in Eqs. (3) and (5) for the bosonic and fermionic cases as in, e.g., Sec. 58 of Ref. . For example, in the case of fermions, the single channel energy and entropy currents can be rewritten as follows
$`\dot{E}={\displaystyle \frac{\pi (k_BT_L)^2}{12\mathrm{}}}[1\left({\displaystyle \frac{T_R}{T_L}}\right)^2`$ $`+`$ $`{\displaystyle \frac{3}{(\pi k_BT_L)^2}}\left(\mu _L^2\mu _R^2\right)+{\displaystyle \frac{6}{\pi ^2}}{\displaystyle _{x_L^0}^{\mathrm{}}}𝑑x\left(\mu _L/k_BT_Lx\right)f(x)`$ (7)
$``$ $`{\displaystyle \frac{6}{\pi ^2}}\left({\displaystyle \frac{T_R}{T_L}}\right)^2{\displaystyle _{x_R^0}^{\mathrm{}}}dx(\mu _R/k_BT_Rx)f(x)]`$ (8)
and
$`\dot{S}={\displaystyle \frac{\pi k_B^2T_L}{6\mathrm{}}}\{1{\displaystyle \frac{T_R}{T_L}}`$ $`+`$ $`{\displaystyle \frac{3}{\pi ^2}}{\displaystyle _{x_L^0}^{\mathrm{}}}𝑑x\left[f\mathrm{ln}f+(1f)\mathrm{ln}(1f)\right]`$ (9)
$``$ $`{\displaystyle \frac{3}{\pi ^2}}{\displaystyle \frac{T_R}{T_L}}{\displaystyle _{x_R^0}^{\mathrm{}}}dx[f\mathrm{ln}f+(1f)\mathrm{ln}(1f)]\}.`$ (10)
Taking the limit $`x_L^0+\mathrm{}`$, with the conditions $`\mu _L=\mu _R`$ and $`T_R=0`$, only the first term remains on the right-hand-sides of Eqs. (8) and (10) and the energy and entropy currents coincide with those for bosons with $`\mu _L=\mu _R=0`$ and $`T_R=0`$.
It is not possible to recover the single channel thermal conductance bound from (6). The best we can do is to derive an upper bound on the rate of heat emission from an isolated reservoir for bosons with zero chemical potential (see also Ref. ). Setting $`\mu _L=\mu _R=0`$, $`T_R=0`$, identifying the heat emission rate with the total energy emission rate $`\dot{Q}_L=\dot{E}_L`$, and using $`\dot{Q}_L/T_L\dot{S}_L`$, bound (6) gives
$$\dot{Q}_L\frac{\pi k_B^2T_L^2}{3\mathrm{}}.$$
(11)
Note that for particles with nonzero chemical potential, the heat emission rate is $`\dot{Q}_L=\dot{E}_L\mu _L\dot{N}_L`$, where $`\dot{N}`$ denotes the number current, and in this case (11) does not follow from (6). If we had equality in (11), then the thermal conductance could be obtained by taking the difference $`\dot{Q}_L\dot{Q}_R=\pi k_B^2(T_L^2T_R^2)/3\mathrm{}=2\pi k_B^2\overline{T}\delta T/3\mathrm{}`$, where $`\overline{T}=(T_L+T_R)/2`$. But this gives the incorrect coefficient ($`2/3`$ instead of $`1/6`$). What is wrong with this argument is the assumption that $`\dot{Q}_{L(R)}/T_{L(R)}=\dot{S}_{L(R)}`$. In fact, $`\dot{Q}_{L(R)}/T_{L(R)}=\dot{S}_{L(R)}/2`$, signalling the irreversible nature of the heat emission.
A conjectured, tighter bound suggested by the form of expressions (8) and (10) which does yield the thermal conductance bound as a special case, is the following
$$\dot{S}^2\frac{\pi k_B^2}{3\mathrm{}}\left(\frac{T_LT_R}{T_L+T_R}\right)\left(\dot{E}\mu \dot{N}\right),$$
(12)
provided $`T_L>T_R`$ and $`\mu _L=\mu _R=\mu `$. Again, we have numerically tested this bound extensively in $`\mu `$ and $`T`$ parameter space, for several rational values of the statistical parameter $`g`$ ranging between between zero and one (see Figs. 2 and 3). In the case of bosons with constant $`\mu =0`$, we obtain equality for all $`T_L>T_R`$. For bosons with non-constant $`\mu 0`$, the bound is approached asymptotically in the degenerate limit $`\mu /k_BT_{L(R)}0^{}`$. For particles with $`g>0`$, the bound is approached asymptotically in the degenerate limit $`\mu /k_BT_{L(R)}+\mathrm{}`$. Note that the heat current $`\dot{Q}=\dot{E}\mu \dot{N}`$ appears instead of the energy current $`\dot{E}`$ on the right-hand-side of bound (12). This replacement is essential: if the energy current is used, then the bound can be violated for bosons with $`\mu <0`$. It is remarkable that the need to recover the thermal conductance and also to satisfy the bound both lead to the replacement of the energy current with the heat current.
## III information bounds
Consider a communication channel, characterised by an input alphabet $`A`$ with letters labeled by an index $`a=1,\mathrm{},𝒜`$ and a set of probabilities $`p_A(a)`$ for transmitting letter $`a`$, an output alphabet $`B`$ labeled by $`b=1,\mathrm{},`$, and a set of conditional probabilities $`p_{B|A}(b|a)`$ for receiving letter $`b`$, given transmission of letter $`a`$. The mutual information gives the measure of the information successfully transmitted from input to output of the communication channel:
$$H(B;A)=\underset{a,b}{}p_{B|A}(b|a)p_A(a)\mathrm{log}_2\left(\frac{p_{B|A}(b|a)}{p_B(b)}\right),$$
(13)
where $`p_B(b)=_ap_{B|A}(b|a)p_A(a)`$.
Suppose the quantum channel medium supports particles for some given rational, statistical parameter value $`g`$, $`0g1`$. Let the input letter $`a`$ be encoded in some quantum state $`\widehat{\rho }_a`$, and the output detection scheme be described, in the most general case, by a set of non-negative, bounded Hermitian operators $`\widehat{F}_b`$ satisfying $`_b\widehat{F}_b=\widehat{1}`$, with $`p_{B|A}(b|a)=\text{tr}(\widehat{\rho }_a\widehat{F}_b)`$. The operators $`\widehat{\rho }_a`$ and $`\widehat{F}_b`$ act on the channel Fock space for statistical parameter value $`g`$. More precisely, since information is transmitted in only one direction, these operators act on the subspace describing right-moving states, say.
Holevo’s theorem provides an upper bound on the mutual information for all possible detection schemes:
$$\underset{\{\widehat{F}_b\}}{\mathrm{max}}H(B;A)S(\widehat{\rho })\underset{a}{}p_A(a)S(\widehat{\rho }_a),$$
(14)
where $`\widehat{\rho }=_ap_A(a)\widehat{\rho }_a`$ and $`S(\widehat{\rho })=\text{t}r(\widehat{\rho }\mathrm{log}_2\widehat{\rho })`$ is the quantum entropy in bits. Note that, while this theorem is usually applied to bosonic communication channels (c.f., Ref. ), it is in fact applicable to channels for arbitrary fractional statistics. All that is required is that the channel obeys the usual rules of quantum mechanics.
Maximizing the mutual information with respect to the output detection scheme and the input states and probabilities gives the optimum capacity $`C`$ of the channel. From (14), we have
$$C=\frac{1}{𝒯}\underset{\widehat{\rho }}{\mathrm{max}}\underset{\{\widehat{F}_b\}}{\mathrm{max}}H(B;A)\frac{1}{𝒯}\underset{\widehat{\rho }}{\mathrm{max}}S(\widehat{\rho })=\frac{S_{\text{max}}}{𝒯},$$
(15)
where $`𝒯`$ is the transmission time. As shown in Sec. IV.B of Ref. , the upper bound $`S_{\text{max}}`$ can in fact be attained: Find a complete, orthonormal set of diagonalizing basis states $`|a`$ for the $`\widehat{\rho }`$ which maximizes $`S(\widehat{\rho })`$, i.e., $`\widehat{\rho }=_aq(a)|aa|`$. Choose $`\widehat{\rho }_a=|aa|=\widehat{F}_a`$ and $`p_A(a)=q(a)`$. Then $`H(B;A)=_ap_A(a)\mathrm{log}_2p_A(a)=S(\widehat{\rho })=S_{\text{max}}`$.
Thus, the optimum capacity is just the maximum quantum entropy in bits divided by the transmission time, subject to the given constraints on the channel. One common constraint is to fix the total energy of the transmitted message; the optimum capacity then gives the maximum information that can be transmitted in a time $`𝒯`$ for a given, allowed signal energy. For a single, wideband channel with longitudinal single-particle energies $`hf_j=hj/𝒯`$, $`j=1,2,\mathrm{}`$, the total longitudinal energy $`E_N`$ of a given Fock state is $`E_N=_jhf_jn_j=Nh/𝒯`$, where $`N=_{j=1}^{\mathrm{}}jn_j`$, and $`n_j`$ is the occupation number of, say, the right-propagating mode $`j`$. The maximum entropy is then $`S_{\text{max}}=\mathrm{log}_2𝒩_N`$, where $`𝒩_N`$ is just the number of different ways the sum $`N`$ can be partitioned. For bosons ($`g=0`$), $`𝒩_N`$ is given by the number of unrestricted partitions, while for the fermions ($`g=1`$), $`𝒩_N`$ is given by the number of partitions into distinct parts. More generally, for particles obeying fractional statistics with $`g=1/n`$, $`n=1,2,\mathrm{}`$, no number can appear more than $`n`$ times in a given partition. Note, however, that for $`0<g<1`$ there are additional constraints on the allowed partitions, as discussed in Ref. .
For long transmission times $`𝒯`$, or equivalently large $`N`$, one obtains the following asymptotic approximation to the optimum capacity of a wideband, bosonic channel for fixed energy :
$$C_{\text{boson}}=\frac{\pi }{\mathrm{ln}2}\sqrt{\frac{2P}{3h}}\frac{1}{𝒯}\mathrm{log}_2\left(\frac{4\sqrt{3}P𝒯^2}{h}\right),$$
(16)
where $`P=E_N/𝒯`$ is the time-averaged power. Caves et al. also derive the bosonic optimum capacity subject to the alternative constraints that the maximum-energy or the message-ensemble-averaged-energy of the channel be fixed. All give the same leading order term as on the right-hand-side of Eq. (16), with $`P`$ appropriately defined in each case.
In order to write down the long transmission-time optimum capacity of a wideband, fermionic channel for fixed energy, we require the asymptotic approximation to the number of distinct partitions of $`N`$ (see, e.g., Sec. 24.2.2 of Ref. ):
$$𝒩_N\frac{1}{43^{1/4}N^{3/4}}e^{\pi \sqrt{1/3}\sqrt{N}}.$$
(17)
This gives
$$C_{\text{fermion}}=\frac{\pi }{\mathrm{ln}2}\sqrt{\frac{P}{3h}}\frac{1}{𝒯}\mathrm{log}_2\left[43^{1/4}\left(\frac{P𝒯^2}{h}\right)^{3/4}\right].$$
(18)
Note that, in the limit $`𝒯\mathrm{}`$, the fermionic optimum capacity is smaller than the bosonic optimum capacity by a factor $`\sqrt{2}`$ for given power $`P`$. Note also that these optimum capacities satisfy the information-theoretic counterpart to bound (6) :
$$C<\frac{\pi }{\mathrm{ln}2}\sqrt{\frac{2P}{3h}}$$
(19)
for finite $`𝒯`$.
Asymptotic approximations to $`C`$, analogous to Eqs. (16) and (18), can no doubt also be written down for certain other rational $`g`$ values. However, rather than attempting to derive $`C`$ through the nontrivial route which involves first obtaining the asymptotic approximation to the number of partitions, we can appeal to the fact that different ensemble derivations of the entropy give the same result in the thermodynamic limit when the ensemble energies coincide. In particular, we can instead use expressions (3) and (5) for $`\mu =0`$ (only the energy is constrained and not the particle number) to derive the leading order term in the asymptotic approximation to $`C`$. For example, in the case of ‘semions’ ($`g=1/2`$), carrying out the integrals in (3) and (5) gives $`C_{\text{semion}}(𝒯\mathrm{})=(\pi /\mathrm{ln}2)\sqrt{2P/5h}`$, which falls between the Fermi and Bose capacities, again satisfying bound (19). Solving numerically Eqs. (3) and (5) for a range of $`g`$-values, we find that $`C_g(𝒯\mathrm{})`$ decreases monotonically as $`g`$ increases from 0 to 1 (Fig. 4). Thus bound (19) holds for all $`0g1`$.
However, unlike the analogous upper bound (6) on the single-channel physical entropy current, the information-theoretic bound (19) cannot be approached arbitrarily closely independently of the particle statistics: only for bosons is the upper bound approached in the limit $`𝒯\mathrm{}`$. But recall, from the form of the conditions for approaching the upper bound (6), and also from the form of Eqs. (8) and (10) for fermions, that it is crucial for both ends of the channel to be connected to reservoirs providing two-way energy and entropy flows in the channel. This suggests that, with a suitable generalization of the communication channel allowing for two-way information flow, the channel capacity will approach the upper bound (19) arbitrarily closely independently of the particle statistics.
Consider, therefore, two sender-receiver ‘stations’, with one station at each end of the single channel, thus sharing the channel. Station $`L`$ at the left end encodes information in right-moving states and decodes information from left-moving states, while station $`R`$ at the right end encodes information in left-moving states and decodes information from right-moving states. A single ‘use’ of the channel involves $`L`$ and $`R`$ each sending and subsequently detecting a message, the whole operation taking place during an interval $`𝒯`$. Station $`L`$ uses an input alphabet $`A_L`$ with letters labeled by an index $`a_L=1,\mathrm{},𝒜_L`$ and a set of probabilities $`p_{A_L}(a_L)`$ for transmitting letter $`a_L`$, and an output alphabet $`B_L`$ labeled by $`b_L=1,\mathrm{},_L`$. Station $`R`$ similarly uses an input alphabet $`A_R`$ with transmission probabilities $`p_{A_R}(a_R)`$, and an output alphabet $`B_R`$. The probability that $`R`$ receives letter $`b_R`$, given that $`L`$ sends letter $`a_L`$ is denoted as $`p_{B_R|A_L}(b_R|a_L)`$, and analogously for the conditional probability $`p_{B_L|A_R}(b_L|a_R)`$. We assume throughout that the joint probabilities for $`L`$ and $`R`$ to send a message are uncorrelated, i.e., $`p_{A_L,A_R}(a_L,a_R)=p_{A_L}(a_L)p_{A_R}(a_R)`$. We also assume to begin with that the left- and right-moving information flows do not interfere with each other.
Using the formula for the single channel entropy current as a guide \[see, e.g., Eqs. (11) and (12) of Ref. \], we define the net information transmitted from the $`L`$ and $`R`$ inputs to the $`R`$ output during a single use of the channel to be
$$H(B_R;A_L,A_R)=H(B_R;A_L)H(A_R),$$
(20)
where $`H(B_R;A_L)`$ is defined as in Eq. (13) and $`H(A_R)=_{a_R}p_{A_R}(a_R)\mathrm{log}_2p_{A_R}(a_R)`$. Similarly, the net information transmitted to the $`L`$ output is $`H(B_L;A_R,A_L)=H(B_L;A_R)H(A_L)`$. Note the asymmetry of the two terms on the right-hand-side of definition (20), reflecting an analogous asymmetry of the left- and right-moving components making up the net entropy current . With the information defined with respect to the receiver at the right end of the channel, it makes more sense to use the information $`H(A_R)`$ rather than the mutual information $`H(B_L;A_R)`$ which takes into account the channel noise and receiver properties at the other end of the channel. As we shall soon see when we generalize (20) to include interfering left- and right-moving information, one reason why it might be a good thing to subtract, rather than to add, the information $`H(A_R)`$ is that it gives reasonable answers in familiar examples such as that of a returned or ‘bounced’ message.
But perhaps the most appealing property of $`H(B_R;A_L,A_R)`$ as defined is that it satisfies a generalized Holevo theorem:
$$\underset{\{\widehat{F}_{b_R}\}}{\mathrm{max}}H(B_R;A_L,A_R)S(\widehat{\rho }_L)S(\widehat{\rho }_R)\underset{a_L}{}p_{A_L}(a_L)S(\widehat{\rho }_{a_L})+\underset{a_R}{}p_{A_R}(a_R)S(\widehat{\rho }_{a_R}),$$
(21)
where equality holds if and only if the left input states $`\widehat{\rho }_{a_L}`$ commute and the right input states $`\widehat{\rho }_{a_R}`$ are orthogonal (see, e.g., Sec. IV.B of Ref. ). Inequality (21) is a consequence both of the Holevo theorem (14), which bounds $`H(B_R;A_L)`$, and also of the inequality $`H(A_R)S(\widehat{\rho }_R)_{a_R}p_{A_R}(a_R)S(\widehat{\rho }_{a_R})`$. Note that the latter inequality goes in the opposite direction to that of (14), so that one must subtract $`H(A_R)`$ in order that $`H(B_R;A_L,A_R)`$ be bounded.
Maximizing the information $`H(B_R;A_L,A_R)`$ with respect to the $`R`$ output detection scheme and the $`L`$ and $`R`$ input states and probabilities gives the optimum capacity of the channel. From (21), we have
$$C=\frac{1}{𝒯}\underset{\widehat{\rho }_L,\widehat{\rho }_R}{\mathrm{max}}\underset{\{\widehat{F}_{b_R}\}}{\mathrm{max}}H(B_R;A_L,A_R)\frac{1}{𝒯}\underset{\widehat{\rho }_L}{\mathrm{max}}S(\widehat{\rho }_L)=\frac{S_{\text{max}}}{𝒯}.$$
(22)
The upper bound $`S_{\text{max}}`$ can in fact be attained: find a complete, orthonormal set of diagonalizing basis states $`|a_L`$ for the $`\widehat{\rho }_L`$ which maximizes $`S(\widehat{\rho }_L)`$, i.e., $`\widehat{\rho }_L=_{a_L}q(a_L)|a_La_L|`$. Choose $`\widehat{\rho }_{a_L}=|a_La_L|=\widehat{F}_{a_L}`$ and $`p_{A_L}(a_L)=q(a_L)`$. Choose any set $`\{\widehat{\rho }_{a_R}\}`$ and probabilities $`p_{A_R}(a_R)=\delta _{a_R,a_R^{}}`$ for some fixed $`a_R^{}`$. Then $`H(B_R;A_L,A_R)=H(B_R;A_L)=_{a_L}p_{A_L}(a_L)\mathrm{log}_2p_{A_L}(a_L)=S(\widehat{\rho }_L)=S_{\text{max}}`$. The choice for $`\widehat{\rho }_{a_R}`$ and $`p_{A_R}(a_R)`$ reflects the obvious fact that, for the definition (20), maximizing $`H(B_R;A_L,A_R)`$ requires that $`H(A_R)`$ be minimized, so that any left-moving message component can be sent, provided it is with probability one so that its information content is zero.
Thus, the optimum capacity is just the maximum quantum entropy in bits for right-moving states divided by the transmission time, subject to the given constraints on the channel. Of particular interest are the constraints for which the optimum capacity in the limit $`𝒯\mathrm{}`$ is independent of the statistical parameter $`g`$. Recalling the conditions for approaching asymptotically the entropy bound (6) for $`g>0`$, namely $`\mu _L/k_BT_L+\mathrm{}`$ with $`\mu _L=\mu _R`$ and $`T_L>T_R=0`$, a little thought establishes that two constraints are: fixed power (i.e., fixed energy current) $`P>0`$ and fixed number current $`\dot{N}=0`$. Again, as for the unidirectional optimum capacity, the choice of ensemble for the definition of $`P`$ and $`\dot{N}`$—microcanonical, grand canonical etc.—is immaterial in the limit $`𝒯\mathrm{}`$. Given that the unidirectional optimum capacity for $`0<g1`$ is strictly less than the bosonic optimum capacity $`(\pi /\mathrm{ln}2)\sqrt{2P/3h}`$ in the limit $`𝒯\mathrm{}`$, it may seem paradoxical that additional constraints have to be imposed (namely, $`\dot{N}=0`$) in order to attain the latter, larger capacity. The resolution lies in the fact that the dimension of the channel Hilbert space accessible for information and energy transmission has been doubled through the accomodation of left-moving states.
The two above constraints, while necessary, are not sufficient. The problem lies in the fact that optimization step (22) places no conditions on the left-moving states $`\widehat{\rho }_{a_R}`$, with the result that it is rather easy to find examples where the power $`P`$ can be made arbitrarily small for given $`S_{\text{max}}`$, while at the same time satisfying the constraint $`\dot{N}=0`$. One possible way to overcome this problem is to introduce the further constraint on the left-moving states $`\widehat{\rho }_{a_R}`$ that they be completely degenerate. Then it is possible to show that $`S_{\text{max}}=(\pi /\mathrm{ln}2)\sqrt{2P/3h}`$ in the limit $`𝒯\mathrm{}`$, i.e., $`S_{\text{max}}`$ coincides with the limiting, unidirectional bosonic optimum capacity independently of $`0<g1`$. Furthermore, in the case of bosons adding a left-moving degenerate state does not change the energy current, so that the above constraints can also be applied to bosons with the unidirectional bosonic optimum capacity again being obtained in the limit.
What we have essentially done both here and in the previous section is cancel part of the right-moving energy current component with a left-moving, degenerate component, leaving the information and entropy currents unchanged, thus increasing the optimum capacity and entropy current bound for given energy current \[Eqs. (8) and (10) show this more explicitly\]. What is remarkable is that the optimum capacity (19) and entropy current bound (6) are attained asymptotically for a common set of constraints independent of the statistics $`0g1`$.
In the final part of this section, we generalize our two-way information definition (20) so as to allow for the possibility of interference between the left- and right-moving information flows. Our definition is motivated by the formula for the single-channel entropy current in the presence of elastic scattering in the channel . We define the net information transmitted from the $`L`$ and $`R`$ inputs to the $`R`$ output during a single use of the channel to be
$`H(B_R;A_L,A_R)`$ $`=`$ $`{\displaystyle \underset{a_L,a_R,b_R}{}}p_{B_R|A_L,A_R}(b_R|a_L,a_R)p_{A_L}(a_L)p_{A_R}(a_R)`$ (23)
$`\times `$ $`\mathrm{log}_2\left(p_{B_R|A_L,A_R}(b_R|a_L,a_R)/p_{B_R}(b_R)\right)+{\displaystyle \underset{a_R}{}}p_{A_R}(a_R)\mathrm{log}_2p_{A_R}(a_R),`$ (24)
with an analogous definition for $`H(B_L;A_R,A_L)`$ and where we again assume that the joint probabilities for $`L`$ and $`R`$ to send a message are uncorrelated. The channel interference is conveniently implemented by a unitary ‘scattering’ operator $`\widehat{𝒮}`$, acting on the states as $`\widehat{𝒮}(\widehat{\rho }_L\widehat{\rho }_R)\widehat{𝒮}^{}\widehat{\rho }_L^{}\widehat{\rho }_R^{}`$, where we restrict ourselves to non-correlating interfering processes. The conditional probabilities are constructed as follows:
$$p_{B_R|A_L,A_R}(b_R|a_L,a_R)=\text{tr}\left[(\widehat{F}_{b_R}\widehat{1})\widehat{𝒮}(\widehat{\rho }_{a_L}\widehat{\rho }_{a_R})\widehat{𝒮}^{}\right]$$
(26)
and
$$p_{B_L|A_L,A_R}(b_L|a_L,a_R)=\text{tr}\left[(\widehat{1}\widehat{F}_{b_L})\widehat{𝒮}(\widehat{\rho }_{a_L}\widehat{\rho }_{a_R})\widehat{𝒮}^{}\right],$$
(27)
where the the right and left detector operators are written as $`\widehat{F}_{b_R}\widehat{1}`$ and $`\widehat{1}\widehat{F}_{b_L}`$, respectively, reflecting the fact that, in the absence of interference, i.e., when $`\widehat{𝒮}`$ is the identity operator, the right(left) detector can only receive left(right) input states. Note that definition (24) reduces to the two-way information definition (20) when there is no interference.
The more general, two-way information definition (24) can be applied to certain situations which are beyond the scope of the unidirectional mutual information (13). As a simple example, consider the situation of a ‘bounced’ message, an all too common occurence with electronic mail. This example can be modeled as follows: let the right letters be encoded in the orthonormal states $`|a_R`$ and sent with probabilities $`p_{A_R}(a_R)`$. Let the states $`|a_L`$ encoding the left letters, and sent with probabilities $`p_{A_L}(a_L)`$, be in one-to-one correspondence with the right states $`|a_R`$, with the mapping achieved simply by reversing the propagation direction. Let the right detector be characterized by projection operators $`\widehat{F}_{a_R}=|a_La_L|`$. Finally, suppose the scattering operator reverses the direction of the propagating states, i.e., $`\widehat{𝒮}|a_{R(L)}=|a_{L(R)}`$. Then evaluating the two-way information (24), we find that the first term on the right-hand-side reduces to the information $`H(A_R)`$, thus cancelling the second term and giving the value $`H(B_R;A_L,A_R)=0`$. This coincides with our common-sense measure: if a message gets bounced back, then no information was sent.
From the Holevo theorem (14) for unidirectional information flow and also the inequality $`H(A_R)S(\widehat{\rho }_R)_{a_R}p_{A_R}(a_R)S(\widehat{\rho }_{a_R})`$, it is straightforward to show that the information $`H(B_R;A_L,A_R)`$ as defined in (24) also satisfies a generalized Holevo theorem:
$$\underset{\{\widehat{F}_{b_R}\},\widehat{𝒮}}{\mathrm{max}}H(B_R;A_L,A_R)S(\widehat{\rho }_L)\underset{a_L}{}p_{A_L}(a_L)S(\widehat{\rho }_{a_L}).$$
(28)
Such a bound enables us to determine the optimum capacity allowing also for interfering left- and right-moving information flows. We shall leave this to a future investigation.
## IV conclusion
We have provided evidence for the validity of two related conjectures which state that the entropy current and optimum capacity for information transmission of a single channel are universally bounded for given energy current/power, independently of the channel materials properties and particle statistics according to Haldane’s definition. What is most notable, is that these bounds can be approached arbitrarily closely no matter the particle statistics. A less general, tighter bound on the entropy current was also conjectured, from which the recently discovered statistics-independent thermal conductance bound follows as a special case. The statistics-independent, limiting bound on the optimum capacity required a generalisation of the definition for the transmitted information, allowing for two-way information flow. The bound then followed from a generalized Holevo theorem, with certain constraints placed on the channel and input states.
The results presented here can be extended in several ways. It would be more satisfying to have an analytic proof of the conjectured bounds, rather than an exhaustive numerical check. The entropy current bound should be tested under more general conditions, for example in the presence of channel scattering. Similarly, the optimum capacity bound should be tested also allowing for interference between two-way information flows.
Finally, we point out the recent demonstration that Holevo’s theorem follows from Landauer’s principle of information erasure . In the light of this, it would be interesting to try to rederive the universal upper bound on the optimum capacity starting from Landauer’s erasure principle.
###### Acknowledgements.
M.P.B. thanks Michael Roukes and his group at Caltech for stimulating discussions and for their hospitality during a visit. We would also like to thank Jay Lawrence and Martin Plenio for helpful discussions and also for suggesting improvements to the manuscript. V.V. was partially funded by a Nuffield Foundation Bursary for Undergraduate Research. |
warning/0001/astro-ph0001029.html | ar5iv | text | # Long-term Variability Properties and Periodicity Analysis for Blazars
## 1 Introduction
The nature of the central regions of quasars and other Active Galactic Nuclei (AGNs) is still an open problem. The study of AGN optical variability can yield valuable information about the mechanisms operating in these sources, with important implications for quasar modeling (see, for instance, Fan et al. 1998). Some particular objects have been claimed to display periodicity in their lightcurves over a variety of timescales (e.g. Jurkevich 1971; Hargen-Thorn et al. 1987; Sillanpaa et al. 1988; Webb et al. 1988; Babadzhanyants & Belokon 1991; Kidger et al. 1992; Liu et al. 1995; Marchenko et al. 1996; Fan et al. 1997, 1998, Fan & Su 1998; Fan 1999a; Zhang et al. 1998; Fan & Lin 1999a; Fan et al. 1999a; Lainela et al. 1999; Lin 1999; Su 2000 and references therein), but, in general, the clear identification of periodic behaviour has been very elusive due to the complexity of the optical lightcurves and the lack of databases large enough as to provide an adequate sampling over large periods.
Blazars are an extremely subclass of AGNs, characterized as rapid and high amplitude variations, high and variable polarization, superluminal radio components, nonthermal continuum. Blazars have been monitored for a long time, compilations are available for optical (Fan & Lin 1999b,c) and infrared (Fan & Lin 1999d; Fan 1999b) bands. The infared data are available in http://xxx.lanl.gov/astro-ph/9908104 and http://xxx.lanl.gov/astro-ph/9910269. The long-term measurements make it possible for one to search for periodicities from the light curves and to discuss the long-term variability properties. Because the characteristic of the astronomical measurements, we adopted the Jurkevich (Jerkevich 1971) and DCF (Discrete Correlation Function) (see Fan & Lin 1999a) methods to the available data in searching for the periods. In the second section we will give the variation properties, in the third section, we present the results and finally in section 4, we give some discussions.
## 2 Variation Property
From the complication (Fan et al. 1999b,c,d; Fan 1999b), we found that the largest variations are comparable in both the optical and infrared bands. Largest variations at different wavelengths increase with decreasing wavelength. The variations are also correlated with the highest observed optical polarization, which can be explained using the beaming model. Some objects show that the spectra flatten when the sources brighten while the opposite cases are for some others. For the largest infrared and optical variations, we listed them in Table 1.
## 3 Periodicity Analysis Methods
### 3.1 Jurkevich Method
The Jurkevich method ( Jurkevich 1971, a lso see Kidger et al. 1992; Fan et al. 1998a; Fan 1999a) is based on the expected mean square deviation and it is less inclined to generate spurious periodicity than the Fourier analysis used by other authors (e.g. Fan et al. 1997). It tests a run of trial periods around which the data are folded. All data are assigned to $`m`$ groups according to their phases around each trial period. The variance $`V_i^2`$ for each group and the sum $`V_m^2`$ of all groups are then computed. If a trial period equals the true one, then $`V_m^2`$ reaches its minimum. So, a “good” period will give a much reduced variance relative to those given by other false trial periods and with almost constant values. A further test is the relationship between the depth of the minimum and the noise in the “flat” section of the $`V_m^2`$ curve close to the adopted period. If the absolute value of the relative change of the minimum to the “flat” section is large enough as compared with the standard error of this “flat” section (say, five times), the periodicity in the data can be considered as significant and the minimum as highly reliable (Fan et al. 1998a).
### 3.2 DCF Method
The DCF (Discrete Correlation Function) method is intended for analyses of the correlation of two data set. It is described in detail by Edelson & Krolik (1988) (also see Fan et al. 1998b). This method can indicate the correlation of two variable temporal series with a time lag, and can be applied to the periodicity analysis of a unique temporal data set. If there is a period, $`P`$, in the lightcurve, then the DCF should show clearly whether the data set is correlated with itself with time lags of $`\tau `$ = 0 and $`\tau `$ = $`P`$ (see Fan & Lin 1999a). We have implemented the method as follows.
Firstly, we have calculated the set of unbinned correlation (UDCF) between data points in the two data streams $`a`$ and $`b`$, i.e.
$$UDCF_{ij}=\frac{(a_i\overline{a})\times (b_j\overline{b})}{\sqrt{\sigma _a^2\times \sigma _b^2}},$$
(1)
where $`a_i`$ and $`b_j`$ are points in the data sets, $`\overline{a}`$ and $`\overline{b}`$ are the average values of the data sets, and $`\sigma _a`$ and $`\sigma _b`$ are the corresponding standard deviations. Secondly, we have averaged the points sharing the same time lag by binning the $`UDCF_{ij}`$ in suitably sized time-bins in order to get the $`DCF`$ for each time lag $`\tau `$:
$$DCF(\tau )=\frac{1}{M}\mathrm{\Sigma }UDCF_{ij}(\tau ),$$
(2)
where $`M`$ is the total number of pairs. The standard error for each bin is
$$\sigma (\tau )=\frac{1}{M1}\{\mathrm{\Sigma }[UDCF_{ij}DCF(\tau )]^2\}^{0.5}.$$
(3)
## 4 Results and Discussion
Very recently, we compiled four data bases for BL Lac objects and OVV/HPQs in the optical and infrared bands (Fan & Lin 1999b,c,d; Fan 1999b). The data bases can be used to discuss the long-term variability properties, the correlated variations and to search for periodicity signatures in the light curves. When the both methods are adopted to those optical data, the both methods give consistent results. The results are presented in Table 2. It is interesting to notice that there is a common period of $``$ 1.0 year, which is likely from the effect of the Sun on the measurements. The detail analysis is presented in the paper by Fan et al. (1999c)
For the long-term periodicity of variations, there are several explanations: the double black hole model, the hectic jet model, the slim disk model, and the effect of external perturbations to the accretion disk. (e.g. Sillanpaa et al. 1988; Meyer & Meyer-Hofmeister 1984; Horiuchi & Kato 1990; Abraham & Romero 1999; Fan et al. 1997, 1998a, 1999; Villata & Raiteri 1999; Romero et al. 2000).
This work has been supported by the National PanDeng Project of China and the National Natural Scientific Foundation of China. I thank Prof. K.S. Cheng to provide me with the opportunity to present this work and Prof. R.G. Lin and C.Y. Su for their help. |
warning/0001/math0001177.html | ar5iv | text | # The module of logarithmic p-forms of a locally free arrangement
## 1 Introduction
If $`X`$ is a complex manifold and $`D`$ a divisor with normal crossings ($`D=D_i`$, $`D_i`$ smooth and meet transversely), then associated to $`D`$ is the sheaf $`\mathrm{\Omega }^1(\text{log }D)`$ of meromorphic one forms with logarithmic poles on $`D`$. Deligne introduced this sheaf in and shows (among other things) that it is locally free. Dolgachev and Kapranov seem to have been the first to examine in depth the case where $`D`$ is a set of hyperplanes in general position in $`𝐏^n`$; one striking result they obtain is that if $`D=_{i=1}^dH_i`$ and $`d2n+3`$, then the hyperplanes can be recovered from $`\mathrm{\Omega }^1(\text{log }D)`$ unless the $`H_i`$ osculate a rational normal curve of degree $`n`$.
We also consider the case when the divisor is a set of hyperplanes in $`𝐏^n`$, but assume only that the hyperplanes are distinct. There are two main themes of this paper. In , Solomon and Terao give a formula for the Poincaré polynomial of an (essential, central) arrangement in terms of the Hilbert series of certain graded modules $`D^i`$ associated to the arrangement. The formula generalizes Terao’s famous freeness theorem : If $`D^1`$ is a free module, then the Poincaré polynomial factors. This suggests a connection to Chern polynomials; motivated by the Solomon-Terao result we prove a formula relating the Chern polynomial of a bundle $``$ on $`𝐏^n`$ to the Hilbert series of the modules $`_{m𝐙}H^0(𝐏^n,^i(m))`$.
Since an arbitrary arrangement does not have normal crossings, $`\mathrm{\Omega }^1(\text{log }D)`$ is in general no longer locally free. Silvotti studies this situation, and remedies the problem by blowing up the arrangement at the non-normal crossings; $`\sigma ^{}D`$ automatically has normal crossings on the blowup X, so yields a locally free sheaf $`\mathrm{\Omega }_X^1(\text{log }\sigma ^{}D)_j`$. Using a vanishing result of Esnault, Schechtman and Viehweg , Silvotti obtains a formula for the coefficients of the Poincaré polynomial in terms of $`\chi (\mathrm{\Lambda }^i_j)`$. However, the computations can be quite complicated; in particular, Silvotti does not recover Terao’s theorem.
The second point of this paper is that even when the hyperplanes are not in general position, there are situations where $`\mathrm{\Omega }^1(\text{log }D)`$ is a vector bundle on $`𝐏^n`$. We relate $`\mathrm{\Omega }^1(\text{log }D)`$ to the modules $`D^i`$ mentioned above, and prove a criterion for the associated sheaves to be locally free. This class of arrangements was studied by Yuzvinsky in ; Yuzvinsky proves that for such arrangements the Hilbert polynomial of $`D^1`$ is a combinatorial invariant. We show that the Chern polynomial of the dual of $`D^1`$ is in fact the Poincaré polynomial of the arrangement (truncated by $`t^{n+1}`$). Hence, the Hilbert polynomial of $`D^1`$ may be obtained from the Poincaré polynomial via Hirzebruch-Riemann-Roch. We close with some results specific to the situation where $`\mathrm{\Omega }^1`$ or $`D^1`$ has projective dimension one. First, we review some facts about arrangements.
## 2 Hyperplane Arrangements
A hyperplane arrangement $`𝒜`$ is a finite collection of codimension one linear subspaces of a fixed vector space V. $`𝒜`$ is central if each hyperplane contains the origin 0 of V. A fundamental invariant of $`𝒜`$ is the Poincaré polynomial $`\pi (𝒜,t)`$. There are various ways of defining $`\pi (𝒜,t)`$; the simplest is from the intersection lattice $`L_𝒜`$ of $`𝒜`$. $`L_𝒜`$ consists of the intersections of the elements of $`𝒜`$, ordered by reverse inclusion. The rank function on $`L_𝒜`$ is given by the codimension in $`V`$. V is the lattice element $`\widehat{0}`$; the rank one elements are the hyperplanes themselves. $`𝒜`$ is called essential if rank $`L_𝒜=`$ dim $`V`$. Henceforth, unless explicitly stated otherwise, all arrangements will be essential and central, and $`V`$ will be $`k^{n+1}`$, with $`k`$ an arbitrary field. We briefly review some fundamental definitions; for more information see Orlik and Terao ().
###### Definition 2.1
The Möbius function $`\mu `$ : $`L_𝒜𝐙`$ is defined by
$$\begin{array}{ccc}\mu (\widehat{0})& =& 1\\ \mu (t)& =& \underset{s<t}{}\mu (s)\text{, if }\widehat{0}<t\end{array}$$
###### Definition 2.2
The Poincaré polynomial $`\pi (𝒜,t)`$ and characteristic polynomial $`\chi (𝒜,t)`$ are defined by:
$$\pi (𝒜,t)=\underset{XL_𝒜}{}\mu (X)(t)^{rank(X)},\text{ }\chi (𝒜,t)=\underset{XL_𝒜}{}\mu (X)t^{dim(X)}$$
The two polynomials are related via $`\chi (𝒜,t)=t^{n+1}\pi (𝒜,t^1)`$. Let $`S=\mathrm{Sym}(V^{})`$, $`\underset{¯}{m}`$ the irrelevant maximal ideal, $`K`$ the fraction field of $`S`$ and suppose $`𝒜`$ consists of $`d`$ distinct hyperplanes in $`V`$. For each hyperplane $`H_i`$ of $`𝒜`$, fix $`l_i`$ a nonzero linear form vanishing on $`H_i`$ and put $`Q=\underset{1}{\overset{d}{}}l_i`$. Denote the module of $`p`$ differentials over $`k`$ of $`S`$ and $`K`$ by $`\mathrm{\Omega }_S^p`$ and $`\mathrm{\Omega }_K^p`$, respectively, and let $`\mathrm{Der}_kS`$ denote the module of $`k`$ derivations of $`S`$.
###### Definition 2.3
$`D^p(𝒜)`$ is the submodule of $`\mathrm{\Lambda }^p(\mathrm{Der}_kS)`$ defined by
$$D^p(𝒜)=\{\theta \mathrm{\Lambda }^p(\mathrm{Der}_kS)\text{ }|\text{ }\theta (Q,f_2,\mathrm{},f_p)(Q),f_iS\}.$$
$`\mathrm{\Omega }^p(𝒜)`$ is the submodule of $`\mathrm{\Omega }_K^p`$ defined by
$$\mathrm{\Omega }^p(𝒜)=\{\omega \mathrm{\Omega }_K^p\text{ }|\text{ }Q\omega \mathrm{\Omega }_S^p\text{ and }Qd\omega \mathrm{\Omega }_S^{p+1}\}.$$
$`D^1(𝒜)`$ is usually called the module of $`𝒜`$ derivations, while $`\mathrm{\Omega }^p(𝒜)`$ is called the module of logarithmic $`p`$ forms with poles along $`𝒜`$. When the arrangement is clear from the context, we will drop it from the notation. Note that we have $`D^0(𝒜)=\mathrm{\Omega }^0(𝒜)=S`$ and we make the convention $`D^p(𝒜)=\mathrm{\Omega }^p(𝒜)=0`$, for $`p<0`$.
If char $`kd`$, then $`D^1(𝒜)`$ has a direct sum decomposition as $`D_0^1S(1)`$, where $`D_0^1`$ is the kernel of the Jacobian matrix of $`Q`$ and $`S(1)`$ corresponds to the Euler derivation. Correspondingly, we have a decomposition $`\mathrm{\Omega }^1\mathrm{\Omega }_0^1S(1)`$.
Since all these modules are graded, we may consider the corresponding sheaves on $`𝐏^n`$, written as usual as $`\stackrel{~}{\mathrm{\Omega }^1}`$ for the sheaf associated to $`\mathrm{\Omega }^1`$. An arrangement is called generic if for every $`H_1,\mathrm{},H_m𝒜`$, with $`mn+1`$, $`rank(H_1\mathrm{}H_m)=m`$. For a generic arrangement $`𝒜`$, the sheaves $`\mathrm{\Omega }^1(\text{log }D)`$ and $`\stackrel{~}{\mathrm{\Omega }^1}`$ are related as follows: after a change of coordinates, we may assume the first $`n+1`$ hyperplanes are the coordinate hyperplanes; for $`i\{1,\mathrm{},dn1\}`$ write $`l_i=_{j=0}^na_{i,j}x_j`$ for the remaining hyperplanes. In , Ziegler gives a free resolution for $`\mathrm{\Omega }^1`$ for a generic arrangement:
$$0S^{dn1}\stackrel{\tau }{}S(1)^d\mathrm{\Omega }^10,$$
where $`\tau `$ is given by
$$\left(\begin{array}{cccc}a_{1,0}x_0& \mathrm{}& \mathrm{}& a_{dn1,0}x_0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_{1,n}x_n& \mathrm{}& \mathrm{}& a_{dn1,n}x_n\\ l_1& 0& \mathrm{}& 0\\ 0& l_2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& 0\\ 0& \mathrm{}& 0& l_{dn1}\end{array}\right)$$
In corollary 3.4 of , Dolgachev and Kapranov present $`\mathrm{\Omega }^1(\text{log }D)`$ as the cokernel of a map
$$V𝒪_{𝐏^n}(1)\stackrel{\tau ^{}}{}W𝒪_{𝐏^n},$$
where $`V`$ is the subspace of $`k^d`$ consisting of relations on the linear forms defining $`𝒜`$, $`W`$ is the subspace of $`k^d`$ orthogonal to $`(1,1,\mathrm{},1)`$, and $`\tau ^{}:(a_1,\mathrm{},a_d)(a_1l_1,\mathrm{},a_dl_d)`$. Thus, the images of $`\tau `$ and $`\tau ^{}`$ are isomorphic, and we have
$`\begin{array}{ccccccccc}0& & S^{dn1}& \stackrel{\tau ^{}}{}& S^{d1}(1)& & \mathrm{\Omega }^1(\text{log }D)(1)& & 0\\ & & & & & & & & \\ 0& & S^{dn1}& \stackrel{\tau }{}& S^d(1)& & \mathrm{\Omega }^1& & 0\\ & & & & & & & & \\ & & 0& & S(1)& & S(1)& & 0\end{array}`$
By the snake lemma, we have $`\mathrm{\Omega }^1(\text{log }D)(1)\stackrel{~}{\mathrm{\Omega }}_0^1`$. This also follows from the local description, but the above makes explicit the different gradings.
For every $`X`$ in the intersection lattice $`L_𝒜`$, the subarrangement $`𝒜_X`$ of $`𝒜`$ is defined by
$$𝒜_X=\{H𝒜\text{ }|\text{ }XH\}.$$
In general this is not an essential arrangement, but we can write it as $`𝒜_X𝒜_X^{}\times \mathrm{\Phi }`$, where $`𝒜_X^{}`$ is essential and $`\mathrm{\Phi }`$ is an empty arrangement.
The functors on the intersection lattice $`XD^p(𝒜_X)`$ and $`X\mathrm{\Omega }^p(𝒜_X)`$ are local (see Orlik and Terao , Chapter 4.6). What we will use is the fact that if $`\underset{¯}{q}`$ is a prime ideal in $`S`$ and $`X=_{\alpha _H\underset{¯}{q}}H`$, then we have a canonical isomorphism:
$$D^p(𝒜)_{\underset{¯}{q}}D^p(𝒜_X)_{\underset{¯}{q}}$$
and a similar one for $`\mathrm{\Omega }^p`$.
By definition, an arrangement $`𝒜`$ is free if $`D^1(𝒜)`$ is a free $`S`$-module. Following Yuzvinsky , we will say that an (essential, central) arrangement $`𝒜`$ is locally free if for every $`XL_𝒜`$ with $`\mathrm{rank}X<\mathrm{dim}V`$, the arrangement $`𝒜_X`$ is free.
For a graded module $`M`$, let $`P(M,X)`$ be its Hilbert series. There is a beautiful relation between the modules $`D^p(𝒜)`$ and the characteristic polynomial:
###### Theorem 2.4
(Solomon and Terao,)
$$\chi (𝒜,t)=(1)^{n+1}lim_{X1}\underset{p0}{}P(D^p(𝒜);X)(t(X1)1)^p.$$
There is a dual version of this theorem, which replaces $`D^p(𝒜)`$ with $`\mathrm{\Omega }^p(𝒜)`$. In certain situations, not all the modules $`D^p(𝒜)`$ are needed to compute $`\pi (𝒜,t)`$; the paradigm for this is the case of free arrangements.
###### Theorem 2.5
(Terao,) If $`D^1(𝒜)`$ is free, then $`\pi (𝒜,t)=(1+a_it)`$, where the $`a_i`$ are the degrees of the generators of $`D^1(𝒜)`$.
For arrangements on $`𝐏^2`$, $`\stackrel{~}{\mathrm{\Omega }}^1`$ is always locally free, and suffices to determine the Poincaré polynomial (). For every coherent sheaf $``$ on $`𝐏^n`$, we denote by $`H_{}^0()`$ the $`S`$-module $`_{m𝐙}H^0(𝐏^n,(m))`$. Motivated by Theorem 2.4 we prove
Theorem For every rank $`r`$ bundle $``$ on $`𝐏^n`$,
$$c_t()=lim_{X1}(1)^rt^r(1X)^{n+1r}\underset{i=0}{\overset{r}{}}P(H_{}^0(^i);X)(\frac{X1}{t}1)^i.$$
As a consequence, we get the following generalization of Theorem 2.5:
Theorem If $`𝒜`$ is an arrangement such that $`\stackrel{~}{\mathrm{\Omega }^1}`$ is locally free and $`\overline{\pi }(𝒜,t)`$ is the class of $`\pi (𝒜,t)`$ in $`𝐙[t]/(t^{n+1})`$, then
$$\overline{\pi }(𝒜,t)=c_t(\stackrel{~}{\mathrm{\Omega }^1}).$$
We characterize those arrangements for which $`\stackrel{~}{\mathrm{\Omega }}^p`$ is a bundle:
Theorem $`\stackrel{~}{\mathrm{\Omega }}^p`$ is a bundle iff for every $`XL_𝒜`$ with $`\mathrm{rank}X<\mathrm{dim}V`$, $`\mathrm{\Omega }^p(𝒜_X)`$ is free. We will use freely results from commutative algebra for which our main reference is Eisenbud , as well as results about Chern classes of vector bundles on projective space, for which we refer to Fulton .
Acknowledgements We are very grateful to David Eisenbud for useful discussions.
## 3 Locally Free Arrangements
We start with a general lemma about the depth of the modules $`\mathrm{\Omega }^p`$ and $`D^p`$.
###### Lemma 3.1
For every central arrangement $`𝒜`$ in $`V`$, with $`\mathrm{dim}V=n+12`$, we have $`\mathrm{depth}\mathrm{\Omega }^p2`$ and $`\mathrm{depth}D^p2`$, for every $`p`$, $`1pn+1`$.
Proof. We consider the case of the modules $`\mathrm{\Omega }^p`$. Recall that if $`K`$ is the quotient field of $`S`$, then
$$\mathrm{\Omega }^p=\{\omega ^p\mathrm{\Omega }_S_SK|Q\omega ^p\mathrm{\Omega }_S,Qd\omega ^{p+1}\mathrm{\Omega }_S\},$$
where $`Q`$ is the product of the linear forms defining the elements of $`𝒜`$. In particular, $`\mathrm{\Omega }^p`$ is torsion-free and therefore $`\mathrm{depth}\mathrm{\Omega }^p1`$.
We have to prove that if $`0aS`$ and $`\omega \mathrm{\Omega }^p`$ are such that $`\underset{¯}{m}\omega a\mathrm{\Omega }^p`$, then $`\omega a\mathrm{\Omega }^p`$. Note that since $`\mathrm{depth}S2`$ and the $`S`$-modules $`\mathrm{\Omega }_S^i`$ are free, if for $`\tau \mathrm{\Omega }_K^i`$ and $`0bS`$ we have $`\underset{¯}{m}\tau b\mathrm{\Omega }_S^i`$, then $`\tau \mathrm{\Omega }_S^i`$.
By definition, we have $`\underset{¯}{m}Q\omega a\mathrm{\Omega }_S^p`$ and the above observation gives $`Q\omega /a\mathrm{\Omega }_S^p`$. For every $`f\underset{¯}{m}`$ we have also $`Qd(f\omega /a)\mathrm{\Omega }_S^{p+1}`$. We use
$$Qd(f\omega /a)=dfQ\omega /a+Qfd(\omega /a).$$
Since we have already seen that $`Q\omega /a\mathrm{\Omega }_S^p`$ we obtain $`\underset{¯}{m}Qd(\omega /a)\mathrm{\Omega }_S^{p+1}`$. One more application of the above observation gives $`Qd(\omega /a)\mathrm{\Omega }_S^{p+1}`$ and therefore, $`\omega /a\mathrm{\Omega }^p`$, which completes the proof.
The proof of the fact that $`\mathrm{depth}D^p2`$ is similar, using the definition of this module. $`\mathrm{}`$
The following Proposition is the generalization of Theorem 4.75 in Orlik and Terao which is the case $`p=1`$. Though the general result seems known to experts, we include the proof, as we could not find a reference in the literature.
###### Proposition 3.2
For every central arrangement $`𝒜`$, each of the modules $`D^p`$ and $`\mathrm{\Omega }^p`$ is dual to the other.
Proof. A standard generalization of the argument in Orlik and Terao , Proposition 4.74, gives a bilinear map of $`S`$-modules $`\mathrm{\Omega }^p\times D^pS`$, which induces morphisms $`\alpha :\mathrm{\Omega }^p\mathrm{Hom}_S(D^p,S)`$ and $`\beta :D^p\mathrm{Hom}_S(\mathrm{\Omega }^p,S)`$. The proofs of the fact that $`\alpha `$ and $`\beta `$ are isomorphisms are similar, so that we will give the proof only for $`\alpha `$. We make induction on $`n1`$, the case $`n=1`$ being straightforward.
Lemma 3.1 gives $`\mathrm{depth}\mathrm{\Omega }^p2`$, while it is an easy exercise to see that since $`n1`$, for every graded $`S`$-module $`M`$, $`\mathrm{depth}\mathrm{Hom}_S(M,S)2`$. Since for a module $`N`$ of depth at least two, $`N_tH^0(\stackrel{~}{N}(t))`$, in order to prove that $`\alpha `$ is an isomorphism, it is enough to prove that it is an isomorphism at the sheaf level i.e. $`\alpha S_{\underset{¯}{q}}`$ is an isomorphism for every prime ideal $`\underset{¯}{q}\underset{¯}{m}`$. Note that if the Proposition is true for an arrangement $`𝒜`$ and every $`p`$, then it is true also for $`𝒜\times \mathrm{\Phi }`$, where $`\mathrm{\Phi }`$ is the empty arrangement in $`k`$. Indeed, if $`S^{}=\mathrm{Sym}(V^{}\times k)`$ is the polynomial ring corresponding to $`𝒜\times \mathrm{\Phi }`$, then by Orlik and Terao , Proposition 4.84 and Solomon and Terao , Proposition 5.8, we have canonical isomorphisms
$$\mathrm{\Omega }^p(𝒜\times \mathrm{\Phi })(\mathrm{\Omega }^p(𝒜)_SS^{})(\mathrm{\Omega }^{p1}(𝒜)(1)_SS^{})$$
$$D^p(𝒜\times \mathrm{\Phi })(D^p(𝒜)_SS^{})(D^{p1}(𝒜)(1)_SS^{}).$$
Therefore, it follows by induction that we may assume $`𝒜`$ to be essential.
For a prime ideal $`\underset{¯}{q}\underset{¯}{m}`$, if we take $`X=_{\alpha _H\underset{¯}{q}}H`$, then $`XL_𝒜`$, and $`\mathrm{rank}X<\mathrm{dim}V`$, since $`𝒜`$ is essential. But since $`\mathrm{\Omega }^p()`$ and $`D^p()`$ are local functors, we have canonical isomorphisms $`D_{\underset{¯}{q}}^pD^p(𝒜_X)_{\underset{¯}{q}}`$ and $`\mathrm{\Omega }_{\underset{¯}{q}}^p\mathrm{\Omega }^p(𝒜_X)_{\underset{¯}{q}}`$. Since $`𝒜_X`$ is not essential, we have seen that it satisfies the conclusion of the Proposition, and therefore we get $`\alpha S_{\underset{¯}{q}}`$ isomorphism. $`\mathrm{}`$
###### Theorem 3.3
For an (essential, central) arrangement $`𝒜`$ and every positive integer $`p`$, the following are equivalent:
1. $`\stackrel{~}{D^p}`$ is locally free on $`𝐏^n`$.
1’. $`\stackrel{~}{\mathrm{\Omega }^p}`$ is locally free on $`𝐏^n`$.
2. For every $`XL_𝒜`$ with $`\mathrm{rank}X<\mathrm{dim}V`$, $`D^p(𝒜_X)`$ is free.
2’. For every $`XL_𝒜`$ with $`\mathrm{rank}X<\mathrm{dim}V`$, $`\mathrm{\Omega }^p(𝒜_X)`$ is free.
Proof. The equivalences $`11^{}`$ and $`22^{}`$ follow from Proposition 3.2.
For the proof of $`12`$, let $`XL_𝒜`$ be a nonzero linear subspace and $`I_XS`$ its ideal. By making a linear change of variables we can assume that $`I_X=(X_0,\mathrm{},X_{r1})`$, where $`r=\mathrm{rank}X`$. Since $`𝒜_X𝒜_{}^{}{}_{X}{}^{}\times \mathrm{\Phi }^{n+1r}`$, if $`S_1=k[X_0,\mathrm{},X_{r1}]`$, then
$$D(𝒜_X)(D(𝒜_{}^{}{}_{X}{}^{})_{S_1}S)S^{n+1r}$$
and more generally
$$D^p(𝒜_X)\underset{0ip}{}(D^i(𝒜_{}^{}{}_{X}{}^{})S^{\left(\genfrac{}{}{0pt}{}{n+1r}{pi}\right)}).$$
Since $`D^i(𝒜_{}^{}{}_{X}{}^{})`$ is a free $`S_1`$ module if and only if $`D^i(𝒜_{}^{}{}_{X}{}^{})_{I_X}`$ is a free $`S_1`$-module, it follows that $`D^p(𝒜_X)`$ is free if and only if $`D^p(𝒜_X)_{I_X}`$ is free. But because $`\stackrel{~}{D^p}`$ is locally free and $`I_X\underset{¯}{m}`$, $`D_{I_X}^p`$ is free over $`S_{I_X}`$. On the other hand, since $`D^p()`$ is a local functor we have
$$D_{I_X}^pD^p(𝒜_X)_{I_X},$$
and therefore $`D^p(𝒜_X)`$ is free.
In order to prove $`21`$, let us consider a prime ideal $`\underset{¯}{q}`$ different from $`\underset{¯}{m}`$. If we take $`X=_{\alpha _H\underset{¯}{q}}H`$, then $`XL_𝒜`$, $`\mathrm{rank}X<\mathrm{dim}V`$ and because $`D^p()`$ is a local functor we have $`D_{\underset{¯}{q}}^pD^p(𝒜_X)_{\underset{¯}{q}}`$, which is free over $`S_{\underset{¯}{q}}`$ by hypothesis. This concludes the proof of the theorem. $`\mathrm{}`$
By taking $`p=1`$ in the above theorem we obtain the following:
###### Corollary 3.4
An arrangement $`𝒜`$ is locally free if and only if $`\stackrel{~}{D^1}`$ is locally free on $`𝐏^n`$.
###### Remark 3.5
A famous conjecture due to Terao asserts that the freeness of an arrangement depends only on the intersection lattice. This is equivalent to the fact that the local freeness of an arrangement depends only on the intersection lattice. Indeed, the fact that the second statement is a consequence of Terao’s conjecture follows immediately by induction on rank. Conversely, if two arrangements $`𝒜_1`$ and $`𝒜_2`$ have isomorphic lattices and $`𝒜_1`$ is free, consider the product arrangements $`𝒜_1^{}=𝒜_1\times B_1`$ and $`𝒜_2^{}=𝒜_2\times B_1`$, where $`B_1`$ is the Boolean arrangement in $`W=k`$. Then we have $`𝒜_1^{}`$ free and in particular locally free, while $`𝒜_2^{}`$ is free if and only if it is locally free if and only if $`𝒜_2`$ is free.
Notice also that since free arrangements are always locally free, Terao’s conjecture becomes a question on the splitting of vector bundles on $`𝐏^n`$.
###### Lemma 3.6
For every arrangement $`𝒜`$, every $`p1`$ and every $`XL_𝒜`$, $`D^p(𝒜_X)`$ is free if and only if $`D^i(𝒜_X^{})`$ is free, for all $`i`$ with $`p\mathrm{dim}Xip`$.
Proof. We have seen in the proof of the above Theorem 3.3 that we can write
$$D^p(𝒜_X)=\underset{0ip}{}D^i(𝒜_X^{})_{S_1}S^{\left(\genfrac{}{}{0pt}{}{n+1r}{pi}\right)},$$
where $`S_1=k[X_0,\mathrm{},X_{r1}]`$ and $`r=n+1\mathrm{dim}X`$. To conclude it is enough to notice that $`S^{\left(\genfrac{}{}{0pt}{}{n+1r}{pi}\right)}0`$ if and only if $`p\mathrm{dim}Xip`$. $`\mathrm{}`$
###### Corollary 3.7
For every arrangement $`𝒜`$ and every $`XL_𝒜`$ such that $`\mathrm{dim}Xp1`$, $`D^p(𝒜_X)`$ is free if and only if $`𝒜_X`$ is free. In particular, if $`\stackrel{~}{D^p}`$ is locally free, then for every $`XL_𝒜`$ with $`\mathrm{dim}X\mathrm{max}\{1,p1\}`$, $`𝒜_X`$ is free.
Proof. The first assertion follows from the above lemma, since we have $`p\mathrm{dim}X1`$. The second assertion follows from the first one and Theorem 3.3. $`\mathrm{}`$
###### Corollary 3.8
For every arrangement $`𝒜`$, $`\stackrel{~}{D^1}`$ is locally free if and only if $`\stackrel{~}{D^2}`$ is locally free.
Proof. The “if” part follows from the previous corollary in the case $`p=2`$, while the “only if” part is a consequence of the more general proposition below. $`\mathrm{}`$
###### Proposition 3.9
If $`\stackrel{~}{D^1}`$ is locally free, then the natural map $`^p(D^1)D^p`$ induces an isomorphism $`^p(\stackrel{~}{D^1})\stackrel{~}{D^p}`$ for every $`p1`$. In particular, $`\stackrel{~}{D^p}`$ is locally free. The similar assertion about the natural morphism $`^p(\mathrm{\Omega }^1)\mathrm{\Omega }^p`$ is also true.
Proof. We have to prove that for every $`\underset{¯}{q}\mathrm{Spec}(S)`$, $`\underset{¯}{q}\underset{¯}{m}`$, the localized morphism:
$$^p(D_{\underset{¯}{q}}^1)D_{\underset{¯}{q}}^p$$
is an isomorphism. Since $`D^1()`$ and $`D^p()`$ are local functors, if $`X=_{\alpha _H\underset{¯}{q}}H`$, we have $`D_{\underset{¯}{q}}^1=D^1(𝒜_X)_{\underset{¯}{q}}`$ and $`D_{\underset{¯}{q}}^p=D^p(𝒜_X)_{\underset{¯}{q}}`$.
Because $`\underset{¯}{q}\underset{¯}{m}`$, $`X(0)`$ and by hypothesis $`𝒜_X`$ is free. But for free arrangements $`D^p^p(D^1)`$ (see Solomon-Terao , Prop.3.4), which concludes the proof of the first assertion.
The proof of the last statement is similar. $`\mathrm{}`$
When $`\mathrm{char}k`$ does not divide $`|𝒜|`$, it is possible to characterize the local freeness of $`𝒜`$ using the Ext modules of the Jacobian ideal.
###### Proposition 3.10
For an arrangement $`𝒜`$ such that $`\mathrm{char}k|𝒜|`$, $`𝒜`$ is locally free if and only if the modules $`\mathrm{Ext}_S^i(S/J,S)`$ are supported only at the maximal ideal $`\underset{¯}{m}`$, for all $`i3`$.
Proof. Recall that we have the module $`D_0`$ defined by the exact sequence:
$$0D_0S^{n+1}J(d1)0,$$
where $`d=|𝒜|`$ and $`J`$ is the Jacobian ideal. Since $`\mathrm{char}kd`$, we have $`DD_0S(1)`$. Using Theorem 3.3 it follows that $`𝒜`$ is locally free if and only if $`\stackrel{~}{D_0}`$ is locally free.
For a finitely generated $`S`$-module $`M`$ it is known that the set
$$S(M)=\{\underset{¯}{q}\mathrm{Spec}(S)|M_{\underset{¯}{q}}\mathrm{is}\mathrm{not}\mathrm{a}\mathrm{free}S_{\underset{¯}{q}}\mathrm{module}\}$$
can be written as
$$S(M)=\underset{i1}{}\mathrm{Supp}\mathrm{Ext}_S^i(M,S)$$
(see, Hartshorne , p. 238, exercise 6.6).
Therefore $`\stackrel{~}{D_0}`$ is locally free if and only if $`\mathrm{Supp}\mathrm{Ext}_S^i(D_0,S)\{\underset{¯}{m}\}`$ for every $`i1`$. Since $`\mathrm{Ext}_S^i(S/J,S)\mathrm{Ext}_S^{i2}(D_0,S)`$ for every $`i3`$, the proof of the proposition is complete. $`\mathrm{}`$
## 4 Chern classes of vector bundles on $`𝐏^n`$
We consider a vector bundle $``$ of rank $`r`$ on $`𝐏^n`$. Motivated by the application in the context of arrangements which will be given in the next section, we introduce $`R(;t,X)𝐙[t]((X))/(t^{n+1})`$, defined by
$$R(;t,X)=(1)^rt^r(1X)^{n+1r}\underset{i=0}{\overset{r}{}}P(H_{}^0(^i);X)(\frac{X1}{t}1)^i.$$
Here $`H_{}^0(^i)`$ is the finitely generated graded module $`_{m𝐙}H^0(𝐏^n,^i(m))`$. Recall that for a finitely generated graded $`S`$-module $`M`$, $`P(M;X)`$ denotes the Hilbert series of $`M`$ which is a Laurent series, but also a rational function in $`X`$. The main result of this section is that $`R(;t,X)`$ can be used to compute the Chern polynomial of $``$. More precisely, we have the following
###### Theorem 4.1
If $``$ is a vector bundle on $`𝐏^n`$, then
$$\underset{X1}{lim}R(;t,X)=c_t().$$
###### Remark 4.2
We will see in the proof that in order to compute the above limit, in the definition of $`R(;t,x)`$ we may replace each $`H_{}^0(^i)`$ with a different finitely generated module $`M_i`$ such that $`^i\stackrel{~}{M_i}`$.
Before proving the theorem we give two lemmas.
###### Lemma 4.3
The assertion of Theorem 4.1 is true in the case of a split vector bundle $``$.
Proof. Suppose that $`𝒪(a_1)\mathrm{}𝒪(a_r)`$. In this case we have
$$R(;t,X)=(1)^rt^r(1X)^{n+1r}\underset{i=0}{\overset{r}{}}P(H_{}^0(_{1k_1<\mathrm{}<k_ir}𝒪(a_{k_1}+\mathrm{}+a_{k_r});X)(\frac{X1}{t}1)^i.$$
Since for every $`a𝐙`$, $`H_{}^0(𝒪(a))=S(a)`$ and $`P(S(a);X)=X^a(1X)^{n1}`$, we get
$$R(;t,X)=(1)^rt^r(1X)^{n+1r}\underset{i=0}{\overset{r}{}}\underset{1k_1<\mathrm{}<k_ir}{}X^{a_{k_1}\mathrm{}a_{k_i}}(1X)^{n1}(\frac{X1}{t}1)^i$$
$$=(1)^rt^r(1X)^r\underset{i=1}{\overset{r}{}}(1+X^{a_i}(\frac{X1}{t}1))=(1)^r\underset{i=1}{\overset{r}{}}(t\frac{1X^{a_i}}{1X}X^{a_i}).$$
It follows from this that the limit exists and
$$\underset{X1}{lim}R(;t,X)=(1)^r\underset{i=1}{\overset{r}{}}(a_it1)=\underset{i=1}{\overset{r}{}}(1+a_it)=c_t().$$
$`\mathrm{}`$
###### Lemma 4.4
If $`rn`$ is fixed and $`P𝐐[X_1,\mathrm{},X_n]`$ is a polynomial such that
$$P(c_1(),\mathrm{},c_n())=0,$$
for every split vector bundle of rank $`r`$, $`=𝒪(a_1)\mathrm{}𝒪(a_r)`$, then $`P=0`$.
Proof. If $`s_i`$, $`1in`$ is the $`i^{\mathrm{th}}`$ symmetric polynomial in $`a_1,\mathrm{},a_r`$, then $`c_i()=s_i(a_1,\mathrm{},a_r)`$, for every split vector bundle $``$ as in the hypothesis. Therefore we get $`P(s_1,\mathrm{},s_n)=0`$.
But the morphism
$$\varphi :𝐐[X_1,\mathrm{},X_n]𝐐[Y_1,\mathrm{},Y_r]$$
given by $`\varphi (X_i)=s_i(Y_1,\mathrm{},Y_r)`$ is the restriction to $`𝐐[X_1,\mathrm{},X_n]`$ of the monomorphism
$$\psi :𝐐[X_1,\mathrm{},X_r]𝐐[Y_1,\mathrm{},Y_r],$$
where $`\psi (X_i)=s_i(Y)`$, for $`1ir`$. Therefore we have $`P=0`$. $`\mathrm{}`$
We are now ready to give the proof of the theorem.
Proof. It is easy to check that if $`^{}=𝒪`$, then $`R(^{};t,X)=R(;t,X)`$. As we have also $`c_t(^{})=c_t()`$, it follows that by taking the direct sum with a large enough number of trivial bundles, we may suppose that $`\mathrm{rank}=rn`$.
In this case, using the above two lemmas, we see that in order to prove the theorem it is enough to show that the existence of the limit and its value can be expressed in terms of some polynomial identities with rational coefficients in the Chern classes of $``$. We have
$$R(;t,X)=(1)^rt^r(1X)^{n+1r}\underset{i=0}{\overset{r}{}}P(H_{}^0(^i);X)\underset{j=0}{\overset{i}{}}(1)^{ij}\left(\genfrac{}{}{0pt}{}{i}{j}\right)(\frac{X1}{t})^j$$
$$=(1)^r\underset{j=0}{\overset{r}{}}t^{rj}(1X)^{n+1r+j}\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)P(H_{}^0(^i);X).$$
Since for $`0ir`$, $`H_{}^0(^i)`$ is a $`S`$-module of dimension $`n+1`$, we can write
$$P(H_{}^0(^i());X)=\frac{Q_i(X)}{(1X)^{n+1}},$$
where $`Q𝐙[X,X^1]`$. Moreover, if we consider the Taylor expansion of $`Q_i(X)`$ around $`X=1`$:
$$Q_i(X)=\underset{l0}{}e_l^{(i)}(X1)^l,$$
then the first $`n+1`$ coefficients of this expansion can be recovered from the Hilbert polynomial of $`H_{}^0(^i)`$, which can be written as:
$$T_i(X)=\underset{l=0}{\overset{n}{}}(1)^{nl}e_{nl}^{(i)}\left(\genfrac{}{}{0pt}{}{X+l}{l}\right).$$
For these results, see for example Bruns and Herzog , Chapter 4.1.
Therefore we have $`e_k^{(i)}=(1)^k\mathrm{\Delta }^{nk}T_i(0)`$. For every graded $`S`$-module $`M`$, its Hilbert polynomial is given by the formula $`T(m)=\chi (\stackrel{~}{M}(m))`$, for every $`m𝐙`$. Using the short exact sequences corresponding to successive hyperplane sections we get
$$e_k^{(i)}=(1)^k\chi (^i|_{H_k}),$$
where $`H_k𝐏^n`$ is a linear subspace of dimension $`k`$, for $`0kn`$. We deduce
$$R(;t,X)=(1)^r\underset{j=rn}{\overset{r}{}}t^{rj}(1X)^{n+1r+j}\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)\frac{\underset{k0}{}(1)^ke_k^{(i)}(1X)^k}{(1X)^{n+1}}.$$
$$R(;t,X)=(1)^r\underset{j=rn}{\overset{r}{}}t^{rj}\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)\underset{k0}{}(1)^k\frac{e_k^{(i)}}{(1X)^{rjk}}.$$
By considering the coefficient of $`t^{rj}`$, for $`rnjr`$, the fact that $`lim_{X1}R(;t,X)`$ exists and is equal to $`c_t()`$ is equivalent to:
$$(1)\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)e_k^{(i)}=0,\mathrm{for}\mathrm{\hspace{0.17em}0}krj1,$$
$$(2)\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)e_{rj}^{(i)}=(1)^rc_{rj}(),$$
for every $`j`$ with $`rnjr`$.
Using the formulas we have for $`e_k^{(i)}`$, these relations become:
$$(1^{})\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)\chi (^i|_{H_k})=0,\mathrm{for}\mathrm{\hspace{0.17em}0}krj1$$
$$(2^{})\underset{i=j}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{j}\right)\chi (^i|_{H_{rj}})=(1)^rc_{rj}(),$$
for every $`j`$, with $`rnjr`$.
For future reference, notice that for $`j=rn`$, $`(2^{})`$ becomes
$$(2^{\prime \prime })\underset{i=rn}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{rn}\right)\chi (^i)=(1)^rc_n().$$
In order to finish the proof of the theorem, it is enough to notice that using Hirzebruch-Riemann-Roch theorem (see Fulton Corollary 15.2.1) all the Euler-Poincaré characteristics can be expressed as polynomials with rational coefficients in the Chern classes of the exterior powers $`^i`$. Indeed, since the Chern classes of these exterior powers can be computed as polynomials in the Chern classes of $``$, we can apply Lemma 4.4 and Lemma 4.3 to conclude the proof of the theorem. $`\mathrm{}`$
###### Corollary 4.5
If $``$ is a vector bundle on $`𝐏^n`$ with $`rank=rn`$, then we have the following formula for the top Chern class:
$$\underset{i=rn}{\overset{r}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{rn}\right)\chi (^i)=(1)^rc_n().$$
In particular, if $`r=n`$ we have
$$\underset{i=0}{\overset{n}{}}(1)^i\chi (^i)=(1)^nc_n().$$
Proof. This is just the identity $`(2^{\prime \prime })`$ in the proof of Theorem 4.1. $`\mathrm{}`$
###### Remark 4.6
If $``$ is a vector bundle on $`𝐏^n`$ such that $`^{}`$ has $`rn+1`$ sections $`\sigma _1,\mathrm{},\sigma _{rn+1}`$ such that the degeneration locus $`\mathrm{\Gamma }`$ is zero dimensional and the degeneration locus $`\mathrm{\Gamma }^{}`$ of $`\sigma _1,\mathrm{},\sigma _{rn}`$ is empty, then the formula in Corollary 4.5 is equivalent to the formula giving the degree of $`\mathrm{\Gamma }`$.
Indeed, $`\sigma _1,\mathrm{},\sigma _{rn+1}`$ define a morphism $`=𝒪_{𝐏^n}^{rn+1}`$ which gives an Eagon-Northcott type complex
$$0^r(S_n)^{}\mathrm{}^{rn+1}^{}^{rn}.$$
It follows from Eisenbud , Theorems A.2.10 and A.2.14 that since $`dim\mathrm{\Gamma }=0`$ and $`\mathrm{\Gamma }^{}=\mathrm{}`$, this complex is exact and moreover, it gives a resolution of $`𝒪_\mathrm{\Gamma }`$. Therefore we obtain
$$\underset{i=rn}{\overset{r}{}}(1)^i\chi (^i)=(1)^{rn}\chi (𝒪_\mathrm{\Gamma })=(1)^{rn}\mathrm{deg}\mathrm{\Gamma }.$$
On the other hand, the Thom-Porteous formula (see Fulton , Theorem 14.4) says that under our hypothesis $`\mathrm{deg}\mathrm{\Gamma }=(1)^nc_n()`$ and we get the formula in Corollary 4.5.
Let $``$ be a vector bundle on $`𝐏^n`$, with $`\mathrm{rank}=n`$. For every $`i`$, we denote by $`Q(^i;X)`$ the Hilbert polynomial of $`^i`$.
###### Corollary 4.7
With the above notation, we have
$$\underset{i=0}{\overset{n}{}}(1)^iQ(^i();iX)=(1)^n\underset{i=0}{\overset{n}{}}c_i()X^{ni}.$$
Proof. We prove that the two polynomials take the same value for every $`a𝐙`$. Indeed, if in Corollary 4.5 we replace $``$ by $`(a)`$, then we have
$$\chi (^i((a)))=\chi (^i)(ai)=Q(^i;ai),$$
while $`c_n((a))=_{i=0}^nc_i()a^{ni}`$ (see Fulton , Remark 3.2.3). $`\mathrm{}`$
## 5 The Characteristic Polynomial
In this section we will apply the general results we have obtained so far to the case of the vector bundle associated to the module of $`𝒜`$ derivations of a locally free arrangement $`𝒜`$. Recall that $`\pi (𝒜,t)`$ is a polynomial of degree $`n+1`$. We will denote by $`\overline{\pi }(𝒜,t)`$ its class in $`𝐙[t]/(t^{n+1})`$. The main result is the following.
###### Theorem 5.1
If $`𝒜`$ is a locally free arrangement, then
$$\overline{\pi }(𝒜,t)=c_t(\stackrel{~}{\mathrm{\Omega }^1}).$$
Proof. From the basic relation between the Poincaré and the characteristic polynomial we get
$$\pi (𝒜,t)=(t)^{n+1}\chi (𝒜,t^1).$$
Combining this with Theorem 2.4 we obtain
$$\pi (𝒜,t)=\underset{X1}{lim}t^{n+1}\underset{i0}{}P(D^i;X)(\frac{1}{t}(X1)1)^i.$$
By Proposition 3.9 we have $`\stackrel{~}{D^i}=^i\stackrel{~}{D^1}`$; so from Remark 4.2 it follows that
$$\underset{X1}{lim}R(\stackrel{~}{D^1};t,X)=\underset{X1}{lim}\underset{i=0}{\overset{n}{}}P(D^i;X)(\frac{1}{t}(X1)1).$$
By Theorem 4.1 this limit is equal to
$$c_t(\stackrel{~}{D^1})=c_t(\stackrel{~}{\mathrm{\Omega }^1}).$$
We therefore obtain
$$\overline{\pi }(𝒜,t)=c_t(\stackrel{~}{\mathrm{\Omega }^1}).$$
$`\mathrm{}`$
###### Remark 5.2
Since it is known that $`\pi (𝒜;1)=0`$ (see Orlik and Terao , Proposition 2.5.1), in order to know $`\pi (𝒜,t)`$ it is enough to know $`\overline{\pi }(𝒜,t)`$.
In fact, when $`\mathrm{char}k`$ does not divide $`|𝒜|`$, then from Theorem 5.1 we can deduce a formula for $`\pi (𝒜,t)`$ involving the vector bundle $`\stackrel{~}{\mathrm{\Omega }_0^1}`$.
###### Corollary 5.3
If $`𝒜`$ is a locally free arrangement such that $`\mathrm{char}k`$ does not divide $`|𝒜|`$, then we have
$$\pi (𝒜,t)=(1+t)c_t(\stackrel{~}{\mathrm{\Omega }_0^1}).$$
Proof. Since $`\pi (𝒜,t)/(1+t)`$ is a polynomial of degree $`n`$, it follows that it is enough to prove that its class in $`𝐙[t]/(t^{n+1})`$ is equal to $`c_t(\stackrel{~}{\mathrm{\Omega }_0^1})`$. But by Theorem 5.1 it follows that this class is equal to
$$\overline{\pi }(𝒜,t)/(1+t)=c_t(\stackrel{~}{\mathrm{\Omega }^1})/(1+t)=c_t(\stackrel{~}{\mathrm{\Omega }_0^1}),$$
since $`\stackrel{~}{\mathrm{\Omega }^1}\stackrel{~}{\mathrm{\Omega }_0^1}𝒪(1)`$. $`\mathrm{}`$ In , Yuzvinsky proves that for a locally free arrangement, the Hilbert polynomial of the module $`D^1(𝒜)`$ depends only on the lattice. The following corollary makes this more precise.
###### Corollary 5.4
If $`𝒜`$ is a locally free arrangement, then giving the Hilbert polynomial of $`D^1(𝒜)`$ is equivalent to giving the Poincaré polynomial $`\pi (𝒜,t)`$ of the arrangement.
Proof. The statement follows from Theorem 5.1 and the Hirzebruch-Riemann-Roch theorem (see Fulton , Corollary 15.2.1). $`\mathrm{}`$
###### Example 5.5
Let $`𝒜`$ be the arrangement in $`𝐏^3`$ defined by the vanishing of the fifteen linear forms $`a_0x_0+a_1x_1+a_2x_2+a_3x_3`$, where $`a_i`$ is either zero or one. This example was constructed by Edelman-Reiner () as a counterexample to a conjecture of Orlik; $`\pi (𝒜,t)=1+15t+80t^2+170t^3+104t^4`$. There are 45 rank three elements of $`L_𝒜`$; we consider the corresponding subarrangements as essential arrangements in $`𝐏^2`$. The subarrangements are of three distinct types, described below ($`\mu (L_2)`$ is the Möbius function of the rank two elements of $`L_{𝒜_X}`$) :
$`\begin{array}{ccc}& & \\ 20& \text{subarrangements on 3 hyperplanes, }& \mu (L_2)=(1,1,1)\\ & & \\ 15& \text{subarrangements on 5 hyperplanes, }& \mu (L_2)=(2,2,1,1,1,1)\\ & & \\ 10& \text{subarrangements on 7 hyperplanes, }& \mu (L_2)=(2,2,2,2,2,2,1,1,1)\end{array}`$
It is easy to check that all of these subarrangements are free, so $`𝒜`$ is locally free. We illustrate Corollary 5.4 for this example. For a rank three bundle $``$ on $`𝐏^3`$ we have:
$$\text{ch}((m))\text{td}(𝒯_{𝐏^3})=\frac{1}{2}m^3+(3+\frac{c_1}{2})m^2+(\frac{11}{2}+2c_1+\frac{c_1^2}{2}c_2)m+3+\frac{11c_1}{6}+c_1^22c_2+\frac{c_1^3}{3}\frac{c_1c_2}{2}+\frac{c_3}{2}.$$
Since we know $`\pi (𝒜,t)`$, we may apply Corollary 5.3 to obtain the Chern classes $`c_i`$, and from Hirzebruch-Riemann-Roch we obtain the Hilbert polynomial:
$$\chi (D_0^1(m))=\frac{1}{2}m^34m^2+\frac{57}{6}m6.$$
The point is that for a locally free arrangement, knowing the combinatorial data (i.e. Poincaré polynomial) means knowing the Hilbert polynomial, which can make the computation of the free resolution much faster; for this example the free resolution is:
$$0S(6)S^4(5)D_0^10.$$
Thus, we see that $`\stackrel{~}{\mathrm{\Omega }_0^1}\mathrm{\Omega }_{𝐏^\mathrm{𝟑}}(6)`$, and $`c_t(\stackrel{~}{\mathrm{\Omega }_0^1})=\frac{(1+5t)^4}{(1+6t)}\text{ mod }t^4=1+14t+66t^2+104t^3`$, as expected.
Remark In , Dolgachev and Kapranov point out that for a generic arrangement, $`\mathrm{\Omega }^1(\text{log }D)`$ is a Steiner bundle (hence stable); in the previous example $`\stackrel{~}{D_0^1}`$ is Steiner. In $`𝐏^2`$ there are many examples of non-generic arrangements for which $`\stackrel{~}{\mathrm{\Omega }_0^1}`$ is indecomposable but not semistable.
## 6 Minimal free resolutions for the modules of logarithmic forms
In this section we give a minimal free resolution for the modules $`\mathrm{\Omega }^p(𝒜)`$ of logarithmic forms in the case of a locally free arrangement $`𝒜`$ with $`\mathrm{pdim}\mathrm{\Omega }^1(𝒜)=1`$. This generalizes the results of Rose and Terao in the case of generic arrangements. The same idea can be used to give a minimal free resolution for the modules $`D^p(𝒜)`$ when $`\mathrm{pdim}D^1(𝒜)=1`$.
We first recall the definition of syzygy modules:
###### Definition 6.1
A module $`M`$ is a $`k^{\mathrm{𝑡ℎ}}\mathrm{𝑠𝑦𝑧𝑦𝑔𝑦}`$ if there exists an exact sequence
$$0MF_1F_2\mathrm{}F_k,$$
with $`F_i`$ free.
###### Lemma 6.2
If $`\mathrm{pdim}M=1`$ and $`Ext^1(M,S)_{\underset{¯}{q}}=0`$ for every prime ideal $`\underset{¯}{q}\underset{¯}{m}`$, then $`M`$ is an $`(n1)^{st}`$ syzygy.
Proof.
From the short exact sequence:
$$0F_1F_0M0,$$
we obtain an exact sequence:
$$0M^{}F_0^{}F_1^{}Ext^1(M,S)0.$$
Since $`Ext^1(M,S)`$ is supported only at $`m`$, $`Ext^1(M,S)`$ is a module of finite length, so has a free resolution of length $`n+1`$; hence $`M^{}`$ has projective dimension $`n1`$. Dualizing once more and using the fact that $`Ext^i(Ext^1(M,S),S)`$ is zero for $`in+1`$, we find that $`M`$ is an $`(n1)^{st}`$ syzygy. $`\mathrm{}`$
Suppose now that $`M`$ is a finitely generated $`S`$-module. Note that if $`\stackrel{~}{M}`$ is locally free, then the second condition in Lemma 6.2 is automatically satisfied. In Lebelt shows that if $`\mathrm{pdim}M=i`$ and if $`M`$ is an $`i(p1)^{st}`$ syzygy, then one can obtain a free resolution for $`\mathrm{\Lambda }^pM`$ from a free resolution of $`M`$. In the situation considered above, i.e. $`i=1`$ and $`M`$ is an $`(n1)^{st}`$ syzygy, we obtain a minimal free resolution of $`\mathrm{\Lambda }^p(M)`$, for $`pn1`$, which is an Eagon-Northcott type complex. Namely, if
$$0F_1F_0M0$$
is a minimal free resolution of $`M`$, then a minimal free resolution of $`^pM`$ is given by:
$$(F_{}^{(p)})(M):0D_pF_1D_{p1}F_1F_0D_{p2}F_1\mathrm{\Lambda }^2F_0\mathrm{}\mathrm{\Lambda }^pF_0\mathrm{\Lambda }^pM0,$$
where $`D_iF_1=(S_i(F_1^{}))^{}`$ denotes the $`i^{th}`$ divided power of $`F_1`$. If $`\mathrm{char}k=0`$, then $`D_iF_1S_iF_1`$ is the usual symmetric power of $`F_1`$.
We can now give the main result of this section:
###### Theorem 6.3
If $`𝒜`$ is a locally free arrangement and $`\mathrm{pdim}\mathrm{\Omega }^1=1`$, then the natural morphism
$$^p\mathrm{\Omega }^1\mathrm{\Omega }^p$$
is an isomorphism, and $`(F_{}^{(p)}(\mathrm{\Omega }^1))`$ gives a minimal free resolution of $`\mathrm{\Omega }^p`$, for every $`p`$, $`pn1`$.
Proof. From Lemma 6.2 and Lebelt’s result cited above, it follows that for every $`p`$, $`pn1`$, $`F_{}^{(p)}(\mathrm{\Omega }^1)`$ is a (minimal) free resolution of $`^p\mathrm{\Omega }^1`$. In particular, we have $`\mathrm{pdim}^p\mathrm{\Omega }^1=p`$. By the Auslander-Buchsbaum formula we obtain $`\mathrm{depth}^p\mathrm{\Omega }^1=n+1p2`$.
We consider the commutative diagram:
$`\begin{array}{ccc}^p\mathrm{\Omega }^1& \stackrel{\alpha }{}& \mathrm{\Omega }^p\\ & & \\ H_{}^0(^p\stackrel{~}{\mathrm{\Omega }^1})& \stackrel{\beta }{}& H_{}^0(\stackrel{~}{\mathrm{\Omega }^p})\end{array}`$
Proposition 3.9 implies that $`\beta `$ is an isomorphism. By Lemma 3.1, $`\mathrm{depth}\mathrm{\Omega }^p2`$ and we also have $`\mathrm{depth}^p\mathrm{\Omega }^12`$ and therefore both the vertical maps in the diagram are isomorphisms. We conclude that $`\alpha `$ is an isomorphism. $`\mathrm{}`$
###### Corollary 6.4
If $`𝒜`$ is a locally free arrangement with $`\mathrm{pdim}\mathrm{\Omega }^1=1`$, then we have
$$\mathrm{pdim}\mathrm{\Omega }^p=p,\mathrm{for}pn1,$$
$$\mathrm{pdim}\mathrm{\Omega }^{n+1}=0.$$
Moreover, if $`\mathrm{char}k`$ does not divide $`|𝒜|`$, then we have also
$$\mathrm{pdim}\mathrm{\Omega }^n=n1,$$
and therefore $`\mathrm{pdim}D^1=n1`$.
Proof. The first assertion follows from Theorem 6.3, while the second one is true for an arbitrary arrangement (see Orlik and Terao , Proposition 4.68).
If $`\mathrm{char}k|𝒜|`$, then we have $`\mathrm{\Omega }^1=\mathrm{\Omega }_0^1S(1)`$. Therefore,
$$^p\mathrm{\Omega }^1=^p\mathrm{\Omega }_0^1^{p1}\mathrm{\Omega }_0^1(1),$$
for every $`p`$. It follows immediately that $`\mathrm{pdim}^p\mathrm{\Omega }_0=p`$, for $`pn1`$. In particular, from the Auslander-Buchsbaum formula we get $`\mathrm{depth}^{n1}\mathrm{\Omega }_0^1(1)=2`$.
From this and the decomposition $`^n\mathrm{\Omega }^1=^n\mathrm{\Omega }_0^1^{n1}\mathrm{\Omega }_0^1(1)`$ it follows that $`^{n1}\mathrm{\Omega }_0^1(1)`$ is a direct summand of $`H_{}^0(\stackrel{~}{^n\mathrm{\Omega }})`$.
By Lemma 3.1 we have $`\mathrm{depth}\mathrm{\Omega }^n2`$. This implies that $`\mathrm{\Omega }^nH_{}^0(\stackrel{~}{\mathrm{\Omega }^n})H_{}^0(^n\stackrel{~}{\mathrm{\Omega }}^1)`$. From the fact that this module has a summand of depth two and has depth at least two, we conclude that the depth is exactly two, and the result follows from one more application of the Auslander-Buchsbaum formula.
The fact that $`\mathrm{pdim}D^1=n1`$ is a consequence of the general fact that $`\mathrm{\Omega }^nD^1(d)`$, where $`d=|𝒜|`$ (see Rose and Terao , Lemma 4.4.1). $`\mathrm{}`$
A generic arrangement $`𝒜`$ is locally free since for every $`XL_𝒜`$, with $`\mathrm{rank}X<\mathrm{dim}V`$, $`𝒜_X`$ is isomorphic to the product between a Boolean arrangement and an empty arrangement. On the other hand, Ziegler , Corollary 7.7 shows that in this case $`\mathrm{pdim}\mathrm{\Omega }^1=1`$ and therefore Theorem 6.3 and Corollary 6.4 apply in this case and give the results in Rose and Terao .
Analogous results with similar proofs hold if we replace the modules $`\mathrm{\Omega }^p`$ with the modules $`D^p`$. Namely, we have
###### Theorem 6.5
If $`𝒜`$ is a locally free arrangement and $`\mathrm{pdim}D^1=1`$, then the natural morphism
$$^pD^1D^p$$
is an isomorphism and $`(F_{}^{(p)}(D^1))`$ is a minimal free resolution of $`D^p`$, for every $`p`$, $`pn1`$. We have
$$\mathrm{pdim}D^p=p,\mathrm{for}pn1,$$
$$\mathrm{pdim}D^{n+1}=0,$$
and if $`\mathrm{char}k`$ does not divide $`|𝒜|`$, then
$$\mathrm{pdim}D^n=n1.$$ |
warning/0001/math0001133.html | ar5iv | text | # Domains of holomorphy with edges and lower dimensional boundary singularities
## 1. Introduction
The positive solution to the classical Levi problem due to Oka \[O42, O53\], Bremermann \[B54\] and Norguet \[No54\] asserts that a domain $`\mathrm{\Omega }^N`$ whose boundary is of class $`C^2`$ is a domain of holomorphy provided the Levi form of the boundary is everywhere positively semidefinite (see e.g. surveys \[S84, Pe94\]). In contrast to this, for domains with singularities on the boundary there seems to be a lack of such geometric conditions in the literature.
### 1.1. Piecewise smooth domains
A natural generalization of smooth domains is given by the class of so-called piecewise smooth domains whose boundaries are pieces of hypersurfaces satisfying suitable transversality conditions (see e.g. \[SH81, Pi82, Na88, K92, F93, MP94\]). However, for the domains of holomorphy, the class of piecewise smooth domains seems to be very restrictive. For instance, an envelope of holomorphy of a domain with real-analytic (even algebraic) boundary does not need to be piecewise smooth as the example
$$\mathrm{\Omega }:=\{(z,w)^2:|w|^2+(|z|^21)^2<2\}$$
shows. Indeed, the Cauchy formula argument implies that the envelope of holomorphy (and also convex, polynomially convex and rationally convex hulls) of $`\mathrm{\Omega }`$ is the union $`\mathrm{\Omega }\{|z|<1,|w|<\sqrt{2}\}`$.
In this paper we consider a larger class of domains $`\mathrm{\Omega }`$, whose smooth boundary pieces may not be extended to closed smooth hypersurfaces in a neighborhood of $`\mathrm{\Omega }`$. We also allow singular subsets in the boundary that are only controlled to have zero $`1`$-codimensional (with respect to the dimension of the boundary) Hausdorff measure.
### 1.2. The class $`L^{2,\mathrm{}}`$.
For an open subset $`U^m`$ ($`m1`$), denote by $`L^{2,\mathrm{}}(U)`$ the space of real continuous functions $`h`$ on $`U`$ that are twice continuously differentiable with bounded derivatives on an open dense subset of $`U`$.
We say that a domain $`\mathrm{\Omega }^n`$ is of class $`L^{2,\mathrm{}}`$ if for every $`a\mathrm{\Omega }`$ there exists a system of local ($`C^2`$-smooth) coordinates $`(x_1,\mathrm{},x_n)=(x^{},x_n)^{n1}\times `$ in a neighborhood $`U=U^{}\times I`$ of $`a`$ and a function $`hL^{2,\mathrm{}}(U^{})`$ such that
$$\mathrm{\Omega }U=\{(x^{},x_n)U:x_n>h(x^{})\}.$$
(1)
It is easy to see that every domain with piecewise $`C^2`$-smooth boundary is of class $`L^{2,\mathrm{}}`$.
### 1.3. Regular and edge points.
Given a subset $`A^n`$, we say that a point $`aA`$ is ($`C^2`$-)regular if $`AU_a`$ is a smooth hypersurface of class $`C^2`$ for some neighborhood $`U_a^n`$ of $`a`$. If $`aA`$ is not regular, we call it a ($`C^1`$-)edge point if there exists a neighborhood $`U_a`$ of $`a`$ and a connected closed $`(n2)`$-dimensional submanifold $`M_a\mathrm{\Omega }U_a`$ of class $`C^1`$, referred to as an edge at $`a`$, that contains all nonregular points of $`AU_a`$.
It $`\mathrm{\Omega }^N`$ is a domain of holomorphy, it is a standard fact that the Levi form (see §2) at every regular point is positively semidefinite. The classical example of Hartogs
$$\mathrm{\Omega }:=\{|z|<1,|w|<1/2\}\{1/2<|z|<1,|w|<1\}^2$$
shows that the converse does not hold in general even for piecewise smooth domains. The edge boundary points of $`\mathrm{\Omega }`$ in this example, where all holomorphic functions extend, are precisely those whose tangent cones are not convex. Recall that the tangent cone (in the sense of Whitney) of $`\mathrm{\Omega }^N`$ at a point $`a\mathrm{\Omega }`$, denoted by $`T_a\mathrm{\Omega }`$, is defined to be the set of all possible limits of $`t_k(a_ka)^N`$, where $`a_k\mathrm{\Omega }`$ and $`t_k_+`$ are sequences with $`a_ka`$ as $`k\mathrm{}`$.
### 1.4. Main results.
It turns out that, together with the Levi form condition for regular points of $`\mathrm{\Omega }`$, the cone convexity for edges points guarantees that $`\mathrm{\Omega }`$ is a domain of holomorphy. No condition on the other points is required provided the set of those points if of Hausdorff $`(2N2)`$-dimensional measure zero. More precisely, we have the following result.
###### Theorem 1.1.
Let $`\mathrm{\Omega }^N`$ be a domain of class $`L^{2,\mathrm{}}`$ and $`E\mathrm{\Omega }`$ be a closed subset of Hausdorff $`(2N2)`$-dimensional measure zero. Suppose that the following hold:
1. every point of $`\mathrm{\Omega }E`$ is either regular or an edge point;
2. at every regular point of $`\mathrm{\Omega }E`$ the Levi form is positively semidefinite;
3. at every edge point of $`\mathrm{\Omega }E`$ the tangent cone of $`\mathrm{\Omega }`$ is convex.
Then $`\mathrm{\Omega }`$ is a domain of holomorphy.
By the well-known fact, $`\mathrm{\Omega }`$ is a domain of holomorphy if and only if the function $`\psi (z):=\mathrm{log}d(z,\mathrm{\Omega })`$ is plurisubharmonic in $`\mathrm{\Omega }`$, where $`d`$ denotes the euclidean distance (see e.g. Theorems 2.6.5 and 4.2.8 in \[H90\]). In particular, it follows from Theorem 1.1 for $`N2`$ that the tangent cone $`C_a`$ of $`\mathrm{\Omega }`$ at any $`a\mathrm{\Omega }`$ cannot be strictly concave (i.e. the interior of $`C_a`$ cannot contain a hyperplane). Indeed, otherwise the function $`\psi `$ would be equal to $`\mathrm{log}za`$ in an open subset of $`\mathrm{\Omega }`$ and hence would not be plurisubharmonic. This shows, on the other hand, that $`\psi `$ cannot be directly used to prove Theorem 1.1 because in Theorem 1.1 there is no convexity condition on the cone at the points from the “exceptional” subset $`E\mathrm{\Omega }`$. In fact we prove the plurisubharmonicity of the function
$$\varphi (z):=\mathrm{log}(x_{2N}h(x^{}))+\lambda z^2$$
(2)
near the boundary for some $`\lambda >0`$ rather than of $`\psi `$, where $`(x_1,\mathrm{},x_{2N})`$ and $`h`$ satisfy (1).
The necessity of the convexity condition (iii) depends on the complex geometry of the edges. We show:
###### Proposition 1.2.
Let $`\mathrm{\Omega }^N`$ be a domain of holomorphy and suppose that for an edge point $`a\mathrm{\Omega }`$, there exists an edge $`M_a`$ which is not a complex hypersurface in any neighborhood of $`a`$. Then the tangent cone of $`\mathrm{\Omega }`$ at $`a`$ is convex.
On the other hand, if an edge can be chosen to be a complex hypersurface, the convexity condition (iii) does not need to hold as the example of $`\mathrm{\Omega }:=D\times ^2`$ shows with $`D`$ a nonconvex polygon. Therefore we have to distinguish between edge points satisfying the assumptions of Proposition 1.2 that we call real edge points and other edge points $`a\mathrm{\Omega }`$, where any edge must be locally a complex hypersurface. In the second case $`a`$ is said to be a complex edge point. Then we impose the convexity condition only at real edge points. In this more general situation the above function $`\varphi `$ given by (2) is not always plurisubharmonic. Nevertheless, we obtain the following necessary and sufficient geometric conditions for domains to be domains of holomorphy as a consequence of Theorem 1.1 and Proposition 1.2.
###### Corollary 1.3.
Let $`\mathrm{\Omega }^N`$ be a domain of class $`L^{2,\mathrm{}}`$ and $`E\mathrm{\Omega }`$ be a closed subset of Hausdorff $`(2N2)`$-dimensional measure zero. Suppose that the following hold:
1. every point in $`\mathrm{\Omega }E`$ is either regular or an edge point;
2. for every $`a\mathrm{\Omega }`$ there exist a neighborhood $`U_a`$ and a complex hypersurface $`N\mathrm{\Omega }`$ that contains all complex edge points in $`\mathrm{\Omega }U_a`$.
Then $`\mathrm{\Omega }`$ is a domain of holomorphy if and only if the Levi form at every regular point $`a\mathrm{\Omega }E`$ is positively semidefinite and the tangent cone at every real edge point $`a\mathrm{\Omega }E`$ is convex.
Finally we would like to mention that the statements of Theorem 1.1 (and of Corollary 1.3) also hold for relatively compact domains in Stein manifolds. Indeed, in this case Theorem 1.1 implies that the domain is locally Stein. Hence it is a domain of holomorphy by a result of Fornaess and Narasimhan (\[FN80\], Theorem 3.1.1).
## 2. The Levi form and plurisubharmonicity
Recall that the Levi form at a point $`a`$ of a real function $`\rho `$ of class $`C^2`$ in an open subset of $`^N`$ is the Hermitian form defined in local holomorphic coordinates $`z=(z_1,\mathrm{},z_N)`$ by
$$L\rho (\xi ,\eta )=L\rho (a)(\xi ,\eta ):=\underset{k,l}{}\frac{^2\rho }{z^k\overline{z^l}}(a)\xi ^k\overline{\eta ^l},\xi ,\eta ^N.$$
We write
$$\rho (a)^{}:=\{\xi ^N:\rho (a)(\xi )=0\}.$$
The Levi form of a domain $`\mathrm{\Omega }`$ at a regular point $`a\mathrm{\Omega }`$ is the restriction $`L\rho |_\rho ^{}`$, where $`d\rho 0`$ and $`\mathrm{\Omega }`$ is locally given by $`\rho <0`$. The norms are defined in the standard way:
$$\rho (a):=\underset{\xi =1}{sup}|\rho (a)(\xi )|,L\rho (a):=\underset{\xi =\eta =1}{sup}|L\rho (a)(\xi ,\eta )|.$$
###### Lemma 2.1.
Let $`\rho <0`$ be a negative function of class $`C^2`$ in an open subset of $`^N`$ and $`\lambda >0`$ be a constant such that the following holds:
1. $`L\rho |_\rho ^{}`$ is positive definite;
2. $`L\rho ^2\lambda (\rho ^2+\rho L\rho )`$.
Then the function $`\varphi (z):=\mathrm{log}(\rho (z))+\lambda z^2`$ is plurisubharmonic.
###### Proof.
We have
$$L\varphi (\xi ,\eta )=\rho ^2\rho (\xi )\overline{\rho (\eta )}\rho ^1L\rho (\xi ,\eta )+\lambda \xi ,\eta ,$$
(3)
where $`\xi ,\eta :=\xi _1\overline{\eta }_1+\mathrm{}+\xi _N\overline{\eta }_N`$. Every vector $`\zeta ^N`$ can be written as $`\zeta =\zeta _1+\zeta _2`$ with
$$|\rho (a)(\zeta _1)|=\rho (a)\zeta _1\text{ and }\rho (a)(\zeta _2)=0.$$
(4)
Applying (3) to $`\xi =\eta =\alpha _1\zeta _1+\alpha _2\zeta _2`$ with $`\zeta _1,\zeta _2`$ satisfying (4) we obtain
$$L\varphi (a)(\alpha _1\zeta _1+\alpha _2\zeta _2,\alpha _1\zeta _1+\alpha _2\zeta _2)=(\alpha _1,\alpha _2)(A+B)\left(\begin{array}{c}\overline{\alpha }_1\\ \overline{\alpha }_2\end{array}\right),$$
where
$$A=\left(\begin{array}{cc}\rho ^2\rho ^2\zeta _1^2\rho ^1L\rho (\zeta _1,\zeta _1)& \rho ^1L\rho (\zeta _1,\zeta _2)\\ \rho ^1L\rho (\zeta _2,\zeta _1)& \lambda \zeta _2^2\end{array}\right)$$
and
$$B=\left(\begin{array}{cc}\lambda \zeta _1^2& 0\\ 0& \rho ^1L\rho (\zeta _2,\zeta _2)\end{array}\right).$$
The matrix $`B`$ is positively semidefinite by (i). It is sufficient to show that $`A`$ is also positively semidefinite, i.e. $`detA0`$ by Sylvester’s criterion. But this follows from (ii):
$$detA\rho ^2(\lambda \rho ^2+\lambda \rho L\rho L\rho ^2)\zeta _1^2\zeta _2^20.$$
## 3. Piecewise plurisubharmonicity
### 3.1. One-dimensional case.
In the following let $`I`$ denote an open interval and $`AI`$ a finite subset. If $`f`$ is continuously differentiable with bounded derivative on $`IA`$, then for every $`aS`$ there exist one-sided limits
$$f(a)_{}:=\underset{xa,x<a}{lim}f(x)\text{ and }f(a)_+:=\underset{xa,x>a}{lim}f(x).$$
By elementary calculus we have
###### Lemma 3.1.
Let $`AI`$ and $`f`$ be as before and suppose that $`f`$ has compact support in $`I`$. Then
$$_If^{}𝑑x+\underset{aA}{}(f(a)_+f(a)_{})=0.$$
###### Corollary 3.2.
Let $`AI`$ be as before, $`\varphi C^0(I)C^2(IA)L^{2,\mathrm{}}(I)`$ be arbitrary and $`\alpha C^2(I)`$ have compact support in $`I`$. Then
$$_I(\varphi \alpha ^{\prime \prime }\varphi ^{\prime \prime }\alpha )𝑑x=\underset{aA}{}(\varphi ^{}(a)_+\varphi ^{}(a)_{})\alpha (a).$$
(5)
The corollary is obtained by applying Lemma 3.1 to $`f:=\varphi \alpha ^{}\varphi ^{}\alpha `$.
### 3.2. Higher-dimensional case.
Now consider an open subset $`\mathrm{\Omega }^n`$ and let $`S\mathrm{\Omega }`$ be a (locally closed) hypersurface of class $`C^1`$. Given a point $`a_0S`$ we fix a neighborhood $`U\mathrm{\Omega }`$ of $`a_0`$ such that $`US`$ has exactly two connected components $`U_+`$ and $`U_{}`$.
###### Lemma 3.3.
Let $`\varphi C^2(US)L^{2,\mathrm{}}(U)`$ be a function in $`U`$. Then for every $`v^n`$ the directional derivatives $`D_v\varphi |_U_{}`$ and $`D_v\varphi |_{U_+}`$ extend Lipschitz-continuously to $`U_{}S`$ and $`U_+S`$ respectively with the one-sided limits
$$D_v\varphi (a)_{}:=\underset{xa,xU_{}}{lim}D_v\varphi (x)\text{ and }D_v\varphi (a)_+:=\underset{xa,xU_+}{lim}D_v\varphi (x)$$
(6)
for $`aS`$. Moreover, if $`\varphi `$ is in addition continuous on $`U`$, one has $`D_v\varphi (a)_{}=D_v\varphi (a)_+`$ whenever $`v`$ is tangent to $`S`$ at $`a`$. In particular, the sign of the expression
$$D_v\varphi (a)_+D_v\varphi (a)_{}$$
(7)
is independent of the choice of a transversal vector $`v`$ pointing into $`U_+`$.
###### Proof.
Boundedness of the first and second derivatives of $`\varphi `$ on $`U_{}`$ and $`U_+`$ implies the one-sided Lipschitz extendibility of $`\varphi `$ and its first derivatives. In particular, if $`\varphi `$ is continuous on $`U`$, then the restriction $`\varphi |_S`$ coincides with both one-sided limits. Hence for $`v`$ tangent to $`S`$, one has $`D_v\varphi _+=D_v(\varphi |_S)=D_v\varphi _{}`$ as required. ∎
We observe that, if we interchange $`U_{}`$ with $`U_+`$, the sign of (7) remains the same, because also $`v`$ (pointing into $`U_+`$) changes the sign. This consideration motivates the following definition.
###### Definition 3.4.
A function $`\varphi C^2(US)L^{2,\mathrm{}}(U)`$, continuous in $`U`$, is said to be transversally convex at $`S`$ if the expression (7) is non-negative for any $`aS`$ and any $`v`$ pointing into $`U_+`$.
In the following we write $`^m`$ for the Hausdorff $`m`$-dimensional measure.
###### Proposition 3.5.
Let $`\mathrm{\Omega }^n`$ be open, $`G`$ be closed in $`\mathrm{\Omega }`$ with $`^{n1}(G)=0`$, $`S\mathrm{\Omega }`$ be a hypersurface of class $`C^1`$ with $`\overline{S}SG`$ and
$$\varphi C^0(\mathrm{\Omega })C^2(\mathrm{\Omega }(SG))L^{2,\mathrm{}}(\mathrm{\Omega })$$
be a function which is transversally convex at $`S`$. Then for every non-negative function $`\alpha C^2(\mathrm{\Omega })`$ with compact support in $`\mathrm{\Omega }`$ the inequality
$$_\mathrm{\Omega }\varphi \frac{^2\alpha }{x_j^2}𝑑x_\mathrm{\Omega }\frac{^2\varphi }{x_j^2}\alpha 𝑑x$$
(8)
holds for every $`j=1,\mathrm{},n`$.
Remark. Since all functions in (8) are measurable and bounded, both integrals exist with respect to the Lebesgue measure on $`^n`$.
###### Proof.
Without loss of generality we may assume $`j=n`$ and, by taking a suitable partition of unity,
$$\mathrm{supp}(\alpha )U^{}\times I\mathrm{\Omega }$$
for an open subset $`U^{}^{n1}`$ and an interval $`I`$. We write $`x=(x^{},x_n)^{n1}\times `$. Since $`^{n1}(G)=0`$, we see from Fubini’s theorem that $`G(\{x^{}\}\times I)=\mathrm{}`$ for all $`x^{}`$ outside a zero measure subset $`G^{}U^{}`$. Furthermore, by Sard’s theorem applied to the projection $`S^{n1}`$, the vector $`v:=(0,1)`$ is not tangent to $`S`$ at the points of $`S(\{x^{}\}\times I)`$ for all $`x^{}`$ outside a zero measure subset $`G^{\prime \prime }U^{}`$. Then for $`x^{}G^{}G^{\prime \prime }`$ the set
$$A_x^{}:=\{x_nI:(x^{},x_n)(S\text{supp}(\alpha ))\}$$
is finite. Thus we can apply Corollary 3.2 to the restriction of $`\varphi `$ and $`\alpha `$ to $`\{x^{}\}\times I`$. According to Definition 3.4, the right-hand side in (5) is non-negative. Then the required inequality is obtained by integrating over $`U^{}(G^{}G^{\prime \prime })`$ the non-negative left-hand side of (5). ∎
###### Corollary 3.6.
Under the assumptions of Proposition 3.5 suppose that $`\mathrm{\Omega }^N`$ and $`\varphi `$ is plurisubharmonic in $`\mathrm{\Omega }(SG)`$. Then $`\varphi `$ is plurisubharmonic in the whole $`\mathrm{\Omega }`$.
###### Proof.
Given a vector $`\xi ^N`$ we can find linear complex coordinates $`z_k=x_k+iy_k`$ ($`k=1,\mathrm{},N`$) such that $`\xi =(0,\mathrm{},0,1)`$. Then
$$L\varphi (a)(\xi ,\xi )=\frac{^2\varphi }{z_N\overline{z_N}}(a)\xi ^N\overline{\xi ^N}=\frac{1}{2}\left(\frac{^2\varphi }{x_N^2}(a)+\frac{^2\varphi }{y_N^2}(a)\right)0$$
(9)
for all $`a\mathrm{\Omega }(SG)`$ by the plurisubharmonicity of $`\varphi `$ there. Since $`SG\mathrm{\Omega }`$ is of zero measure, (9) implies
$$_\mathrm{\Omega }\frac{^2\varphi }{x_N^2}\alpha 𝑑z\overline{dz}+_\mathrm{\Omega }\frac{^2\varphi }{y_N^2}\alpha 𝑑z\overline{dz}0$$
(10)
for every non-negative function $`\alpha C^2(\mathrm{\Omega })`$ with compact support in $`\mathrm{\Omega }`$. By Proposition 3.5,
$$2_\mathrm{\Omega }\varphi L\alpha (a)(\xi ,\xi )𝑑z\overline{dz}=_\mathrm{\Omega }\varphi \frac{^2\alpha }{x_N^2}𝑑\mathrm{\Omega }+_\mathrm{\Omega }\varphi \frac{^2\alpha }{y_N^2}𝑑\mathrm{\Omega }0.$$
(11)
Since $`\xi `$ is arbitrary, we conclude that $`\varphi `$ has a non-negative Levi form in distributional sense. By the continuity of $`\varphi `$, the last fact is equivalent to the plurisubharmonicity of $`\varphi `$ on the whole $`\mathrm{\Omega }`$ (see \[H90\], Theorem 1.6.11). ∎
## 4. Proof of Theorem 1.1
Suppose that conditions (i) and (ii) hold. Since a domain in $`^N`$ is a domains of holomorphy if it is pseudoconvex (see \[H90\], Theorem 4.2.8) and due to the local characterization of pseudoconvexity (\[H90\], Theorem 2.6.10) it is sufficient to show that every point $`x_0\mathrm{\Omega }`$ has a neighborhood $`U`$ such that $`\mathrm{\Omega }U`$ is pseudoconvex. We use the standard identification $`^N^n`$ with $`n:=2N`$. Since $`\mathrm{\Omega }`$ is of class $`L^{2,\mathrm{}}`$, we can choose a neighborhood $`U=U^{}\times I^{n1}\times `$ of $`x_0`$ and a continuous function $`h:U^{}I`$ of class $`L^{2,\mathrm{}}`$ such that (1) holds. Define a continuous function $`\rho :U`$ with $`\mathrm{\Omega }U=\{\rho <0\}`$ by
$$\rho (x^{},x_n):=h(x^{})x_n.$$
Set $`\stackrel{~}{\mathrm{\Omega }}:=\mathrm{\Omega }U`$ and denote by $`\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$ the subset of all $`x=(x^{},x_n)\stackrel{~}{\mathrm{\Omega }}`$ such that $`(x^{},h(x^{}))\mathrm{\Omega }`$ is ($`C^2`$-)regular. We wish to apply Lemma 2.1 to the restriction of $`\rho `$ to $`\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$. To check the assumption (i) in Lemma 2.1 we observe that
$$L\rho (x^{},x_n)=L\rho (x^{},h(x^{})),\rho (x^{},x_n)=\rho (x^{},h(x^{}))$$
(12)
for all $`x\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$. Hence it follows by condition (ii) in Theorem 1.1 that $`L\rho (x)|_{\rho (x)^{}}`$ is positively semidefinite whenever $`x\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$. In order to satisfy the assumption (ii) in Lemma 2.1 we shrink the neighborhood $`U`$ of $`x_0`$ such that
$$|\rho (x)|L\rho (x)1/8$$
(13)
holds for all $`x\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$. This is possible because the second derivatives of $`\rho `$ are bounded on $`\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$ and $`\rho `$ is continuous with $`\rho (x_0)=0`$. Since $`\rho (x)1/2`$ for $`x\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$, (13) implies
$$\rho ^2+\rho L\rho 1/41/8=1/8$$
on $`\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$ and hence the existence of $`\lambda >0`$ such that the assumption (ii) in Lemma 2.1 in satisfied. By Lemma 2.1, the function
$$\varphi :\stackrel{~}{\mathrm{\Omega }},\varphi (z):=\mathrm{log}(\rho (z))+\lambda z$$
is plurisubharmonic in $`\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$.
Next we wish to apply Corollary 3.6 to $`\varphi `$ in $`\stackrel{~}{\mathrm{\Omega }}`$. For this we first construct a $`C^1`$ hypersurface $`S\stackrel{~}{\mathrm{\Omega }}`$ satisfying the assumptions. By using the definition, for every edge point $`a\mathrm{\Omega }U`$, we can choose a neighborhood $`U_a=U_a^{}\times I^{n1}\times `$ with $`U_a^{}^{n1}`$ a euclidean ball and a closed $`(n2)`$-dimensional real submanifold $`M_a\mathrm{\Omega }U_a`$ of class $`C^1`$ such that all nonregular points of $`\mathrm{\Omega }U_a`$ are contained in $`M_a`$. Furthermore we can shrink $`U_a`$ and $`M_a`$ such that all nonregular points of $`\mathrm{\Omega }\overline{U_a}`$ are contained in $`\overline{M_a}`$ and $`^{n2}(\overline{M_a}M_a)=0`$. By the choice of the coordinates $`(x^{},x_n)`$, the projection
$$M_a^{}:=\{x^{}^{n1}:x_n,(x^{},x_n)M_a\}$$
is a closed submanifold of $`U_a^{}`$. In view of condition (i) any nonregular point of $`\mathrm{\Omega }E`$ is an edge point. By considering a sequence of all rational points in $`^N`$ we may choose a sequence $`a_m`$ ($`1m<\mathrm{}`$) of edge points such that the union $`_m(V_{a_m}^{}\times I)`$ covers all nonregular points in $`(\mathrm{\Omega }E)U`$, where $`V_{a_m}^{}\times I`$ is an open neighborhood of $`a_m`$ and $`\overline{V_{a_m}^{}}U_{a_m}^{}`$. Then by applying Fubini’s theorem, we may shrink the balls $`U_{a_m}^{}`$ and their submanifolds $`M_{a_m}^{}`$ to obtain the additional property
$$^{n2}(U_{a_m}^{}M_{a_k}^{})=0\text{ for all }m,k1.$$
(14)
Define $`U_m^{}:=_{jm}U_{a_j}^{}`$.
We claim that there exist increasing sequences of $`C^1`$ submanifolds $`S_m^{}U_m^{}`$ and closed subsets $`E_m^{}U_m^{}`$ with $`^{n2}(E_m^{})=0`$ such that
1. $`S_k^{}U_m^{}=S_m^{}`$ and $`E_k^{}U_m^{}=E_m^{}`$ for all $`km`$,
2. all nonregular points in $`(\mathrm{\Omega }E)(\overline{U_m^{}}\times I)`$ outside $`E_m^{}\times I`$ are contained in $`\overline{S_m^{}}\times I`$,
3. $`^{n2}(\overline{S_m^{}}S_m^{})=0`$.
We construct the sequences $`S_m^{}`$ and $`E_m^{}`$ by induction on $`m`$. For this we set $`S_1^{}:=M_{a_1}^{}`$, $`E_1^{}:=\mathrm{}`$. If $`S_{m1}^{}`$ and $`E_{m1}^{}`$ are already constructed, define
$$S_m^{}:=S_{m1}^{}(M_{a_m}^{}(U_{a_m}^{}\overline{U_{m1}^{}}))$$
and $`E_m^{}:=E_{m1}(\overline{M_{a_m}^{}}U_{m1}^{})`$. It is easy to see from our construction that (i)-(iii) are satisfied.
Define
$$S:=((_mS_m^{})\times I)\stackrel{~}{\mathrm{\Omega }}\text{ and }F:=((_mE_m^{})\times I)\stackrel{~}{\mathrm{\Omega }}.$$
Then $`^{n1}(F)=0`$ and $`F`$ is closed in $`\stackrel{~}{\mathrm{\Omega }}`$ by (i). It follows from (i) that $`S`$ is a $`C^1`$ hypersurface in $`\stackrel{~}{\mathrm{\Omega }}`$. Furthermore, $`^{n1}(\overline{S}S)=0`$ by (iii).
Denote by $`E^{}U^{}`$ the projection of the subset $`E\mathrm{\Omega }`$. Since $`\mathrm{\Omega }U`$ is the graph of $`h`$, $`E^{}`$ is closed in $`U^{}`$. Next we subtract $`E^{}\times I`$ from $`S`$ and denote the remainder again by $`S`$. Then for every $`x=(x^{},x_n)S`$, $`(x^{},h(x^{}))\mathrm{\Omega }`$ is either $`C^2`$-regular or an edge point. Finally we define
$$G:=F(\overline{S}S)((E^{}\times I)\stackrel{~}{\mathrm{\Omega }}),$$
where the closure of $`S`$ is taken in $`\stackrel{~}{\mathrm{\Omega }}`$. By (ii), $`\stackrel{~}{\mathrm{\Omega }}(SG)\stackrel{~}{\mathrm{\Omega }}_{\mathrm{reg}}`$ and therefore $`\varphi `$ is plurisubharmonic in $`\stackrel{~}{\mathrm{\Omega }}(SG)`$.
In order to apply Corollary 3.6 it remains to show that $`\varphi `$ is transversally convex at $`S`$ (see Definition 3.4). But this is a direct consequence of condition (iii) in Theorem 1.1 on the convexity of tangent cones. We conclude that $`\varphi `$ is plurisubharmonic on $`\stackrel{~}{\mathrm{\Omega }}`$. Then, for a sufficiently small euclidean ball $`B(x_0,\epsilon )`$ centered at $`x_0`$,
$$\mathrm{max}(\varphi (z),\mathrm{log}(\epsilon zx_0))$$
is a plurisubharmonic exhaustion function for $`\mathrm{\Omega }B(x_0,\epsilon )`$. This shows pseudoconvexity of $`\mathrm{\Omega }B(x_0,\epsilon )`$ and hence completes the proof.
## 5. Proof of Proposition 1.2
Let $`a\mathrm{\Omega }`$ be a real edge point satisfying the assumptions of Proposition 1.2. As before we choose a neighborhood $`U=U^{}\times I^{n1}\times `$ of $`a`$ and a continuous function $`h:U^{}I`$ of class $`L^{2,\mathrm{}}`$ such that (1) holds. Let $`M_a\mathrm{\Omega }`$ be an edge and denote by $`M_a^{}`$ its projection on $`U^{}`$. As in §3.2 choose $`U^{}`$ sufficiently small such that $`M_a^{}`$ divides $`U^{}`$ into two parts $`U_{}^{}`$ and $`U_+^{}`$. Then the tangent cone is given by
$$T_a\mathrm{\Omega }=\{(v^{},v_n)^n:v_nD_v^{}h(a)_\pm \text{ for }v^{}T_aU_\pm ^{}\},$$
(15)
where the one-sided limits exist by Lemma 3.3.
We first suppose that $`M_a`$ is generic at $`a`$, i.e. $`T_aM_a+iT_aM_a=^N`$. We prove the statement by contradiction assuming that the tangent cone of $`\mathrm{\Omega }`$ at $`a`$ is not convex. Then there exists a linear disc $`A:\mathrm{\Delta }T_a\mathrm{\Omega }`$, $`ttv`$ with $`vT_aM_a`$ and $`A(\mathrm{\Delta }\{1,1\})`$ in the interior of $`T_a\mathrm{\Omega }`$. Here $`\mathrm{\Delta }:=\{|t|<1\}`$. Furthermore, for $`\zeta :=(0,1)^{n1}\times `$ and $`\epsilon >0`$ sufficiently small, the “deformed disc” $`A^{}(t):=tv+t^2\epsilon \zeta `$ sends the whole boundary $`\mathrm{\Delta }`$ to the interior of $`T_a\mathrm{\Omega }`$. Finally for $`\mu >0`$ sufficiently small the “rescaled disc” $`A^{\prime \prime }(t):=a+\mu A^{}(t)`$ sends the boundary $`\mathrm{\Delta }`$ into $`\mathrm{\Omega }`$. Then the Cauchy formula argument for $`A^{\prime \prime }`$ and its translations shows that all holomorphic functions in $`\mathrm{\Omega }`$ extend holomorphically to a neighborhood of $`a=A^{\prime \prime }(0)`$ which contradicts the assumption that $`\mathrm{\Omega }`$ is a domain of holomorphy.
Now consider the general case. Let $`M_a`$ be an edge at $`a`$ satisfying the assumptions. Then $`M_a`$ contains generic points arbitrarily close to $`a`$. Hence the tangent cones at those points are convex by the above argument. But then the explicit formula (15) shows that the cone $`T_a\mathrm{\Omega }`$ is also convex completing the proof.
## 6. Proof of Corollary 1.3
The necessity of the Levi form condition is well-known. The necessity of the convexity follows from Proposition 1.2. For the converse it is sufficient to show the local pseudoconvexity at every point $`a\mathrm{\Omega }`$ as in §4. Let $`NU_a`$ be given by condition (ii). By a coordinate change we may assume that $`N`$ is locally given by $`z_N=0`$. Then it follows from (15) that, if $`U_a`$ is sufficiently small polydisc centered at $`a`$, the intersection $`\mathrm{\Omega }U_a`$ can be mapped via
$$(z_1,\mathrm{},z_{N1},z_N)(z_1,\mathrm{},z_{N1},z_N^\alpha )$$
biholomorphically onto a domain satisfying assumptions of Theorem 1.1. Here $`\alpha `$ is a sufficiently small positive number. The required conclusion follows now from Theorem 1.1. |
warning/0001/astro-ph0001116.html | ar5iv | text | # Magnetic ionization fronts II: Jump conditions for oblique magnetization
## 1 Introduction
The sequence of evolution of the H$`\mathrm{II}`$ region around a star with a strong ionising radiation field which turns on rapidly is well known (Kahn 1954; Goldsworthy 1961; see also Spitzer 1978, Osterbrock 1989 or Dyson & Williams 1997). Initially, the ionization front (IF) between ionized and neutral gas moves outwards at a speed limited by the supply of ionizing photons, but it begins to decelerate as the ionizing flux at the front surface is cut by geometrical divergence and absorption by recombined atoms within the front. Eventually, when the speed of the front decreases to roughly twice the sound speed in the (now highly overpressured) ionized gas, a shock is driven forwards into the neutral gas ahead of the front. Before this stage, the front is referred to as R-type, while subsequently it is referred to as D-type. The shock propagates outwards, gradually weakening, until, in principle, the H$`\mathrm{II}`$ region eventually reaches pressure equilibrium with its surroundings.
Where the external medium has an ordered magnetic field, the obvious critical flow speeds are the fast, Alfvén and slow speeds rather than the isothermal sound speed. Redman et al. \[Paper I, 1998, hereafter\] studied IFs with the magnetic field vector in the plane of the front, and found that the fast-mode speed plays the same role as the sound speed does in the hydrodynamic case. In this paper, we extend their work to treat the case of an IF moving into a medium in which the magnetic field is oblique to the direction of propagation of the front (note also we here follow the more conventional usage in which the magnetic fields are termed parallel or perpendicular with respect to the front-normal).
Jump conditions for IFs with oblique magnetization have previously been studied by Lasker . Here we consider a wider range of upstream magnetic fields, since observations have shown that higher magnetic fields are found around H$`\mathrm{II}`$ regions than once thought likely \[1993, 1995, e.g.\]. We determine the properties of the jumps as functions of upstream conditions, using the velocity of the front as a parameter rather than as the variable for which we solve. We use evolutionary conditions to isolate the stable IF solutions, and verify these conclusions for a simple model of the internal structure of the fronts and using numerical simulations. We find that rather weak parallel magnetic fields can lead to a substantial decrease in the D-critical (i.e. photoevaporation) velocity from dense clumps except where the magnetic field is exactly parallel to the IF, and also find additional solutions to the jump conditions in the range of front velocities forbidden by the hydrodynamical jump conditions, which were not considered by Lasker.
In the following sections, we present the jump conditions for MHD shocks (Section 2), and discuss the regions for which evolutionary conditions suggest that these solutions are stable (Section 3). We verify that the evolutionary solutions are those with resolved internal structures for one simple model for the internal structure of the fronts (Section 4). In the context of these results, we discuss the development of an IF over time (Section 5) and illustrate this development using numerical models (Section 6).
Finally (in Section 7), we summarize our results, and provide an example of their application to observations of the H$`\mathrm{II}`$ region S106. Our physical interpretation of the development of MHD IFs will, we hope, facilitate the future application of these results.
## 2 Jump conditions
We orient axes so that $`\widehat{𝒛}`$ is normal to the front, and (without loss of generality) that the upstream velocity and magnetic field are in the $`(x,z)`$ plane. We use the usual MHD jump conditions:
$`[\rho v_z]`$ $`=`$ $`0`$ (1)
$`[\rho v_z^2+p+B_x^2/8\pi ]`$ $`=`$ $`0`$ (2)
$`[\rho v_zv_xB_zB_x/4\pi ]`$ $`=`$ $`0`$ (3)
$`[B_z]`$ $`=`$ $`0`$ (4)
$`[v_xB_zv_zB_x]`$ $`=`$ $`0,`$ (5)
except that instead of using the energy flux condition, we adopt the isothermal equation of state $`p=\rho c_\mathrm{s}^2`$ where the sound speed, $`c_\mathrm{s}`$, increases across the front but is constant on either side of it. We use subscripts 1 and 2 to denote upstream and downstream parameters, respectively, and write $`v_x=u_{1,2}`$, $`v_z=v_{1,2}`$, $`c_\mathrm{s}=c_{1,2}`$ and $`B_x=B_{1,2}`$. Hence
$`\rho _1/\rho _2=v_2/v_1`$ $``$ $`\delta `$ (6)
$`\rho _2(v_2^2+c_2^2)+B_2^2/8\pi `$ $`=`$ $`\rho _1(v_1^2+c_1^2)+B_1^2/8\pi `$ (7)
$`\rho _2v_2u_2B_zB_2/4\pi `$ $`=`$ $`\rho _1v_1u_1B_zB_1/4\pi `$ (8)
$`u_2B_zv_2B_2`$ $`=`$ $`u_1B_zv_1B_1`$ (9)
Equations (6), (8) and (9) give
$$B_2=\frac{m_1^22\eta _1}{\delta m_1^22\eta _1}B_1,$$
(10)
where $`m=v/c_\mathrm{s}`$ and we define $`\eta =B_z^2/8\pi \rho c^2`$ and $`\xi =B_x^2/8\pi \rho c^2`$ (the $`z`$ and $`x`$ contributions to the reciprocal of the plasma beta). The dependence on the upstream transverse velocity has disappeared, as expected as a result of frame-invariance.
In equation (7), we substitute with (10) for $`B_2`$ and use equation (6) to eliminate $`\rho _2`$ and $`v_2`$ to find that, so long as $`\delta 0`$ and $`\delta m_1^22\eta _1`$, the dilution factor $`\delta `$ is given by the quartic equation \[1966, see also\]
$`m_1^6\delta ^4m_1^4(1+m_1^2+4\eta _1+\xi _1)\delta ^3`$
$`+m_1^2(\alpha m_1^2+4\eta _1(1+m_1^2+\eta _1+\xi _1))\delta ^2`$
$`+(\xi _1m_1^44\eta _1(\eta _1+m_1^2(\eta _1+\xi _1+\alpha )))\delta `$
$`+4\alpha \eta _1^2`$ $`=`$ $`0,`$ (11)
where $`\alpha =(c_2/c_1)^2`$ ($`100`$ is a typical value). It is easily verified that this equation has the correct form in the obvious limiting cases (of isothermal MHD shocks where $`\alpha =1`$, and perpendicular-magnetized IFs where $`\eta _1=0`$, see Paper I).
## 3 Solutions
Equation (11) relates the dilution factor, $`\delta `$, to the upstream properties of the flow. Its roots can be found by standard techniques. If the flow is to have a unique solution based only on initial and boundary conditions, then some of these roots must be excluded. It is possible to exclude roots on the basis that they correspond to flows that do not satisfy an evolutionary condition, analogous to that long used in the study of MHD shocks. The evolutionary condition is based on the requirement that the number of unknowns in the jump conditions matches the number of constraints applied to the flow \[1964, e.g.\]. For MHD shocks, the number of characteristics entering the shock must be two greater than the number leaving it (since the number of independent shock equations is equal to the number of conserved variables and if there are no internal constraints on the front structure).
The applicability of the evolutionary conditions to MHD shocks has been the subject of controversy in the recent past \[1988, 1990, e.g.\], with various authors suggesting that intermediate shocks (i.e. shocks which take the flow from super- to sub-Alfvén speeds) may be stable. However, Falle & Komissarov have shown that the non-evolutionary solutions are only stable when the symmetry is artificially constrained, so that the magnetic fields ahead of and behind the shocks are precisely coplanar. In any cases in which the boundary conditions differed from this special symmetry, the solutions including non-evolutionary shocks were found to be unstable.
The equations governing the dynamics across an IF are no longer a system of hyperbolic conservation laws, since the ionization source term cannot be neglected on the scale of the front. It seems reasonable to apply analogous evolutionary and uniqueness conditions, but the mathematical proofs for hyperbolic conservation laws with dissipation no longer apply. The ‘strong’ evolutionary conditions suggest that for IF the number of incoming characteristics is equal to the number of outgoing characteristics, since applying the ionization equation means that there is an additional constraint on the velocity of the front.
Solutions which are under-specified by the external characteristic constraints, termed ‘weakly evolutionary’, may occur if there are internal constraints, as is the case for strong D-type IFs in hydrodynamics. Stable weakly evolutionary solutions only occur for limited classes of upstream parameters which depend on the internal structure of the fronts. Where the number of characteristics entering the front is greater than suggested by the evolutionary conditions, the solution can be realized as an MHD shock leading or trailing an evolutionary IF.
In the limit in which $`\alpha `$ tends to unity from above, the phase change through the IF has no dynamical consequences, and the IF jump conditions approach those for isothermal MHD shocks. This can be seen if we rewrite equation (11) as
$`(\delta 1)[m^6\delta ^3m^4(1+4\eta +\xi )\delta ^2`$
$`+m^2(4\eta (1+\eta +\xi )m^2\xi )\delta 4\eta ^2]`$
$`=(\alpha 1)[m^2\delta 2\eta ]^2.`$ (12)
The IF solution which satisfies the strong evolutionary conditions becomes the trivial ($`\delta =1`$) solution of the MHD jump condition. The other solutions to the IF jump conditions become non-evolutionary or evolutionary isothermal MHD shocks. The analysis of Falle & Komissarov rules out the former as physical solutions. The latter are treatable as separate discontinuities, which will propagate away from the IF when the flow is perturbed, since the coincidence between the speeds of the shock and of the IF will be broken. This argument by continuity supports the use of the evolutionary conditions for IFs.
As a result of this discussion, we will proceed for the present to isolate solutions in which the number of characteristics entering an IF is equal to the number leaving it (and discuss in Section 4 the weakly evolutionary solutions for a simplified model of the internal structure of IFs). The velocity of fronts obeying the strong evolutionary conditions must be between the same critical speeds in the upstream and downstream gas (somewhat confusingly, the fronts which obey the strong evolutionary conditions are termed ‘weak’ in the standard nomenclature of detonations and IFs). We follow the usual classification of flow speeds relative to the fast, Alfvén and slow mode speeds $`1>v_\mathrm{f}>2>v_\mathrm{a}>3>v_\mathrm{s}>4`$, so the allowed fronts are $`11`$, $`22`$, $`33`$ and $`44`$. By analogy with the nomenclature of non-magnetized IF, we call these fast-R, fast-D, slow-R and slow-D type IF, respectively.
The panels of Fig. 1 show regions of $`m_1`$ and $`\xi _1`$ space corresponding to evolutionary MHD IFs for several values of $`\eta _1`$. In the figures, we see regions corresponding to the two distinct classes of R- and D-type solutions. The flows into the R-type fronts are super-fast or super-slow, while those into the D-type fronts are sub-fast or sub-slow. At the edges of the regions of solutions either the velocity into the front is the Alfvén speed or the exit velocity from the front is equal to a characteristic speed (i.e. the fast mode speed at the edge of the fast-R region, etc.). For comparison, the solid lines on these plots show the edges of the forbidden region for perpendicular magnetization ($`\eta _1=0`$, see Paper I): these lines reach $`\xi _1=0`$ at the edges of the forbidden region for unmagnetized IF, $`0.05m_120`$.
Since equation (11) is a cubic in $`m_1^2`$, quadratic in $`\eta _1`$ and linear in $`\alpha `$ and $`\xi _1`$, there is no simple analytic form for the boundaries of the regions. However, certain critical values can be determined analytically. For $`\eta _12\alpha `$, the slow-R-critical line terminates where it hits the Alfvén speed at $`\xi _1=\alpha 1`$, while the position at which the fast D-critical line terminates is given by
$$(1+2\eta _1+\xi _1)^2=8\alpha \eta _1.$$
(13)
These points are linked, respectively, by steady switch-off and switch-on shocks to the points on the limiting slow-D and fast-R critical loci at which these loci hit the axis $`\xi _1=0`$. The intercept between the slow-D critical locus and the axis is at
$$m_1^2=2\eta _1\frac{12\eta _1}{\alpha 2\eta _1},$$
(14)
for $`2\eta _1(\alpha \sqrt{\alpha ^2\alpha })`$, beyond which the limiting value is $`m_1^21/(4\alpha )`$ as for D-critical hydrodynamic IF. To illustrate the reason for this change in solution, we rewrite equation (11) for $`\xi _1=0`$ as
$$(m_1^2\delta 2\eta )^2(m_1^2\delta ^2(1+m_1^2)\delta +\alpha )=0.$$
(15)
The flow leaving a front with no upstream perpendicular component of magnetic field can be either at the Alfvén speed or at the velocity of the corresponding non-magnetized IF. Where the flow is in the region beyond the edge of the slow-D-critical region shown in Figure 1 (a) or (b), the root of equation (15) for flow out at the Alfvén speed is not a real solution, since satisfying equation (7) would require that $`B_2^2<0`$.
Even for as small a ratio between magnetic and thermal energy upstream of the front as implied by $`\eta _1=0.01`$, the effect of parallel magnetization on the slow-D-critical velocity is dramatic. Only once $`\eta _11/(8\alpha )`$ (so the Alfvén speed in the upstream gas is below the unmagnetized D-critical speed) does the fast-critical locus reach $`\xi _1=0`$, so that the parallel magnetization may be ignored. As $`\eta _1`$ increases, the vertical line at the Alfvén speed moves across the plot (see Figure 1), decreasing the region of fast-mode IFs and increasing that of slow-mode IFs. When $`\eta _1\mathrm{}`$ (a very strong parallel magnetic field), the (slow-mode) forbidden region is identical to that in the unmagnetized case (independent of $`\xi _1`$).
If the upstream flow is at the Alfvén velocity, $`m_1^2=2\eta _1`$, then the physical solution (when due care is taken with the singularity of equation (10)) is often at $`\delta =1`$, i.e. the ionization of the gas does not change the flow density and it remains at the Alfvén speed. Both classes of D-type front have $`\delta >1`$ (rarefy the gas), while both classes of R-type have $`\delta <1`$ (compress it). The perpendicular component of the magnetic field, $`B_x`$, increases in a fast-R- or slow-D-type IF, while it decreases in a fast-D- or slow-R-type.
We will now study the internal structure of the fronts for one simple model.
## 4 Resolved fronts
Up to now we have assumed, by investigating the jump conditions, that the processes within the IF take place on scales far smaller than those of interest for the global flow problem. In fact, the flow structure within an ionization front will vary smoothly on scales comparable to the ionization distance in the neutral gas. The flow may take several recombination lengths behind the front to relax to equilibrium. Here we describe the internal structure of MHD IF in one simple approximation, that the temperature of the gas varies smoothly through the front but the flow obeys the inviscid MHD equations throughout \[1961, as used in the study of hydrodynamic IF structure by\]. We shall see that, in this approximation, only fronts obeying the evolutionary conditions can have smooth structures. For the structures to be generic, the jumps across them must satisfy the strong evolutionary conditions, although singular classes of weakly evolutionary fronts with internal constraints on their flow structures are also possible.
In this model the form of an IF is given by the variation of the roots of equation (12) with the temperature of the gas, specified by $`\alpha `$. The left hand side of this equation is a quartic independent of $`\alpha `$, which is positive at $`\delta =0,\mathrm{}`$ and zero at $`\delta =1`$, and has two or four positive roots \[1963, since the condition that its value is zero is the normal MHD shock condition for isothermal gas, cf.\]. The right hand side is a quadratic, which is zero at $`m_1^2\delta =2\eta `$ and depends on $`\alpha `$ through its (negative) scale. If the value of $`\alpha `$ increases steadily through the front, the manner in which the solutions vary can be followed, as illustrated by the schematic plot, Fig. 2.
For the upstream conditions ($`\alpha =1`$), there will be either two or four solutions where the quartic curve crosses the axis. Of these, one is the trivial solution, $`\delta =1`$, and at most one corresponds to an evolutionary shock.
When internal heating occurs in an IF, the solutions at a given $`\alpha `$ correspond to where the solid curve in Fig. 2 crosses the corresponding dashed curve, $`(\alpha 1)(m_1^2\delta 2\eta _1)`$. If gas enters the IF at the point marked $``$, in region 2 between the Alfvén speed and the fast mode speed, then it can move following the arrow as $`\alpha `$ increases. For sufficiently large $`\alpha `$, the dashed curve becomes tangent to the solid curve: when this occurs, the gas is moving at the slow- or fast-mode speed, and the IF is called a critical front. In between the initial and final solutions, smooth, steady IF solutions must remain between the same characteristic speeds as they were when they started. The strong evolutionary conditions give just those cases in which a front structure calculated for a smoothly varying, monotonically increasing $`\alpha `$ (and non-zero perpendicular magnetic field) has a continuous solution from the upstream to the downstream case.
When the upstream flow is at the fast or the slow critical speed, the l.h.s. of equation (12) has a second root $`\delta =1`$. Thus for $`\alpha >1`$, this pair of roots disappears, and a forbidden region is generated. For $`\alpha 1`$ there is no forbidden region.
For any $`\alpha `$, the points where the solid curve crosses the dashed curve in Fig. 2 are related to each other by the isothermal MHD shock jump conditions, so a steady shock can form anywhere within the IF structure for identical upstream and downstream conditions. However, since such a flow is over-specified, the equality of the speed between the shock and IF is a coincidence which will be broken by any perturbation of the flow (in which case the shock will escape from one or other side of the IF). This situation is in direct analogy with strong R-type IFs in unmagnetized flows.
Strong D-type fronts, for which the exhaust leaves the front rather above the critical speed, can occur where the heating is not monotonic. For these, the highest temperature is attained when the flow is at the critical speed (i.e. the curves in Figure 2 become tangent exactly at the highest value of $`\alpha `$), and as it subsequently cools the solution can move back up the other branch. These fronts will form a more restrictive limiting envelope on the allowed weak solutions than that given by the critical solutions. The evolutionary conditions are necessary but not sufficient, so this behaviour would be expected when more detailed physics was included. The actual envelope will correspond to the case for critical fronts with exhausts at the highest temperature attained within the front. Since unmagnetized IF models suggest that any overshoot in the temperature of the gas is likely to be small, the envelope will probably not differ greatly from that found for critical solutions, although the transonic nature of these fronts can be important in determining the structure of global flows.
Equation (10) suggests that no MHD flow can pass through the Alfvén velocity (where $`\delta m_1^2=2\eta _1`$) in a front unless it has zero perpendicular magnetic field. Heating the gas will generally move the flow in regions 2 and 3 away from the Alfvén speed in any case (see Figures 2 and 3). However, where the perpendicular magnetic field is zero, a smooth, weakly evolutionary, front structure can be found (the internal constraint being zero perpendicular magnetic field). The internal structure will be identical to an unmagnetized IF. Indeed, the strong-D hydrodynamic front can become, by analogy, an ‘extra-strong’ front which passes through both the Alfvén and sound speeds (for zero perpendicular field, the slow and fast velocities are each equal to one of these).
By analogy with equation (10), the $`y`$-components of velocity and magnetic field are zero everywhere if the MHD equations apply throughout the front (except if it were to pass through the Alfvén velocity). Components in these directions can be generated if the velocity coupling between different components of the fluid – electrons, ions, neutrals or dust – is not perfect \[1994, 1998, as in shock structures,\]. These components must, however, damp at large enough scales: far beyond the front the magnetic field must be in the same plane and of the same sign as the upstream field, from the evolutionary conditions. Exactly this behaviour has been found to occur in time-dependent multifluid models of MHD shocks (Falle, private communication). A full treatment of ionization fronts in multicomponent material is, however, beyond the scope of the present paper.
## 5 Time development
In this section, we discuss the development of the IF in a magnetized H$`\mathrm{II}`$ region, by combining the well-understood development of IF in unmagnetized environments with the classes of physical roots of equation (11) found above.
An IF driven into finite density gas from a source which turns on instantaneously will start at a velocity greater than the fast-R-critical velocity. Unless the density decreases rapidly away from the source, the speed of the front decreases so that eventually the ionized gas exhaust is at the fast-mode speed (at the fast-R-critical velocity). When this occurs, two roots of equation (11) merge, and become complex for smaller $`m_1`$. As a result, the front will then have to undergo a transition of some sort. As in the unmagnetized case , if the size of the ionized region is large compared to the lengthscales which characterise the internal structures of shocks and IFs, the IF will evolve through emitting (one or more) shocks. The evolution of an initial IF discontinuity can be treated as a modified Riemann problem because the speed of the IF is determined by the flow properties on either side of it, together with the incident ionizing flux which we assume varies slowly. The development of this modified Riemann problem will be self-similar, just as for conventional Riemann problems. One complication is that the IF may be located within a rarefaction wave, but this does not occur for the circumstances we discuss in the present section.
The simplest possibility for a slowing fast-R-critical IF is that it will become fast-D-type by emitting a single fast-mode shock. This has obvious limits to the cases where the magnetization is zero (where the shock is a normal hydrodynamic shock), and where it is purely perpendicular. If the speed of the IF is specified by the mass flux through it, then the leading shock driven into the surrounding neutral gas must be a fast-mode shock, since the upstream neutral gas must still be advected into the combined structure at more than the Alfvén speed after the transition, and so only a fast-mode shock can escape.
While there may be no fast-D-type solutions at the value of $`\xi _1`$ which applied for the fast-R-type front, the fast-mode shock moving ahead will act to increase the value of $`\xi `$ upstream of the IF, since the (squared) increase in the perpendicular component of the magnetic field dominates over the increase in the gas density after the shock. At fast-R-criticality, there is in general a second solution with the same upstream and downstream states in which a fast-mode shock leads a fast-D-critical IF, because the l.h.s. of equation (12) (for which a zero value implies an isothermal shock solution) is positive for large $`\delta `$ and for post shock flow at the Alfvén speed (i.e. $`m_1^2\delta =2\eta _1`$), unless the flow into the front is at the Alfvén velocity, or $`\eta _1`$ or $`\xi _1`$ is zero. For example, for a downstream state $`\eta _2=\xi _2=5\times 10^3`$, there are two evolutionary solutions: a $`11`$ front with $`\eta _1=0.994`$, $`\xi _1=0.248`$, $`m_1=20.0`$ and $`\delta =0.503`$, and a $`22`$ front with $`\eta _1=3.11\times 10^2`$, $`\xi _1=11.1`$, $`m_1=0.625`$ and $`\delta =16.1`$. Figure 1 parts (c) and (b), respectively, contain points which correspond to these solutions. These two upstream states are linked by a fast-mode shock with the same velocity as the IF. Thus the emission of a single fast-mode shock is a valid solution of the modified Riemann problem which occurs as the flow passes through criticality whatever the internal structure of the shock and IF, so long as the evolutionary shocks and IFs exist. (Since the l.h.s. of equation (12) is greater than zero so long as $`\eta _10`$ for $`\delta =0`$, an equivalent argument holds for slow-critical transitions.)
It is possible that further waves may be emitted at the transition, for instance a slow-mode shock into the neutral gas together with a slow-mode rarefaction into the ionized gas. For the model resolved IF, this seems unlikely to occur unless the flow velocity reaches the slow-mode speed somewhere within the fast-D-critical IF. These other solutions are not seen in our numerical simulations below. Additional solutions would also make the development of the IF non-unique, if the simpler possibility is allowed.
The fast-critical transition may be followed using the jump conditions for a front with its exhaust at the fast-mode speed (i.e. a fast-critical front). In the smoothly-varying $`\alpha `$ model of Section 4, the upstream and downstream states can be joined by a front in which an isothermal MHD shock is at rest in the IF frame anywhere within the IF structure, since the quantities conserved through the front are also conserved by the shock. The evolution of an IF through criticality will occur by an infinitesimally-weak fast-mode wave at the exhaust of the IF moving forward through its structure and strengthening until it eventually escapes into the neutral gas as an independent shock \[1996, as illustrated for recombination fronts by\].
The escaping fast-mode shock leads to a near-perpendicular field configuration upstream of the IF. This boost in the perpendicular component is required if the transition is to proceed though the fast-D type solutions, which, as Figure 1 illustrates, are near-perpendicular (large $`\xi _1`$) except where $`B_z`$ is very large or very small. This will result in a rapid change in downstream parameters across a front where the upstream field is nearly parallel to the IF. As an example, the flow downstream of the IF will either converge onto or diverge from lines where the upstream magnetic field is parallel to the ionization front, and as a result may produce inhomogeneities in the structure of H$`\mathrm{II}`$ regions on various scales (from bipolarity to clumps).
As the velocity of the front decreases further, eventually it will approach the Alfvén speed. We find that the evolution depends qualitatively on whether $`\alpha 1`$ is smaller than $`\xi _1`$. In Fig. 3, we plot $`\delta `$ as a function of $`\alpha `$ for a range of values of $`m_1`$ (i.e. the internal structure of the fronts in the simple model of Section 4). We take $`\eta _1=3\times 10^2`$ and $`\xi _1=11`$, corresponding to our numerical example above, so the ratio of the upstream Alfvén speed to the upstream isothermal sound speed is $`0.245`$. For $`m_1`$ slightly larger than $`0.245`$ (i.e. just super-Alfvén), the curves remain flat until $`\alpha `$ is close to $`\xi _1+1`$ and then turn upwards when they reach this value. If the value of $`\alpha `$ in the fully ionized gas were less than $`\xi _1+1`$, the solutions would move through a case where the flow is of uniform density and moves at the Alfvén speed throughout the front before becoming slow-R-critical for some $`m_1^2<2\eta `$, as can be seen in the leftmost part of Fig. 3, for fronts in which the maximum $`\alpha `$ is smaller than 12.
For larger values of $`\alpha `$, however, the solutions develop a gradient discontinuity in their structure when the inflow is at the Alfvén speed, $`m_1^2=2\eta `$. At this discontinuity, the transverse component of the magnetic field becomes zero (i.e. switch-off occurs). Once the propagation speed of the IF drops below the Alfvén speed, there is no form of steady evolutionary structure with a single wave in addition to the IF which is continuous with that which applied before. Non-evolutionary $`32`$ type solutions do exist for fronts just below this limit, but in numerical simulations (see Section 6) these break up. A slow-mode switch-off shock moves into the neutral gas and a slow-mode switch-on rarefaction is advected away into the ionized gas. Between them, these waves remove the parallel component of magnetic field at the D-type IF between them. A precursor for the rarefaction is apparent in internal structure of the critical front, Figure 3. We find in numerical simulations that the IF which remains is trans-Alfvénic. Note that a steady resolved structure is possible for such (weakly evolutionary) fronts only because it has exactly zero parallel field throughout.
Analogous processes must occur in an accelerating fast-mode IF with weak parallel magnetization when $`\eta _1`$ is very large: in Fig. 1 (d), it is the fast-D-critical line rather than the slow-R-critical line which meets the Alfvén locus at finite $`\xi _1`$.
## 6 Numerical solutions
In this section, we present some numerical examples to illustrate the processes discussed in the preceding sections. We have implemented linear scheme A as described by Falle, Komissarov & Joarder in one dimension, and added an extra conserved variable corresponding to the flow ionization. To study the local development of the IFs we have neglected recombination terms and just chosen to set the mass flux through the ionization flux as a function of time.
Figure 4(a) shows the propagation of an IF into gas with density $`1`$, $`B_z=0.3\sqrt{4\pi }`$, $`B_x=2\sqrt{4\pi }`$. We reset the flow temperature at the end of each time step so that $`p=(0.1+0.9x)\rho `$ (the temperature ratio between ionized and neutral gas is rather small so that any shells of shocked neutral gas are more easily resolved). For these values, the characteristic speeds in the neutral gas are $`v_\mathrm{s}=0.046`$, $`v_\mathrm{a}=0.3`$ and $`v_\mathrm{f}=2`$. In the first figure, the incident flux varies as $`80/(20+t)`$, and the transitions through fast-R, fast-D and slow-D are clearly visible, while the slow-R stage (which is often narrow in the parameter space of Fig. 1) is less clear.
In Figure 4(b) the upstream conditions are the same, but the flux was set to a constant value of $`0.26`$ in order to isolate the slow-R transition. This is between the upstream slow-mode and Alfvén speeds, and equation (11) predicts that a slow-R transition exists which will increase the flow density from $`1`$ to $`1.18`$ (for comparison, the $`34`$ jump relations require the density to increase to 1.75). In the simulation, a weak fast-mode wave propagates off from the IF initially, but does not greatly change the upstream conditions. The slow-R IF which follows it increases the density, and is backed by a rarefaction because of the reflective boundary condition applied at the left of the grid. The small overshoot within the front is presumably due to numerical viscosity, and can be removed by broadening the IF \[1999, e.g. by the method described in\]. The plateau between the rarefaction and the IF has density $`1.22`$, which is in adequate agreement with the jump conditions (particularly when account is taken of the slight perturbation of the conditions upstream of the front).
These numerical solutions illustrate the orderly progression of a magnetized IF through the various transitions described in the previous section. In a realistic model, however, the density perturbations generated by the transitions will have important effects on the evolution, as the consequent changes in recombination rates alter the flux incident on the IF. These processes should ideally be studied in the context of a two- or three-dimensional global model for the evolution of magnetized H$`\mathrm{II}`$ regions, which is beyond the scope of the present paper.
In Figure 5, we illustrate the development of the flow from initial conditions in which a non-evolutionary IF is stationary in the grid. The upstream (neutral) gas has density $`1`$, $`B_z=0.3\sqrt{4\pi }`$, and $`B_x=\sqrt{4\pi }`$ and moves into the front at $`v_z=0.253`$ (with no transverse velocity), while the downstream (fully ionized) gas has $`\rho =0.423`$, $`v_z=0.598`$, $`v_x=1.688`$ $`B_z=0.3\sqrt{4\pi }`$, and $`B_x=0.424\sqrt{4\pi }`$ so the IF is of $`32`$ type. We set the pressure as above, and the value of the incident flux as $`0.253`$ so the initial IF would remain steady in the grid. The IF breaks up immediately, driving a slow-mode shock away to the right, into the neutral gas, while a slow-mode rarefaction moves away to the left, into the ionized gas. The D-type IF which remains is marginally trans-Alfvénic, but has zero transverse magnetic field.
To study this further, in Fig. 6, we illustrate a simulation of an IF close to the (hydrodynamic) D-critical condition, with a parallel magnetic field which makes it trans-Alfvénic. When perturbed with small but significant perpendicular field components, this IF again switches off these fields by emitting slow-mode waves. It remains stable, satisfying the hydrodynamic jump conditions as a weakly evolutionary solution of the MHD jump conditions.
## 7 Conclusions
We have presented the jump conditions for obliquely-magnetized ionization fronts. We have determined the regions of parameter space in which physical IF solutions occur, and have discussed the nature of the interconversions between the types of front. Fast-D and slow-R solutions with high transverse fields are found in the region of front velocities forbidden by the hydrodynamic jump conditions: in the evolution of an H$`\mathrm{II}`$ region, the fast-mode shock sent into the neutral gas by the fast-critical transition will act to generate these high transverse fields. In the obliquely-magnetized case, the fronts are significantly perturbed as long as the Alfvén speed in the neutral gas is greater than $`c_1^2/2c_2`$. However, the stability of parallel-magnetized weakly evolutionary IF means that the flow may still leave at the isothermal sound speed in the ionized gas over much of the surface of magnetized globules exposed to ionizing radiation fields.
Large ($`100\mu \mathrm{G}`$), highly ordered magnetic fields have been observed in the molecular gas surrounding some H$`\mathrm{II}`$ regions \[1993, 1995, e.g.\]. Roberts et al. suggested that the highest observed magnetizations in S106 are associated with unshocked rather than shocked gas, as a result of the relatively low density and velocity shift observed for the strongly magnetized gas. They suggested that the magnetic field becomes tangled close to the IF, leading to the decrease in detectable magnetization.
It is interesting to compare these results with the example fast-critical front we discuss in Section 5. Scaling the parameters of this front to an exhaust hydrogen density of $`10^4\mathrm{cm}^3`$, typical of an ultracompact H$`\mathrm{II}`$ region, a mean mass per hydrogen nucleus of $`10^{24}\mathrm{g}`$, and a sound speed in the ionized gas of $`10\mathrm{km}\mathrm{s}^1`$, the limiting $`11`$ front takes the flow from an upstream state with $`n_\mathrm{H}=5.03\times 10^3\mathrm{cm}^3`$, $`𝒗=(0.0497,0,20.0)\mathrm{km}\mathrm{s}^1`$ and $`𝑩=(17.7,0,35.4)\mu \mathrm{G}`$ to an exhaust with $`n_\mathrm{H}=10^4\mathrm{cm}^3`$, $`𝒗=(0,0,10.1)\mathrm{km}\mathrm{s}^1`$ and $`𝑩=(35.4,0,35.4)\mu \mathrm{G}`$ (while the fast-R-critical speed is almost exactly twice the exit speed of the front, this ratio decreases in more strongly magnetized fronts). The limiting $`22`$ front takes the flow from $`n_\mathrm{H}=1.61\times 10^5\mathrm{cm}^3`$, $`𝒗=(1.78,0,0.625)\mathrm{km}\mathrm{s}^1`$ and $`𝑩=(671,0,35.4)\mu \mathrm{G}`$ to the same final state. A fast-mode shock ahead of the front and at rest with respect to it would change the upstream state from that of the limiting $`11`$ front to that of the limiting $`22`$ front.
This fast-mode shock, which precedes a limiting fast-weak-D type front, boosts the $`x`$-component of the magnetic field from $`18\mu \mathrm{G}`$ to $`670\mu \mathrm{G}`$, with a $`20\mathrm{km}\mathrm{s}^1`$ change in the $`z`$-velocity and only a $`2\mathrm{km}\mathrm{s}^1`$ change to the $`x`$-velocity component. A succeeding slow-mode shock would further increase the $`z`$-velocity component and gas density while weakening the $`x`$-component of magnetic field, without recourse to field-tangling. If we tentatively identify OH component B of Roberts et al. as fast-shocked material and component A as doubly-shocked material, component B is more edge-brightened and has a smaller blueshift than component A as would be expected. Component A is kinematically warmer and most blue shifted towards the centre of the region. The strong line-of-sight magnetic fields in S106 are seen at the edges of the region, in a ‘toroidal’ distribution. The (poorly resolved) line-of-sight velocity of component B has little gradient in the equatorial plane of the region, but this might result in part from a combination of flow divergence and the value of $`B_{\mathrm{los}}`$ (which is measured close to the centre of the region) being rather larger than that of $`B_z`$ in our example. While the qualitative properties of this assignment are attractive, it remains to calculate a proper model tuned to the properties of the region, in particular its geometry. Nevertheless, the present discussion at least illustrates how both fast and slow shocks should be considered in the analysis of regions with well-ordered magnetization.
In future work, we will model in detail the global structure of magnetized H$`\mathrm{II}`$ regions and the local structure of photoevaporated magnetized clumps.
### Acknowledgements.
We thank Sam Falle and Serguei Komissarov for helpful discussions on evolutionary conditions, and the referee for constructive comments which brought several issues into sharper focus. RJRW acknowledges support from PPARC for this work. |
warning/0001/cond-mat0001158.html | ar5iv | text | # Casimir amplitudes in critical quantum systems
## Casimir amplitudes in critical quantum systems
Let us consider a critical quantum system with a film geometry $`L\times \mathrm{}^{d1}\times L_\tau `$, where $`L_\tau =\mathrm{}/(k_BT)`$ is the “finite-size” in the temporal (imaginary time) direction and let us suppose that periodic boundary conditions are imposed across the finite space dimensionality $`L`$ (in the remainder we will set $`\mathrm{}=k_B=1`$).
The confinement of critical fluctuations of an order parameter field induces long-ranged force between the boundary of the plates . This is known as “statistical-mechanical Casimir force”. The Casimir force in statistical-mechanical systems is characterized by the excess free energy due to the finite-size contributions to the free energy of the bulk system. In the case it is defined as
$$F_{\mathrm{Casimir}}(T,g,L|d)=\frac{f^{\mathrm{ex}}(T,g,L|d)}{L},$$
(1)
where $`f^{\mathrm{ex}}(T,g,L|d)`$ is the excess free energy
$$f^{\mathrm{ex}}(T,g,L|d)=f(T,g,L|d)Lf(T,g,\mathrm{}|d).$$
(2)
Here $`f(T,g,L|d)`$ is the full free energy per unit area and per $`k_BT`$, and $`f(T,g,\mathrm{}|d)`$ is the corresponding bulk free energy density.
Then, near the quantum critical point $`g_c`$, where the phase transition is governed by the non thermal parameter $`g`$, one could state that ( see, )
$$\frac{1}{L}f^{\mathrm{ex}}(T,g,L|d)=\left(TL_\tau \right)L^{(d+z)}X_{\mathrm{ex}}^\mathrm{u}(x_1,\rho |d),$$
(3)
with scaling variables
$$x_1=L^{1/\nu }\delta g,\text{and}\rho =L^z/L_\tau .$$
(4)
Here $`\nu `$ is the usual critical exponent of the bulk model, $`\delta ggg_c`$, and $`X_{\mathrm{ex}}^\mathrm{u}`$ is the universal scaling function of the excess free energy. According to the definition (1), one gets
$$F_{\mathrm{Casimir}}^d(T,g,L)=\left(TL_\tau \right)L^{(d+z)}X_{\mathrm{Casimir}}^\mathrm{u}(x_1,\rho |d),$$
(5)
where $`X_{\mathrm{Casimir}}^\mathrm{u}(x_1,\rho |d)`$ is the universal scaling functions of the Casimir force.
It follows from Eq. (5) that depending on the scaling variable $`\rho `$ one can define Casimir amplitudes
$$\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{u}\left(\rho |d\right):=X_{\mathrm{Casimir}}^\mathrm{u}(0,\rho |d).$$
(6)
In addition to the “usual” excess free energy and Casimir amplitudes, denoted by the superscript “$`u`$”, one can define, in a full analogy with what it has been done above, “temporal excess free energy density” $`f_\mathrm{t}^{\mathrm{ex}}`$,
$$f_\mathrm{t}^{\mathrm{ex}}(T,g,|d)=f(T,g,\mathrm{}|d)f(0,g,\mathrm{}|d).$$
(7)
If the quantum parameter $`g`$ is in the vicinity of $`g_c`$, then one expects
$$f_\mathrm{t}^{\mathrm{ex}}(T,g|d)=TL_\tau ^{d/z}X_{\mathrm{ex}}^\mathrm{t}\left(x_1^t|d\right),$$
(8)
i.e. instead of $`X_{\mathrm{ex}}^\mathrm{u}(x_1,\rho |d)`$. one has a scaling function $`X_{\mathrm{ex}}^\mathrm{t}\left(x_1^t|d\right)`$ which is the corresponding analog with scaling variables
$$x_1^t=L^{1/\nu z}\delta g.$$
(9)
Obviously one can define the ”temporal Casimir amplitude”
$$\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{t}\left(d\right):=X_{\mathrm{ex}}^\mathrm{t}\left(0|d\right).$$
(10)
Whereas the “usual” amplitudes characterize the leading $`L`$ corrections of a finite size system to the bulk free energy density at the critical point, the “temporal amplitudes” characterize the leading temperature-dependent corrections to the ground state energy of an infinite system at its quantum critical point $`g_c`$.
For the universality class under consideration the following exact results are obtained:
(i)For the ”usual” Casimir amplitudes
$$\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{u}\left(0|2,2\right)=\frac{2\zeta (3)}{5\pi }0.1530,$$
(11)
here $`\zeta (x)`$ is the Riemann zeta function, and
$$\mathrm{\Delta }_{\mathrm{Casimir}}^u(0|1,1)=0.3157.$$
(12)
(ii)For the ”temporal” Casimir amplitudes in the case ($`0<\sigma 2`$)
$$\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{t}(\sigma ,\sigma )=\frac{16}{5\sigma }\frac{\zeta (3)}{(4\pi )^{\sigma /2}}\frac{1}{\mathrm{\Gamma }(\sigma /2)}.$$
(13)
Note that the defined ”temporal Casimir amplitude” $`\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{t}(\sigma ,\sigma )`$ reduces for $`\sigma =2`$ to the ”normal” Casimir amplitude $`\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{u}\left(0|2,2\right)`$, given by Eq. (11). This reflects the existence of a special symmetry in that case between the ”temporal” and the space dimensionalities of the system.
When $`\sigma 2`$ it is easy to verify that the following general relation
$$\frac{\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{t}(\sigma ,\sigma )}{\mathrm{\Delta }_{\mathrm{Casimir}}^\mathrm{t}(2,2)}=\frac{8\pi }{\sigma (4\pi )^{\sigma /2}\mathrm{\Gamma }(\sigma /2)}$$
(14)
between the temporal amplitudes holds. The r.h.s. of (14) is a decreasing function of $`\sigma `$.
### Relation with the Zamolochikov’s $`C`$-function
Let us note that if $`z=1`$ the temporal excess free energy introduced above coincides, up to a (negative) normalization factor, with the proposed by Neto and Fradkin definition of the non-zero temperature generalization of the $`C`$-function of Zamolodchikov (see e.g. Ref. ).
For $`z1`$ a straightforward generalization of this definition can be proposed at least in the case of long-range power-low decaying interaction
$$C(T,g|d,z)=T^{(1+d/z)}\frac{v^{d/z}}{n(d,z)}f_{\mathrm{ex}}^\mathrm{t}(T,g|d),$$
(15)
where $`z=\sigma /2`$, $`v=TL_\tau `$ and
$$n^t(d,\sigma )=\frac{4}{\sigma }\frac{\zeta \left(1+2d/\sigma \right)}{(4\pi )^{d/2}}\frac{\mathrm{\Gamma }(2d/\sigma )}{\mathrm{\Gamma }(d/2)}.$$
(16)
The quantity $`\stackrel{~}{c}_0(d,\sigma ):=C(T,g_c|d,z)`$ is an important universal characteristic of the theory. The behavior of $`\stackrel{~}{c}_0(d,\sigma )`$ is calculated numerically for dimensions between the lower critical dimension $`\sigma /2`$ and upper critical dimension $`3\sigma /2`$ for arbitrary values of $`0<\sigma 2`$. The results are universal as function of $`d/\sigma `$ as it is presented on Fig.1. In the particular case $`d/\sigma =1`$, one can obtain analytically
$$\stackrel{~}{c}_0(\sigma ,\sigma )=4/5.$$
(17)
This generalizes the result obtained for $`d=\sigma =2`$ to the case of long-range interaction.
To shed some light to what extend the amplitudes presented above are close to that one of more realistic models we present a comparison of the scaling functions of the excess free energy of the Ising, XY, Heisenberg and spherical model (limit $`n\mathrm{}`$) in FIG. 2. The results for the spherical model are exact while that ones for the Ising, XY and Heisenberg models are obtained by $`ϵ`$-expansion technique up to the first order in $`ϵ`$. The Monte Carlo results for he 3d Ising model give $`0.1526\pm 0.0010`$ , which is surprisingly close to the exact value (11). This makes difficult to resolve the question how $`X^{ex}/n`$ approaches the corresponding result for the spherical model when $`n\mathrm{}`$. Note that all the curves practically overlap for $`L>2\xi `$, where $`\xi `$ is the correlation length.
## Other amplitudes
Other important universal critical amplitudes, in finite-size scaling, depend upon the geometry $`L_{dd^{}}\times \mathrm{}^d^{}\times L_\tau ^z`$ as well as the range of the interaction. One of the most important quantities for a numerical analysis is the Binder’s cumulant ratio. For the quantum 2d spherical model with $`\sigma =2`$ at the critical point it is
$$B=\frac{2\pi }{\sqrt{5}\mathrm{ln}^3\tau }25.21657,$$
(18)
where $`\tau =(1+\sqrt{5})/2`$ is the ”golden mean” value.
In what follows we will list a number of results obtained in the framework of the quantum spherical model and the $`𝒪(n)`$ quantum $`\phi ^4`$ model .
(i) Finite system at zero temperature:
$$d=\sigma =1:\frac{L}{\xi }=0.624798\mathrm{for}d^{}=0.$$
(19)
$$d=\sigma =2:\frac{L}{\xi }=\{\begin{array}{cc}1.511955\hfill & \mathrm{for}d^{}=0,\hfill \\ 0.962424\hfill & \mathrm{for}d^{}=1.\hfill \end{array}$$
(20)
(ii) Bulk system at finite temperature:
$$d=\sigma :\frac{L_\tau }{\xi }=0.962424.$$
(21)
This result is a just a point in graph presented in FIG. 3, where we show the behaviour of $`L_\tau /\xi `$ as a universal function of the ratio $`d/\sigma `$. The point corresponding to $`(\frac{d}{\sigma }=1,y_0=0.962424)`$ can be obtained analytically .
The above result are obtained for the case when the quantum parameter controlling the phase transition is fixed at its critical value. Now we will present results obtained when the quantum parameter is fixed by “running” values corresponding the shifted critical quantum parameter. We are limited to the case $`d=\sigma =2`$
$$\frac{L}{\xi }=\{\begin{array}{cc}7.061132\hfill & \mathrm{for}d^{}=1,\hfill \\ 4.317795\hfill & \mathrm{for}d^{}=0,\hfill \end{array}$$
(22)
for finite system at zero temperature and
$$\frac{L_\tau }{\xi }=\{\begin{array}{cc}7.061132\hfill & \mathrm{for}d^{}=1,\hfill \\ 6.028966\hfill & \mathrm{for}d^{}=0,\hfill \end{array}$$
(23)
for the bulk system at finite temperature .
This work is supported by The Bulgarian Science Foundation (Projects F608/96 and MM603/96). |
warning/0001/math0001040.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Wess-Zumino-Novikov-Witten (WZNW) models provide important examples of a two-dimensional conformal field theory. They can roughly be described as follows. The gauge algebra of the theory is the affine algebra associated to a finite-dimensional gauge algebra (i.e. a simple finite-dimensional Lie algebra). The geometric data consists of a compact Riemann surface (with complex structure) of genus $`g`$ and a finite number of marked points on this surface. Starting from representations of the gauge algebra the space of conformal blocks can be defined. It depends on the geometric data. Varying the geometric data should yield a bundle over the moduli space of the geometric data. In Knizhnik and Zamolodchikov considered the case of genus 0 (i.e. the Riemann sphere). There, changing the geometric data consists in moving the marked points on the sphere. The space of conformal blocks could completely be found inside the part of the representation associated to the finite-dimensional gauge algebra. On this space an important set of equations, the Knizhnik-Zamolodchikov (KZ) equations, was introduced. In a general geometric setting, solutions are the flat sections of the bundle of conformal blocks over the moduli space with respect to the KZ connection.
For higher genus it is not possible to realize the space of conformal blocks inside the representation space associated to the finite-dimensional algebra. There exists different attacks to the generalization. Some of them add additional structure on these representation spaces (e.g. twists, representations of the fundamental group,..). Here I do not have the place to pay proper reference to all these approaches. Let me only give a few names: Bernard ,, Felder and Wieczerkowski , , Hitchin , and Ivanov .
An important approach very much in the spirit of the original Knizhnik-Zamolodchikov approach was given by Tsuchia, Ueno and Yamada . In Section 2 I will present a very short outline of their theory. The main point in their approach is that at the marked points, after choosing local coordinates, local constructions are done. In this setting the well-developed theory of representations of the traditional affine Lie algebras (Kac-Moody algebras of affine type) can be used. It appears a mixture between local and global objects and considerable effort is necessary to extend the local constructions to global ones.
Oleg Sheinman and myself propose a different approach to the WZNW models which uses consequently only global objects. These objects are the Krichever-Novikov (KN) algebras and their representations, respectively their multi-point generalizations given by me . An outline of this approach is presented in Section 3. In the remaining sections more details on the construction are given. Here I only want to point out that a subspace of the KN algebra of vector fields is identified with tangent directions on the moduli space of the geometric data. Conformal blocks can be defined. We are able to incorporate a richer theory because in our set-up we are able to deal with more general representations of the global algebras.
Finally, with the help of the global Sugawara construction for the higher genus and multi-point situation proven in , it is possible to define the higher genus multi-point Knizhnik-Zamolodchikov equations (see Definition 7.3). The complete proofs of the results appeared in . A detailed study of the KZ equations, resp. of the connection is work in progress .
## 2 Outline of the Tsuchya-Ueno-Yamada approach
Let me first recall a few steps of the approach of Tsuchia, Ueno and Yamada to the WZNW models. I will concentrate on the steps which are of relevance in our approach. More details can be found in , resp. in the more pedagogical introduction .
(1) A finite-dimensional complex simple Lie algebra $`𝔤`$ with $`(.|.)`$ a symmetric, nondegenerate, invariant bilinear form (e.g. the Killing form) is fixed. This Lie algebra is the finite-dimensional gauge algebra of the theory. Some authors prefer to call this algebra the horizontal algebra.
(2) Fix a natural number $`N`$ and a compact Riemann surface $`M`$ of arbitrary genus $`g`$ (resp. a smooth projective curve over $``$). Take a tuple of $`N`$ distinct points $`(P_1,P_2,\mathrm{},P_N)`$ on $`M`$ and around every such point $`P_i`$ a coordinate $`\xi _i`$. This defines the geometric data
$$\mathrm{{\rm Y}}^{}=(M,(P_1,P_2,\mathrm{},P_N),(\xi _1,\xi _2,\mathrm{},\xi _N)).$$
(1)
(3) The affine algebra (or Kac-Moody algebra of affine type) associated to $`𝔤`$ is given as
$$\widehat{𝔤}=𝔤[t^1,t]]c$$
(2)
with Lie structure
$$[xt^n,yt^m]=[x,y]t^{n+m}+(x|y)n\delta _m^nc,n,m,$$
(3)
$$[xt^n,c]=0,n.$$
(4)
Note that in this approach one has to consider Laurent series instead of the usual Laurent polynomials. In our approach we will return to Laurent polynomials. By ignoring the central element $`c`$ one sees that $`\overline{𝔤}=𝔤[t^1,t]]`$ carries also a Lie algebra structure. This algebra is called the loop algebra (associated to $`𝔤`$). There is a well-defined theory of highest weight representations $`H_\lambda `$ of the affine algebra $`\widehat{𝔤}`$ associated to a level and certain weights $`\lambda `$ of the finite dimensional Lie algebra $`𝔤`$, see . Recall that the central element $`c`$ operates as level$`\times id`$ on $`H_\lambda `$.
Another Lie algebra appearing in this context is the Virasoro algebra $`V`$ which is the Lie algebra with basis $`\{l_n,n\}\{c_1\}`$ and Lie structure
$$[l_n,l_m]=(mn)l_{n+m}+\delta _m^n\frac{n^3n}{12}c_1,n,m,$$
(5)
$$[l_n,c_1]=0,n.$$
(6)
Starting from a highest weight representation of the affine algebra the Sugawara construction defines also a representation of the Virasoro algebra on the representation space. Further down I will describe the Sugawara construction in a general geometric setting. This setting will incorporate also the classical Sugawara construction. Note that in the theory a graded structure is implicitly given and is employed. For example the Virasoro algebra becomes a graded affine Lie algebra by defining $`\mathrm{deg}(l_n):=n`$ and $`\mathrm{deg}(c_1):=0`$.
(4) After fixing the level globally, one assigns to every marked point $`P_i`$ a heighest weight representation $`H_{\lambda _i}`$ of the gauge algebra. Set $`\lambda :=(\lambda _1,\lambda _2,\mathrm{},\lambda _N)`$ and
$$H_\lambda :=H_{\lambda _1}H_{\lambda _2}\mathrm{}H_{\lambda _N}.$$
(7)
The space $`H_\lambda `$ is a representation space for the Lie algebra $`\widehat{𝔤_{(N)}}`$, where this algebra is defined as the one-dimensional central extension of $`N`$ copies of the loop algebra. Here the $`i`$-th copy of the loop algebra operates on the $`i`$-th factor in the tensor product and is associated to the $`i`$-th point $`P_i`$. To distinguish the different copies we use $`t_i`$ for the affine parameter $`t`$ in the loop algebra.
(5) Up to now the geometry was not really involved. This changes in this step. The affine parameter $`t_i`$ corresponding formally to $`P_i`$ via the assignment of the representation $`H_{\lambda _i}`$ to this point, will be identified with the coordinate $`\xi _i`$ at this point. One considers the algebra $`A(\mathrm{{\rm Y}}^{})`$ of meromorphic functions on $`M`$ which have poles at most at the points $`\{P_1,P_2,\mathrm{},P_N\}`$, and sets $`𝔤(\mathrm{{\rm Y}}^{}):=𝔤A(\mathrm{{\rm Y}}^{})`$. This algebra is also called the block algebra. By taking the Laurent expansion of $`fA(\mathrm{{\rm Y}}^{})`$ at the point $`P_i`$ with respect to the coordinate $`\xi _i`$ there, we get an embedding $`𝔤(\mathrm{{\rm Y}}^{})\overline{𝔤}_N`$ by assigning to $`f`$ the corresponding Laurent series in the affine parameter $`t_i`$ in the $`i`$-th copy of $`\overline{𝔤}`$. The cocycle defining the central extension vanishes on $`𝔤(\mathrm{{\rm Y}}^{})`$. Hence it can be considered as a subalgebra of $`\widehat{𝔤_{(N)}}`$.
(6) The space of conformal blocks (also called chiral blocks) are defined as the coinvariants
$$V_\lambda :=H_\lambda /𝔤(\mathrm{{\rm Y}}^{})H_\lambda .$$
(8)
In some context is is better to work with the dual objects $`V_\lambda ^{}=\mathrm{Hom}_{}(V_\lambda ,)`$. This space can be described as the space of linear forms on $`H_\lambda `$ vanishing on $`𝔤(\mathrm{{\rm Y}}^{})H_\lambda `$. The vector spaces $`V_\lambda `$ turn out to be finite-dimensional. Their dimension is given by the Verlinde formula.
(7) One of the motivations of was to supply a proof of the Verlinde formula. For this the authors pass to the moduli space $`_{g,N}^{\mathrm{}}`$ of the data $`\mathrm{{\rm Y}}^{}`$ (with the obvious identifications under isomorphisms). Clearly, not only the goal to proof the Verlinde formula forces one to consider moduli spaces but also the concept of quantization requires to consider all possible configurations. Now everything has to be sheavified. One obtains the sheaf of conformal blocks over moduli space. On this sheaf the Knizhnik-Zamolodchikov connection is constructed. The sheaf is indeed a vector bundle, hence the dimension of the conformal blocks will be constant along the moduli space, only depending on the genus $`g`$, the number $`N`$ of marked points, and the associated weights $`\lambda `$ (of course the Lie algebra $`𝔤`$ and the level will be fixed). The construction of the connection involves the Sugawara construction which is done locally on the Riemann surface. By introducing a projective connection on the Riemann surface it is possible to show that it globalizes.
(8) It is shown that the essential data can be given in terms of the moduli space $`_{g,N}`$ of smooth projective curves of genus $`g`$ with $`N`$ marked points, by forgetting the coordinate systems at the marked points.
(9) Finally, by passing to the Deligne-Mumford-Knudsen boundary of $`_{g,N}`$ corresponding to stable singular curves and by proving factorization rules (e.g. behaviour of the conformal blocks under normalization of the singular curves) it is possible to express the dimension of the space of conformal blocks for the $`(g,N,\lambda )`$ situation in dimensions of spaces of conformal blocks for lower genera.
## 3 Outline of the global operator approach
The approach presented in Section 2 is very successful and a beautiful piece of mathematics. Nevertheless it has some problems.
(a) It is necessary to choose coordinates around the points $`P_i`$. The authors have to work over infinite-dimensional moduli spaces. Finally, a complicated reduction process is needed to reduce the relevant data to the moduli space of curves with marked points. Note that by Tsuchimoto some simplifications with respect to this problem has been given.
(b) The proof that the local Sugawara construction globalizes (with the help of a projective connection) on the Riemann surface is difficult.
(c) Some objects of the theory are only defined locally, other objects globally.
(d) The representation $`H_\lambda `$ of the affine algebra $`\widehat{𝔤_{(N)}}`$ does not see the geometry.
We (Oleg Sheinman and myself) propose a different approach to the WZNW models which still stays in the algebraic-geometric setup. In this section I will outline what has been done so far. Some more details will be given in the following sections. First note that the geometric data we consider consist of a compact Riemann surface $`M`$ (or a smooth projective curve over $``$ – I will use the terms interchangeable) and $`N`$ (distinct) marked points:
$$\mathrm{{\rm Y}}=(M,(P_1,P_2,\mathrm{},P_N)).$$
(9)
In contrast to the data (1) it does not contain coordinates. Again, we fix a simple finite-dimensional complex Lie algebra $`𝔤`$.
(1) We replace all local algebras by algebras of Krichever-Novikov(KN)-type and their multi-point generalizations. They consists of meromorphic objects which might have poles only at the points $`\{P_1,P_2,\mathrm{},P_N\}`$ and a fixed reference point $`P_{\mathrm{}}`$. This can be done for the function algebra, the affine algebra, the vector field algebra, the differential operator algebra. Central extensions can be defined. They carry an almost-graded structure (see the definition below) induced by the vanishing order at $`\{P_1,P_2,\mathrm{},P_N\}`$, resp. at $`P_{\mathrm{}}`$. In particular, all these algebras can be decomposed into subalgebras of elements of positive degree corresponding roughly to elements holomorphic at $`\{P_1,P_2,\mathrm{},P_N\}`$ and subalgebras of elements of negative degree corresponding roughly to elements holomorphic at $`P_{\mathrm{}}`$ and a “critical” finite-dimensional subspace spanned by other elements.
(2) In the next step we consider highest weight representations of the higher genus affine algebras. Examples are given again by representations $`\widehat{W}_\lambda `$ which are constructed from weights assigned to the marked points. For highest weight representations the Sugawara construction works and gives a representation of the centrally extended vector field algebra (see , ).
(3) The objects make direct sense over a dense subset of $`_{g,N}`$. We obtain sheaf versions of our objects.
(4) Note that the block algebra $`𝔤(\mathrm{{\rm Y}})`$ considered above is naturally a subalgebra of our higher genus multi-point affine algebra. There is no need to choose coordinates. It is possible to define again
$$\widehat{V}_\lambda :=\widehat{W}_\lambda /𝔤(\mathrm{{\rm Y}})\widehat{W}_\lambda .$$
(10)
as the space of conformal blocks.
(5) At a fixed point $`b`$ of the moduli space $`_{g,N}`$ the basis elements $`l_k`$ of the “critical subspace” of the vector field algebra can be identified with the tangent directions $`X_k`$ to the moduli space at $`b`$. Additionally the elements operate as Sugawara operators $`T[l_k]`$ on the representations $`\widehat{W}_\lambda `$ defined above $`b`$. This allows us to define the formal Knizhnik-Zamolodchikov(KZ) equations to be
$$(_k+T[l_k])\mathrm{\Phi }=0,k=1,\mathrm{},3g3+N.$$
(11)
Here $`\mathrm{\Phi }`$ is assumed to be a section of $`\widehat{W}_\lambda `$ , resp. of $`\widehat{V}_\lambda `$, and $`_k`$ is an action of $`X_k`$ on $`\widehat{W}_\lambda `$ which fulfills suitable conditions.
In the following sections I will give additional information on these steps. Details appeared in . Further work is in progress . In particular what has to be done is to construct a projectively flat connection on the sheaf of conformal blocks and to extend the construction to the whole moduli space. Especially it should be extended its boundary to obtain again a proof of the Verlinde formula.
## 4 The Krichever-Novikov objects
Let $`M`$ be a compact Riemann surface of genus $`g`$. Let $`A`$ be a fixed set of finitely many points on $`M`$ which is splitted into two non-empty subsets $`I`$ and $`O`$. This is the general situation dealt with in . Here it is enough to consider
$$I:=\{P_1,P_2,\mathrm{},P_N\},O:=\{P_{\mathrm{}}\}.$$
(12)
The classical (Virasoro) situation is: $`M=^1`$ with $`I=\{z=0\}`$ and $`O=\{z=\mathrm{}\}`$. The Krichever-Novikov situation is: $`M`$ an arbitrary Riemann surface with $`I=\{P_+\}`$ and $`O=\{P_{}\}`$.
Let $`𝒦`$ be the canonical bundle, i.e. the bundle whose local sections are the holomorphic 1-differentials, and $`𝒦^\lambda `$ (for $`\lambda `$) its tensor powers, resp. for $`\lambda `$ negative, the tensor power of its dual. Denote by $`F^\lambda (A)`$ the space of meromorphic sections of $`𝒦^\lambda `$ consisting of those elements which are holomorphic outside of $`A`$. We set $`𝒜(A):=F^0(A)`$ for the associative algebra of functions and $`(A):=F^1(A)`$ for the Lie algebra of vector fields (with the usual Lie bracket as Lie structure). If the set $`A`$ is clear from the context, we will drop it in the notation. By multiplication with the elements of $`𝒜`$ the space $`F^\lambda `$ becomes a module over $`𝒜`$. By taking the Lie derivative with respect to the vector fields it becomes also a Lie module over $``$.
###### Theorem 4.1.
(, ) There exists a decomposition
$$F^\lambda =\underset{n}{}F_n^\lambda ,$$
(13)
where the $`F_n^\lambda `$ are subspaces of dimension $`N`$. By defining the elements of $`F_n^\lambda `$ to be the homogeneous elements of degree $`n`$, the algebras $`𝒜`$ and $``$ are almost-graded (Lie) algebras and the $`F_n^\lambda `$ are almost-graded (Lie) modules over $`𝒜`$, respectively over $``$.
For the convenience of the reader let me recall the definition of an almost-grading at the example of the Lie algebra $``$. A Lie algebra $``$ is called almost-graded if it can be decomposed (as vector space)
$$=\underset{n}{}_n,dim_n<\mathrm{},$$
(14)
such that there exists $`L,K`$ with
$$[_n,_m]\underset{h=n+mL}{\overset{n+m+K}{}}_h,n,m.$$
(15)
This definition (suitable modified) works for the associative algebra $`𝒜`$ and the modules $`F^\lambda `$. In our situation we always can do with lower shift $`L=0`$.
The almost-grading is important to obtain a decomposition of the algebras into positive and negative parts and to define the concept of highest weight representations (see below). It is also necessary to obtain an embedding of the algebras $`𝒜`$ and $``$ (and more generally also of $`𝒟`$, the Lie algebra of differential operators, see ) into $`\overline{gl}(\mathrm{})`$ via their action on the modules $`F^\lambda `$. Here $`\overline{gl}(\mathrm{})`$ denotes the algebra of (both-sided) infinite matrices with finitely many diagonals.
The degree is introduced by the order of $`fF^\lambda `$ at $`I`$. More precisely, $`F_n^\lambda `$ is given by exhibiting a basis $`f_{n,p}^\lambda ,p=1,\mathrm{},N`$. For the generic situation ($`\lambda 0,1`$, $`g1`$, and a generic choice of $`A`$) the element is fixed up to multiplication with a scalar by
$$\mathrm{ord}_{P_i}(f_{n,p}^\lambda )=\{\begin{array}{cc}n\lambda ,\hfill & i=p\hfill \\ n\lambda +1,\hfill & ip,i=1,\mathrm{},N\hfill \\ N(n+1\lambda )+(2\lambda 1)(g1),\hfill & i=\mathrm{}.\hfill \end{array}$$
(16)
For the non-generic situation for finitely many $`n`$ the prescription at $`i=\mathrm{}`$ has to be adjusted. A detailed prescription is given in and . To fix the element $`f_{n,p}^\lambda `$ uniquely it is necessary to fix a coordinate $`\xi _i`$ centered at $`P_i`$, for $`i=1,\mathrm{},N`$ and to require $`f_{n,p}^\lambda (\xi _p)_|=\xi _p^{n\lambda }(1+O(\xi _p))d\xi _p^\lambda `$. Note that the fixing does not really depend on the full information in the coordinate. It depends only on the first infinitesimal neighbourhood. In detail: Let $`\xi _p^{}=a_1\xi _p+_{j2}a_j\xi _p^j`$ (clearly with $`a_10`$) be another coordinate centered at $`P_p`$ then the normalization will only depend on the value $`a_1`$.
We have the important duality
$$F^\lambda \times F^{1\lambda },f,g=\frac{1}{2\pi \mathrm{i}}_Cfg=\mathrm{res}_P_{\mathrm{}}(fg),$$
(17)
where $`C`$ is any separating cycle for $`(P_1,P_2,\mathrm{},P_N)`$ and $`P_{\mathrm{}}`$ which is cohomologous to (-1) $`\times `$ the circle around $`P_{\mathrm{}}`$. For the scaled basis elements we obtain
$$f_{n,p}^\lambda ,f_{m,r}^{1\lambda }=\delta _n^m\delta _p^r.$$
(18)
Again note that the duality relation (18) will not depend on the coordinates chosen. For special values of $`\lambda `$ we set
$$A_{n,p}:=f_{n,p}^0,e_{n,p}:=f_{n,p}^1,\omega ^{n,p}:=f_{n,p}^1,\mathrm{\Omega }^{n,p}:=f_{n,p}^2.$$
(19)
In particular $`(A_{n,p},\omega ^{n,p})`$ and $`(e_{n,p},\mathrm{\Omega }^{n,p})`$ are dual systems of basis elements.
Next we decompose our algebras. First we consider $`𝒜`$, the algebra of functions:
$$𝒜=𝒜_{}^{}𝒜_{(0)}^{}𝒜_+,$$
$$𝒜_{}^{}=\underset{nK1}{}𝒜_n,𝒜_+=\underset{n1}{}𝒜_n,𝒜_{(0)}^{}=A_{n,p}1pN,Kn0.$$
Here $`𝒜_n`$ is the subspace of $`𝒜`$ consisting of the elements of degree $`n`$ and $`K`$ is the constant appearing in the definition of the almost-grading (15) for the algebra $`𝒜`$. The constant depends on the genus $`g`$ and the number of points $`N`$. From the almost-grading it follows that $`𝒜_+`$ and $`𝒜_{}^{}`$ are subalgebras. In general $`𝒜_{(0)}^{}`$ is only a subspace. For our purpose here it is more convenient to take as $`𝒜_{}𝒜_{}^{}`$ the subalgebra of meromorphic functions vanishing at $`P_{\mathrm{}}`$, make a change of basis in the lowest degree part $`𝒜_K`$ in $`𝒜_{(0)}`$, and take $`𝒜_{(0)}𝒜_{(0)}^{}`$ correspondingly smaller. This yields the decomposition
$$𝒜=𝒜_{}𝒜_{(0)}𝒜_+.$$
(20)
A completely analogous decomposition exists for the vector field algebra $``$
$$=_{}_{(0)}_+.$$
(21)
The vector fields in $`_+`$ are vanishing of order 2 at the points $`\{P_1,P_2,\mathrm{},P_N\}`$, and the vector fields in $`_{}`$ are vanishing of order 2 at the point $`P_{\mathrm{}}`$. The space $`_{(0)}`$ is $`(3g3+2N+2)`$-dimensional. We call it the “critical strip”.
For further reference let me identify the basis elements of the critical strip. First we have $`N`$ elements $`e_{0,p},p=1,\mathrm{},N`$ vanishing at all $`P_i,i=1,\mathrm{},N`$, but not vanishing of 2nd order at every point. Next we have $`N`$ elements $`e_{1,p},p=1,\mathrm{},N`$ regular at all $`P_i,i=1,\mathrm{},N`$, but not vanishing at every point. Later we will see that they will be responsible for moving the points $`P_p,p=1,\mathrm{},N`$. In the middle we have $`3g3`$ elements which are neither regular at $`\{P_1,P_2,\mathrm{},P_N\}`$ nor at $`P_{\mathrm{}}`$. They will correspond to deformations of the complex structure of $`M`$. Finally we have two elements regular of order 1, resp. of order 0 at $`P_{\mathrm{}}`$.
For the algebras we can construct central extensions $`\widehat{𝒜}`$ and $`\widehat{}_R`$ defined via the following geometric cocycles
$`𝒜:`$ $`\gamma (g,h):={\displaystyle \frac{1}{2\pi \mathrm{i}}}{\displaystyle _C}g𝑑h,`$ (22)
$`:`$ $`\chi _R(E,f):={\displaystyle \frac{1}{24\pi \mathrm{i}}}{\displaystyle _C}\left({\displaystyle \frac{1}{2}}\left(e^{\prime \prime \prime }fef^{\prime \prime \prime }\right)R\left(e^{}fef^{}\right)\right)𝑑z.`$ (23)
Here $`R`$ is a projective connection on $`M`$. The cocycles are local in the sense that there exists $`T`$ and $`S`$ such that for all $`n,m`$: $`\gamma (𝒜_n,𝒜_m)0Tm+n0`$ and $`\chi _R(_n,_m)0Sm+n0`$. Again the $`T`$ and $`S`$ can be explicitly calculated . Restricted to the $`(+)`$ and $`()`$ subalgebras in (20) and (21) the cocycles are vanishing. In particular, the subalgebras can be considered in a natural way as subalgebras of the central extensions $`\widehat{𝒜}`$ and $`\widehat{}_R`$. By setting $`\mathrm{deg}(t):=0`$ for the central elements the almost-grading can be extended to the central extensions.
Different connections $`R`$ yield cohomologous cocycles $`\chi _R`$, hence equivalent central extensions (see for details). In the classical situation of $`^1`$ with two marked points one obtains the Virasoro algebra.
## 5 The affine multi-point Krichever-Novikov algebras
Let $`𝔤`$ be a finite-dimensional reductive Lie algebra (e.g. semi-simple or abelian) and $`(.|.)`$ a symmetric, nondegenerate invariant bilinear form on $`𝔤`$. The (higher genus multi-point) loop algebra is defined to be
$$\overline{𝔤}:=𝔤𝒜\text{with}[xg,yh]:=[x,y](gh).$$
(24)
Its elements can be considered as $`𝔤`$-valued meromorphic functions on $`M`$ which are holomorphic outside of $`\{P_1,P_2,\mathrm{},P_N,P_{\mathrm{}}\}`$. A central extension $`\widehat{𝔤}`$ is given as vector space $`\widehat{𝔤}=\overline{𝔤}`$ with Lie structure ($`\widehat{x}:=(0,x),t:=(1,0)`$)
$$[\widehat{xf},\widehat{yg}]:=\widehat{[x,y](fg)}(x,|y)\gamma (f,g)t,[t,\widehat{𝔤}]=0,$$
(25)
and $`\gamma `$ from (22). We call this algebra higher genus multi-point affine algebra. For $`g=0`$ and 2 points the classical affine algebras are obtained (with respect to Laurent polynomials). For higher genus and two points it was introduced by Krichever-Novikov and extensively studied by Sheinman ,,. Its generalization to higher genus and an arbitrary number of marked points was given in .
By the decomposition (20) of $`𝒜`$ we obtain a decomposition
$$\widehat{𝔤}=\widehat{𝔤}_{}\widehat{𝔤}_{(0)}\widehat{𝔤}_+.$$
(26)
In particular, $`\widehat{𝔤}_{(0)}=𝒜_{(0)}𝔤t`$. Using the duality (18) we see that $`1=_{p=1}^NA_{0,p}`$ and hence that via $`𝔤𝔤1\widehat{𝔤}`$ the finite dimensional Lie algebra $`𝔤`$ can be naturally considered as subalgebra of $`\widehat{𝔤}_0\widehat{𝔤}_+`$, of $`\overline{𝔤}`$, and of $`\widehat{𝔤}`$.
## 6 Verma modules and Sugawara construction
In this section we assume $`𝔤`$ to be a simple complex Lie algebra. Choose $`𝔥𝔤`$ a Cartan subalgebra, $`𝔟`$ a corresponding Borel subalgebra, and $`𝔫`$ a corresponding upper nilpotent subalgebra. We choose $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _N),\lambda _p𝔥^{}`$, take $`N`$ copies of everything, and label them by $`p=1,\mathrm{},N`$. We set
$$𝔤_{(N)}:=𝔤_1\mathrm{}𝔤_N,𝔟_{(N)}:=𝔟_1\mathrm{}𝔟_N.$$
(27)
Let $`W_p`$ be a one-dimensional vector space over $``$ with basis $`w_p`$ for $`p=1,\mathrm{},N`$. On $`W_p`$ a one-dimensional representation of $`𝔟_p`$ is defined via
$$h_p\omega _w:=\lambda _p(h_p)w_p,n_pw_p:=0,h_p𝔥_p,n_p𝔫_p.$$
(28)
This yields a representation of $`𝔟_{(N)}`$ on $`W=_{p=1}^NW_p`$. As usual the Lie algebra acts via “Leibniz rule” on the tensor product, and on the $`p`$-th tensor factor only the $`p`$-th summand operates non-trivially. We call this representation $`\tau _\lambda `$.
###### Remark.
It is possible to allow twists to get a richer theory. They can be obtained via automorphisms $`\phi _p:𝔤𝔤=𝔤_p`$ given by $`\phi _p=Ad\gamma _p,\gamma _pG`$, where $`G`$ is the Lie group associated to $`𝔤`$.
###### Lemma 6.1.
The subspace $`\overline{𝔤}_0t\widehat{𝔤}_+`$ is a Lie subalgebra of $`\widehat{𝔤}`$ and
$`\psi :\overline{𝔤}_0t\widehat{𝔤}_+\widehat{𝔤_{(N)}}`$ defined by
$$\psi (\underset{p=1}{\overset{N}{}}x_pA_{0,p}):=(x_1,\mathrm{},x_N),\psi (t):=0,\psi (\widehat{𝔤}_+):=0$$
(29)
is a Lie homomorphism.
###### Proof.
We calculate
$$[x_pA_{0,p},y_rA_{0,r}]=[x_p,y_r](A_{0,p}A_{0,r})=[x_p,y_r](A_{0,r}\delta _p^r+\underset{h>0}{}\underset{s=1}{\overset{N}{}}\alpha _{(0,p),(0,r)}^{(h,s)}A_{h,s}),$$
with $`\alpha _{(0,p),(0,r)}^{(h,s)}`$. First, this shows that the subspace is indeed closed under the Lie bracket. Note that the defining cocycle for the central extension vanishes for the subalgebra. Next, we see that all mixing terms are of higher degree. They will be annulated under $`\psi `$. Only the term with $`p=r`$ will survive. This shows that $`\psi `$ is a Lie homomorphism. ∎
Denote by $`\widehat{𝔟}`$ the subalgebra $`\overline{𝔟}_0t\widehat{𝔤}_+`$ (recall $`\overline{𝔟}_0=𝔟𝒜_0`$). By restricting $`\psi `$ to $`\widehat{𝔟}`$ we obtain from the lemma that $`\psi :\widehat{𝔟}𝔟_{(N)}`$ is a Lie homomorphism. Now we can define for $`\delta `$ a representation of $`\widehat{𝔟}`$ on $`W`$ by
$$\tau _{\lambda ,\delta }(x\alpha tx_+):=\tau _\lambda (\psi (x))+\alpha \delta id.$$
(30)
###### Definition 6.2.
The linear space
$$\widehat{W}_{\lambda ,\delta }:=U(\widehat{𝔤})_{U(\widehat{𝔟})}W$$
(31)
with its natural structure of a $`\widehat{𝔤}`$-module is called the Verma module of the Lie algebra $`\widehat{𝔤}`$ corresponding to $`(\lambda ,\delta )`$ where $`\lambda (𝔥^{})^N`$ is the weight of the Verma module and $`\delta `$ is the level of the Verma module.
As usual $`U(\widehat{𝔤})`$ denotes the universal enveloping algebra.
Now the question arises can the finite-dimensional representations of $`𝔤`$ be recovered. The answer is yes. It is given by
###### Proposition 6.3.
The $`\widehat{𝔤}`$-module $`\widehat{W}_{\lambda ,\delta }`$ is under the natural embedding of $`𝔤`$ into $`\widehat{𝔤}`$ also a $`𝔤`$-module and contains the $`𝔤`$-module
$$\stackrel{~}{W}_\lambda =\stackrel{~}{W}_{\lambda _1}\mathrm{}\stackrel{~}{W}_{\lambda _N},$$
(32)
where the $`\stackrel{~}{W}_{\lambda _i}`$ are the heighest weight modules of $`𝔤`$ of weight $`\lambda _i`$.
The proof is straightforward. The module $`\stackrel{~}{W}_\lambda `$ lies in the degree zero part of $`\widehat{W}_{\lambda ,\delta }`$, but there are other elements of degree zero not corresponding to $`\stackrel{~}{W}_\lambda `$.
Denote by $`\widehat{𝔤}_{}^{}=\overline{𝔤}_{}^{}\widehat{𝔤}`$ the subalgebra of $`𝔤`$-valued meromorphic functions which are regular at $`P_{\mathrm{}}`$, then we can define the space of conformal blocks as the space of coinvariants
$$\widehat{V}_{\lambda ,\delta }:=\widehat{W}_{\lambda ,\delta }/\widehat{𝔤}_{}^{}\widehat{W}_{\lambda ,\delta }.$$
(33)
Verma modules are examples of admissible modules $`\widehat{W}`$ of $`\widehat{𝔤}`$ in the following sense:
1. the central element $`t`$ operates as $`cid`$ with $`c`$, ($`c`$ is called the level),
2. for all $`w\widehat{W}`$ we have $`\widehat{𝔤}_n.w=0`$, for all $`n0`$.
For admissible modules we showed (see also ) that there exists a Sugawara construction which yields a representation of the centrally extended vector field algebra. Here I can only give the main ideas. Choose $`(u_i)_{i=1,\mathrm{},dim𝔤}`$ and $`(u^i)_{i=1,\mathrm{},dim𝔤}`$ a system of dual basis of $`𝔤`$ with respect to the bilinear form $`(.|.)`$. The current associated to the basis element $`u_i`$ is given as
$$J_i(Q)=\underset{n}{}\underset{p=1}{\overset{N}{}}u_i(n,p)\omega ^{n,p}(Q),QM.$$
(34)
Here we used $`u_i(n,p)`$ to denote the operator corresponding to $`u_iA_{n,p}`$ on $`\widehat{W}`$. The current is a formal operator-valued $`1`$-differential on $`M`$. The similar expression is used for the current $`J^i(Q)`$ associated to $`u^i`$, resp. for the current associated to an arbitrary element $`x𝔤`$. The energy-momentum tensor is defined as the formal sum
$$T(Q):=\frac{1}{2}\underset{i=1}{\overset{dim𝔤}{}}:J_i(Q)J^i(Q):,QM.$$
(35)
In this expression $`:\mathrm{}:`$ denotes some normal ordering, which moves the positive degree elements to the right. Using the admissibility and the normal ordering we can conclude that the energy momentum tensor is indeed a well-defined formal operator-valued $`2`$-differential on $`M`$ written as
$$T(Q)=\underset{k}{}\underset{r=1}{\overset{N}{}}L(k,r)\mathrm{\Omega }^{k,r}(Q),QM.$$
(36)
Let $`k^{}`$ be the dual Coxeter number (i.e. $`2k^{}`$ is the eigenvalue of the Casimir operator in the adjoint representation), then
###### Theorem 6.4.
Assume $`𝔤`$ to be simple or abelian and $`k^{}`$ the dual Coxeter number, resp. $`k^{}=0`$ for the abelian case. Assume $`\widehat{W}`$ to be an admissible module of $`\widehat{𝔤}`$ with level $`c`$, then the $`L_{k,r}`$ are well-defined operators on $`\widehat{W}`$. If $`c+k^{}0`$ then the rescaled operators
$$L_{k,r}^{}=\frac{1}{c+k^{}}L_{k,r}$$
(37)
define a representation of a centrally extended vector field algebra $`\widehat{}`$.
It is shown in that a different normal ordering yields an equivalent central extension.
## 7 Moduli spaces and formal <br>Knizhnik-Zamolodchikov equations
The construction done in the previous sections globalizes over a dense subset of the moduli space $`_{g,N}`$ of smooth projective curves of genus $`g`$ with $`N`$ marked points. Recall that two configurations $`m=(M,(P_1,\mathrm{},P_N))`$ and $`m^{}=(M^{},(P_1^{},\mathrm{},P_N^{}))`$ describe the same point in moduli space if there is an algebraic isomorphism $`\phi :MM^{}`$ which maps $`P_iP_i^{}`$ for $`i=1,\mathrm{},N`$. “Bad” points in moduli space correspond to configurations which admit nontrivial automorphisms. At such points singularities can occur. For $`N1`$ and $`g2`$ the generic moduli point does not admit any nontrivial automorphism. Over the locus $`Y`$ of these points there exists a universal family of curves with marked points (a versal family would suffice). The situation for $`g=0`$ and $`g=1`$ is a little bit different. Here one should better work with the configuration space. This situation will not be covered here, see .
Beside $`_{g,N}`$ we have also to consider $`_{g,N+1}`$. By forgetting the last marked point we obtain a morphism $`f:_{g,N+1}_{g,N}`$. Denote by $`Y`$ the dense open subset of $`_{g,N}`$, where the universal family exists and set $`\stackrel{~}{Y}=f^1(Y)`$ then there exists also a universal family over $`\stackrel{~}{Y}`$. Recall that a universal family over $`\stackrel{~}{Y}`$ consists of a suitable well-behaved morphism (i.e. proper and flat) $`\pi :𝒰\stackrel{~}{Y}`$ for
$`\stackrel{~}{b}=[(M,P_1,P_2,\mathrm{},P_N,P_{\mathrm{}})]\stackrel{~}{Y}`$ with $`\pi ^1(\stackrel{~}{b})M`$, and $`N+1`$ sections
$$\sigma _i:\stackrel{~}{Y}𝒰,\text{with}\sigma _i(\stackrel{~}{b})=P_i\pi ^1(\stackrel{~}{b}),i=1,\mathrm{},N,\mathrm{}.$$
(38)
Note that for $`g2`$ the family $`𝒰`$ can be obtained by pulling back the universal curve defined over an open dense subset of $`_{g,0}`$ via the morphism corresponding to forgetting the points. For $`g=2`$ we have to start with $`_{2,1}`$.
Now we fix a reference section $`\widehat{\sigma }_{\mathrm{}}`$. This corresponds to choosing a reference point for every $`M`$ depending algebraically on varying $`M`$. Inside of $`\stackrel{~}{Y}`$ we have the subset
$$Y^{}=\{[(M,P_1,P_2,\mathrm{},P_N,P_{\mathrm{}})]\stackrel{~}{Y}P_{\mathrm{}}=\widehat{\sigma }_{\mathrm{}}(M)\}.$$
(39)
By genericity the forget morphism restricted to $`Y^{}`$ is 1:1 onto $`Y`$. We will identify in the following $`Y^{}`$ with $`Y`$. But note that this identification will not be canonical. It depends on the reference section chosen.
Let $`\pi :𝒰Y`$ be the projection. For every open subset $`U`$ of $`Y`$ set
$$\stackrel{~}{S}_U:=\underset{i=1}{\overset{N}{}}\sigma _i(U)+\widehat{\sigma }_{\mathrm{}}(U)$$
(40)
for the divisor of sections on $`𝒰_{|\pi ^1(U)}`$. The sheaf $`𝒜_Y`$ is defined as the sheaf over $`Y`$ given by defining $`𝒜_Y(U)`$ to be the space of functions on $`\pi ^1(U)`$ with possible poles along $`\stackrel{~}{S}_U`$. The same works for the central extension $`\widehat{𝒜}_Y`$, the loop algebra $`\overline{𝔤}_Y=𝒜_Y𝔤`$ and the affine algebra $`\widehat{𝔤}_Y=\overline{𝔤}_Y𝒪_Yt`$. Here $`𝒪_Y`$ denotes the structure sheaf of $`Y`$. The central extensions are given in a natural globalization of (22), resp. (25).
###### Definition 7.1.
A sheaf $`𝒲`$ of $`𝒪_Y`$-modules is called a sheaf of representations for the affine algebra sheaf $`\widehat{𝔤}_Y`$ if the $`𝒲(U)`$ are modules over $`\widehat{𝔤}_Y(U)`$ for every $`U`$ fulfilling the obvious compatibility conditions on the sheaf level.
Admissible representation sheaves are defined in a similar way. The Verma module construction can be made on the sheaf level. This yields admissible representation sheaves $`\widehat{W}_{(\lambda ,\delta ),Y}`$. The sheaf of conformal blocks is the quotient sheaf
$$\widehat{V}_{(\lambda ,\delta ),Y}:=\widehat{W}_{(\lambda ,\delta ),Y}/\widehat{𝔤}_{,Y}^{}\widehat{W}_{(\lambda ,\delta ),Y}.$$
(41)
Next we want to describe the tangent space at a moduli point.
Let $`b=[(M,P_1,P_2,\mathrm{},P_N)]`$ be the moduli point and set $`S=_{i=1}^NP_i`$, then the Kodaira-Spencer map gives an isomorphism of the tangent space to the moduli space with certain cohomology spaces
$$\mathrm{T}_{[M]}_{g,0}\mathrm{H}^1(M,T_M),\mathrm{T}_b_{g,N}\mathrm{H}^1(M,T_M(S)).$$
(42)
Here $`T_M`$ is the holomorphic tangent line bundle on $`M`$, i.e. $`T_M=𝒦^1`$. The following theorem proven in gives an identification of the tangent space at the moduli point $`b`$ with a certain subspace of the critical strip of the vector field algebra associated to the geometric data $`(M,(P_1,\mathrm{},P_N,\widehat{\sigma }_{\mathrm{}}(M))`$.
###### Theorem 7.2.
There exists a natural isomorphism
$$(_{k2}_{k3}\mathrm{}_1)_{(0)}^{}\mathrm{H}^1(M,T_M(kS)),k0.$$
(43)
Here $`_{(0)}^{}`$ are the elements of the reduced critical strip generated by the basis elements with poles at $`\{P_1,P_2,\mathrm{},P_N,P_{\mathrm{}}\}`$. For $`g2`$ its dimension is $`3g3`$ and it starts with $`_2`$ from above.
Let me only indicate the construction of these isomorphisms. It is based on the calculation of the Cech cohomology of $`T_M(kS)`$ with respect to the affine (resp. Stein) covering of $`M`$ given as follows. Let $`U_{\mathrm{}}`$ be a disk around $`P_{\mathrm{}}`$ and set $`U_1:=MS`$. Note that $`U_1`$ is affine. Then $`U_1U_{\mathrm{}}=U_{\mathrm{}}^{}`$ is the disc with $`P_{\mathrm{}}`$ removed. Hence the Cech 1-cocycles can be given as sections on $`U_{\mathrm{}}^{}`$. In this way $`ff_{|U_{\mathrm{}}^{}}`$ gives a linear map from the vector field algebra to the Cech 1-cocycles, and further to the cohomology group $`\mathrm{H}^1(M,T_M(kS))`$. The restrictions of elements coming from outside of the strip given in the formulation of the theorem can by their very definition be extended either to $`U_{\mathrm{}}`$ or to $`U_1`$ with the required zero-order at $`S`$. Hence they are coboundaries. A closer examination shows that the basis of the rest stays linearly independent in cohomology. By dimension count it follows that the map is an isomorphism.
To give an example: for the element $`e_{1,p}`$ (with $`p=1,\mathrm{},N`$) the lowest order term of its expansion at the point $`P_p`$ has the form $`/\xi _p`$. At all other marked points it has a zero. Under the isomorphism it corresponds to moving the marked point $`P_p`$ on $`M`$.
Now take $`X_k\mathrm{T}_b_{g,N}_{(0)}^{}_1`$ a tangent vector corresponding to an element $`l_k`$ of the critical strip. We assume that $`X_k`$ operates linearly as operator $`_k`$ on the space of sections of a representation sheaf $`𝒲`$. For $`\mathrm{\Phi }`$ a section we set
$$_k\mathrm{\Phi }:=(_k+T[l_k])\mathrm{\Phi },\text{with}T[l]:=\frac{1}{(c+k^{})\mathrm{\hspace{0.17em}2}\pi \mathrm{i}}_CT(Q)l(Q)$$
(44)
the Sugawara operator, which corresponds to $`l`$.
###### Definition 7.3.
The formal Knizhnik-Zamolodchikov equations are defined to be the set of equations
$$_k\mathrm{\Phi }=0,k=1,\mathrm{},3g3+N.$$
(45)
Note that these equations can be expressed in terms of the geometric data of the curve and the points which can be moved. In the equations have been explicitely calculated for genus 0 and 1.
## Acknowledgments
First, let me point out that the new results presented here are joint work with Oleg K. Sheinman. I would like to thank him for fruitful cooperation continuing now over several years. For the financial support of this cooperation I like to thank the DFG and the RFBR. It is a pleasure for me to thank Jürgen Fuchs and Christoph Schweigert for many discussions and useful hints on the subject. |
warning/0001/hep-th0001126.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Motivated, to a large extent, by the Hořava-Witten scenario of the strongly coupled heterotic string , the cosmology of a five dimensional universe bounded by 3-branes has recently been the subject of many investigations. If one estimates the parameters in the Hořava-Witten scenario using, as input, the phenomenological values of Newton’s constant, $`M_{GUT}`$ and $`\alpha _{GUT}`$ one finds that there exists a regime below $`M_{GUT}`$ (the inverse size of a Calabi-Yau manifold, on which 6 dimensions are compactified) in which a five dimensional description of our universe is appropriate.
The topology of the five dimensional universe is the one of a $`S^1/Z_2`$ orbifold: There are two distinct 3-branes, and the fields in the bulk between the branes are either symmetric or anti-symmetric with respect to $`Z_2`$, i.e. reflections at the branes. The field content in the five dimensional bulk is the one of $`N=1`$ (gauged) supergravity, plus “matter” originating from a 3-form in 11 dimensions, and from internal degrees of freedom of the metric (moduli) \[4 – 7\]. On the branes we have, in addition, matter originating from $`E_8`$ Yang-Mills theories in 10 dimensions.
Solutions to the equations of motion in the five dimensional bulk have been obtained in \[4 – 12\]. Due to the presence of potentials for the moduli on the branes (with opposite signs) one finds that the moduli configurations vary with $`y`$, the fifth coordinate, leading to a $`y`$-dependent energy-momentum tensor in the bulk.
Independently from the Hořava-Witten theory, scenarios for the cosmology of a five dimensional universe with branes have been developed. Assuming a $`y`$-independent cosmological constant in the bulk, and fine-tuned cosmological constants on the branes (with opposite signs), Randall and Sundrum (RSI) found a static solution of the Einstein equations in the bulk with an exponential dependence of the metric on $`y`$. They argued that this phenomenon could solve the hierachy problem, provided our four dimensional universe lives on a brane with negative tension (negative cosmological constant). Alternatively, if we live on a brane with positive tension surrounded by a bulk with a $`y`$-independent cosmological constant (RSII), the physical distance between the branes may go off to infinity without affecting the validity of Newton’s law on our brane (due to additional massless graviton Kaluza-Klein modes).
Interestingly, once one studies the cosmological evolution induced by additional matter on the branes, the conventional (and successful) relation $`3H^2=8\pi G_N\rho `$ between the Hubble parameter $`H`$, Newton’s constant $`G_N`$ and the energy density $`\rho `$ is not always obtained. Consequently the (non)-validity of this relation, at least at temperatures below $``$ 1 MeV, serves to eliminate certain scenarios.
First, the cosmological evolution induced by matter on the branes has been studied in the simplest case of vanishing cosmological constants in the bulk and on the branes by Binétruy, Deffayet and Langlois in . The astonishing result was that the conventional relation $`H^2\rho `$ does not hold, rather one obtains $`H\rho `$ (in contradiction, e.g., with the otherwise successful big bang nucleo-synthesis).
Subsequently, the matter-induced cosmological evolution in the RSI scenario was investigated in with the “negative” result $`H^2\rho `$ (if we live on the brane with negative tension, which is a necessary condition for the solution of the hierarchy problem in this case).
On the other hand, the cosmological evolution induced by matter on a brane with positive tension agrees with the conventional evolution (at late times) \[17 – 23\] if one assumes that the physical distance between the branes is stabilized, and the relative values of the cosmological constants in the bulk and on the branes are fine-tuned.
The purpose of the present paper is the study, in the (compactified) Hořava-Witten theory, of the cosmological evolution induced by matter on the branes. We allow for an arbitrary equation of state, i.e., a relation $`p=w\rho `$ among the pressure $`p`$ and the energy density $`\rho `$. As fields in the bulk we consider, apart from the components of the graviton, a scalar field $`\phi `$ which corresponds to the universal Calabi-Yau modulus. The presence of this field is model independent and, due to the non-vanishing “internal” components of the field strength associated to the 3-form, one obtains potentials $`V(\phi )`$ both in the bulk and on the branes \[4 – 7\]. If one introduces a dimensionful parameter $`\gamma `$ with $`\gamma 𝒪(M_{GUT}^4/M_{11}^3)`$, where $`M_{11}`$ is scale of the gravitational coupling in eleven dimensions, one finds for the potential in the bulk $`V^{(bulk)}(\phi )𝒪(\gamma ^2)`$, whereas the potentials $`V^{(n)}`$ on the branes satisfy $`V^{(n)}(\phi )𝒪(\gamma )`$.
Subsequently we will assume $`\rho ,p\gamma \kappa _5^2`$ where $`\kappa _5^2`$ is the inverse gravitational coupling in five dimensions with mass dimension $`\kappa _5^2=M_5^3`$. This assumption is certainly realistic with respect to “Standard Model” matter on our brane, but our subsequent results will also apply to “non-standard” matter (e.g. associated with gaugino condensation on the hidden brane) as long as the above inequality is satisfied. Notably we will allow for an arbitrary dependence of the “matter” Lagrangians on the branes, $`^{(n)}`$, on the field $`\phi `$: $`^{(n)}=^{(n)}(\phi )`$ (with $`^{(n)}𝒪(\rho ,p)𝒪(\gamma \kappa _5^2))`$ which affects the junction conditions of the field $`\phi `$ on the branes, see below.
The inflationary evolution in this set-up, induced by matter potentials on the branes (corresponding to an equation of state $`p=\rho `$), has previously been discussed in . The results obtained in correspond either to the case $`\gamma R_51`$, which allows for a linearized approximation of the $`y`$ dependence of the fields, or to “matter” Lagrangians $`^{(n)}`$ on the branes which are $`\phi `$ independent.
Below, we will construct solutions in the bulk which are exact in $`y`$, since we will not assume $`\gamma R_51`$. Notably, we will not assume that the physical distance between the branes (or the $`yy`$-component of the metric in the bulk) is stabilized in an ad hoc fashion (corresponding mechanisms are discussed in \[23, 25 – 30\]): One of our aims is to see, under which general circumstances solutions with a static size of the fifth dimension exist. In fact we will obtain, depending on the $`\phi `$ dependence of $`^{(n)}(\phi )`$, both solutions corresponding to a “rolling radius”, as well as solutions with a static $`yy`$-component of the metric. Not astonishingly, only the latter case allows for a conventional cosmological evolution on a brane.
In the next section we will present the general set-up of our approach: the action in the bulk and on the branes, the Einstein equations and the equation of motion of $`\phi `$, the junction conditions on the branes, and the exact (static) solution of . In section 3 we introduce our ansatz for the time-dependent fields in the bulk, which is motivated by the search for a conventional cosmological evolution on the branes: we will assume, that the non-trivial time-dependence is induced by the presence of matter on the branes. A self-consistent ansatz is seen to be $`_t𝒪((\rho \gamma \kappa _5^2)^{1/2})`$, where $`_t`$ denotes the time derivative of any field in the bulk or on the branes. (Given the existence of time dependent solutions with $`_t𝒪(\gamma )`$ , this ansatz is certainly not the most general one. It represents, on the other hand, the “minimal” time dependence which is induced by the presence of matter on the branes). Then the Einstein equations and the equation of motion of $`\phi `$ can be solved exactly in $`y`$ and $`t`$, but neglecting terms of relative order $`\rho \kappa _5^2\gamma ^1`$. In section 4 we discuss the physical properties of our solutions and conclude.
## 2 Equations of motion and junction conditions
Our starting point is a five dimensional action in the bulk which depends, apart from the gravitational sector, on a scalar field $`\phi `$ which is related to the universal modulus of the internal Calabi-Yau manifold \[4 – 7\]. In the notation of , our field $`\phi `$ is related to the field $`V`$ in by $`V=\mathrm{exp}(2\phi )`$. Likewise, our parameter $`\gamma `$ appearing in the potentials of $`\phi `$ in the bulk and on the branes is related to the parameter $`\alpha `$ in through $`\gamma =\sqrt{2}\alpha `$. In any case $`\gamma `$ is of $`𝒪(M_{GUT}^4/M_{11}^3)`$. The exact value of $`\gamma `$ depends on the shape of the Calabi-Yau manifold . The bulk action then reads
$$S^{bulk}=\frac{1}{\kappa _5^2}\sqrt{g_5}d^5x\left\{\frac{1}{2}R+_\mu \phi ^\mu \phi +\frac{\gamma ^2}{12}e^{4\phi }\right\}.$$
(2.1)
Subsequently we denote the fifth dimension (across the bulk) by $`y`$ and assume, that all fields depend on $`t`$ and $`y`$, but not on the 3 spatial coordinates $`x_i`$. A convenient (diagonal) ansatz for the metric is then
$$ds^2=e^{2\nu }dt^2+e^{2\alpha }(d\stackrel{}{x})^2+e^{2\beta }dy^2.$$
(2.2)
Denoting time derivatives by points, and derivatives with respect to $`y`$ by primes, the Einstein equations in the bulk become
$`e^{2\nu }\left(\dot{\alpha }^2+\dot{\alpha }\dot{\beta }\right)+e^{2\beta }\left(\alpha ^{\prime \prime }+2\alpha ^2\alpha ^{}\beta ^{}\right)={\displaystyle \frac{\kappa _5^2}{3}}T_0^0,`$ (2.3a)
$`e^{2\nu }\left(2\ddot{\alpha }+3\dot{\alpha }^22\dot{\alpha }\dot{\nu }+\dot{\beta }(2\dot{\alpha }\dot{\nu })+\ddot{\beta }+\dot{\beta }^2\right)`$
$`+e^{2\beta }\left(2\alpha ^{\prime \prime }+\nu ^{\prime \prime }+3\alpha ^2+\nu ^2+2\alpha ^{}\nu ^{}\beta ^{}(2\alpha ^{}+\nu ^{})\right)=\kappa _5^2T_S,`$ (2.3b)
$`\dot{\alpha }\nu ^{}\dot{\alpha }^{}\dot{\alpha }\alpha ^{}+\alpha ^{}\dot{\beta }={\displaystyle \frac{\kappa _5^2}{3}}T_{05},`$ (2.3c)
$`e^{2\nu }\left(\ddot{\alpha }+2\dot{\alpha }^2\dot{\alpha }\dot{\nu }\right)+e^{2\beta }(\alpha ^2+\alpha ^{}\nu ^{})={\displaystyle \frac{\kappa _5^2}{3}}T_5^5.`$ (2.3d)
Here $`T_S`$ denotes the diagonal element of the spatial components of the energy momentum tensor,
$$T_i^j=T_S\delta _i^j.$$
(2.4)
Given the action (2.1), and recalling that $`\phi `$ depends only on $`t`$ and $`y`$, the non-vanishing components of the energy momentum tensor in the bulk (away from the branes) are
$`\kappa _5^2T_0^0=e^{2\nu }\dot{\phi }^2e^{2\beta }\phi ^2{\displaystyle \frac{\gamma ^2}{12}}e^{4\phi },`$ (2.5a)
$`\kappa _5^2T_S=e^{2\nu }\dot{\phi }^2e^{2\beta }\phi ^2{\displaystyle \frac{\gamma ^2}{12}}e^{4\phi },`$ (2.5b)
$`\kappa _5^2T_5^5=e^{2\nu }\dot{\phi }^2+e^{2\beta }\phi ^2{\displaystyle \frac{\gamma ^2}{12}}e^{4\phi },`$ (2.5c)
$`\kappa _5^2T_{05}=2\dot{\phi }\phi ^{}.`$ (2.5d)
Finally we have to consider the equation of motion of the field $`\phi `$, which takes the form (away from the branes)
$`e^{2\nu }\left(\ddot{\phi }+\dot{\phi }(3\dot{\alpha }\dot{\nu }+\dot{\beta })\right)+e^{2\beta }\left(\phi ^{\prime \prime }+\phi ^{}(3\alpha ^{}+\nu ^{}\beta ^{})\right)={\displaystyle \frac{\gamma ^2}{6}}e^{4\phi }.`$
Now we consider the actions on the branes, which are situated at $`y^{(n)}=0,\pi R_5`$. We parametrize them as
$$S^{(n)}=\sqrt{g_4}d^4x\left\{_{matter}^{(n)}(\phi )V^{(n)}(\phi )\right\}.$$
(2.7)
Here $`_{matter}^{(n)}(\phi )`$ depends, of course, on many additional matter fields $`\varphi _i`$. Below, however, only the possible dependence on the bulk field $`\phi `$ will play a role. Depending on the precise form of $`_{matter}^{(n)}`$, and on the configurations of the fields $`\varphi _i`$, this part of the brane actions will contribute to the energy momentum tensors on the branes in the form of energy densities $`\rho ^{(n)}`$ and pressures $`p^{(n)}`$, see below.
The potentials $`V^{(n)}(\phi )`$ on the two branes are known to be of the form
$$V^{(n)}(\phi )=\frac{\gamma }{\kappa _5^2}e^{2\phi }.$$
(2.8)
Here the minus sign applies to the brane at $`y=0`$, and the plus sign to the brane at $`y=\pi R_5`$. Since we assume
$$_{matter}^{(n)}\rho ^{(n)}p^{(n)}\frac{\gamma }{\kappa _5^2},$$
(2.9)
the brane at $`y=0`$ is thus the brane with negative tension, if $`\gamma `$ is positive. We can easily interchange the role of the two branes (negative vs. positive tension) by changing the sign of $`\gamma `$.
As is well known, the presence of branes leads to additional singular terms (proportional to $`\delta `$-functions in $`y`$) on the right-hand sides of the Einstein equations (2.3a), (2), and the equation of motion (2), which have to be matched by singularities in the second derivatives in $`y`$ on the left-hand side. Since all fields under consideration are symmetric under the orbifold symmetry $`Z_2`$, these jumps in the first derivatives in $`y`$ fix these first derivatives completely at $`y^{(n)}=0,\pi R_5`$. Here, these junction conditions read
$`\alpha ^{(n)}={\displaystyle \frac{\kappa _5^2}{6}}e^\beta T_0^{0(n)},`$ (2.10a)
$`\nu ^{(n)}={\displaystyle \frac{\kappa _5^2}{6}}e^\beta \left(3T_S^{(n)}2T_0^{0(n)}\right),`$ (2.10b)
$`\phi ^{(n)}={\displaystyle \frac{\kappa _5^2}{4}}e^\beta {\displaystyle \frac{\delta }{\delta \phi }}\left(V^{(n)}(\phi )_{matter}^{(n)}(\phi )\right),`$ (2.10c)
with
$`T_0^{0(n)}=V^{(n)}(\phi )\rho ^{(n)},`$ (2.11a)
$`T_S^{(n)}=V^{(n)}(\phi )+p^{(n)}`$ (2.11b)
and $`V^{(n)}(\phi )`$ as in eq. (2.8) As discussed above, $`\rho ^{(n)}`$ and $`p^{(n)}`$ originate from the matter Lagrangian on the branes, which we will not specify any further at this stage.
Static solutions to the Einstein equations (2), the scalar equation of motion (2) and the junction conditions (2) - in the absence of matter on the branes - have been given in . We will denote these solutions by $`\stackrel{~}{\alpha }`$, $`\stackrel{~}{\nu }`$, $`\stackrel{~}{\beta }`$ and $`\stackrel{~}{\phi }`$, which read
$`\stackrel{~}{\alpha }(y)=\stackrel{~}{\nu }(y)={\displaystyle \frac{1}{2}}\mathrm{}nH+A,`$ (2.12a)
$`\stackrel{~}{\beta }(y)=2\mathrm{}nH+B,`$ (2.12b)
$`\stackrel{~}{\phi }(y)={\displaystyle \frac{3}{2}}\mathrm{}nH+{\displaystyle \frac{1}{2}}B,`$ (2.12c)
$`H={\displaystyle \frac{\gamma }{3}}|y|+C\mathrm{for}\pi R_5y\pi R_5,H\left(y+2\pi R_5\right)=H(y).`$ (2.12d)
$`A`$, $`B`$ and $`C`$ are arbitrary constants; in terms of a supergravity Lagrangian in four dimensions $`B`$ and $`C`$ correspond to linear combinations of the arbitrary vevs of the standard fields $`Re(S)`$ and $`Re(T)`$ .
Time dependent solutions have been obtained in , but in the next sections we seek for cosmological solutions, where the time dependence is induced by the presence of matter on the branes. Hence the static solutions (2) will serve as a starting point.
## 3 Cosmological evolution induced by matter on the branes
In this section we will construct solutions for $`\alpha `$, $`\nu `$, $`\beta `$ and $`\phi `$ in the presence of matter on the branes. Matter on the branes contributes to the junction conditions (2) which affects, by continuity in $`y`$, the solutions across the bulk for all $`y`$. We will assume that the “amount” of matter is small compared to the fundamental scales $`M_{11}`$, $`M_{GUT}`$ (which are quite close to each other) and $`M_5`$, i.e.
$$\rho ^{(n)}p^{(n)}_{matter}^{(n)}\gamma \kappa _5^2.$$
(3.1)
On the other hand we will need no assumption on $`\gamma R_5`$, since we will obtain solutions exact in $`y`$.
First we recall that, in the case of an “empty” bulk, the Hubble constant $`H`$ was shown to satisfy $`H\rho `$ , i.e. time derivatives $`_t`$ are of the order $`_t𝒪(\rho )`$. A corresponding result was also obtained in in the case where $`\rho `$ represents a $`\phi `$ independent potential and $`\rho R_5\kappa _5^21`$. Conventional cosmology, on the other hand, corresponds to $`H^2\rho `$, i.e. $`_t𝒪(\sqrt{\rho })`$. This latter behaviour can be obtained in the case of additional cosmological constants in the bulk and on the branes \[17 – 23\], and we will also obtain corresponding solutions in the present case.
Our ansatz for solutions in the presence of matter on the branes will be generalizations of the static solutions (2) in two respects: First, we add $`y`$ and $`t`$ dependent contributions to $`\stackrel{~}{\alpha }`$, $`\stackrel{~}{\nu }`$, $`\stackrel{~}{\beta }`$ and $`\stackrel{~}{\phi }`$, which are of $`𝒪(\rho \kappa _5^2\gamma ^1)`$ where $`\rho `$ denotes the order of magnitude of all terms on the left-hand side of eq. (3.1). Second, we promote the constants $`A`$, $`B`$ and $`C`$ in (2) to time dependent parameters assuming, however, that all time derivatives are of the order
$$_t𝒪\left((\rho \gamma \kappa _5^2)^{1/2}\right)$$
(3.2)
or smaller. (This ansatz allows, of course, for a still “weaker” time dependence with $`_t𝒪(\rho )`$). Thus we write
$`\alpha (y,t)={\displaystyle \frac{1}{2}}\mathrm{}nH(y,t)+\alpha _0(t)+\overline{\alpha }(y,t),`$ (3.3a)
$`\nu (y,t)={\displaystyle \frac{1}{2}}\mathrm{}nH(y,t)+\nu _0(t)+\overline{\nu }(y,t),`$ (3.3b)
$`\beta (y,t)=2\mathrm{}nH(y,t)+\beta _0(t)+\overline{\beta }(y,t),`$ (3.3c)
$`\phi (y,t)={\displaystyle \frac{3}{2}}\mathrm{}nH(y,t)+\phi _0(t)+\overline{\phi }(y,t),`$ (3.3d)
$`\beta _0(t)=2\phi _0(t),`$ (3.3e)
$`H(y,t)={\displaystyle \frac{\gamma }{3}}|y|+C(t),`$ (3.3f)
with
$$\overline{\alpha },\overline{\nu },\overline{\beta },\overline{\phi }𝒪\left(\rho \kappa _5^2\gamma ^1\right).$$
(3.4)
Equation (3.3e) can be obtained from the dominant terms of the Einstein equations, and agrees with the static limit (2) after $`B2\phi _0(t)`$.
Let us first use this ansatz in the (05) component of the Einstein equations (2.3c). Subsequently we replace $`\beta _0`$ by $`2\phi _0`$ everywhere according to eq. (3.3e), which simplifies several expressions. Using (3.2), (3.4), and keeping all terms up to $`𝒪(\rho ^{3/2})`$, eq. (2.3c) can be brought into the form (with (2.5d) for $`T_{05}`$)
$`\left(\dot{\alpha }_0+{\displaystyle \frac{\dot{C}}{2H}}\right)\left(\overline{\nu }^{}\overline{\alpha }^{}\right)\dot{\overline{\alpha }}^{}+\left(2\dot{\phi }_0+3{\displaystyle \frac{\dot{C}}{H}}\right)\left(\overline{\alpha }^{}{\displaystyle \frac{\overline{\phi }^{}}{3}}\right)`$
$`+{\displaystyle \frac{\gamma }{6H}}\left(\dot{\overline{\beta }}2\dot{\overline{\phi }}\right)+{\displaystyle \frac{\gamma }{6H^2}}\dot{C}=0.`$ (3.5)
One notes that all terms are of $`𝒪(\rho ^{3/2})`$, except for the last term on the left-hand side, which is a priori (from (3.2)) of $`𝒪(\rho ^{1/2})`$. Thus $`\dot{C}`$ is exceptionally at most of $`𝒪(\rho ^{3/2})`$, and will not contribute to the dominant orders in $`\rho `$ in the following equations.
Next we insert our ansatz into the remaining Einstein equations (2) as well as (2). We use eqs. (2) for the components of the energy momentum tensor, and expand each expression in $`\rho `$ using (3.2) and (3.4). The dominant terms of $`𝒪(\rho ^0)`$ cancel as they should, and subsequently we display all terms of $`𝒪(\rho )`$. The subsequent equations follow from eqs. (2.3a), (2), (2.3d) and (2), after moving all time derivatives to the left (and with $`\beta _0=2\phi _0`$):
$`\dot{\alpha }_0^2+2\dot{\alpha }_0\dot{\phi }_0{\displaystyle \frac{1}{3}}\dot{\phi }_0^2=H^3e^{2\nu _04\phi _0}(\overline{\alpha }^{\prime \prime }+{\displaystyle \frac{\gamma }{6H}}(2\overline{\phi }^{}\overline{\beta }^{})`$
$`+{\displaystyle \frac{\gamma ^2}{18H^2}}(\overline{\beta }2\overline{\phi })),`$ (3.6a)
$`2\ddot{\alpha }_0+3\dot{\alpha }_0^22\dot{\nu }_0\left(\dot{\alpha }_0+\dot{\phi }_0\right)+4\dot{\alpha }_0\dot{\phi }_0+2\ddot{\phi }_0+5\dot{\phi }_0^2`$
$`=H^3e^{2\nu _04\phi _0}\left(2\overline{\alpha }^{\prime \prime }+\overline{\nu }^{\prime \prime }+{\displaystyle \frac{\gamma }{2H}}\left(2\overline{\phi }^{}\overline{\beta }^{}\right)+{\displaystyle \frac{\gamma ^2}{6H^2}}\left(\overline{\beta }2\overline{\phi }\right)\right),`$ (3.6b)
$`\ddot{\alpha }_0+2\dot{\alpha }_0^2\dot{\alpha }_0\dot{\nu }_0+{\displaystyle \frac{1}{3}}\dot{\phi }_0^2=H^3e^{2\nu _04\phi _0}({\displaystyle \frac{\gamma }{6H}}(3\overline{\alpha }^{}+\overline{\nu }^{}2\overline{\phi }^{})`$
$`+{\displaystyle \frac{\gamma ^2}{18H^2}}(\overline{\beta }2\overline{\phi })),`$ (3.6c)
$`\ddot{\phi }_0+\dot{\phi }_0(3\dot{\alpha }_0\dot{\nu }_0+2\dot{\phi }_0)=H^3e^{2\nu _04\phi _0}(\overline{\phi }^{\prime \prime }+{\displaystyle \frac{\gamma }{2H}}(3\overline{\alpha }^{}+\overline{\nu }^{}\overline{\beta }^{})`$
$`+{\displaystyle \frac{\gamma ^2}{3H^2}}(\overline{\beta }2\overline{\phi })).`$ (3.6d)
These four equations fix the $`y`$ dependence of $`\overline{\alpha }`$, $`\overline{\nu }`$, $`\overline{\beta }`$ and $`\overline{\phi }`$:
$`\overline{\alpha }(y,t)=\widehat{\alpha }(t)H^5+{\displaystyle \frac{1}{3}}F(y),`$ (3.7a)
$`\overline{\nu }(y,t)=\widehat{\nu }(t)H^5+{\displaystyle \frac{1}{3}}F(y),`$ (3.7b)
$`\overline{\beta }(y,t)=\widehat{\beta }(t)H^5+2F(y)+{\displaystyle \frac{2H}{\gamma }}F^{}(y),`$ (3.7c)
$`\overline{\phi }(y,t)=\widehat{\phi }(t)H^5+F(y)`$ (3.7d)
with $`H`$ as in eq. (2.12d) and $`F(y)`$ arbitrary.
The four time-dependent parameters $`\widehat{\alpha }`$, $`\widehat{\nu }`$, $`\widehat{\beta }`$ and $`\widehat{\phi }`$ are constrained by the three junction conditions (2) in terms of the matter on the branes. Plugging our ansätz (3) as well as (3) into eqs. (2), one finds again that the leading terms of $`𝒪(\rho ^0)`$ cancel, and the terms of $`𝒪(\rho )`$ can be brought into the form
$`10\gamma \epsilon ^{(n)}\widehat{\alpha }=\left(H^{(n)}\right)^2\kappa _5^2\rho ^{(n)}e^{2\phi _0}+\gamma (\widehat{\beta }2\widehat{\phi }),`$ (3.8a)
$`10\gamma \epsilon ^{(n)}\widehat{\nu }=\left(H^{(n)}\right)^2\kappa _5^2\left(3p^{(n)}+2\rho ^{(n)}\right)e^{2\phi _0}+\gamma (\widehat{\beta }2\widehat{\phi }),`$ (3.8b)
$`10\gamma \epsilon ^{(n)}\widehat{\phi }={\displaystyle \frac{3}{2}}\left(H^{(n)}\right)^2\kappa _5^2{\displaystyle \frac{\delta _{matter}^{(n)}}{\delta \phi }}e^{2\phi _0}+3\gamma (\widehat{\beta }2\widehat{\phi }),`$ (3.8c)
where $`H^{(n)}`$ denotes $`H(y^{(n)})`$ with $`H(y)`$ as in eq. (2.12d), and $`\epsilon ^{(n)}=+1`$ for $`y^{(n)}=0`$, $`\epsilon ^{(n)}=1`$ for $`y^{(n)}=\pi R_5`$. The dependence on the function $`F(y)`$ of eqs. (3) cancels in the junction conditions (2).
Altogether we thus have six equations, three junction conditions at the brane $`n=1`$, and three junction conditions at the brane $`n=2`$. Let us first discuss, to what extent these six equations restrict the properties of the matter on the different branes. We recall that, e.g., in the case of an “empty” bulk considered in , $`\rho ^{(2)}`$ is fixed in terms of $`\rho ^{(1)}`$, and $`p^{(2)}`$ in terms of $`p^{(1)}`$. In the present case the three equations which involve the properties of the matter on the brane 2 can be brought into the form
$`\left(H^{(1)}\right)^2p^{(1)}+\left(H^{(2)}\right)^2p^{(2)}=\left(H^{(1)}\right)^2\rho ^{(1)}\left(H^{(2)}\right)^2\rho ^{(2)},`$ (3.9a)
$`\left(H^{(1)}\right)^2{\displaystyle \frac{\delta _{matter}^{(1)}}{\delta \phi }}+\left(H^{(2)}\right)^2{\displaystyle \frac{\delta _{matter}^{(2)}}{\delta \phi }}=2\left(H^{(1)}\right)^2\rho ^{(1)}+2\left(H^{(2)}\right)^2\rho ^{(2)},`$ (3.9b)
$`\gamma \kappa _5^2\left[\left(H^{(1)}\right)^2\rho ^{(1)}+\left(H^{(2)}\right)^2\rho ^{(2)}\right]=2\left(\widehat{\beta }2\widehat{\phi }\right).`$ (3.9c)
At this stage we have to recall that $`H`$, defined in eq. (2.12d), depends on an as yet arbitrary parameter $`C`$ (with negligible time dependence): we have $`H^{(1)}=C`$, $`H^{(2)}=\frac{\gamma }{3}\pi R_5+C`$. Hence eqs. (3) do not necessarily constrain the matter on the brane 2 in terms of the matter on the brane 1, but can rather serve to fix $`C`$.
Of particular interest is the case where the matter fields $`\varphi _i`$ on the branes are constant (at the minima of their potentials), and the only role of $`_{matter}^{(n)}`$ is thus to provide additional (possibly constant) potentials for the modulus field $`\phi `$ on the branes. Then we can write
$`_{matter}^{(n)}=\widehat{V}^{(n)}(\phi ),`$ (3.10a)
$`\rho ^{(n)}=\widehat{V}^{(n)}(\phi ),`$ (3.10b)
$`p^{(n)}=\widehat{V}^{(n)}(\phi ).`$ (3.10c)
Hence eq. (3.9a) is satisfied identically. Eq. (3.9b) serves to fix $`C`$ in terms of the potentials on the branes and $`\phi `$. However, in order not to contradict the previous result $`\dot{C}<𝒪(\rho ^{3/2})`$, one must have either $`\dot{\phi }_00`$, or $`\widehat{V}^{(1)}(\phi )=\mathrm{const}.\widehat{V}^{(2)}(\phi )`$. Eq. (3.9c) fixes the combination $`(\widehat{\beta }2\widehat{\phi })`$.
Herewith we conclude the discussion on the relation between the matter on the different branes, and concentrate subsequently on the physics on brane 1 assuming that eqs. (3) are satisfied. Note, however, that the two branes are physically distinct in Hořava-Witten theory, since the dominant potentials $`V^{(n)}(\phi )`$ on the branes differ in their sign (cf. eq. (2.8)). On the other hand we find from eqs. (3) that we can exchange the role of the two branes by changing the sign of $`\gamma `$, and redefining the parameter $`\widehat{\beta }`$:
$$\gamma \gamma ,\widehat{\beta }4\widehat{\phi }\widehat{\beta }.$$
(3.11)
In order to solve eqs. (3) on brane 1 it is convenient to define the parameters $`w`$, $`d`$ as follows:
$`p^{(1)}=w\rho ^{(1)},`$ (3.12a)
$`{\displaystyle \frac{\delta _{matter}^{(1)}}{\delta \phi }}=d\rho ^{(1)}.`$ (3.12b)
In general the parameter $`d`$ will depend on $`\phi `$, unless $`_{matter}^{(1)}`$ and $`\rho ^{(1)}`$, $`p^{(1)}`$ happen to be related as, e.g., in eqs. (3) with
$$\widehat{V}^{(n)}(\phi )e^{d\phi }.$$
(3.13)
With eqs. (3) and $`H^{(1)}=C`$, the general solution of eqs. (3) on brane 1 can be written as
$`\widehat{\alpha }={\displaystyle \frac{1}{16}}\widehat{\beta }{\displaystyle \frac{\kappa _5^2\rho ^{(1)}e^{2\phi _0}}{160\gamma C^2}}(3d+16),`$ (3.14a)
$`\widehat{\nu }={\displaystyle \frac{1}{16}}\widehat{\beta }+{\displaystyle \frac{\kappa _5^2\rho ^{(1)}e^{2\phi _0}}{160\gamma C^2}}\left(3d+16(3w+2)\right),`$ (3.14b)
$`\widehat{\phi }={\displaystyle \frac{3}{16}}\widehat{\beta }+{\displaystyle \frac{3\kappa _5^2\rho ^{(1)}e^{2\phi _0}}{32\gamma C^2}}d,`$ (3.14c)
with $`\widehat{\beta }`$ arbitrary. (Eventually $`\widehat{\beta }`$ can be fixed, combining eqs. (3.14c) and (3.9c), in terms of $`\rho ^{(2)}`$).
Combining eqs. (3) and eqs. (3) we have thus obtained the general solutions for $`\overline{\alpha }`$, $`\overline{\nu }`$, $`\overline{\beta }`$, $`\overline{\phi }`$ in terms of the properties of the matter on the branes and for all values of $`y`$.
Let us recall the ansatz (3) for the fields $`\alpha `$, $`\nu `$, $`\beta `$ and $`\phi `$. Since we obtained $`\dot{C}<𝒪(\rho ^{3/2})`$, the dominant contributions to these fields decompose into a sum of $`y`$ dependent and $`t`$ dependent terms. The cosmological evolution on any brane is thus determined by $`\alpha _0(t)`$, $`\nu _0(t)`$ and $`\phi _0(t)`$. The knowledge of the subdominant contributions $`\overline{\alpha }`$, $`\overline{\nu }`$, $`\overline{\beta }`$ and $`\overline{\phi }`$ to eqs. (3) is required in order to obtain the $`t`$ dependence of $`\alpha _0`$, $`\nu _0`$ and $`\phi _0`$ from eqs. (3), which we exploit in the following.
Once we plug our solutions (3) for $`\overline{\alpha }`$, $`\overline{\nu }`$, $`\overline{\beta }`$ and $`\overline{\phi }`$, together with eqs. (3), into eqs. (3), we first observe that all dependence on the arbitrary function $`F(y)`$ in (3) as well as on $`\widehat{\beta }`$ cancels out. Hence the time dependence of $`\alpha _0`$, $`\nu _0`$ and $`\phi _0`$ can be related to the properties of the matter on the brane 1 only, which are encoded in our parameters $`\rho ^{(1)}`$, $`w`$ and $`d`$. Second, we may choose the gauge $`\nu _0(t)=0`$, such that $`t`$ is proportional - up to small corrections - to the cosmic time: now we have $`_t\nu (y,t)<𝒪(\rho ^{3/2})`$ (cf. eq. 3.3b), since $`\dot{C}`$ and $`\dot{\overline{\nu }}`$ are at most of $`𝒪(\rho ^{3/2})`$.
Since one finds that one of eqs. (3), e.g. eq. (3), is redundant, we just give our results for eqs. (3), (3) and (3):
$`\dot{\alpha }_0^2+2\dot{\alpha }_0\dot{\phi }_0{\displaystyle \frac{1}{3}}\dot{\phi }_0^2={\displaystyle \frac{2}{9}}{\displaystyle \frac{\gamma \kappa _5^2\rho ^{(1)}}{C^2}}e^{2\phi _0}`$ (3.15a)
$`\ddot{\alpha }_0+2\dot{\alpha }_0^2+{\displaystyle \frac{1}{3}}\dot{\phi }_0^2={\displaystyle \frac{\gamma \kappa _5^2\rho ^{(1)}}{36C^2}}e^{2\phi _0}(3w13d),`$ (3.15b)
$`\ddot{\phi }_0+\dot{\phi }_0\left(3\dot{\alpha }_0+2\dot{\phi }_0\right)={\displaystyle \frac{\gamma \kappa _5^2\rho ^{(1)}}{12C^2}}e^{2\phi _0}(3w1+d).`$ (3.15c)
Eqs. (3) justify, a posteriori, our initial ansatz (3.2) for the order of $`_t`$, provided $`C^2\mathrm{exp}(2\phi _0)𝒪(1)`$.
Taking the time derivative of eq. (3.15a) and using all eqs. (3), the analog of the standard energy conservation condition can be obtained:
$$\dot{\rho }^{(1)}=3\rho ^{(1)}\dot{\alpha }_0(1+w)+d\rho ^{(1)}\dot{\phi }_0.$$
(3.16)
First, if we insert our solution (3) for $`\overline{\alpha }`$, $`\overline{\nu }`$, $`\overline{\beta }`$ and $`\overline{\phi }`$, together with eqs. (3) and eq. (3.16), into eq. (3), we obtain a trivial time dependence of $`C`$:
$$\dot{C}=0.$$
(3.17)
Second, eq. (3.16) differs from the standard condition for energy conservation due to the last term. This term describes the “disappearance” of energy into the fifth dimension, if both $`\delta _{matter}^{(1)}/\delta \phi d0`$ and $`\dot{\phi }_00`$. In fact, if $`_{matter}^{(1)}`$ is of the form of a potential $`\widehat{V}^{(1)}(\phi )`$ (as in eqs. (3)) we have $`w=1`$ and, using $`\dot{\phi }\dot{\phi }_0`$, eq. (3.10b) and eq. (3.12b),
$$\dot{\rho }^{(1)}=_t\widehat{V}^{(1)}(\phi )=\frac{\delta \widehat{V}^{(1)}}{\delta \phi }\dot{\phi }\frac{\delta _{matter}^{(1)}}{\delta \phi }\dot{\phi }_0=d\rho ^{(1)}\dot{\phi }_0$$
(3.18)
in agreement with (3.16).
For $`d0`$, a “standard” cosmological evolution on any brane, which requires standard energy conservation, is only possible for $`\dot{\phi }_0=\dot{\beta }_0=0`$. From eq. (3.15c) this situation is feasable only for
$$d=13w,$$
(3.19)
which requires a particular dependence of $`_{matter}^{(1)}`$ on $`\phi `$. If, again, $`_{matter}^{(1)}`$ is of the form of a potential (hence $`w=1`$), we need $`d=4`$ or
$$\widehat{V}^{(1)}(\phi )=\mathrm{const}.e^{4\phi }.$$
(3.20)
Furthermore, from eq. (3.15a), $`\dot{\phi }_0=0`$ implies immediately
$$\gamma \rho ^{(1)}0.$$
(3.21)
Thus, if we want to accomodate a standard positive energy density $`\rho ^{(1)}`$, we have to choose $`\gamma <0`$, i.e. brane 1 has to be the one with a positive tension (since now the dominant contribution $`V^{(1)}(\phi )`$, from eq. (2.8), is positive). This result coincides with the one obtained in the case of cosmological constants in the bulk and on the branes \[17 – 23\].
If eq. (3.19) is satisfied, eqs. (3) allow for $`\dot{\phi }_0=0`$. Since we have $`\dot{\alpha }\dot{\alpha }_0`$ and, from our definition (2.2) of the metric, $`\dot{\alpha }`$ corresponds to the Hubble parameter, eqs. (3.15a) and (3.15b) are easily seen to correspond to the ordinary Freedman equations for $`\gamma <0`$. (A constant term in $`\nu `$ allows for a constant rescaling of $`t`$ such that the right-hand sides of eqs. (3.15a) and (3.15b), for $`\dot{\phi }_0=0`$, can always be brought into standard form).
If eq. (3.19) is not satisfied, we have necessarily $`\dot{\phi }_00`$ and, from eq. (3.3e), $`\dot{\beta }_00`$, i.e. solutions corresponding to a “rolling radius”: the physical distance between the two branes, given by $`e^{\beta (y)}𝑑y`$ with $`y=\{0,\pi R_5\}`$, varies with $`t`$.
Explicit solutions to eqs. (3) can be given if $`w=`$ const. and $`d=`$ const. From eq. (3.12b) the second condition is satisfied if
$$_{matter}^{(1)}(\phi )\rho ^{(1)}(\phi ,t)=\widehat{\rho }^{(1)}(t)e^{d\phi }.$$
(3.22)
Then eqs. (3) are solved for
$`\alpha _0(t)=\mathrm{const}.+r\mathrm{ln}t,\phi _0(t)=\mathrm{const}.+s\mathrm{ln}t,`$
$`\rho ^{(1)}\mathrm{const}.t^{2s2}(\mathrm{or}\widehat{\rho }^{(1)}\mathrm{const}.t^{(2d)s2}),`$ (3.23)
with
$$r=\frac{2(w3d)}{3w^26wdd^211},s=\frac{2(3w1+d)}{3w^26wdd^211},$$
(3.24)
provided the denominators are non-zero. In the case where $`_{matter}^{(1)}(\phi )`$ is just an additional potential in $`\phi `$, as in eqs. (3) with $`n=1`$, we have $`w=1`$ and thus
$$r=\frac{2(4+d)}{(4d)(2d)},s=\frac{2}{2d}\mathrm{if}d2,4.$$
(3.25)
(The particular case $`d=2`$, where $`\widehat{V}^{(n)}(\phi )`$ in eqs. (3) has the same functional dependence on $`\phi `$ as $`V^{(n)}(\phi )`$ in eq. (2.8), has already been considered in .) For $`d=4`$ we are back in the situation where eq. (3.19) is valid (and where standard inflation on the branes is obtained for $`w=1`$), whereas for $`d=2`$ we obtain inflationary evolution both on the branes and across the bulk, i.e. in $`\beta (t)2\phi _0(t)`$:
$`a(t)=e^{\alpha (t)}e^{\alpha _0(t)}e^{const.t},`$
$`\dot{\phi }_0={\displaystyle \frac{1}{3}}\dot{\alpha }_0,\rho (t)=\overline{\rho }e^{2\phi _0(t)}`$ (3.26)
with $`\dot{\overline{\rho }}=0`$, and const. $`\pm \sqrt{\overline{\rho }}`$.
Herewith we conclude the different solutions to eqs. (3), which will be discussed in the next section.
## 4 Discussion and conclusions
In the previous section we have constructed various cosmological solutions in (compactified) Hořava-Witten theory with additional matter on the branes, assuming $`\rho \gamma \kappa _5^2`$, cf. eq. (3.1). Actually, whenever $`\rho `$ varies with $`t`$, this assumption restricts the validity of the solutions to corresponding regimes in $`t`$, typically to sufficiently late time once $`\rho `$ decays in $`t`$.
Generally compactified Hořava-Witten theory suffers from the usual moduli problem, i.e. there are scalar degrees of freedom with vanishing potentials. In the present framework this phenomenon corresponds to the presence of arbitrary constants in the static solutions (2) (in the absence of additional matter), and to the existence of time dependent solutions with $`_t𝒪(\gamma )`$ . Since we assumed a much weaker time dependence, cf. eq. (3.2), our solutions require particular initial conditions in the form of nearly static (and homogeneous) configurations. However, once matter on the branes is present, the fields cannot vary slower with $`t`$ than indicated by our solutions.
One aspect of the moduli problem is the fact that generically the physical distance between the branes, i.e. the $`yy`$ component of the metric (parametrized by $`\beta `$), varies in time. Our results show that, once matter is present only on the branes, this phenomenon is unavoidable, unless the matter couples to the universal Calabi-Yau modulus field $`\phi `$ in a particular way, cf. eq. (3.19). Possibly the five dimensional radius can be stabilized by means of an additional potential for $`\phi `$ in the bulk \[23, 25 – 30\]. Then our cosmological solutions would be relevant at times $`t`$, at which potentials $`V^{(n)}(\phi )\gamma \kappa _5^2`$ on the branes are present (e.g. due to gaugino condensation), but where additional bulk-potentials are not yet switched on.
Let us summarize our results in the case where the additional matter Lagrangian on the branes corresponds to potentials for $`\phi `$, i.e. where the equation of state corresponds to $`w=1`$. Exact solutions have been obtained for potentials of the form
$$\widehat{V}^{(n)}(\phi )\left(=\rho ^{(n)}\right)=\widehat{\rho }^{(n)}e^{d\phi }.$$
(4.1)
Generically we obtain solutions with a power law behaviour in $`t`$ for the scale factor $`a`$ and the physical distance $`R_{phys}`$ between the branes (cf. eqs. (3) and below; note that $`\widehat{\rho }^{(n)}`$ is time independent for $`w=1`$, where eqs. (3.25) are valid):
$`a(y,t)=e^{\alpha (y,t)}\mathrm{const}.(y)e^{\alpha _0(t)}t^r,`$ (4.2a)
$`R_{phys}(t)e^{\beta _0(t)}=e^{2\phi _0(t)}t^{2s},`$ (4.2b)
with $`r`$, $`s`$ as in eqs. (3.25) (for $`d=0`$, e.g., one has $`r=s=1`$). In the particular case $`d=2`$ one obtains, from eqs. (3), inflationary (i.e. exponentially increasing or decreasing) evolutions for $`a`$, $`R_{phys}`$ and $`\widehat{\rho }^{(n)}`$. In the other particular case $`d=4`$ the ordinary Freedman equations hold on the branes: Now $`R_{phys}`$ can be time independent (since $`\dot{\beta }_0=2\dot{\phi }_0=0)`$, and inflationary evolution happens only “parallel” to the branes in terms of an exponential $`t`$ dependence of $`a`$.
Once the equation of state of the matter on the branes differs from $`w=1`$ we can still obtain solutions with constant $`R_{phys}`$, if eq. (3.19) is satisfied. In these cases the cosmological evolution is of the standard form. For radiation dominated matter ($`w=\frac{1}{3}`$), e.g., standard cosmological evolution is obtained for $`d=0`$, i.e. when the (dominant part of the) matter action is independent of $`\phi `$. For nonrelativistic matter ($`w=0`$) standard cosmological evolution is obtained once the matter action satifies eq. (3.12b) with $`d=1`$. It remains to be seen whether these solutions constitute an alternative to the (presently ad hoc) radius fixation by potentials in the bulk, i.e. whether corresponding couplings of the Calabi-Yau modulus $`\phi `$ to matter on the branes and sufficiently well-behaved initial conditions can be obtained.
In any case we have seen that matter induced cosmological evolution in Hořava-Witten theory (albeit in its simplest version with just the universal Calabi-Yau modulus in the bulk) differs considerably from simpler scenarios as an empty bulk or a cosmological constant in the bulk. |
warning/0001/astro-ph0001254.html | ar5iv | text | # CO detection of the extremely red galaxy HR10
## 1 Introduction
The recent detections of CO emission at cosmological distances provide hints about the physical structure of newly formed objects (Combes, Maoli & Omont (1999) and reference therein). A measure of the total gas and dust mass is indeed a very useful indicator of the object evolutionary status, because it provides an estimation of the fraction of the galaxy which has yet to be turned into stars at the epoch of observation. At high redshifts such measurements, therefore, provide hints about the occurrence of active star-formation processes and help in investigating models of galaxy formation (Silk & Spaans (1997)).
Only a handful of distant objects was detected so far in CO, and most of them appear to be magnified by gravitational lenses (table 1). For these objects the mass of molecular gas inferred from the CO intensities – corrected for gravitational amplification and using the Galactic CO to H<sub>2</sub> conversion factor – turns out to be $`30÷80`$ % of the total dynamical mass, with typical values of M(H<sub>2</sub>) $`10^{10}÷10^{11}`$ M. These masses could be somewhat smaller if the CO to H<sub>2</sub> conversion factor were higher than in the Galaxy. Large quantities of CO are expected when large FIR luminosities and dust content are detected. Most of the objects detected until now in CO were indeed selected because they were dust–rich systems.
HR10 belongs to the class of objects with very red colours ($`RK>6`$). Their faintness at optical/NIR wavelengths makes the redshift determination and the investigation of their nature difficult even with 4m class telescopes. HR10 is so far the only one with a measured spectroscopic redshift (Graham & Dey (1996)). One of the main issues regarding ERGs is whether they are young and starbursting galaxies hidden in the optical by a large amount of dust or whether they are old passively evolving galaxies at $`z1`$. Recent results indicate that both classes contribute to the population of ERGs (Cimatti et al (1999)). HR10 was first detected in the submm/mm continuum with the IRAM 30m equipped with the MPIfR bolometer and with the JCMT equipped with the SCUBA double arrays (Cimatti et al (1998)). Subsequent observations confirmed the submm detection of this galaxy (Dey et al. (1999)). The inferred properties of this object show its extraordinary nature: its dust mass is $`4÷8\times 10^8`$ M and its total FIR luminosity in the range 10 –2000 $`\mu `$m rest–frame is $`2÷2.5\times 10^{12}`$ L (H$`{}_{0}{}^{}=50`$ Mpc/km/s). This places HR10 in the class of the ultra-luminous infrared galaxies and suggests the presence of a star–forming object with a SFR of $`200÷500`$ M yr<sup>-1</sup> or even higher.
In this paper we present CO(2-1) and CO(5-4) observations of HR10 made with the IRAM Plateau de Bure interferometer. The observations are described in Sect. 2, the resulting detection of both lines are presented in Sect. 3, while implications of these measurements are reported in Sect. 4. Throughout the paper, we adopt H$`{}_{0}{}^{}=50`$ Mpc/km/s and q$`{}_{0}{}^{}=0.5`$.
## 2 Observations
The observations were done partly during the winter 1998-99 and partly during the summer of 1999 with the IRAM Plateau de Bure Interferometer. A dual SIS 3mm/1.3mm receiver was used to observe simultaneously the CO(2-1) (redshifted in the 3mm band) and the CO(5-4) (redshifted in the 1.3mm band) lines. In both cases, the 3mm receiver was connected to two units of the Correlator, providing a velocity coverage of about 850 km s<sup>-1</sup>, while 4 units were connected to the 1.3mm receiver providing a velocity coverage of about 650 km s<sup>-1</sup>. The observations made during the winter 1998-99 assumed a redshift of $`z=1.443`$ deduced from the H$`\alpha `$ line. This yields frequencies of $`\nu =94.405`$ GHz for the CO(2-1) line and $`\nu =235.982`$ GHz for the CO(5-4) line. Both lines were detected, but blue-shifted by about 400 km s<sup>-1</sup> with respect to the H$`\alpha `$ line and thus truncated. The observations were repeated during the summer 1999 with a new frequency setup: $`\nu =94.531`$ GHz and $`\nu =236.297`$ GHz for the CO(2-1) and CO(5-4) lines respectively. The observations were made in the CD configuration of the interferometer. The phase drifts were calibrated by observing the nearby quasars 1633+382 and 1732+389 every 20 minutes throughout the observations. The amplitude was calibrated with the quasars 3C273, 3C345 and 2145+067, and the compact H ii region MWC349 at the beginning and/or the end of each transit. The passband of the system was calibrated using 3C273 or 2145+067. At 1.3mm, the total useful integration time on source was about 9 hours during the winter 1998/99, and 12 hours during the summer 1999. More integration time was available at 3mm because some additional data were obtained, marginal at 1.3mm but of good quality at 3mm. The 3mm receiver was tuned in single side band, while the 1.3mm was tuned in double side band, with the CO(5-4) in the lower side band. The upper side band was then available to measure the continuum.
Three visibility tables were produced from the calibrated data: two spectral tables containing the CO(2-1) and CO(5-4) data respectively and one continuum table with the 1.3mm continuum. The spectral data were smoothed in frequency to improve their signal-to-noise ratio: the velocity resolution of the CO(2-1) data was smoothed to 50 km s<sup>-1</sup> and that of the CO(5-4) to 75 km s<sup>-1</sup>. The UV tables were imaged, CLEANed and restored with elliptical gaussian beams of $`5.6^{\prime \prime }\times 4.6^{\prime \prime }`$ ($`PA=56^{}`$) and $`2.9^{\prime \prime }\times 2.0^{\prime \prime }`$ ($`PA=65^{}`$) for the 3mm and 1.3mm data respectively.
## 3 Results
Signal is detected in the central channels of both the CO(2-1) and CO(5-4) spectral data sets at 10 and 5 sigmas respectively. The averages of the channels with detected emission were used to produce the integrated maps shown on Figs. 1a and 1c. The CO(2-1) integrated flux is 1.4 Jy km s<sup>-1</sup>. There is no indication for any extension of the emission at this resolution. The peak of the source is offset by about $`1^{\prime \prime }`$ from the position of the optical source as measured by HST (Dey et al. (1999)). This offset is well within the astrometric accuracy of HST but if real it may be indicative of a spatial extinction.
The corresponding CO(2-1) profile (Fig 1b) is roughly gaussian, with a FWHM of almost exactly 400 km s<sup>-1</sup>, and a central frequency that corresponds to a redshift of $`z=1.439\pm 0.001`$. The appearance of a flat-top or even double peak profile around $`150`$ km s<sup>-1</sup> (see Fig1b) cannot be checked with the present data and only observations at higher signal-to-noise ratio (those exploiting the full spectral resolution) can settle its reality. The line-width observed is fairly large, although not atypical for this kind of sources (e.g. SMMJ02399, Frayer et al. 1998) and could be due either to an edge-on system or to various separate components. The redshift deduced from the CO line is apparently shifted from that deduced from the H$`\alpha `$ line ($`z=1.443`$), but corresponds to the redshift deduced from the \[O ii\] line at 3727 Å ($`z=1.439`$). This shift is still within the uncertainties in the optical redshift but if real it would be different from what is typically found for low redshift luminous galaxies (Sanders & Mirabel (1996)), where systematic blue-ward offsets of optical lines from the CO redshift are attributed to outflows with dust obscuration (see e.g., Gonzalez-Delgado et al. (1998)). The question remains, however, completely open since a recent analysis by McIntosh et al. (1999) of a sample of quasars shows how high-redshift objects present H$`\beta `$ lines with a systematic mean red-ward shift of $``$ 500 km/s with respect to the systemic redshift of the objects (that defined by the narrow line region). Even though the comparison with quasars may not be fair since the line emitting regions could be different, HR10 seems to show similar properties with the CO redshift corresponding to that of the narrow forbidden lines and coinciding with the centre of mass of the system, while the H$`\alpha `$ line is shifted with respect to that.
The integrated map corresponding to the CO(5-4) map (Fig. 1c) shows a source at the same position as the CO(2-1) integrated map. The integrated intensity of that source (corrected for the contribution of the continuum) is 1.35 Jy km s<sup>-1</sup>. The corresponding CO(5-4) profile (Fig. 1d) is roughly gaussian, with a FWHM of 380 km s<sup>-1</sup> similar to that of the CO(2-1) profile. The central frequency corresponds to a redshift of $`z=1.440\pm 0.001`$ similar to that deduced from the CO(2-1) line. Although the CO $`J=5`$ level is $`J(J+1)2.77\mathrm{K}=83\mathrm{K}`$ above the ground state the integrated flux (in Jy km s<sup>-1</sup>) of CO(5-4) is equal to that of the CO(2-1) line, $`(54)/(21)1`$. CO luminosities, in solar units, are $`1.5`$ and $`3.7\times 10^7`$ L for the CO(2-1) and CO(5-4) line respectively, while the total line luminosities L$`{}_{\mathrm{CO}}{}^{}{}_{}{}^{}`$ are $`4\times 10^{10}`$ and $`6\times 10^9`$ K km s<sup>-1</sup> pc<sup>2</sup> $`h_{50}^2`$, for the CO(2-1) and CO(5-4) line respectively. When expressed in these latter units, the ratio between the CO(2-1) and CO(5-4) luminosities of the same source is proportional to the line intrinsic brightness (Rayleigh-Jeans) temperature ratio integrated over the area of the source:
$$\mathrm{}=\frac{T_\mathrm{b}[\mathrm{CO}(54)]}{T_\mathrm{b}[\mathrm{CO}(21)]}=\frac{L_{\mathrm{CO}}^{}(54)}{L_{\mathrm{CO}}^{}(21)}\frac{\mathrm{\Omega }_\mathrm{s}(21)}{\mathrm{\Omega }_\mathrm{s}(54)}$$
(1)
where $`\mathrm{\Omega }_\mathrm{s}`$ is the source solid angle. If the spatial extent of the CO(5-4) emission region is similar to that of the CO(2-1) – a plausible hypothesis since both transitions have same line-width and profile – $`\mathrm{}=0.57`$ and corresponds to a value of the excitation temperature of $`T_{\mathrm{ex}}18`$ K (see e.g. Maloney, & Black (1988)). If gas and dust are in thermodynamic equilibrium the kinetic temperature $`T_{\mathrm{kin}}`$ would equal $`T_{\mathrm{dust}}`$, but it is usually found that $`T_{\mathrm{kin}}<T_{\mathrm{dust}}`$. In HR10 the dust temperature was estimated to be $`40\mathrm{K}`$ (Cimatti et al (1998); Dey et al. (1999)) and this value can be taken as the upper limit to the gas temperature. If we assume that $`T_{\mathrm{kin}}20`$ K the gas density implied by this ratio is less than $`10^3\mathrm{cm}^3`$. As an example, at $`T_{\mathrm{rad}}=2.77(1+z)=6.76`$ K and $`T_{\mathrm{kin}}=20`$ K the estimated excitation temperature is $`T_{\mathrm{ex}}=11`$ K for a CO density of $`300\mathrm{c}\mathrm{m}^3`$. The present data cannot distinguish between a picture where the dominant component of the ISM in this system is a diffuse ($`n(H_2)10^210^3\mathrm{cm}^3`$) gas or whether the medium is clumpy. The spatial shift between the CO and H$`\alpha `$ lines would be more compatible with this latter picture.
If the ratio between L$`{}_{\mathrm{CO}}{}^{}{}_{}{}^{}`$ and the mass of molecular gas is similar in HR10 and Arp 220 (Scoville et al. 1997b ), the molecular gas mass in HR10, using the CO(2-1) line, is $`\mathrm{M}(\mathrm{H}_2)=1.6\times 10^{11}\mathrm{h}_{50}^2\mathrm{M}_{}`$, larger than what is usually found in local ULIRGs (Solomon et al. (1997); Braine & Dumke (1998)) but similar to that of other detected high-z sources (Frayer et al. (1998)).
The 1.3mm continuum map shows only a marginal ($`2\sigma `$) detection at the position of the source, with an integrated flux of $`2.2\pm 0.9`$ mJy beam<sup>-1</sup>. The flux detected with the IRAM 30m telescope was $`4.9\pm 0.8`$ mJy (Cimatti et al (1998)) at 240 GHz. Scaling the flux as $`\nu ^4`$, the PdBI detection would correspond to an expected flux of $`2.4\pm 1.0`$ mJy at the observed frequency of the 30m. We exclude that this discrepancy is due to an extended component; it is more likely due to the difference in the calibration of the two instruments since the two values are consistent within the error bars.
## 4 Implications
The analysis of the spectral energy distribution of HR10 shows the presence of thermal emission at rest-frame $`\lambda >60\mu m`$ with a range of dust temperatures between 30 and 45 K. The implied total dust mass is $`84\times 10^8h_{50}^2`$ M (for a dust emissivity index $`\beta `$ of 2, Cimatti et al (1998)). Therefore the resulting gas-to-dust mass ratio for HR10 ranges between 200 and 400, as local spirals (Andreani, Casoli & Gerin (1995)), ULIRGs (Solomon et al. (1997)) and also sub-mm selected luminous sources show (Frayer et al. (1999)). The total rest-frame far-IR luminosity in the range $`102000\mu m`$ is $`22.5\times 10^{12}h_{50}^2`$ L (Cimatti et al (1999)) as estimated taking into account the ISO upper limits at 90 and 170 $`\mu m`$ (Ivison et al (1997)). When these latter are not considered and the 450 $`\mu m`$ detection is included the luminosity turns out to be a factor of 3 larger (Dey et al. (1999)). The ratio $`\frac{L_{\mathrm{FIR}}}{L_{\mathrm{CO}}^{}}`$ lies therefore in the range $`60175L_{}`$ (K km s<sup>-1</sup> pc<sup>2</sup>)<sup>-1</sup>, which agrees with the relation found for nearby luminous galaxies (Sanders & Mirabel (1996)), whose emission is mainly powered by star-formation. Objects whose FIR emission is dominated by an AGN – as the hyperluminous Infrared Galaxies – show much larger $`\frac{L_{\mathrm{FIR}}}{L_{\mathrm{CO}}^{}}`$ and do not even show up in CO (e.g., Evans et al. (1998)). This indicates that the overall FIR emission by HR10 is dominated by star formation. Assuming that most of the FIR luminosity is due to recent OB star formation activity, the star formation rate turns out to be $`\mathrm{SFR}=\mathrm{\Psi }10^{10}L_{\mathrm{FIR}}200500h_{50}^2`$ M/yr. Star formation efficiency is usually measured by the ratios $`\frac{L_{\mathrm{FIR}}}{M_{\mathrm{H}_2}}`$ and $`\frac{L_{\mathrm{H}\alpha }}{M_{\mathrm{H}_2}}`$ (see e.g. Young (1999)). While the former shows indeed quite a high value ($`1644`$) similar to that of merging local systems (Young (1999)), the latter is of only 0.007 and very likely indicates a large extinction affecting the H$`\alpha `$ emission.
With the values above for molecular mass and SFR this active phase of gas depletion lifetime should have lasted at least:
$$t_{\mathrm{gas}}=4\times 10^{10}(\frac{L_{\mathrm{FIR}}}{L_{\mathrm{CO}}^{}})^1=(0.2÷0.6)10^9\mathrm{yr}$$
(2)
The large value of gas conversion into stars (with respect to local galaxies) could be consistent with two possible scenarios: either a genuinely young galaxy in the process of active star-formation (and the detected amount of gas seems enough to feed it), or the presence of a large amount of gas could be the result of a merging process of two discs (in this latter case the resulting galaxy will have a mass of a present-day massive elliptical).
Most of the properties of HR10 suggest that the ‘locus’, which best characterizes it, is that of local ULIRGs (Hughes, Dunlop, Rawlings (1997)). HR10 follows also the expected tight correlation between the infrared flux and the radio continuum: in fact the logarithmic ratio of FIR (60 $`\mu m`$) and radio (1.5GHz) continuum flux density (in HR10 rest-frame) $`q=\mathrm{log}\frac{f_{\mathrm{FIR}}}{f_{\mathrm{radio}}}3`$ again falls within the value of nearby starbursts (Sanders & Mirabel (1996)). Furthermore, the ratio between the line (2-1) and (5-4) and the FIR luminosities, $`\frac{L_{\mathrm{CO}(21)}}{L_{\mathrm{FIR}}}=(26)\times 10^6`$ and $`\frac{L_{\mathrm{CO}(54)}}{L_{\mathrm{FIR}}}=(0.51.5)\times 10^5`$, agree with a model of CO emission in high redshift galaxies, based on an extrapolation of the properties of local ULIRGs (Blain et al. (1999)).
With the detected line width and the upper limit on the CO source size, given by the effective beam width $`\theta <5^{\prime \prime }`$, the upper limit to the total dynamical mass contained within the CO emitting region is:
$$M_{\mathrm{dyn}}=\frac{R}{2G}(\frac{\mathrm{\Delta }V}{\mathrm{sin}i})^2<7.7\times 10^{11}(\mathrm{sin}i)^2h_{50}^1M_{}$$
(3)
where $`\mathrm{\Delta }V`$ is the observed deconvolved line width, $`i`$ is the inclination and $`R`$ is the linear diameter of the source ($`R<43h_{50}^1kpc`$). The resulting dynamical mass is a factor of 5 larger than the estimation of the molecular mass. The two values would coincide if the CO emission were concentrated within the inner 10 kpc.
###### Acknowledgements.
We are grateful to Roberto Neri for his help during the data reduction and analysis, all the IRAM staff for providing us with excellent data and the referee, Francoise Combes, for helping in the interpretation of these observations. Part of this work was supported by the CNR contract progetto strategico ”Astronomia Submillimetrica” del Comitato Scienze Fisiche. P.A. thanks MPE for hospitality during 1999 when this work was done. We dedicate this paper to the memory of the 25 people who lost their lives in the tragedies of the Plateau de Bure cable-car and helicopter. |
warning/0001/gr-qc0001091.html | ar5iv | text | # On Stationary Vacuum Solutions to the Einstein Equations
## 0. Introduction.
A stationary space-time $`(M,g)`$ is a 4-manifold $`M`$ with a smooth Lorentzian metric $`g`$, of signature $`(,+,+,+),`$ which has a smooth 1-parameter group $`G`$ of isometries whose orbits are time-like curves in $`M`$. We assume throughout the paper that $`M`$ is a chronological space-time, i.e. $`M`$ admits no closed time-like curves, c.f. §1.1 for further discussion.
Let $`S`$ be the orbit space of the action $`G`$. Then $`S`$ is a smooth 3-manifold and the projection
$$\pi :MS$$
is a principle $``$-bundle, with fiber $`G`$. The chronology condition implies that $`S`$ is Hausdorff and paracompact, c.f. \[Ha\] for example. The infinitesimal generator of $`G`$ is a time-like Killing vector field $`X`$ on $`M`$, so that
$$_Xg=0.$$
The metric $`g=g_M`$ restricted to the horizontal subspaces of $`TM`$, i.e. the orthogonal complement of $`<X>TM`$ then induces a Riemannian metric $`g_S`$ on $`S`$. Since $`X`$ is non-vanishing on $`M`$, $`X`$ may be viewed as a time-like coordinate vector field, i.e. $`X=/t,`$ where $`t`$ is a global time function on $`M`$. The time function $`t`$ gives a global trivialization of the bundle $`\pi `$ and so induces a diffeomorphism from $`M`$ to $`\times S`$. The metric $`g_M`$ on $`M`$ may be then written globally in the form
(0.1)
$$g_M=u^2(dt+\theta )^2+\pi ^{}g_S,$$
where $`\theta `$ is a connection 1-form for the $``$-bundle $`\pi `$ and
(0.2)
$$u^2=X,X>0.$$
The 1-form $`\xi `$ dual to $`X`$ is thus given by $`\xi =u^2(dt+\theta ).`$ The 1-form $`\theta `$ is uniquely determined by $`g_M`$ and the time function $`t`$, but of course changes by an exact 1-form if the trivialization of $`\pi `$ is changed. We point out that $`(M,g_M)`$ is geodesically complete as a Lorentzian manifold if and only if $`(S,g_S)`$ is complete as a Riemannian manifold, c.f. Lemma 1.1.
The vacuum Einstein field equations on the space-time $`(M,g)`$ are
(0.3)
$$r_M=0,$$
where $`r_M`$ is the Ricci curvature of $`(M,g_M).`$
Stationary vacuum space-times are usually considered as the possible final, i.e. time-independent, states of evolution of a physical system, in particular isolated physical systems such as isolated stars or black holes, outside regions of matter. The most important non-trivial example is the Kerr metric, c.f. \[W\], modeling the time-independent gravitational field outside a rotating star.
It is easy to see from the field equations, c.f.(1.4) below, that there are no non-flat stationary vacuum solutions of the field equations (0.3) whose orbit space is a closed 3-manifold $`S`$. Hence, we will always assume that $`S`$ is an open 3-manifold.
Next, it is natural to consider the class of stationary vacuum space-times which are geodesically complete. In this respect, Lichnerowicz \[L, §90\] proved that any such solution $`(M,g)`$ for which the 3-manifold $`(S,g_S)`$ is complete and asymptotically flat is necessarily flat Minkowski space.
The assumption that $`S`$ is asymptotically flat is very common in general relativity in that such space-times serve as natural models for isolated physical systems, e.g. stars or black holes. The reasoning here is that as one moves further and further away from an isolated gravitational source, the corresponding gravitational field should decay as it does in Newtonian gravity, giving in the limit of infinite distance the empty Minkowski space-time.
However, mathematically the requirement that $`S`$ is asymptotically flat is a very strong assumption on both the topology and geometry of $`S`$ outside large compact sets. Further, the reasoning above is not at all rigorous. It presupposes that a geodesically complete stationary solution of the vacuum equations, i.e. a stationary solution without sources, is necessarily empty, and so in particular flat.
Consider the fact that there are geodesically complete, non-stationary vacuum space-times consisting of gravitational waves, c.f. \[MTW, §35.9\] or \[R, §8.8\] for example. Again, physically, such space-times can be considered as idealized limiting configurations at infinite distance from radiating sources. Similarly, if there does in fact exist a complete non-flat stationary vacuum solution, say $`(M_{\mathrm{}},g_{\mathrm{}}),`$ then there could well exist models $`(M,g)`$ for isolated physical systems which are asymptotic to $`(M_{\mathrm{}},g_{\mathrm{}})`$ at space-like infinity. For instance, it is not even clear apriori that the curvature of a stationary space-time, vacuum outside a compact source region, should decay anywhere at infinity.
The first main result of this paper is that in fact there are no such non-trivial stationary space-times; this of course places the physical reasoning above on stronger footing.
###### Theorem 0.1.
Let $`(M,g)`$ be a geodesically complete, chronological, stationary vacuum space-time. Then $`(M,g)`$ is the flat (i.e. empty) Minkowski space $`(^4,\eta ),`$ or a quotient of Minkowski space by a discrete group $`\mathrm{\Gamma }`$ of isometries of $`^3,`$ commuting with $`G`$. In particular, $`M`$ is diffeomorphic to $`S\times `$, $`d\theta =`$ 0 and $`u=`$ const.
This result, together with Lemma 1.1 below implies that if $`(M,g)`$ is a non-flat stationary vacuum space-time, then the orbit space $`S`$ must have a non-empty metric boundary. More precisely, since $`(S,g_S)`$ is Riemannian, let $`\overline{S}`$ denote the metric, (equivalently the Cauchy), completion of $`S`$ and let $`S=\overline{S}S.`$ Hence
(0.4)
$$\mathrm{\Sigma }=S\mathrm{},$$
if $`(M,g)`$ is not flat.
In order to avoid trivial ambiguities, we will only consider maximal stationary quotients $`S`$. For example any domain $`\mathrm{\Omega }`$ in $`^3`$ with the flat metric, $`u`$ a positive constant, and $`\theta =`$ 0 generates a stationary vacuum solution, (namely a domain in Minkowski space). In this case, the metric boundary $`\mathrm{\Omega }`$ is artificial, and has no intrinsic relation with the geometry of the solution. The solution obviously extends to a larger domain, i.e. all of Minkowski space. Thus, we only consider maximal solutions $`(S,g_S,u,\theta ),`$ in the sense that the data $`(S,g_S,u,\theta )`$ does not extend to a larger domain $`(S^{},g_S^{},u^{},\theta ^{})(S,g_S,u,\theta )`$ with $`u^{}>`$ 0 on $`M^{}.`$ It follows that in any neighborhood of a point $`q\mathrm{\Sigma }=S,`$ either the metric $`g_S`$ or the connection 1-form $`\theta `$ degenerates in some way, or $`u`$ approaches 0 in some way, or both.
Without any further restrictions, the behavior of the data near $`S`$ can be quite complicated; numerous concrete examples of this can be found among the axi-symmetric stationary, or even axi-symmetric static, i.e. Weyl, solutions; c.f. \[A1\] for further discussion. In particular, singularities, both of curvature type and of non-curvature type, may form at the boundary. The horizon $`H=\{u=0\}`$, viewed as a subset of $`S`$, may or may not be well-defined in this generality; of course it corresponds to the locus in $`M`$ where the Killing vector $`X`$ becomes null. Even when $`H`$ is well-defined and smooth, in general there may be other, possibly singular, parts to $`S.`$
Theorem 0.1 leads to the following apriori estimate on the norm of the curvature of a stationary vacuum solution away from the boundary of $`S`$, and on the rate of curvature blow-up on approach to the boundary.
###### Theorem 0.2.
There is a constant $`K<\mathrm{}`$ such that if $`(M,g)`$ is any chronological stationary vacuum solution, (not geodesically complete), then
(0.5)
$$|R_M|[x]K/\rho ^2[x],$$
where $`R_M`$ is the curvature tensor of $`(M,g)`$, $`[x]`$ is the Killing orbit through $`xM`$ and $`\rho (x)=dist_{g_S}([x],S).`$ The constant $`K`$ is independent of the data $`(M,g)`$.
Note that Theorem 0.2 implies Theorem 0.1 by letting $`\rho \mathrm{}.`$ On the other hand, Theorem 0.2 requires Theorem 0.1 for its proof. In particular, this result shows that if $`S`$ is compact in the completion $`\overline{S}`$, then the curvature of $`(M,g)`$ decays at least quadratically w.r.t. the distance from $`S`$.
The contents of the paper are as follows. We discuss some background information and preliminary results in §1, needed for the work to follow. Theorem 0.1 is proved in §2 and Theorem 0.2 is proved in §3.
I would like to thank Piotr Chrusciel and Jim Isenberg for useful discussions, the referee for pointing out some needed clarifications and Grisha Perelman for pointing out an error in a previous version of the paper.
## 1. Background and Preliminary Results.
§1.1. A stationary space-time $`(M,g)`$ uniquely determines the orbit data $`(S,g_S,u,\mathrm{\Omega })`$ described in §0, where $`\mathrm{\Omega }=d\theta `$ is the curvature 2-form of the bundle $`\pi `$ on $`S`$. Conversely, given arbitrary orbit data $`(S,g_S,u,\mathrm{\Omega })`$, $`u>`$ 0, satisfying certain equations, (c.f. (1.3)-(1.6) below), there is a unique stationary space-time $`(M,g)`$ in the sense of §0, i.e. a chronological space-time with a global isometric $``$-action with the given orbit data.
Of course, if $`(M,g)`$ is not chronological, then it will not be uniquely determined by the orbit data. One may for instance take a $``$-quotient of $`(M,g)`$, preserving the orbit data. More importantly, if $`(M,g)`$ is not chronological, then the orbit space $`S`$ may not be a manifold; even if $`S`$ is a manifold, it may not be Hausdorff, c.f. \[Ha\]. Since the arguments to follow are global on $`S`$, we require that $`S`$ is globally well-behaved, which is ensured by the chronology condition. It is not known for instance if Theorem 0.1 is valid without this assumption.
Recall that a space-time $`(M,g)`$ is geodesically complete if all geodesics in $`(M,g)`$, parametrized by an affine parameter $`s`$, are defined for all $`s.`$ The vertical subspace of $`TM`$ is the subspace spanned by the Killing field $`X`$ and the horizontal distribution $``$ is its orthogonal complement in $`TM`$, defined by the metric $`g_M.`$
###### Lemma 1.1.
A stationary space-time $`(M,g_M)`$ is geodesically complete if and only if the orbit space $`(S,g_S)`$ is geodesically complete.
Proof: Suppose $`(M,g_M)`$ is geodesically complete. Let $`\gamma `$ be a geodesic in $`S`$. Since the projection $`\pi :MS`$ is a principle fiber bundle, with horizontal spaces $`TM`$, the geodesic $`\gamma `$ may be lifted to a horizontal geodesic $`\overline{\gamma }`$ in $`(M,g_M),`$ with the same parametrization. Since $`(M,g_M)`$ is complete, $`\overline{\gamma }`$ is defined for all values of the parameter, and hence so is $`\gamma .`$
Conversely, suppose $`(S,g_S)`$ is geodesically complete, and hence complete as a metric space. Let $`\gamma `$ be a geodesic in $`M`$, with affine parameter $`s`$ and tangent vector $`T`$. Then the projection $`\sigma =\pi \gamma `$ is a curve in $`S`$, whose acceleration is given by
(1.1)
$$_VV=\frac{1}{2}\kappa ^2u^2\frac{1}{2}\kappa L(V).$$
Here $`V=d\sigma /ds=\pi _{}T,`$ is the covariant derivative in $`(S,g_S),\kappa =<X,T>=`$ const and $`L`$ is the linear map defined by $`<L(A),B>X=`$ \[A, $`B]^v`$ where $`A,B`$ are horizontal vector fields on $`M`$ and $`v`$ is the vertical projection, c.f. \[T, Ch.18.3\] for example. Conversely, any curve $`\sigma `$ satisfying (1.1) lifts to a geodesic in $`(M,g)`$.
The equations (1.1) form a $`2^{\mathrm{nd}}`$ order system of ODE w.r.t. the parameter $`s`$; note that $`L(V)`$ is linear in $`V`$, while $`\kappa `$ is a constant in $`s`$, depending linearly on $`V`$. By local existence and uniqueness, there exist locally defined solutions $`\sigma `$ for arbitrary initial data $`(x,V(x))TS.`$ Since $`S`$ is complete, it follows that $`\sigma `$ exists for all values of $`s.`$ Hence $`(M,g)`$ is geodesically complete.
###### Remark 1.2.
It is easy to verify that if $`(M,g)`$ is a stationary, (strongly) globally hyperbolic space-time, in the sense that $`(M,g)`$ admits a geodesically complete Cauchy surface $`L`$, (w.r.t. the induced metric), then $`(M,g)`$ is geodesically complete. The converse issue however, i.e. whether a chronological, stationary and geodesically complete space-time is necessarily globally hyperbolic, is not clear to the author, at least without further assumptions on $`u`$ and $`\theta `$.
For brevity, we will often say that $`(M,g_M)`$ or $`(S,g_S)`$ is complete instead of geodesically complete.
§1.2. Let $`\xi =u^2(dt+\theta )`$ be the 1-form dual to the Killing vector $`X`$, as in §0. The twist potential $`\omega `$ is the 1-form on $`M`$ defined by
(1.2)
$$\omega =\frac{1}{2}(\xi d\xi ),$$
It is easily verified that $`\omega `$ is $`G`$-invariant, and that it descends to a 1-form $`\omega `$ on the base space $`S`$. The form $`\omega `$ represents the obstruction to integrability of the horizontal distribution in $`TM`$, and so is related to the curvature 2-form $`\mathrm{\Omega }`$ of the connection 1-form $`\theta .`$ In fact, one easily verifies that
$$2\omega =u^4d\theta =u^4\mathrm{\Omega },$$
on $`(S,g_S).`$
The vacuum Einstein equations (0.3) on $`(M,g)`$ are $`G`$-invariant, and so also descend to equations on $`S`$. The vacuum equations are equivalent to the following equations on $`(S,g_S):`$
(1.3)
$$r=\frac{1}{u}D^2u+2u^4(\omega \omega |\omega |^2g),$$
(1.4)
$$\mathrm{\Delta }u=2u^3|\omega |^2,$$
(1.5)
$$div\omega =3dlogu,\omega ,$$
(1.6)
$$d\omega =0.$$
Here $`r=r_S`$ is the the Ricci curvature of $`(S,g_S),D^2u`$ is Hessian of $`u`$ on $`(S,g_S)`$, $`\mathrm{\Delta }u=tr_{g_S}D^2u`$ and $`log`$ is the natural logarithm; we refer for instance to \[Kr, Ch. 16\] for a derivation of these equations, (but note that \[Kr\] does not use the factor $`\frac{1}{2}`$ in (1.2)). The equation (1.3) comes from the pure space-like (or horizontal) part of $`r_M,`$ the equation (1.4) from the vertical part of $`r_M,`$ i.e. $`r_M(X,X),`$ while the equations (1.5)-(1.6) come from the mixed directions. The equation (1.6) implies that $`\omega `$ is locally exact, i.e. there exists $`\varphi ,`$ the twist potential, such that
(1.7)
$$2\omega =d\varphi $$
locally. On the universal cover $`\stackrel{~}{S}`$ of $`S`$, (1.7) holds globally.
Observe that these equations are invariant under the substitutions
(1.8)
$$u\lambda u,\omega \lambda ^2\omega ,$$
corresponding to $`\xi \lambda \xi ,`$ and $`\theta \lambda ^2\theta `$.
§1.3. To prove Theorems 0.1 and 0.2, we will need to study sequences of stationary (vacuum) solutions, where all the data $`(S,g_S,u,\omega )`$ are allowed to vary. Thus, in effect, we need to understand aspects of the moduli space of stationary solutions. For this, we will frequently use the following two Lemmas, which will be proved together.
###### Lemma 1.3.
(Convergence). Let $`(\mathrm{\Omega }_i,g_i,u_i,\omega _i)`$ represent data for a sequence of solutions to the stationary vacuum equations (0.1). Suppose on the domains $`(\mathrm{\Omega }_i,g_i)`$,
(1.9)
$$|r_i|\mathrm{\Lambda },diam\mathrm{\Omega }_iD,vol\mathrm{\Omega }_i\nu _o,$$
and
(1.10)
$$dist(x_i,\mathrm{\Omega }_i)\delta ,$$
for some $`x_i\mathrm{\Omega }_i`$ and positive constants $`\nu _o,\mathrm{\Lambda },D,\delta .`$ Then, for any $`\epsilon =\epsilon (\delta )>`$ 0 sufficiently small, there are domains $`U_i\mathrm{\Omega }_i,`$ with $`\epsilon /2dist(U_i,\mathrm{\Omega }_i)\epsilon ,`$ and $`x_iU_i`$ such that a subsequence of the Riemannian manifolds $`(U_i,g_i,x_i)`$ converges, in the $`C^{\mathrm{}}`$ topology, modulo diffeomorphisms, to a limit manifold $`(U,g,x)`$, with limit base point $`x=`$ lim $`x_i.`$
Further, the potentials $`u_i`$ and 1-forms $`\omega _i`$ may be renormalized by scalars $`\lambda _i,`$ as in (1.8), so that they converge smoothly to limit potential $`u`$ and 1-form $`\omega .`$ The limit $`(U,g,x,u,\omega )`$ represents a smooth solution to the stationary vacuum equations.
###### Lemma 1.4.
(Collapse). Let $`(\mathrm{\Omega }_i,g_i,u_i,\omega _i)`$ represent data for a sequence of solutions to the stationary vacuum equations (0.1). Suppose on the domains $`(\mathrm{\Omega }_i,g_i)`$,
(1.11)
$$|r_i|\mathrm{\Lambda },diam\mathrm{\Omega }_iD,vol\mathrm{\Omega }_i0$$
and
(1.12)
$$dist(x_i,\mathrm{\Omega }_i)\delta ,$$
for some $`x_i\mathrm{\Omega }_i`$ and constants $`\mathrm{\Lambda },D,\delta .`$ Then, for any $`\epsilon =\epsilon (\delta )>`$ 0 sufficiently small, there are domains $`U_i\mathrm{\Omega }_i,`$ with $`\epsilon /2dist(U_i,\mathrm{\Omega }_i)\epsilon `$ with $`x_iU_i,`$ such that $`U_i`$ is either a Seifert fibered space or a torus bundle over an interval. In both cases, the $`g_i`$-diameter of any fiber $`F`$, (necessarily a circle $`S^1`$ or torus $`T^2),`$ goes to 0 as $`i\mathrm{},`$ and $`\pi _1(F)`$ injects in $`\pi _1(U_i).`$
Consequently, the universal cover $`\stackrel{~}{U}_i`$ of $`U_i`$ does not collapse and hence has a subsequence converging smoothly to a limit $`(\stackrel{~}{U},g,x)`$, with $`x=`$ lim $`x_i^{},x_i^{}`$ a lift of $`x_i`$ to $`\stackrel{~}{U}_i.`$ In addition, the limit $`(\stackrel{~}{U},g,x)`$ admits a free isometric $``$-action.
As above, the potentials $`u_i`$ and 1-forms $`\omega _i,`$ after possible renormalization by scalars, converge smoothly to limits $`u`$ and $`\omega .`$ The limit $`(\stackrel{~}{U},g,x,u,\omega )`$ is a smooth solution of the stationary vacuum equations, and all data are invariant under a free isometric $``$-action on $`\stackrel{~}{U}.`$
Proofs: The proofs of the first parts of Lemmas 1.3 and 1.4 are essentially immediate consequences of the well-known Cheeger-Gromov theory on convergence and collapse of Riemannian manifolds with bounded curvature, c.f. \[CG1,2\], \[Ka\], \[A3,§2\] for example. We note that we are implicitly using the fact, special to dimension 3, that the full curvature is determined by the Ricci curvature.
More precisely, under the bounds (1.9)-(1.10), one obtains convergence of a subsequence of $`\{g_i\}`$ to a $`C^{1,\alpha }`$ limit metric $`g`$ on the domain $`U`$; the convergence is in the $`C^{1,\alpha ^{}}`$ topology, for any $`\alpha ^{}<\alpha <`$ 1. For a clear introduction to this theory, c.f. \[P, Ch. 10\]. In particular, the bounds (1.9) imply a lower bound on the injectivity radius of every point in $`U_i`$; this is Cheeger’s lemma, c.f. \[C\], \[P, 10.4.5\]
Under the bounds (1.11)-(1.12), the sequence of domains collapses with bounded curvature in the sense that the injectivity radius at every point in $`U_i`$ tends to 0. This implies that the domains $`U_i`$ admit an F-structure, \[CG1,2\]. In dimension 3, this means that $`U_i`$ is topologically a graph manifold, i.e. a union of Seifert fibered spaces ($`S^1`$ fibrations over a surface) or torus bundles over an interval, glued together along toral boundary components of such, c.f. \[Ro, §3\]. A result of Fukaya, c.f. \[F, Ch.11,12\] and references therein, implies that on domains of bounded diameter, i.e. under (1.11)-(1.12), for $`i`$ sufficiently large, the F-structure may be chosen to be pure, so that $`U_i`$ itself is either a Seifert fibered space or a torus bundle over an interval. The collapse takes place by shrinking the fibers, (circles or tori), to points. From the theory of Seifert fibered spaces, c.f. \[O\] or \[Ro, Thm. 4.3\], the fibers inject in $`\pi _1`$ whenever $`U_i`$ is not covered by $`S^3`$. But this is necessarily the case here, since $`U_i`$ is an open domain, (c.f. the remark following (0.3). Thus, one may unwrap the collapse by passing to covers, for instance the universal cover, that unwind the fibers. This ability to unwrap collapse on domains of controlled diameter is special to dimension 3.
It remains to show that the convergence is actually smooth $`(C^{\mathrm{}}),`$ and that the limit, in either case of Lemma 1.3 or 1.4, is a smooth solution to the stationary vacuum equations. This is done by showing that the equations (1.3)-(1.6) form essentially an elliptic system and using elliptic regularity.
By taking the trace of (1.3) and using (1.4), one derives that
(1.13)
$$s=6u^4|\omega |^2,$$
where $`s`$ is the scalar curvature of $`(S,g_S)`$, so that (1.4) is equivalent to
(1.14)
$$\mathrm{\Delta }u=\frac{s}{3}u.$$
Since, by hypothesis, the Ricci curvature is uniformly bounded on $`(\mathrm{\Omega }_i,g_i),`$ so is the scalar curvature $`s_i.`$ Now the potential functions $`u_i`$ may be unbounded, or converge to 0, in neighborhoods of the base points $`x_i.`$ Thus, we renormalize $`u_i`$ by setting
(1.15)
$$\overline{u}_i=u_i/u(x_i),$$
so that $`\overline{u}_i(x_i)=`$ 1. The equation (1.14) is of course invariant under this renormalization. Moreover, since $`u_i>`$ 0 everywhere, and since the local geometry of $`(\mathrm{\Omega }_i,g_i)`$ is uniformly controlled in $`C^{1,\alpha }`$ away from $`\mathrm{\Omega }_i`$, i.e. within $`U_i`$, the Harnack inequality, (c.f. \[GT, Thm. 8.20\]), applied to the elliptic equation (1.14) implies that there is a constant $`\kappa >`$ 0, independent of $`i`$, such that
(1.16)
$$\kappa \frac{sup\overline{u}_i}{inf\overline{u}_i}\kappa ^1;$$
here the $`sup`$ and $`inf`$ are taken over $`U_i`$, or more precisely over an $`\epsilon /4`$ thickening of $`U_i`$. Of course the diameter bound in (1.9) or (1.11) is being used here. It then follows from $`L^2`$ elliptic theory, c.f. \[GT, Thm. 9.11\], that the functions $`\overline{u}_i`$ are uniformly bounded in $`L^{2,p}(U_i),p<\mathrm{}`$. Next, as in (1.15), we renormalize the twist 1-forms $`\omega _i`$ by
(1.17)
$$\overline{\omega }_i=\omega _i/(u(x_i))^2,$$
c.f. (1.8). It then follows from (1.13), (1.15), (1.17) and the uniform $`L^{\mathrm{}}`$ bound on $`s_i`$ that the forms $`\overline{\omega }_i`$ are uniformly bounded in $`L^{\mathrm{}}`$ on $`U_i`$.
Next, to obtain higher regularity, consider the equations (1.5)-(1.6)
$$\mathrm{\Delta }\varphi _i=3dlog\overline{u}_i,d\varphi _i,$$
locally, i.e. in neighborhoods where the twist potential $`\varphi =\varphi _i`$ is defined; (we omit the overbar from the notation for $`\varphi `$). We may add a constant to $`\varphi _i`$ and assume $`\varphi _i(x_i)=`$ 0. By the bound on $`\overline{\omega }_i`$ above, $`|d\varphi _i|`$ is uniformly bounded, as is $`|dlog\overline{u}_i|,`$ so by elliptic regularity, $`\varphi _i`$ is bounded locally in $`L^{2,p},`$ and hence $`\overline{\omega }_i`$ is uniformly bounded locally in $`L^{1,p}`$ everywhere in $`U_i`$. By (1.13) again, this implies $`s_i`$ is bounded in $`L^{1,p},`$ and so by elliptic regularity applied to (1.14), $`\overline{u}_i`$ is uniformly bounded locally in $`L^{3,p}.`$ Hence, the right side of (1.3) is bounded in $`L^{1,p},`$ and so the Ricci curvature $`r_i`$ is uniformly controlled locally in $`L^{1,p}`$ everywhere in $`U_i`$. This implies that the metrics $`g_i`$ are uniformly controlled in $`L^{3,p}`$ in local harmonic coordinates, c.f. \[A3. §3\] for example. Hence, by the Sobolev embedding theorem, the sequence $`\{g_i\}`$ is uniformly bounded in $`C^{2,\alpha }`$, $`\alpha <1`$.
This process may now be iterated inductively to give uniform $`C^k`$ control on $`\{g_i\},`$ for any $`k<\mathrm{},`$ away from the boundary, as well as uniform $`C^k`$ control on $`\{\overline{u}_i\}`$ and on $`\{\overline{\omega }_i\}.`$ This proves that the convergence to the limit is in the $`C^{\mathrm{}}`$ topology, as well as $`C^{\mathrm{}}`$ convergence to limits $`\overline{u}`$ and $`\overline{\omega }.`$ Since the metrics $`g_i`$ are stationary vacuum solutions, it is obvious that the limit $`(U,g,\overline{u},\overline{\omega })`$ is also.
As an application of these results, we prove the following Lemma, which shows that a given complete stationary vacuum solution gives rise to another one with uniformly bounded curvature.
###### Lemma 1.5.
Let $`(S,g,u,\omega ),g=g_S,`$ represent data for a complete non-flat stationary vacuum solution. Then there exists another complete non-flat stationary vacuum solution given by data $`(S^{},g^{},u^{},\omega ^{})`$, $`g^{}=g_S^{}^{},`$ obtained as a geometric limit at infinity of $`(S,g)`$, which has uniformly bounded curvature, i.e.
(1.18)
$$|r_g^{}|1\mathrm{and}|r_g^{}|(y)>0,$$
for some $`yS^{}`$.
Proof: We may assume that $`(S,g)`$ itself has unbounded curvature, for otherwise there is nothing to prove since (1.18) can then be obtained by a fixed rescaling of $`(S,g)`$ if necessary. Let $`\{x_i\}`$ be a sequence in $`S`$ such that
(1.19)
$$|r|(x_i)\mathrm{},\mathrm{as}i\mathrm{}.$$
Let $`B_i=B_{x_i}(1)`$ and let $`d_i(x)=dist(x_i,B_i).`$ Consider the scale-invariant ratio $`(d_i^2|r|)(x),`$ for $`xB_i,`$ and choose points $`y_iB_i`$ realizing the maximum value of $`(d_i^2|r|)(x)`$ on $`B_i.`$ Since $`(d_i^2|r|)(x)`$ is 0 on $`B_i,y_i`$ is in the interior of $`B_i.`$ By (1.19), we have
(1.20)
$$d_i^2(y_i)|r|(y_i)\mathrm{},\mathrm{as}i\mathrm{}$$
and so in particular $`|r|(y_i)\mathrm{}.`$
Now consider the pointed rescaled sequence $`(B_i,g_i,y_i),`$ where
$$g_i=|r|(y_i)g.$$
By construction, $`|r_i|(y_i)=`$ 1, where $`r_i`$ is the Ricci curvature of $`g_i.`$ This, together with (1.20) and its scale-invariance, implies that $`\delta _i(y_i)dist_{g_i}(y_i,B_i)\mathrm{}.`$ Further, by the maximality property of $`y_i,`$
(1.21)
$$|r_i|(x)|r_i|(y_i)\frac{\delta _i(x)}{\delta _i(y_i)}=\frac{\delta _i(x)}{\delta _i(y_i)}.$$
It follows from (1.20) that $`|r_i|(x)`$ 2, at all points $`x`$ of uniformly bounded $`g_i`$-distance to $`y_i,`$ (for $`i`$ sufficiently large, depending on $`dist_{g_i}(x,y_i)).`$
If the pointed sequence $`(B_i,g_i,y_i),`$ (or a subsequence), is not collapsing at $`y_i,`$ i.e. the volume of the unit $`g_i`$-ball at $`y_i`$ is bounded below as $`i\mathrm{},`$ then by Lemma 1.3, $`\{(B_i,g_i,y_i)\}`$ has a subsequence converging, smoothly and uniformly on compact subsets, to a limit $`(U^{},g^{},y)`$, $`y=`$ lim $`y_i.`$ The limit is a complete stationary vacuum solution, (since $`\delta _i(y_i)\mathrm{}),`$ and by the smooth convergence, $`|r_g^{}|`$ 2 everywhere and $`|r_g^{}(y)|=1`$, where $`y`$ = lim$`y_i`$. A further bounded rescaling then gives (1.18). The limit potential $`u`$ and twist form $`\omega `$ are obtained as in Lemma 1.3.
On the other hand, suppose this sequence is collapsing at $`y_i,`$ so that the volume of the unit $`g_i`$-ball at $`y_i`$ converges to 0, (in some subsequence). Then by Lemmas 1.3 and 1.4, it is collapsing everywhere within $`g_i`$-bounded distance to $`y_i`$, i.e. within $`(B_{y_i}(R),g_i)`$, for any fixed $`R<\mathrm{}`$. For any such $`R`$, if $`i`$ is sufficiently large, there are domains $`U_i(R)B_{y_i}(R)`$, with $`U_i(R)`$ near $`B_{y_i}(R)`$ w.r.t. $`g_i`$, which are highly collapsed along an injective Seifert fibered structure or torus bundle structure on $`U_i(R)`$. Hence the universal cover $`(\stackrel{~}{U}_i(R),\stackrel{~}{g}_i)`$ is not collapsing. For any sequence $`R_j\mathrm{}`$, there is then a suitable diagonal subsequence $`U_{i_j}`$ such that the covers $`\stackrel{~}{U}_{i_j}`$ converge smoothly, as above, to a complete stationary vacuum solution; again a bounded rescaling then gives (1.18).
## 2. Proof of Theorem 0.1.
Let $`(M,g_M)`$ be a complete stationary vacuum solution. As above in §1.2 and §1.3, we will work exclusively on the 3-manifold quotient $`S`$, with data $`u`$, $`\omega `$ and $`g`$ satisfying the field equations (1.3)-(1.6). By passing to the universal cover, we may and will assume for this section that $`S`$ is simply connected.
It is very useful to rewrite the metric $`g_M`$ in (0.1) in the form
(2.1)
$$g_M=u^2(dt+\theta )^2+\frac{1}{u^2}\overline{g}_S,$$
where $`\overline{g}_S`$ is the conformally equivalent metric
(2.2)
$$\overline{g}_S=u^2g_S$$
on $`S`$. Using standard formulas for behavior under conformal changes, c.f. \[B, Ch. 1J\], w.r.t this metric the field equations (1.3)-(1.5) are equivalent to:
(2.3)
$$\overline{r}=2(dlogu)^2+2u^4(\omega )^2,$$
(2.4)
$$\overline{\mathrm{\Delta }}logu=2u^4|\omega |^2,$$
(2.5)
$$div\omega =4dlogu,\omega ,$$
c.f. also \[Kr, Ch. 16\]. All metric quantities in (2.3)-(2.5) are w.r.t. the $`\overline{g}=\overline{g}_S`$ metric.
There are two reasons for preferring $`\overline{g}`$ to $`g=g_S.`$ First, it is apparent from (2.3) that
(2.6)
$$\overline{r}0,$$
so that $`(S,\overline{g})`$ has non-negative Ricci curvature. Second, the field equations (2.3)-(2.5) are exactly the Euler-Lagrange equations for the functional
(2.7)
$$S_{eff}=_S(s\frac{1}{2}(\frac{|du^2|^2+|d\varphi |^2}{u^4}))𝑑V.$$
Here we are using the fact that $`S`$ is simply connected, so that the relation (1.7) holds globally on $`S`$. This functional is the Einstein-Hilbert functional on $`G`$-invariant metrics on $`M`$, dimensionally reduced to a functional on data $`(\overline{g},u,\varphi )`$ on $`S`$, when $`g_M`$ is expressed in the form (2.1). It corresponds to a coupling of 3-dimensional gravity to the energy (or $`\sigma `$-model) of the mapping $`E=(\varphi ,u^2)`$ from $`S`$ to the hyperbolic plane. The mapping $`E`$ is called the Ernst potential and the Euler-Lagrange equations (2.3)-(2.5) imply that
(2.8)
$$E:(S,\overline{g}_S)(H^2(1),g_1)$$
is a harmonic map. Here $`H^2(1)`$ is the hyperbolic plane, given as the upper half-plane $`(^2)^+=\{(x,y):y>0\}`$, with metric
(2.9)
$$g_1=\frac{dx^2+dy^2}{y^2}.$$
We refer for instance to \[H1\] or \[H2\] for further details and discussion on $`S_{eff}.`$
From the equation (2.3), we see that
(2.10)
$$\overline{r}=\frac{1}{2}E^{}(g_1).$$
In particular, the energy density of $`e(E)`$ of $`E`$, given by
$$e(E)=\frac{1}{2}|E_{}|^2$$
satisfies
(2.11)
$$\overline{s}=e(E)=\frac{1}{2}tr_{\overline{g}}E^{}(g_1).$$
For clarity, we break the proof up at this stage into two steps.
Step I. Assume the metric $`(S,\overline{g}_S)`$ is complete.
The space $`(S,\overline{g}_S)`$ may or may not have uniformly bounded curvature, i.e. possibly after a bounded rescaling,
(2.12)
$$|\overline{r}|1,$$
everywhere on $`S`$, where the norm is taken w.r.t. $`\overline{g}_S.`$ If (2.12) holds, then the arguments below are applied to $`(S,\overline{g}_S)`$. If instead the curvature of $`(S,\overline{g}_S)`$ is unbounded, (and hence $`(S,\overline{g}_S)`$ is not flat), we apply Lemma 1.5 to obtain a new non-flat stationary space-time $`(S^{},\overline{g}_S^{},u^{},\omega ^{})`$ satisfying (2.12). The arguments below are then applied to $`(S^{},\overline{g}_S^{})`$.
With this understood, we drop the prime from the notation and assume that $`(S,\overline{g}_S)`$ satisfies (2.12).
We now apply the well-known Bochner formula, c.f. \[EL, (3.12)\], to the harmonic Ernst map $`E`$, to obtain
(2.13)
$$\overline{\mathrm{\Delta }}e(E)=|\overline{}DE|^2+r_M,E^{}(g_1)\underset{i,j=1}{\overset{3}{}}(E^{}R_1)(e_i,e_j,e_j,e_i).$$
Here the sign of the curvature tensor for the last term is such that $`R_1(X,Y,Y,X)`$ is the sectional curvature of $`g_1`$ for an orthonormal pair $`(X,Y)`$. We claim that the last two terms in (2.13) are given by
(2.14)
$$\overline{r},E^{}(g_1)=2|\overline{r}|^2,$$
(2.15)
$$(E^{}R_{g_1})(e_i,e_j,e_j,e_i)=4(\overline{s}^2|\overline{r}|^2)0.$$
The equation (2.14) follows immediately from (2.10). For (2.15), using the fact that $`g_1`$ is of constant sectional curvature $`1`$, we have $`(E^{}R_1)(e_i,e_j,e_j,e_i)`$ = $`g_1(E_{}e_i,E_{}e_i)g_1(E_{}e_j,E_{}e_j)g_1(E_{}e_i,E_{}e_j)^2`$. Choosing $`\{e_i\}`$ to be an orthonormal basis in $`(S,\overline{g}_S)`$ diagonalizing the Ricci curvature $`\overline{r}`$, and using (2.10), gives (2.15).
In particular, the equations (2.13)-(2.15) show that the energy density $`e(E)`$ is a subharmonic function on $`(S,\overline{g}_S)`$.
Since (2.12) holds on $`(S,\overline{g}_S)`$, (2.11) implies that $`e(E)`$ is uniformly bounded above on $`(S,\overline{g}_S)`$. Thus, let $`\{x_i\}`$ be a maximizing sequence for $`e(F)`$, i.e.
(2.16)
$$e(F)(x_i)supe(F)<\mathrm{}.$$
Since the curvature of $`(S,\overline{g}_S)`$ is bounded, and this space is complete, it follows from elementary properties of the Laplacian that
$$\mathrm{\Delta }e(F)(x_i)\epsilon _i,$$
where $`\epsilon _i`$ 0, as $`i\mathrm{}.`$ However, (2.13)-(2.15) then imply that
$$|\overline{r}|^2(x_i)\epsilon _i0.$$
This of course forces $`e(E)(x_i)=\overline{s}(x_i)0`$. Since $`x_i`$ is a maximizing sequence, this is only possible if
$$e(E)0,$$
i.e. $`E`$ is a constant map. This means that $`u=`$ const $`>`$ 0, $`\varphi =`$ const, and hence $`(M,g)`$ is flat. Thus $`(M,g)`$ is Minkowski space, (since $`S`$ is simply connected).
Observe that this argument now implies that the passage to the geometric limit $`(S^{},\overline{g}_S^{})`$ at the beginning of Step I was not in fact necessary.
Step II. We now remove the assumption that $`\overline{g}`$ is complete, by transfering the estimates above back to the complete manifold $`(S,g_S)`$.
Exactly as in the beginning of Step I however, since $`(S,g_S)`$ is complete, if necessary we use Lemma 1.5 first to pass to a non-flat geometric limit $`(S^{},g_S^{})`$ with uniformly bounded $`g^{}`$-curvature, i.e. satisfying (1.18). As before, we drop the prime from the notation below.
Since $`\overline{g}_S=u^2g_S,`$ we have the following relation between the Laplacians of $`g_S`$ and $`\overline{g}_S`$, c.f. \[B, Ch. 1J\] for example:
$$\overline{\mathrm{\Delta }}f=u^2\mathrm{\Delta }f+u^3du,df,$$
for any function $`f`$, where metric quantities on the right are w.r.t. $`g_S`$. Setting $`f=\overline{s}`$ then gives
(2.17)
$$\mathrm{\Delta }\overline{s}=u^2\overline{\mathrm{\Delta }}\overline{s}dlogu,d\overline{s}.$$
Now the function $`\overline{s}`$ may well be an unbounded function on $`(S,g_S)`$; (in fact the unboundedness may cause the incompleteness of $`\overline{g}_S).`$ However, in terms of the metric $`g`$, we have
(2.18)
$$\overline{s}=u^2(2|dlogu|^2+\frac{1}{2}u^4|d\varphi |^2)u^2h,$$
where the last inequality defines $`h`$ and the norms on the right are w.r.t. $`g_S`$. This follows by taking the trace of (2.3).
Since the curvature of $`g_S`$ is uniformly bounded, i.e. (1.18) holds, the same arguments as in the proof of Lemma 1.3-1.4 imply that
(2.19)
$$|dlogu|^2+u^4|d\varphi |^2C,$$
for some $`C<\mathrm{}.`$ The estimate (2.19) can also be deduced directly from (1.13) and (1.3)-(1.7). Hence, $`h`$ is uniformly bounded above on $`(S,g_S)`$.
Returning to (2.17), we then have
(2.20)
$$\mathrm{\Delta }\overline{s}=\mathrm{\Delta }u^2h=u^2\mathrm{\Delta }h+h\mathrm{\Delta }u^2+2du^2,dh.$$
Now
(2.21)
$$\mathrm{\Delta }u^2=2u^3\mathrm{\Delta }u+6u^4|du|^2=u^6|d\varphi |^2+6u^4|du|^2,$$
where the last equality uses (1.4) and (1.7). Hence, combining (2.20)-(2.21), we obtain
$$\mathrm{\Delta }h=u^2\mathrm{\Delta }\overline{s}(u^4|d\varphi |^2+6u^2|du|^2)h2u^2du^2,dh.$$
Substituting (2.17) gives
(2.22)
$$\mathrm{\Delta }h=u^4\overline{\mathrm{\Delta }}\overline{s}(u^4|d\varphi |^2+6u^2|du|^2)h2u^2du^2,dhu^2dlogu,d\overline{s}.$$
Since $`\overline{s}=u^2h,d\overline{s}=2u^3hdu+u^2dh,`$ and so (2.22) becomes
$$\mathrm{\Delta }h=u^4\overline{\mathrm{\Delta }}\overline{s}(u^4|d\varphi |^2+6u^2|du|^2)h+4dlogu,dh+2u^2h|du|^2dlogu,dh,$$
i.e.
$$\mathrm{\Delta }h=u^4\overline{\mathrm{\Delta }}\overline{s}(u^4|d\varphi |^2+4u^2|du|^2)h+3dlogu,dh.$$
By (2.18) again, the middle term on the right above equals $`2h^2=2u^4\overline{s}^2`$. Hence, we have
(2.23)
$$\mathrm{\Delta }h3dlogu,dh=u^4\overline{\mathrm{\Delta }}\overline{s}2u^4\overline{s}^2.$$
On the other hand, from the Bochner formula (2.13) and (2.14)-(2.15), we have
$$\overline{\mathrm{\Delta }}\overline{s}=|\overline{}DE|^2+2|\overline{r}|^2+4(\overline{s}^2|\overline{r}|^2),$$
where all quantities are w.r.t. the $`\overline{g}`$ metric. Substituting this in (2.23) then gives
(2.24)
$$\mathrm{\Delta }h3dlogu,dh=u^4|\overline{}DE|^2+2u^4(\overline{s}^2|\overline{r}|^2)0,$$
where the terms on the left are in the $`g`$ metric while those on the right are in the $`\overline{g}`$ metric.
We now basically repeat the argument above in Step I to prove that
(2.25)
$$h0.$$
Thus, recalling from (2.19) that $`h`$ is bounded on $`(S,g_S)`$, let $`\{x_i\}`$ be a maximizing sequence for $`h`$. It follows as before that $`\mathrm{\Delta }h(x_i)ϵ_i,|dh|(x_i)ϵ_i`$ while $`|dlogu|(x_i)`$ remains uniformly bounded.
To prove (2.25), it is most convienient to pass to the limit of the pointed sequence $`(S,g_S,x_i)`$ by use of Lemmas 1.3-1.4. Thus, a subsequence of $`\{(S,g_S,x_i)\}`$ converges smoothly, (passing to covers if necessary in the case of collapse), to a complete stationary vacuum solution $`(S_{\mathrm{}},g_{\mathrm{}},x_{\mathrm{}})`$. Here the limit potentials $`u_{\mathrm{}}`$ and $`\varphi _{\mathrm{}}`$ are limits of the renormalized potentials $`u_i=u/u(x_i),\varphi _i=\varphi /u(x_i)^2`$. Observe that $`h`$ and $`dlogu`$ are invariant under such renormalizations, as is the right side of (2.24) under the changes $`uu_i,\overline{g}_S\overline{g}_i=u_i^2g_S`$.
It follows from these estimates and (2.24), together with the maximum principle, that the limit $`(S_{\mathrm{}},g_{\mathrm{}},x_{\mathrm{}},u_{\mathrm{}},\varphi _{\mathrm{}})`$ satisfies
(2.26)
$$hh_{\mathrm{}}=const,|\overline{}DE|=0,|\overline{r}|^2\overline{s}^2=0,$$
where $`\overline{g}_{\mathrm{}}=u_{\mathrm{}}^2g_{\mathrm{}}`$ and
(2.27)
$$h_{\mathrm{}}=sup_Sh.$$
To see that $`h_{\mathrm{}}=0`$, (2.26) and (2.10) imply that $`\overline{}\overline{r}=0`$, i.e. the Ricci curvature $`\overline{r}_{\mathrm{}}`$ of $`\overline{g}_{\mathrm{}}`$ is parallel. By the Bianchi identity this implies that the scalar curvature $`\overline{s}_{\mathrm{}}`$ of $`\overline{g}_{\mathrm{}}`$ is constant. Since $`h=h_{\mathrm{}}`$ is constant, (2.18) shows that $`u_{\mathrm{}}`$ is also constant on $`(S_{\mathrm{}},g_{\mathrm{}})`$. Hence by (2.4) on $`(S_{\mathrm{}},g_{\mathrm{}})`$, it follows that $`d\varphi _{\mathrm{}}=0`$. By the definition of $`h`$ in (2.18), this of course gives $`h_{\mathrm{}}0`$, which by (2.27) gives (2.25).
The equation (2.25) means that $`u`$ is a constant function and $`\omega =`$ 0, so that $`d\theta =`$ 0. It follows that $`(S,g_S)`$ and $`(M,g_M)`$ are both flat, which proves the result.
## 3. Proof of Theorem 0.2.
The following result gives Theorem 0.2 essentially as an immediate corollary. The proof is a standard consequence of the global result in Theorem 0.1, together with the control on moduli of stationary vacuum solutions given in Lemmas 1.3 and 1.4.
###### Theorem 3.1.
Let $`(M,g_M)`$ be a stationary vacuum solution, with orbit data $`(S,g_S,u,\theta ),`$ and $`US`$ a domain with smooth boundary, so that $`u>`$ 0 on $`\overline{U}.`$ Then there is an (absolute) constant $`K<\mathrm{},`$ independent of $`(M,g_M)`$ and $`U`$, such that for all $`xU,`$
(3.1)
$$|r_S|(x)\frac{K}{\rho (x)^2},$$
where $`\rho (x)=dist_{g_S}(x,U).`$
Proof: The proof is by contradiction. Thus, assume that (3.1) does not hold. Then there are stationary vacuum solutions $`(M_i,g_{M_i}),`$ with orbit data $`(S_i,g_{S_i},u_i,\omega _i),`$ smooth domains $`U_iS_i`$ on which $`u_i>`$ 0 and points $`x_iU_i`$ such that
(3.2)
$$\rho ^2(x_i)|r_i|(x_i)\mathrm{},\mathrm{as}i\mathrm{}.$$
Let $`\rho _i=\rho (x_i).`$ Since it may not be possible to choose the points $`x_i`$ so that they maximize $`|r_i|`$ (over large domains), we shift the base points $`x_i`$ as follows; compare with the proof of Lemma 1.5. Choose $`t_i[0,\rho _i)`$ such that
(3.3)
$$t_i^2sup_{B_{x_i}(\rho _it_i)}|r_i|=sup_{t[0,\rho _i)}t^2sup_{B_{x_i}(\rho _it)}|r_i|\mathrm{},\mathrm{as}i\mathrm{},$$
where the last estimate follows from (3.2), (set $`t=\rho _i).`$ Let $`y_iB_{x_i}(\rho _it_i)`$ be points such that
(3.4)
$$|r_i|(y_i)=sup_{B_{x_i}(\rho _it_i)}|r_i|.$$
Further, setting $`t=t_i(1\frac{1}{k}),k>`$ 1, in (3.3), one obtains the estimate
(3.5)
$$t_i^2|r_i|(y_i)t_i^2(1\frac{1}{k})^2sup_{B_{x_i}(\rho _it_i(1\frac{1}{k}))}|r_i|t_i^2(1\frac{1}{k})^2sup_{B_{y_i}(t_i/k)}|r_i|,$$
so that
(3.6)
$$sup_{B_{y_i}(t_i/k)}|r_i|(1\frac{1}{k})^2|r_i|(y_i),$$
Now rescale or blow-up the metric so that $`|\stackrel{~}{r}_i|(y_i)=`$ 1 by setting $`\stackrel{~}{g}_i=|r_i|(y_i)g,`$ and consider the pointed sequence $`(U_i,\stackrel{~}{g}_i,y_i).`$ We have
(3.7)
$$|\stackrel{~}{r}_i|(y_i)=1,$$
and by (3.3) and scale invariance,
(3.8)
$$dist_{\stackrel{~}{g}_i}(y_i,U_i)\mathrm{},\mathrm{as}i\mathrm{}.$$
Also, (compare with (1.21)), it follows from (3.6) that
(3.9)
$$|\stackrel{~}{r}_i|(x)C(dist_{\stackrel{~}{g}_i}(x,y_i)).$$
We also normalize $`u`$ by setting
(3.10)
$$\stackrel{~}{u}_i(x)=\frac{u(x)}{u(y_i)},$$
and note that $`\stackrel{~}{u}_i>`$ 0 on $`U_i.`$
We may now apply Lemmas 1.3 and 1.4, exactly as in the proof of Lemma 1.5 to conclude that a subsequence of the pointed sequence $`(U_i,\stackrel{~}{g}_i,\stackrel{~}{u}_i,\stackrel{~}{\omega }_i,y_i)`$ converges in the $`C^{\mathrm{}}`$ topology on compact subsets, to a limit stationary vacuum solution $`(U_{\mathrm{}},\stackrel{~}{g}_{\mathrm{}},\stackrel{~}{u}_{\mathrm{}},\stackrel{~}{\omega }_{\mathrm{}},y)`$, which is complete and satisfies $`\stackrel{~}{u}_{\mathrm{}}>`$ 0 everywhere. Here, one must pass to the universal cover in case of collapse, as in Lemma 1.4, and the potential $`\stackrel{~}{u}_i`$ and 1-form $`\stackrel{~}{\omega }_i`$ are normalized so that $`\stackrel{~}{u}_i(y_i)=`$ 1 and $`|\stackrel{~}{\omega }_i(y_i)|`$ is bounded.
By Theorem 0.1, $`\stackrel{~}{g}_{\mathrm{}}`$ must be flat, $`\stackrel{~}{u}_{\mathrm{}}`$ constant and $`d\stackrel{~}{\omega }_{\mathrm{}}=`$ 0. However, the smooth convergence of the sequence $`(U_i,\stackrel{~}{g}_i)`$ guarantees that the equality (3.7) passes to the limit, contradicting the fact that $`\stackrel{~}{g}_{\mathrm{}}`$ is flat.
As in the proof of Lemmas 1.3 and 1.4, it follows from (3.1) that
(3.11)
$$|dlogu|(x)\frac{K}{\rho (x)},$$
and
(3.12)
$$u^2|\omega |(x)\frac{K}{\rho (x)}.$$
Combining the estimates (3.1) and (3.11)-(3.12), one obtains the same bound on the full curvature tensor $`R_M`$ of $`(M,g)`$.
Note that since $`K`$ is independent of the domain $`U`$, (3.1) holds for $`\rho `$ the distance to the boundary $`\mathrm{\Sigma }`$ of $`S`$, even if $`\mathrm{\Sigma }`$ is singular. To see this, just apply Theorem 3.1 to a smooth exhaustion $`U_j`$ of $`S`$, with $`U_j`$ converging to $`S`$ in the Hausdorff metric on subsets of $`(S,g_S)`$. In particular, these results together prove Theorem 0.2.
We note that elliptic regularity further implies that, for any $`j`$ 1,
(3.13)
$$|^jR_M|(x)\frac{K(j)}{\rho ^{2+j}(x)},|^jlogu|(x)\frac{K(j)}{t^j(x)}.$$
Theorem 0.2, when combined with Lemmas 1.3 and 1.4, shows that the moduli space of stationary vacuum solutions is apriori well-controlled away from the boundary $`\mathrm{\Sigma }=S`$. Thus, away from the boundary, sequences of such metrics either have a smoothly convergent subsequence, or they collapse, in which case the universal covers have a convergent subsequence.
Theorems 0.1 and 0.2 give new proofs of similar results for static vacuum solutions in \[An2, Thm. 3.2\]. Similarly, in work to follow, we plan to generalize the results on the asymptotic structure of static vacuum space-times in \[A1\] to stationary space-times as well as consider the Riemannian analogues of these questions.
October 1999/March 2000 |
warning/0001/hep-ph0001321.html | ar5iv | text | # CALCULATION OF DIRECT CP VIOLATION IN B DECAYS
## 1 Introduction
The sourse of CP violation is one of the unsolved problem in the standard model (SM). The richness of charmless hadronic decays of B meson provides a good place for study of CP violation . When $`B^0`$ and $`\overline{B}^0`$ decay to a common final state $`f`$, $`B^0`$-$`\overline{B}^0`$ mixing plays a crucial role. It interferences with the direct CP asymmetries. For other decays, $`B^0`$ and $`\overline{B}^0`$ decay to different final states, for example $`B^0K^+\pi ^{}`$, $`\overline{B}^0K^{}\pi ^+`$. No mixing is involved here. They are similar to charged $`B^\pm `$ decays. CP asymmetry has no time dependence. The direct CP asymmetry is important even if for neutral B meson decays.
If there is only one amplitude contributing to the decay, both the strong phase and weak phase can be factored out. we have $`\mathrm{\Gamma }=\overline{\mathrm{\Gamma }}`$. So there is no direct CP asymmetry. That is the case for D meson decays and B meson going to heavy meson decays, like $`BJ/\psi K_S`$. If there are two amplitudes, the decay rate of $`\mathrm{\Gamma }`$ and $`\overline{\mathrm{\Gamma }}`$ may be different. If the strong phase difference between the two amplitudes $`M_1`$ and $`M_2`$ is not zero ($`\delta _{12}0`$) and the weak phase difference of the two amplitudes is also non zero ($`\varphi _{12}0`$), we have $`\mathrm{\Gamma }\overline{\mathrm{\Gamma }}`$. The direct CP asymmetry is
$`A_{CP}={\displaystyle \frac{2r\mathrm{sin}\delta _{12}\mathrm{sin}\varphi _{12}}{1+r^2+2r\mathrm{cos}\delta _{12}\mathrm{cos}\varphi _{12}}},`$ (1)
where $`r=|M_2|/|M_1|`$. $`A_{CP}`$ depends on $`2r/(1+r^2)`$, $`\mathrm{sin}\delta _{12}`$ and $`\mathrm{sin}\varphi _{12}`$. If one of the three parameters is small, then $`A_{CP}`$ is small. In many decays, we do not have all these conditions, then there is no sizable direct CP violation. However, most charmless decays have large values for $`2r/(1+r^2)`$, where $`M_1`$ is tree amplitude and $`M_2`$ is penguin amplitude. Furthermore, the CKM parameters for the tree diagram and penguin diagrams are different providing weak phase differences.
Direct CP asymmetries require an interference between two amplitudes involving both a CKM phase and a final state strong interaction phase difference. The weak phase difference arises from the superposition of penguin contributions and the tree diagrams. The strong-phase difference arises through the perturbative penguin diagrams (hard final state interaction), or non-perturbatively (soft final state interaction). The soft part is not important which is shown in some model calculations . There are also some other contributions, such as annihilation diagram and Soft final state interaction. Mostly, their contributions to branching ratios are small . Probably their contribution to $`A_{CP}`$ is also small. This is also shown in some model calculations .
The method of Isospin or SU(3) symmetry which requires a set of measurements to solve the uncertainties is sometimes difficult for experiments. We estimate these strong phases in specific models, such like the generalized factorization approach, which can be tested by experiments.
## 2 CP Violation Classification and Formulae
For charged $`B^\pm `$ decays the CP-violating asymmetries are defined as
$$A_{CP}=\frac{\mathrm{\Gamma }(B^+f^+)\mathrm{\Gamma }(B^{}f^{})}{\mathrm{\Gamma }(B^+f^+)+\mathrm{\Gamma }(B^{}f^{})}.$$
(2)
The charged modes are self-tagged decay channels for experiments. They are easy to be measured. For $`B^0`$ decays, more complication is from the $`B^0\overline{B}^0`$ mixing. The CP-asymmetries may be time-dependent, if the final states are the same for $`B^0`$ and $`\overline{B^0}`$
$`A_{CP}(t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(B^0(t)f)\mathrm{\Gamma }(\overline{B}^0(t)f)}{\mathrm{\Gamma }(B^0(t)f)+\mathrm{\Gamma }(\overline{B}^0(t)f)}}`$ (3)
$``$ $`a_ϵ^{}\mathrm{cos}(\mathrm{\Delta }mt)+a_{ϵ+ϵ^{}}\mathrm{sin}(\mathrm{\Delta }mt).`$ (5)
Here the direct CP violation parameter $`a_ϵ^{}`$ is defined as
$`a_ϵ^{}`$ $`=`$ $`A_{CP}^{dir}={\displaystyle \frac{\mathrm{\Gamma }(B^0f)\mathrm{\Gamma }(\overline{B}^0f)}{\mathrm{\Gamma }(B^0f)+\mathrm{\Gamma }(\overline{B}^0f)}},`$ (6)
which is the same defination as the charged B decays. And $`a_{ϵ+ϵ^{}}`$ is mixing-induced CP violation . In this note we will concentrate on the direct CP asymmetries.
### 2.1 $`bs`$ ($`\overline{b}\overline{s}`$), transitions
First we parametrize the decay amplitude like this way
$``$ $`=`$ $`T\xi _uP_t\xi _te^{i\delta _t}P_c\xi _ce^{i\delta _c}P_u\xi _ue^{i\delta _u},`$
$`\overline{}`$ $`=`$ $`T\xi _u^{}P_t\xi _t^{}e^{i\delta _t}P_c\xi _c^{}e^{i\delta _c}P_u\xi _u^{}e^{i\delta _u},`$ (7)
where $`\xi _i=V_{ib}V_{is}^{}`$. $`T`$ and $`P_i`$ are the tree and $`i`$ ($`i=u,c,t`$) quark penguin contributions, respectively. Working in SM, we can use the unitarity relation $`\xi _c=\xi _u\xi _t`$ to simplify the above equation,
$``$ $`=`$ $`T\xi _uP_{tc}\xi _te^{i\delta _{tc}}P_{uc}\xi _ue^{i\delta _{uc}},`$
$`\overline{}`$ $`=`$ $`T\xi _u^{}P_{tc}\xi _t^{}e^{i\delta _{tc}}P_{uc}\xi _u^{}e^{i\delta _{uc}},`$ (8)
where we define
$`P_{tc}e^{i\delta _{tc}}`$ $`=`$ $`P_te^{i\delta _t}P_ce^{i\delta _c},`$
$`P_{uc}e^{i\delta _{uc}}`$ $`=`$ $`P_ue^{i\delta _u}P_ce^{i\delta _c}.`$ (9)
Thus, the direct CP-violating asymmetry is
$$A_{CP}a_ϵ^{}=\left(|\overline{}|^2||^2\right)/\left(||^2+|\overline{}|^2\right)=A^{}/A^+,$$
(10)
where
$`A^{}`$ $`=`$ $`2TP_{tc}|\xi _u^{}\xi _t|\mathrm{sin}\varphi _3\mathrm{sin}\delta _{tc}+2P_{tc}P_{uc}|\xi _u^{}\xi _t|\mathrm{sin}\varphi _3\mathrm{sin}(\delta _{uc}\delta _{tc}),`$ (11)
$`A^+`$ $`=`$ $`(T^2+P_{uc}^2)|\xi _u|^2+P_{tc}^2|\xi _t|^22P_{tc}P_{uc}|\xi _u^{}\xi _t|\mathrm{cos}\varphi _3\mathrm{cos}(\delta _{uc}\delta _{tc})`$ (12)
$`2TP_{uc}|\xi _u|^2\mathrm{cos}\delta _{uc}+2TP_{tc}|\xi _u^{}\xi _t|\mathrm{cos}\varphi _3\mathrm{cos}\delta _{tc}.`$
First, we note that $`|\xi _u||\xi _t||\xi _c|`$, with an upper bound $`|\xi _u|/|\xi _t|0.025`$. In some channels, such as $`B^+K^+\pi ^0`$, $`K^+\pi ^0`$, $`K^+\rho ^0`$, $`B^0K^+\pi ^{}`$, $`K^+\pi ^{}`$, $`K^+\rho ^{}`$, $`|P_{tc}/T|`$ is of $`O(0.1)`$, $`|P_{uc}/P_{tc}|=O(0.3)`$. The CP-violating asymmetry in this case is
$$A_{CP}\frac{2z_{12}\mathrm{sin}\delta _{tc}\mathrm{sin}\varphi _3}{1+2z_{12}\mathrm{cos}\delta _{tc}\mathrm{cos}\varphi _3+z_{12}^2},$$
(13)
where $`z_{12}=|\xi _u/\xi _t|\times T/P_{tc}`$. We show the CP asymmetry of $`BK^+\pi ^{}`$ as an example in Figure 1. It is easy to see that, there may be large CP asymmetries in this decay channel. Besides the CKM parameter $`\rho `$ and $`\eta `$, the CP asymmetry is also sensitive to the gluon momentum $`k^2`$, which is related to the size of strong phase. If $`k^2`$ is known, the strong phase is predictable, we may use $`A_{CP}`$ to determine $`\mathrm{sin}\varphi _3`$. The first 6 channels of Table 1 are this kind of decays. Two of them are reported from CLEO Collaboration with large error-bars . The central values are far away from the theoretical predictions. If more data suport this, we may expect new physics signals here.
There are some decays with vanishing tree contributions ($`T=0`$), such as $`B^+\pi ^+K_S^0`$, $`\pi ^+K^0`$, $`\rho ^+K^0`$. Then for these decays, the CP-violating asymmetry is
$$A_{CP}2\frac{P_{uc}}{P_{tc}}\left|\frac{\xi _u}{\xi _t}\right|\mathrm{sin}(\delta _{uc}\delta _{tc})\mathrm{sin}\varphi _3.$$
(14)
Without the tree contribution, the suppression due to both $`P_{uc}/P_{tc}`$ and $`|\xi _u/\xi _t|`$ is much stronger. The CP-violating asymmetries are only around $``$($`1`$-$`2`$)%.
Some estimates of the channel $`B^+\pi K_S^0`$ show that even including the annihilation and soft final state interaction, the CP asymmetry of this decay is still small . This means that this channel is clean for new physics to show up. In table 1, we can see that CLEO’s central value of this decay indicates a large CP asymmetry maybe possible.
### 2.2 $`bd`$ ($`\overline{b}\overline{d}`$) transitions
Similarly to the $`bs`$ transition, we can define the CP asymmetry as
$$A_{CP}=A^{}/A^+,$$
(15)
where
$`A^{}`$ $`=`$ $`2TP_{tc}|\zeta _u^{}\zeta _t|\mathrm{sin}\varphi _2\mathrm{sin}\delta _{tc}2P_{tc}P_{uc}|\zeta _u^{}\zeta _t|\mathrm{sin}\varphi _2\mathrm{sin}(\delta _{uc}\delta _{tc}),`$ (16)
$`A^+`$ $`=`$ $`(T^2+P_{uc}^2)|\zeta _u|^2+P_{tc}^2|\zeta _t|^22P_{tc}P_{uc}|\zeta _u^{}\zeta _t|\mathrm{cos}\varphi _2\mathrm{cos}(\delta _{uc}\delta _{tc})`$ (17)
$`2TP_{uc}|\zeta _u|^2\mathrm{cos}\delta _{uc}+2TP_{tc}|\zeta _u^{}\zeta _t|\mathrm{cos}\varphi _2\mathrm{cos}\delta _{tc},`$
with $`\zeta _i=V_{ib}V_{id}^{}`$, and again we have used CKM unitarity relation $`\zeta _c=\zeta _t\zeta _u`$.
For the tree-dominated decays, such as $`B^+\pi ^+\eta ^{()}`$, $`\rho ^+\eta ^{()}`$, $`\rho ^+\omega `$, the relation $`P_{uc}<P_{tc}T`$ holds. The CP asymmetry is
$$A_{CP}\frac{2z_1\mathrm{sin}\delta _{tc}\mathrm{sin}\varphi _2}{1+2z_1\mathrm{cos}\delta _{tc}\mathrm{cos}\varphi _2},$$
(18)
with $`z_1=|\zeta _t/\zeta _u|\times TP_{tc}/T^2`$, and $`T^2T^22TP_{uc}\mathrm{cos}\delta _{uc}`$. The CP asymmetries are proportional to $`\mathrm{sin}\varphi _2`$. They are large enough for the experiments to detect them. The theoretical predictions of these decays are shown in table 1.
For the decays with a vanishing tree contribution ($`T=0`$), such as $`B^+K^+K_S^0`$, $`K^+\overline{K}^0`$, $`K^+\overline{K}^0`$, the CP-violating asymmetry is approximately proportional to $`\mathrm{sin}\varphi _2`$ again,
$$A_{CP}=\frac{2z_3\mathrm{sin}(\delta _{uc}\delta _{tc})\mathrm{sin}\varphi _2}{12z_3\mathrm{cos}(\delta _{uc}\delta _{tc})\mathrm{cos}\varphi _2+z_3^2},$$
(19)
with $`z_3=|\zeta _u/\zeta _t|\times P_{uc}/P_{tc}`$. As the suppressions from $`|\zeta _u/\zeta _t|`$ and $`|P_{uc}/P_{tc}|`$ are not very big, the CP-violating asymmetry can again be of order $`(10`$-$`20)\%`$. Unfortunately, these channels have smaller branching ratios .
More charmless decay channels are discussed in ref.. Some of them are more complicated than the ones we discussed above. There are also some other interesting decays like $`BK^{}\gamma `$, $`BD\pi \mathrm{}\nu `$, $`B\pi \pi \mathrm{}\nu `$ , etc. They have small CP asymmetries in SM. They are sensitive to new physics.
## 3 Models of Calculation
In the Factorization Approach, the two body $`B`$ meson decays can be factorized as two products:
$$C_iP_1P_2|O_i|B=C_iP_1|J_\mu |0P_2|J^\mu |B,$$
where $`C_i`$ is the corresponding Wilson coefficients. The second factor on the right side of the equation is proportional to the meson decay constant. The last term is the corresponding form factors.
The strong-phase differences arise through Bander-Silverman-Soni Mechanism (BSS) . In this picture, the perturbative penguin diagrams involving charm and up quark loops, where the light quarks can be on mass shell, providing the strong phases. They are mostly sensitive to the gluon momentum $`k^2`$. For numerical calculations, we use $`k^2=m_b^2/2\pm 2GeV^2`$.
In the perturbative QCD approach (pQCD) , we need one hard gluon connecting the spectator quark. Strong phases are from the non-factorizable diagram and annihilation diagram, where the innner quark or gluon propagator can be on mass shell. The pQCD approach is based on factorization, and goes one step further. In pQCD, we can calculate annihilation diagrams and also the non-factorizable contributions. The $`k^2`$ of gluon is well defined in this approach. We have calculated the $`B\pi \pi `$ decays in this approach, and the results compared with the factorization approach in Table 2.
## 4 Summary
The recently measured direct CP asymmetries for $`B^0K^+\pi ^{}`$, $`B^+K^+\eta ^{}`$, $`B^0K^+\pi ^{}`$, $`B^+\pi ^+K^0`$, and $`B^+\omega \pi ^+`$ are encouraging news for direct CP violation in B decays, although the signal is not significantly excess background.
CP-asymmetries of $`A_{CP}(K^\pm \eta ^{})`$, $`A_{CP}(\pi ^\pm K_S^0)`$ and ACP(ρ±
( )
[-.7ex] K0
)subscript𝐴𝐶𝑃superscript𝜌plus-or-minus
( )
[-.7ex] K0
A_{CP}(\rho^{\pm}\shortstack{{\tiny(\rule[0.86108pt]{5.0pt}{0.28453pt})}
\\
[-.7ex] $K^{*0}$}) are small, but stable against variation in $`N_c`$, $`k^2`$ and $`\mu `$. CP-asymmetries well over $`10\%`$ in these decay modes will be a sign of new physics. The decay channels of $`BK^\pm \pi ^{}`$, $`K^\pm \pi ^0`$, $`K^\pm \eta `$, $`K^\pm \eta ^{}`$, $`K^\pm \rho ^{}`$ and $`K^\pm \rho ^0`$, have measurably large CP-violating asymmetries. A good measurement of the CP-asymmetry in any one of these decays could be used to determine $`k^2`$. Such that the theoretical predictions of all other channels make sense. We also hope that the perturbative QCD approach could solve the remaining uncertainties in the factorization approach. With the two B factories and other hadronic machines, a number of decays is going to be measured soon.
## Acknowledgments
The author thanks A. Ali and G. Kramer for collaboration on the main topic discussed in this report. We thank the organizer H.Y. Cheng and W.S. Hou for a fruitful BCP3 conference. We acknowledge the Grant-in-Aid for Scientific Research on Priority Areas (Physics of CP violation) and JSPS for support.
## References |
warning/0001/hep-ph0001329.html | ar5iv | text | # Lambda polarization and single-spin left-right asymmetry in diffractive hadron-hadron collisions
We discuss Lambda polarization and single-spin left-right asymmetry in diffractive hadron-hadron scattering at high energies. We show that the physical picture proposed in a recent Letter is consistent with the experimental observation that $`\mathrm{\Lambda }`$ polarization in the diffractive process, $`pp\mathrm{\Lambda }K^+p`$, is much higher than that in the inclusive reaction, $`pp\mathrm{\Lambda }X`$. We make predictions for the left-right asymmetry, $`A_N`$, and for the spin transfer, $`D_{NN}^\mathrm{\Lambda }`$, in the single-spin process $`p()p\mathrm{\Lambda }K^+p`$ and suggest further experimental tests of the proposed picture.
It has been known for a long time that hyperons produced in hadron-hadron or hadron-nucleus collisions are polarized transversely to the production plane, although neither the projectile nor the target is polarized before the collision. Significant hyperon polarizations (up to 40%) in inclusive production processes have been observed in the fragmentation regions for moderately large transverse momenta. The effect has been confirmed in past years by a large number of similar experiments at different energies and/or using different projectiles and/or targets and for the production of different hyperons. However, the physical origin of these striking polarizations is still a puzzle. The basic difficulty is that helicity of the almost massless quarks is conserved in perturbative QCD (pQCD) in leading twist and at leading order , but the existence of large hyperon polarization requires significant helicity flip at the hadron level. Different attempts have been made to overcome this difficulty in the context of pQCD, other theoretical models have also been proposed \[References-References\]. It is clear that understanding of this striking spin effect would also shed light on the spin structure of hadron.
While most of the experimental data on hyperon polarization are in inclusive processes, the CERN R608 Collaboration has carried out an experiment in the diffractive dissociation process $`pp\mathrm{\Lambda }K^+p`$. One of the obvious advantages of this exclusive experiment is that, here, one concentrates on a much simpler final state than in inclusive processes. In this way, one hopes to gain deeper insight into the mechanisms of Lambda polarization. The results of this experiment show that $`\mathrm{\Lambda }`$ produced in this process is also transversely polarized, and that the polarization has the same sign (negative) as that observed in the inclusive process $`pp\mathrm{\Lambda }X`$ but the magnitude is very large $`(62\%\pm 4\%)`$ — much larger than those observed in the inclusive process. Since the diffractive process $`pp\mathrm{\Lambda }K^+p`$ is the simplest channel for the inclusive process $`pp\mathrm{\Lambda }X`$, the observation that $`\mathrm{\Lambda }`$ polarization $`P_\mathrm{\Lambda }`$ in this channel has a much larger value than that in $`pp\mathrm{\Lambda }X`$, which is the average over all the different channels, suggests that this process plays indeed a special role.
In a recent Letter , we argued that there is a close relation between the above mentioned hyperon polarization ($`P_H`$) observed in inclusive production processes in unpolarized hadron-hadron collisions and the left-right asymmetry ($`A_N`$) observed in inclusive hadron-hadron collisions using transversely polarized projectile or target. Theoretical arguments and experimental observations have been presented which strongly suggest that the two phenomena have the same origin and should be considered together. If this is indeed the case, it offers a new starting point to understand the origin of $`P_H`$.
In this note, we apply the picture to $`\mathrm{\Lambda }`$ production in diffractive hadron-hadron collisions to make further tests of the picture by comparing the obtained results with the available data. We show in particular that the much larger value of $`P_\mathrm{\Lambda }`$ in $`pp\mathrm{\Lambda }K^+p`$ should be considered as a further strong evidence for the picture. We make suggestions for future experiments.
We now start by recollecting the key points of the picture proposed in \[References\]. The basic idea of the picture is that there is a close relation between $`P_H`$ in unpolarized hadron-hadron collisions and $`A_N`$ in single-spin hadron-hadron collisions. Hence, if we extract the essential points encoded in the $`A_N`$ data, we should be able to understand $`P_H`$ based on these points. The following two points, (a) and (b), have therefore been derived from the existing $`A_N`$ data, and used as inputs to study $`P_H`$.
(a) If a hadron is produced by an upwards polarized valence quark of the transversely polarized projectile, it has a large probability to have a transverse momentum pointing to the “left-hand side” looking down stream. This production mechanism gives rise to positive left-right asymmetry. Here, “left-hand side” is defined by the requirement that the scalar product, $`=\stackrel{}{S}(\stackrel{}{p}_B\times \stackrel{}{p})`$, be positive, where $`\stackrel{}{S}`$ is the polarization vector of the transversely polarized beam, $`\stackrel{}{p}_B`$ and $`\stackrel{}{p}`$ are the momenta of the beam and the produced hadron, respectively. Positive $`A_N`$ measures the excess of hadrons produced to the left-hand side over those produced to the right-hand side. The above point has been derived directly from the data on $`A_N`$ in meson production. We denote the difference of the probabilities for $``$ to be positive and to be negative by $`C`$. $`C`$ should lie in the region $`0<C<1`$.
(b) If a hadron is produced by two valence quarks (valence diquark) of the projectile, the remaining valence quark produces an associated hadron. The left-right asymmetry due to this production mechanism is opposite to that of the associated hadron. This is consistent with the data on $`A_N`$ in $`\mathrm{\Lambda }`$ production. It explains in particular the surprising experimental result that $`\mathrm{\Lambda }`$ produced by a spin zero $`(ud)`$-diquark from the polarized projectile also exhibits left-right asymmetry!
These two points are supplemented by the following two points, which are consistent with the recent ALEPH and OPAL data on longitudinal polarization of $`\mathrm{\Lambda }`$ in $`e^+e^{}Z^0q\overline{q}\mathrm{\Lambda }+X`$, in order to give a description of $`P_H`$ in hadron-hadron collisions.
(1) Quark polarization is not destroyed in fragmentation.
(2) The SU(6) wave-function can be used to describe the relation between the spin of the fragmenting quark and that of the hadron produced in the fragmentation process.
We recall that the SU(6) wave-functions have been widely used in studying hyperon polarization in hadron-hadron collisions in the literature, (see, e.g., \[References\], \[References\], and \[References\]), i.e. point (2) has been assumed to be true. From this point, we obtain immediately that $`\mathrm{\Lambda }`$ polarization is completely determined by its $`s`$ valence quark. This result is quite different from that suggested by the recent polarized deep-inelastic lepton-nucleon scattering (DIS) data . The recent DIS data suggest that, at large $`Q^2`$, the quarks and antiquarks carry only a small fraction of the nucleon spin. When applied to $`\mathrm{\Lambda }`$, it suggests that $`\mathrm{\Lambda }`$ spin is not completely determined by its $`s`$ quark. This initiated the discussions (see, e.g. \[References-References\]) whether such a picture of the spin structure should also be used for describing the relation between the spins of the fragmenting quarks and the polarization of the hyperon produced in the fragmentation processes. It has been pointed out in \[References\] that measurements of the longitudinal polarization of $`\mathrm{\Lambda }`$ in $`e^+e^{}Z^0\mathrm{\Lambda }X`$ can provide some hints to answer this question. The available data are still far from accurate and enormous enough to provide a conclusive judgement. But these available data seem to favor the SU(6) description.
We note that the four points (a), (b), (1) and (2) form the basis of the picture in \[References\]. They are consistent with the data now available. Whether, if yes, how they can be derived from QCD are questions which have been discussed frequently in recent years in literature. Many models have been proposed to understand in particular point (a) and (b) in terms of quark-parton model in the framework of QCD. Since the purpose of \[References\] and that of the present paper are to discuss the close relationship between the left-right asymmetries observed in single-spin hadron-hadron collisions with hyperon polarization in unpolarized hadron-hadron collisions, we use these points as input. In this sense, the results obtained in the following and those in \[References\] are consequences of a phenomenological model which is consistent with the basic principles and the empirical facts from other experiments. Using this picture, we showed that various data on hyperon polarization in inclusive production processes can be understood provided that the $`s`$ and $`\overline{s}`$ taking part in the associated production of the $`\mathrm{\Lambda }`$ and $`K^+`$ have opposite transverse spins. We found in particular that, in this picture, $`\mathrm{\Lambda }`$ polarization comes mainly from the $`\mathrm{\Lambda }`$’s which contain two valence-quarks of the proton projectile and are associated with a spin zero meson such as $`K^+`$ which contains the remaining valence quark of the proton.
We stress that, although the model is consistent with that proposed by DeGrand and Miettinen on some points, it differs very much from the latter. In \[References\], polarization of hyperon originates from a semi-classical effect, Thomas precession, which leads to the “slow-partons-spin-down-fast-partons-spin-up” rule. This would in particular predict that hyperons produced in $`e^+e^{}`$ annihilations should also be transversely polarized, which contradicts the data. The model in \[References\] relates hyperon polarization $`P_H`$ in unpolarized collisions with left-right asymmetries $`A_N`$ observed in single-spin inclusive hadron production processes. Data on $`A_N`$ suggest a correlation between the polarization of the valence quark and the transverse moving direction of the produced hadron \[Points (a) and (b) mentioned above\]. The picture in \[References\] shows that the same correlations (a) and (b) lead also to hyperon polarization in unpolarized collisions.
We now come to the application of the picture to the diffractive process $`pp\mathrm{\Lambda }K^+p`$. We note that this is the simplest channel for the inclusive process $`pp\mathrm{\Lambda }X`$ and it has the following peculiarities: First, unlike in $`pp\mathrm{\Lambda }X`$, $`\mathrm{\Lambda }`$ has to contain two of the valence quarks from the colliding proton in this process. The associated spin zero $`K^+`$ contains the other valence quark. This means that, in this channel, we have only $`\mathrm{\Lambda }`$ produced by mechanism (b) which, according to the picture in \[References\], provides the largest polarization. Second, there is no contribution from hyperon decay. This means that there is no contamination from such decay processes to the $`\mathrm{\Lambda }`$ polarization. In fact, this is the only channel where these conditions are completely fulfilled. We therefore expect that $`\mathrm{\Lambda }`$ polarization should take its maximum in this process. This is consistent with the R608 observation that $`P_\mathrm{\Lambda }`$ in this process is much larger than that in $`pp\mathrm{\Lambda }X`$.
In order to estimate $`P_\mathrm{\Lambda }`$ in this process quantitatively, we recall that $`P_\mathrm{\Lambda }`$ is defined with respect to the production plane (see Fig. 1) and $`P_\mathrm{\Lambda }<0`$ means that the $`\mathrm{\Lambda }`$ has a large probability to be polarized in the $`\stackrel{}{n}`$ direction, where $`\stackrel{}{n}\stackrel{}{p}_B\times \stackrel{}{p}_\mathrm{\Lambda }`$ is the normal of the production plan. According to point (b), if $`\stackrel{}{p}_\mathrm{\Lambda }`$ is in the direction as shown in the figure (denoted by “left"), $`K^+`$ should be going right, thus $`\stackrel{}{p}_B\times \stackrel{}{p}_{K^+}`$ should be in the opposite direction as $`\stackrel{}{n}`$. According to point (a), the $`u`$-valence quark should have a large probability to be polarized in $`\stackrel{}{n}`$ (downwards) direction. The difference of the probability for this $`u_v`$ to be polarized in $`\stackrel{}{n}`$ and that to be polarized in $`\stackrel{}{n}`$ direction is given by the constant $`C`$ mentioned in point (a). Since $`K^+`$ is a spin zero object, the $`\overline{s}`$ in $`K^+`$ should be polarized in $`\stackrel{}{n}`$ direction if $`u_v`$ is polarized in $`\stackrel{}{n}`$ direction. Hence, the $`s`$ quark thus the $`\mathrm{\Lambda }`$ should be polarized in $`\stackrel{}{n}`$ direction if $`s`$ and $`\overline{s}`$ have opposite transverse spins. (See Fig.2). We see, in this case, the polarization of $`\mathrm{\Lambda }`$ is the same as that for the remaining $`u`$-valence quark which is contained in the associatively produced $`K^+`$. This implies that,
$$P_\mathrm{\Lambda }(pp\mathrm{\Lambda }K^+p)=C.$$
(1)
Using the value $`C=0.6`$ determined (see, e.g. \[References\] and the references given there) by fitting the $`A_N`$ data, we obtain that $`P_\mathrm{\Lambda }(pp\mathrm{\Lambda }K^+p)=0.6`$ which is in good agreement with the R608 data $`P_\mathrm{\Lambda }(pp\mathrm{\Lambda }K^+p)=0.62\pm 0.04`$.
This result is rather encouraging. We therefore made a detailed analysis for the related spin effects in this processes. We found that single-spin reaction $`p()p\mathrm{\Lambda }K^+p`$ with transversely polarized proton $`p()`$ is an ideal place to test the picture proposed in \[References\] and its applicability to diffractive processes. We obtained the following:
($`\alpha `$) There should be a large left-right asymmetry $`A_N`$ for $`\mathrm{\Lambda }`$ as well as for $`K^+`$ in $`p()p\mathrm{\Lambda }K^+p`$ in the fragmentation region of the transversely polarized proton $`p()`$, and,
$$A_N^\mathrm{\Lambda }[p()p\mathrm{\Lambda }K^+p]=A_N^{K^+}[p()p\mathrm{\Lambda }K^+p]=C.$$
(2)
This is because $`|\mathrm{\Lambda }^{}>=(ud)_{0,0}s^{}`$, thus only the configuration $`(ud)_{0,0}u^{}`$ in the projectile proton $`p()`$ contributes to the process $`p()p\mathrm{\Lambda }K^+p`$. Hence, the $`u`$ valence quark contained in $`K^+`$ is upwards polarized. According to the points (a) and (b) mentioned above, we obtain the results shown in Eq.(2).
($`\beta `$) The “spin transfer parameter” $`D_{NN}^\mathrm{\Lambda }`$ for the produced $`\mathrm{\Lambda }`$ should be positive and large in $`p()p\mathrm{\Lambda }K^+p`$ in the fragmentation region of $`p()`$. We recall that $`D_{NN}^\mathrm{\Lambda }`$ is defined as the probability for the produced $`\mathrm{\Lambda }`$ to be polarized in the same transverse direction as the projectile. Although the $`ud`$ diquark which comes from the projectile to form the produced $`\mathrm{\Lambda }`$ has to be in the spin-zero state thus carries no information of polarization of the projectile, the remaining $`u`$-valence quark completely determines the polarization of the projectile. They are polarized in the same direction. Hence, the $`\overline{s}`$ has to be polarized in the opposite direction as the projectile since $`K^+`$ has spin zero. The $`s`$ quark, which has opposite transverse spin as the $`\overline{s}`$, thus the $`\mathrm{\Lambda }`$ containing this $`s`$-quark should therefore be polarized in the same transverse direction of the projectile. Hence, we obtain,
$$D_{NN}^\mathrm{\Lambda }[p()p\mathrm{\Lambda }K^+p]=1.$$
(3)
Both ($`\alpha `$) and ($`\beta `$) are predictions which can be tested in future experiments. The predictions for the process $`pp\mathrm{\Lambda }K^+p`$ are summarized in Table I.
Here, it should be mentioned that $`\mathrm{\Lambda }`$ polarization has recently been measured in another exclusive process $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$ at incident momentum of 27.5 GeV/c. The results show the following: (i) The magnitude of $`P_\mathrm{\Lambda }`$ in this channel is much smaller than that in $`ppp\mathrm{\Lambda }K^+`$. (ii) $`P_\mathrm{\Lambda }`$ is approximately the same for events where $`K^+`$ and $`\mathrm{\Lambda }`$ are produced in the same hemisphere and for those where they are in the opposite hemispheres. We show that both (i) and (ii) are consistent with the picture mentioned above. First, unlike that in $`ppp\mathrm{\Lambda }K^+`$, $`\mathrm{\Lambda }`$ in this channel can be produced by two or one of the three valence quarks of the colliding proton. But, according to the picture , only the $`\mathrm{\Lambda }`$’s produced by two valence quarks are polarized, those produced by one are not. Second, while $`\mathrm{\Lambda }`$ in $`pp\mathrm{\Lambda }K^+p`$ is definitely associated with $`K^+`$, in $`ppp\mathrm{\Lambda }K^+\pi ^+\pi ^{}\pi ^+\pi ^{}`$, $`\mathrm{\Lambda }`$ can be associated with vector meson such as $`K^0`$ which subsequently decays into $`K^+\pi ^{}`$. As has been emphasized in \[References\], the correlation between the $`u`$-valence quark in the associated meson and the $`s`$-quark in the produced $`\mathrm{\Lambda }`$ will not be destroyed if more spin-zero mesons are associatively produced, (See Fig.3.), but it will be destroyed if (spin one) vector meson(s) is (are) involved. Hence, we expect that point (i) is true. Furthermore, since $`P_\mathrm{\Lambda }`$ is not changed if more spin zero mesons are produced in between (Fig.3), it is therefore unimportant whether $`K^+`$ and $`\mathrm{\Lambda }`$ are produced in the same or in the opposite hemispheres. This implies (ii) should also be true.
In summary, we have successfully applied the picture proposed in \[References\] to $`\mathrm{\Lambda }`$ production in diffractive processes. The obtained results are in agreement with the data now available. Predictions have been given which can be tested in future experiments.
We are indebted to Professor Meng Ta-chung who initiated this research and participated in the early stage of this work. Part of the results in this paper are taken from an unpublished note (Ref.\[References\]) written together with him. We thank R. Rittel, Wang Qun and Xie Qu-bing for helpful discussions. This work was supported in part by the National Natural Science Foundation of China (NSFC), the State Education Commission of China, and the Australian Research Council. |
warning/0001/cond-mat0001063.html | ar5iv | text | # Interchain conductivity of coupled Luttinger liquids and organic conductors
\[
## Abstract
We reconsider the theory of dc and ac interchain conductivity in quasi-one dimensional systems. Our results are in good agreement with the measured c-axis optical conductivity of $`(\text{TMTSF})_2\text{ClO}_4`$ and suggest that the c-axis dc-conductivity of $`(\text{TMTSF})_2\text{PF}_6`$ in the $`150K<T<300K`$ range is dominated by precursor effects of Mott localization. The crossover from a Luttinger liquid at high energy to a Fermi liquid at low energy is also addressed, within a dynamical mean-field theory. Implications for the inter-chain resistivity and Drude weight in the Fermi liquid regime are discussed.
\]
Inter-chain (or inter-plane) coherence is a key issue for the physics of anisotropic strongly correlated systems. It is of special importance for the quasi-one dimensional organic conductors TMTSF<sub>2</sub>X and TMTTF<sub>2</sub>X which are systems of weakly coupled chains . Since the one dimensional (1D) interacting electron gas has non Fermi liquid (Luttinger liquid) properties , one expects in these compounds a crossover from Luttinger liquid behavior to a coherent behavior as the temperature or frequency is lowered. The latter can be either a metallic three-dimensional conductor (presumably Fermi-liquid like), or a phase with long-range order.
Although there is general agreement on this qualitative scenario, there is considerable debate on how precisely this crossover takes place . As was pointed out early on , the interactions renormalize the effective interchain hopping of weakly-coupled Luttinger chains. This leads to a crossover scale $`E^{}t(t_{}/t)^{1/(1\alpha )}=t_{}(t_{}/t)^{\alpha /(1\alpha )}`$ for the onset of coherence along one of the transverse directions with interchain hopping $`t_{}`$. In this expression, $`t`$ is the intra-chain hopping and $`\alpha \frac{1}{4}(K_\rho +1/K_\rho )\frac{1}{2}`$ is the Fermi surface exponent. The scale $`E^{}`$ can be much smaller than the bare interchain hopping $`t_{}`$. With this idea in mind, analysis of NMR data led to believe that the 1D Luttinger liquid (LL) to Fermi liquid (FL) crossover took place in the $`(\text{TMTSF})_2\text{X}`$ family at a low temperature, of the order of $`1015\text{K}`$. This scale seemed consistent with the above estimate of $`E^{}`$ since the interchain hopping along the $`b`$-axis is $`t_b300K`$, an order of magnitude smaller than the hopping $`t_a`$ along the chains. Recent analysis of the transport and optical properties have led to a reexamination of this commonly accepted view. Measurement of the optical conductivity along the chains agrees with a LL picture at high enough frequency and provides a determination of the LL exponent $`K_\rho `$ (of order $`K_\rho 0.23`$). This value is consistent with the temperature dependence of $`\rho _a(T)`$ in the $`100\text{K}300\text{K}`$ range (suitably corrected for thermal expansion ). However, these studies and most notably the recent measurements of the interchain dc conductivity revealed that the 1D regime only extends down to $`100\text{K}`$, a much higher temperature scale than the early estimates based on the NMR. Furthermore, it is clear from the optical data that another important energy scale is present in these compounds, associated with a Mott gap signaled by the onset of optical absorption at $`100\text{cm}^1150\text{K}`$. Transverse hopping is actually responsible for the metallic character of the $`(\text{TMTSF})_2\text{X}`$ compounds, which otherwise would be Mott insulators . These findings raise a number of important and largely unanswered questions such as the precise nature of the dimensional crossover, the origin of the anomalies observed for $`10K<T<100K`$ (calling in particular for a reinterpretation of the NMR in this range), and the actual range of the LL behavior.
In this letter, we focus on inter-chain transport and optical response, which are ideal probes of these issues. We first reconsider the theoretical analysis of the $`\omega `$ and $`T`$\- dependence of the inter-chain conductivity for the simplest case of weakly coupled metallic Luttinger liquids. We find a power law of the temperature or the frequency $`\sigma _{}(T,\omega )(T,\omega )^{2\alpha 1}`$. The exponent is at variance with previous theoretical predictions , but is in reasonable agreement with recent measurements of c-axis optical conductivity of the $`(\text{TMTSF})_2\text{ClO}_4`$ compound . We then show that understanding the observed T-dependence of c-axis transport in $`(\text{TMTSF})_2\text{PF}_6`$ requires to take into account the physics associated with the Mott gap apparent in optical measurements. Finally, we consider a dynamical mean-field theory (DMFT) for the description of the dimensional crossover, and draw some consequences for the inter-chain conductivity in the low-temperature FL regime.
We consider the hamiltonian of coupled chains:
$$H=\underset{m}{}H_{1D}^{(m)}\underset{m,m^{}}{}t_{}\underset{i\sigma }{}(c_{im\sigma }^+c_{im^{}\sigma }+\text{h.c})$$
(1)
in which $`H_{1D}^{(m)}`$ stands for each individual chain $`m`$, and the inter-chain hopping $`t_{}`$ will be first considered, for simplicity, to be identical for all transverse directions. Using the DMFT approach detailed below, a general formula (Eq. 9) for the transverse conductivity can be established. When the inter-chain hopping $`t_{}`$ can be treated in perturbation theory to lowest order, this formula reduces to that established in Ref. within the tunneling approximation :
$`\text{Re}\sigma _{}(\omega ,T)^{pert}`$ $`t_{}^2{\displaystyle }dx{\displaystyle }d\omega ^{}A_{1D}(x,\omega +\omega ^{})\times `$ (3)
$`\times A_{1D}(x,\omega ^{}){\displaystyle \frac{f(\omega ^{})f(\omega ^{}+\omega )}{\omega }}`$
where $`A_{1D}(x,\omega )`$ is the one-electron spectral function of a single chain, and $`f`$ the Fermi function. Let us consider first the simplest case in which each decoupled chain is a Luttinger liquid characterized by a given value of $`K_\rho `$, and all intra-chain umklapp processes can be neglected. For the values of $`K_\rho `$ of interest here, the transverse hopping is then a relevant perturbation in the renormalization group sense, so that it is legitimate to treat it perturbatively only at sufficiently high energy or temperature. Eq. (3) thus applies only when $`kTE^{}`$ or $`\mathrm{}\omega E^{}`$, where $`E^{}=t(t_{}/t)^{1/(1\alpha )}`$ is the single-particle crossover scale mentioned above . In this high-temperature/high-energy regime, the system can be considered to be quasi one-dimensional. Using known properties of the spectral functions in a Luttinger liquid, we then obtain from (3):
$$\text{Re}\sigma _{}(\omega ,T)^{pert}\left(\frac{t_{}}{t}\right)^2\left(\frac{T}{t}\right)^{2\alpha 1}F_\alpha \left(\frac{\omega }{T}\right)$$
(4)
In this expression, $`F_\alpha `$ is a scaling function such that $`F_\alpha (0)=\text{const}`$ and $`F_\alpha (x1)x^{2\alpha 1}`$. Thus, we find $`\text{Re}\sigma _{}(T\omega )T^{2\alpha 1}`$ for $`kTE^{}`$, while at high-frequency ($`\mathrm{}\omega E^{}`$), we obtain $`\sigma _{}(\omega T)\omega ^{2\alpha 1}`$.
This behavior of $`\sigma _{}(T)`$ differs from the conclusions of previous authors. In , the authors used a somewhat different approach than the direct Kubo formula and concluded that $`\sigma _{}(T)T^{2\alpha }`$. In , it was noted that this temperature dependence also results from (3) provided that the spatial integral is cutoff at a scale of the order of the lattice spacing. According to these authors, this cutoff is justified by the incoherent nature of the propagation of a physical electron within each chain (in which the true quasiparticles are spinons and holons). In our opinion however, this information is already encoded in the spectral function of the physical electron $`A_{1D}(x,\omega )`$ and no extra cutoff has to be introduced in (3). This conclusion would of course be changed in the presence of specific forms of intra-chain disorder such as a random forward scattering, which would spoil momentum conservation and bring in a natural cutoff of the order of the mean-free path. Note however that more realistic forms of disorder would be signaled by localisation effects.
We now turn to a comparison with experiments on $`(\text{TMTSF})_2\text{X}`$. It is increasingly recognized that the properties of these compounds result from the commensurate quarter-filling of the band . Transport and optical measurements lead to consistent values of the LL parameter $`K_\rho 0.23`$. Note that $`K_\rho <0.25`$ is the condition upon which the quarter-filled umklapp process becomes relevant. This suggests that if the transverse hopping was entirely suppressed in these compounds, they would actually be one-dimensional (quarter filled) Mott insulators. Fortunately, optical experiments extend to high frequencies compared to the Mott gap. Thus the complications associated with the proximity to the Mott insulating state can be treated perturbatively, making the comparison to Eq.(4) relevant. Recently, optical studies of $`(\text{TMTSF})_2\text{ClO}_4`$ have been performed , and the high-frequency optical conductivity $`\sigma _c(\omega )`$ along the least conducting c-axis was found to be a broad, very slowly increasing function of frequency in the range $`50\text{cm}^1<\omega <10^5\text{cm}^1`$. This slowly increasing continuum is essentially independent of temperature for $`10\text{K}<T<300\text{K}`$, and is in good agreement with our result $`\sigma _c(\omega )\omega ^{2\alpha 1}`$. The exponent $`2\alpha 1=(K_\rho +1/K_\rho )/22`$ depends quite sensitively on the value of $`K_\rho `$ ($`2\alpha 1=0.12`$ for $`K_\rho =0.25`$ and $`2\alpha 1=0.6`$ for $`K_\rho =0.2`$). The latter value $`2\alpha 10.6`$ provides a reasonable fit to the experimental data of . Note that values of $`K_\rho `$ in the $`0.2`$-$`0.25`$ range are quite consistent with the other experimental findings discussed above. In contrast, the $`\omega ^{2\alpha }`$ law advocated in would require $`K_\rho 0.35`$ to explain these data, a value quite inconsistent with that obtained from the longitudinal optical conductivity. Finally, an important issue in view of the success of this simple theory is whether it can also account for the behavior of the optical conductivity along the b-axis. Since $`t_b/t_a0.1`$, an incoherent regime should also apply there at those high frequencies. Preliminary analysis of the $`\sigma _b`$ data of seems in agreement with such a behavior, but more detailed comparison with the data is needed.
We now discuss the measurements of the c-axis transport in $`(\text{TMTSF})_2\text{PF}_6`$ . Because of the large thermal expansion, these were performed at several pressures so that a volume correction could be made. This yields data for the c-axis resistivity $`\rho _c(T)`$ which can be interpreted as corresponding to a compound with constant lattice parameters. This volume-corrected $`\rho _c(T)`$ displays a marked increase by more than a factor of $`3`$ upon cooling from $`T=350\text{K}`$ down to $`T=100\text{K}`$ (followed by a rapid downturn at lower temperature which presumably marks the onset of inter-chain coherence). Such a large variation of $`\rho _c(T)`$ cannot be explained by the dependence $`\rho _c1/T^{2\alpha 1}`$ found above with acceptable values of $`K_\rho `$ (which correspond to rather small values of $`2\alpha 1`$). In , it was fitted to a $`1/T^{1.4}`$ power law, which happens to be in reasonable agreement with the $`1/T^{2\alpha }`$ dependence advocated by these authors. However, the temperature range in which these data are obtained is such that $`kT`$ is never significantly larger than the “would-be” 1D Mott gap which can be extracted from the optical measurements ($`100200\text{cm}^1`$) leading to transport gaps (half the optical gap) in the range $`\mathrm{\Delta }75150\text{K}`$. As a result, the proximity to the incipient Mott insulator cannot be neglected, and we propose that this effect is the one responsible for the increase of $`\rho _c(T)`$. Indeed, we have attempted a phenomenological fit by an activated behavior $`\rho _c(T)1/T^{2\alpha 1}\mathrm{exp}\mathrm{\Delta }/T`$. and found that it provides a fit of the data of Ref. which is as satisfactory as the power-law fit used in . The value of $`\mathrm{\Delta }`$ used in this fit is of the same order of magnitude than the Mott gap seen in optical measurements. The discrepancy between the theoretical result, Eq.(4) and the experimentally observed behavior of $`\rho _c(T)`$ is, in our opinion, a strong indication that an explanation based on incoherent tunneling between LL chains in a purely metallic regime is insufficient and that the physics of Mott localisation plays an important role in the temperature range $`150\text{K}<T<300\text{K}`$.
A more refined theory of the interplay between interchain hopping and Mott localisation is clearly needed, including both the transverse hopping and the quarter- filling umklapp scattering . Due to the umklapp, the effective $`K_\rho `$ will be renormalized towards smaller values as temperature is decreased, hence leading to larger values of $`\rho _c`$, as expected. The main issue is whether the intra-chain resistivity is also strongly affected by these renormalizations. Indeed, no sign of an upturn in $`\rho _a(T)`$ is observed experimentally for the $`\text{TMTSF}_2X`$ compounds. Since $`\rho _a(T)g_{1/4}^2T^{16K_\rho 3}`$ , it is conceivable that the decrease in $`K_\rho `$ is almost compensated by an associated decrease of the effective $`g_{1/4}`$ due to $`t_{}`$. This however, remains to be verified and is probably the most challenging issue for the picture that we propose. Transport in the $`\text{TMTTF}_2X`$ family is interesting in this respect. These compounds are insulators with a rather large Mott gap, and indeed both $`\rho _a(T)`$ and $`\rho _c(T)`$ display an activated behavior at low temperature. Above some temperature scale however, $`\rho _a(T)`$ becomes metallic, while $`\rho _c(T)`$ still displays insulating behaviour (albeit slower than activated). Studies of the transport properties of these compounds over a more extended range of pressure and temperature (and in particular a detailed comparison of the energy scales appearing in $`\sigma _a(\omega )`$, $`\rho _a(T)`$ and $`\rho _c(T)`$) would certainly prove very interesting.
Let us finally address the issue of the crossover between a Mott-Luttinger liquid at high energy/temperature and a coherent (presumably Fermi-liquid) state at low energy. The main theoretical difficulty is that this crossover is associated with the breakdown of perturbative expansions in $`t_{}`$. On the other hand, perturbative approaches in the interaction strength could be used in the FL regime, but fail in the LL regime. The necessity of treating $`t_{}`$ in a non-perturbative manner was recently noted by Arrigoni , who used a limit of infinite transverse dimensionality as a guidance for resumming the perturbation series in $`t_{}`$. Here, we show that this limit provides a generalized dynamical mean-field theory (DMFT) framework for the description of the crossover, and we draw some consequences for inter-chain conductivity. The limit of infinite dimensions has been used extensively in recent years to formulate a controlled DMFT of strongly-correlated systems . We consider each one-dimensional chain to be coupled to $`z_{}\mathrm{}`$ nearest neighbor ones, with the transverse hopping scaled as $`t_{}=\stackrel{~}{t}_{}/\sqrt{z_{}}`$, so that the non-interacting density of states in the transverse direction $`D(ϵ_{})_k_{}\delta [ϵ_{}ϵ(k_{})]`$ remains well-defined. It is easily shown that in this limit the self-energy becomes independent of transverse momentum $`k_{}`$: $`\mathrm{\Sigma }=\mathrm{\Sigma }(k,i\omega _n)`$ ($`k`$ is the momentum along the chains). Furthermore, the calculation of $`\mathrm{\Sigma }`$ reduces to the solution of an effective one-dimensional problem specified by the effective action:
$`S_{\text{eff}}=`$ $`{\displaystyle _0^\beta 𝑑\tau 𝑑\tau ^{}\underset{ij,\sigma }{}c_{i\sigma }^+(\tau )𝒢_0^1(ij,\tau \tau ^{})c_{j\sigma }(\tau ^{})}`$ (6)
$`+{\displaystyle _0^\beta }𝑑\tau H_{1D}^{int}[\{c_{i\sigma },c_{i\sigma }^+\}]`$
where $`H_{1D}^{int}`$ is the interacting part of the on-chain hamiltonian in (1), while $`𝒢_0`$ is an effective bare propagator which is determined from a self-consistency condition. The latter imposes that the Green’s function $`G(ij,\tau \tau ^{})c(i,\tau )c^+(j,\tau ^{})_{\text{eff}}`$ calculated from $`S_{\text{eff}}`$ coincides with the on-chain Green’s function of the original problem, with a self-energy $`\mathrm{\Sigma }=𝒢_0^1G^1`$. This condition reads:
$$G(k,i\omega _n)=𝑑ϵ_{}\frac{D(ϵ_{})}{i\omega _n+\mu ϵ_k\mathrm{\Sigma }(i\omega _n,k)ϵ_{}}$$
(7)
The DMFT equations (6,7) fully determine the self-energy and Green’s function of the coupled chains. Furthermore, inter-chain transport properties simplify in this approach. Because the current vertex is an odd function of the transverse momentum, vertex corrections drop out of the inter-chain conductivity, which reads:
$`\text{Re}\sigma _{}(\omega ,T)`$ $`t_{}^2{\displaystyle 𝑑ϵ_{}D(ϵ_{})\frac{dk}{2\pi }𝑑\omega ^{}A(ϵ_{},k,\omega ^{})}`$ (9)
$`\times A(ϵ_{},k,\omega +\omega ^{}){\displaystyle \frac{f(\omega ^{})f(\omega ^{}+\omega )}{\omega }}`$
Here, $`A(ϵ_{},k,\omega )=\frac{1}{\pi }\text{Im}G(ϵ_{},k,\omega )`$ is the single-particle spectral function of the coupled chains system. Eq.(9) has a much wider range of applicability than its small-$`t_{}`$ limit, Eq.(3), since it holds all the way from the LL to the FL regime.
Solving quantitatively these DMFT equations is still a challenging problem however. In the single-site DMFT (corresponding to the usual $`d=\mathrm{}`$ limit ) the mapping is onto a single-impurity Anderson model with a self-consistent bath, and several techniques have been developed to handle this problem. Here however, the mapping is onto a self-consistent 1D-chain, and the problem is an order of magnitude more difficult. Nevertheless, several conclusions can be drawn from the above equations even in the absence of a full solution. We shall consider again the simplest case of metallic 1D LL chains, neglecting umklapp scattering and Mott physics. In the 1D regime, the self-energy behaves as $`\mathrm{\Sigma }t((\omega ,k)/t)^{1/(1\alpha )}`$. From (7), it is clear that the inter-chain hopping becomes relevant when $`t_{}>\mathrm{\Sigma }`$, and we recover the crossover scale $`E^{}t(t_{}/t)^{1/(1\alpha )}`$. At energies smaller than $`E^{}`$, the DMFT approach leads to Fermi-liquid behavior. This is clear from viewing the effective 1D model as a set of Anderson impurity models coupled together by the non-local part of the “bare” propagator $`𝒢_0`$. Each of the decoupled Anderson model is a Fermi-liquid at low-energy and the non-local hybridization is a smooth perturbation on the FL ground state. In other words, the physics below $`E^{}`$ is smoothly connected from small $`t_{}/t`$ to large $`t_{}/t`$. When $`t_{}t`$ and for energy scales much smaller than the bandwidth (of order $`t`$), all single-particle quantities should obey a scaling behavior in the variables $`\stackrel{~}{\omega }=\omega /E^{}`$, $`\stackrel{~}{x}=xE^{}/t`$, $`\stackrel{~}{T}=T/E^{}`$. In particular $`t\mathrm{\Sigma }(x,\omega ,T)=E^{}t_{}\stackrel{~}{\mathrm{\Sigma }}(\stackrel{~}{x},\stackrel{~}{\omega },\stackrel{~}{T})`$, $`tG(x,\omega ,T)=(E^{}/t_{})\stackrel{~}{G}(\stackrel{~}{x},\stackrel{~}{\omega },\stackrel{~}{T})`$ with $`\stackrel{~}{\mathrm{\Sigma }}`$ and $`\stackrel{~}{G}`$ universal scaling functions associated with the crossover. A low-frequency expansion of the scaling form of $`\mathrm{\Sigma }`$ in the FL regime yields a finite quasiparticule residue $`Z(t_{}/t)^{\frac{\alpha }{1\alpha }}E^{}/t_{}`$ (shown in to have only weak k-dependence along the Fermi surface). $`Z`$ can be much smaller than unity, while the effective mass (specific heat) ratio $`m^{}/m`$ does not become large because $`\mathrm{\Sigma }/\omega `$ and $`\mathrm{\Sigma }/k`$ scale in the same manner. Using the FL form of $`\mathrm{\Sigma }`$ in (9), we obtain the $`T=0`$ inter-chain conductivity as: $`\sigma _{}(\omega ,T=0)=\frac{\omega _D^2}{8}\delta (\omega )+\sigma _{}^{reg}(\omega )`$ with $`\omega _{D}^{}{}_{}{}^{2}/8Z^2t_{}^2/t=t(t_{}/t)^{2/(1\alpha )}`$. The total spectral weight $`𝑑\omega \sigma _{}(\omega )=\omega _P^2/8`$ is, from the f-sum rule, proportional to the inter-chain kinetic energy, and hence to $`t_{}^2/t`$, so that the relative weight carried by the Drude peak is small, of order $`\omega _D^2/\omega _P^2(t_{}/t)^{2\alpha /(1\alpha )}=Z^2`$ This is in reasonable agreement with the optical data on the b-axis reflectivity of $`(\text{TMTSF})_2\text{PF}_6`$ , for which a Drude peak is seen at low temperature, with $`\omega _D^2/\omega _P^20.03`$ (while $`t_b/t_a0.1`$). Scaling considerations also lead to an inter-chain resistivity $`\rho _{}(T)/\rho _0=(t/E^{})R(T/E^{})`$, with $`R(x1)x^{12\alpha }`$ , $`R(x1)x^2`$ and $`\rho _0=hV_m/(e^2a_{}^2)`$ . This implies that $`\rho _{}/\rho _0=A(T/t)^2`$ in the FL regime, with an enhanced value of the coefficient $`A(t/t_{})^{\frac{3}{1\alpha }}`$. Note that $`\rho _{}`$ is typically much bigger than the Mott limit $`\rho _0`$, as observed experimentally.
In conclusion, we have reconsidered in this paper the theory of interchain conductivity in quasi-1D systems. When applied to organic compounds, good agreement is found with the frequency dependence of $`\sigma _c(\omega )`$. while c-axis dc transport appears to be dominated by the proximity of the Mott gap $`\mathrm{\Delta }`$. This calls for additional theoretical work in the regime where $`E^{}t(t_{}/t)^{1/(1\alpha )}<\mathrm{\Delta }<t_{}`$. It would also be valuable to see whether the regime $`\mathrm{\Delta }<E^{}`$, which is simpler to analyze theoretically, can be realized experimentally in some family of compounds (e.g by applying uniaxial stress to TMTTF<sub>2</sub>X).
We thank D. Jerome, J. Moser, L. Degiorgi, G. Grüner, and V. Vescoli for sharing and discussing their data with us. Discussions with M. Gabay and A. M Tremblay are also acknowledged. A.G. and T.G. would like to thank the ITP, Santa Barbara for support (under NSF grant No. PHY94-07194) and hospitality during part of this work. |
warning/0001/math0001142.html | ar5iv | text | # Vanishing theorems on toric varieties
## Introduction
Our main goal in this article is to give a characteristic free approach to vanishing results on arbitrary toric varieties. We prove that the vanishing of a certain cohomology group depending on a Weil divisor is implied by the vanishing of the analogous cohomology group involving a higher multiple of that divisor. When the variety is complete and the divisor is $``$-Cartier, one recovers in this setting a theorem due to Kawamata and Viehweg. We apply these results to prove a strong form of Fujita’s conjecture on a smooth complete toric variety.
Let $`X`$ be a toric variety and $`D_1,\mathrm{},D_d`$ the invariant Weil divisors on $`X`$, so that $`\omega _X𝒪_X(D_1\mathrm{}D_d)`$. In the first part of the paper we deduce the following generalization of the theorem of Kawamata and Viehweg, for toric varieties.
###### Theorem 0.1.
Let $`D`$ be an invariant Weil divisor on $`X`$ as above. Suppose that we have $`E=_{j=1}^da_jD_j`$, with $`a_j`$ and $`0a_j1`$ such that for some integer $`m1`$, $`m(D+E)`$ is integral and Cartier. If for some $`i0`$ we have $`H^i(𝒪_X(D+m(D+E)))=0`$, then $`H^i(𝒪_X(D))=0`$. In particular, if $`X`$ is complete and there is $`E`$ aforementioned such that $`D+E`$ is $``$-ample, then $`H^i(𝒪_X(D))=0`$ for every $`i1`$.
As a particular case of this theorem, we see that if for some $`m1`$ and $`L\mathrm{Pic}(X)`$ we have $`H^i(L^m(D_{j_1}\mathrm{}D_{j_r}))=0`$, then $`H^i(L(D_{j_1}\mathrm{}D_{j_r}))=0`$. The cases $`r=0`$ and $`r=d`$ of this assertion were known to hold by reduction to a field of positive characteristic. Over such a field $`X`$ is Frobenius split and one concludes using arguments in Mehta and Ramanathan in \[MR\]. The fact that $`X`$ is Frobenius split will follow also from our results.
Our method yields other vanishing results as well. For example, we prove that if $`X`$ is a smooth toric variety and $`L\mathrm{Pic}(X)`$ is such that for some $`m1`$ and $`i0`$ we have $`H^i(\mathrm{\Omega }_X^jL^m)=0`$, then $`H^i(\mathrm{\Omega }_X^jL)=0`$. By taking $`X`$ complete and $`L`$ ample, we thus recover a theorem of Bott, Steenbrink and Danilov.
In the second part of the paper we give some applications of vanishing theorems on toric varieties. Our main result is a proof in the toric case of a strong version of the following conjecture due to Fujita (see \[La\] for discussion and related results).
###### Conjecture 0.2.
Let $`X`$ be a smooth projective variety of dimension $`n`$ and $`L\mathrm{Pic}(X)`$ an ample line bundle. Then $`\omega _XL^{n+1}`$ is globally generated and $`\omega _XL^{n+2}`$ is very ample.
When the ample line bundle $`L`$ is generated by global sections, an argument of Ein and Lazarsfeld \[EL\] based on vanishing results proves the conjecture over a field of characteristic zero. Under the same hypothesis on $`L`$, the first assertion of the conjecture is proved in arbitrary characteristic by Smith \[Sm\]. On a smooth projective toric variety every ample line bundle is very ample (see\[De\]), so these results prove the conjecture in this setting.
We give a direct proof of a strengthened form of the conjecture in the case of a toric variety, where instead of making a conclusion about a power $`L^m`$, we make a statement about any line bundle $`L`$ satisfying conditions on the intersection numbers with the invariant curves. We are able to replace $`\omega _X`$ by the negative of the sum of any set of $`D_i`$, and also improve the bound by one in the case when $`X`$ is not the projective space. More precisely, we prove:
###### Theorem 0.3.
Let $`X`$ be an $`n`$-dimensional complete smooth toric variety, $`L\mathrm{Pic}(X)`$ a line bundle and $`D_1,\mathrm{},D_m`$ distinct prime invariant divisors.
1. If $`(LC)n`$ for every invariant integral curve $`CX`$, then $`L(D_1\mathrm{},D_m)`$ is globally generated, unless $`X\text{P}^n`$, $`L𝒪(n)`$ and $`m=n+1`$.
2. If $`(LC)n+1`$ for every invariant integral curve $`CX`$, then $`L(D_1\mathrm{}D_m)`$ is very ample, unless $`X\text{P}^n`$, $`L𝒪(n+1)`$ and $`m=n+1`$.
To obtain these results we use Cox’s notion of homogeneous coordinate ring of $`X`$. When the fan defining $`X`$ is nondegenerate (i.e., it is not contained in a hyperplane), this is a polynomial ring $`S=k[Y_1,\mathrm{},Y_d]`$ together with a reduced monomial ideal $`B`$ and with a grading in the class group of $`X`$ which is compatible with the $`^d`$-grading by monomials. In general we need to slightly adjust this definition, but we leave this generalization for the core of the paper. As in the case of projective space, each graded $`S`$-module $`P`$ gives a quasi coherent sheaf $`\stackrel{~}{P}`$ on $`X`$ and for each $`i1`$, the Zariski cohomology $`H^i(X,\stackrel{~}{P})`$ is the degree zero part of the local cohomology module $`H_B^{i+1}(P)`$. This idea has been used in \[EMS\] to give an algorithm for the computation of cohomology of coherent sheaves on a toric variety.
Our basic result says that if $`P`$ is in fact $`^d`$-graded and if the multiplication by $`Y_j`$ is an isomorphism in certain $`^d`$-degrees, then the same is true for $`H_B^i(P)`$. The main example is $`P=S`$ in which we get that the multiplication
$$\nu _{Y_j}:H_B^i(S)_\alpha H_B^i(S)_{\alpha +e_j}$$
is an isomorphism for every $`\alpha =(\alpha _j)^d`$ such that $`\alpha _j1`$. In particular, $`H_B^i(S)_\alpha `$ depends only on the signs of the components of $`\alpha `$. This case was used in \[EMS\] in order to describe the support of $`H_B^i(S)`$. Similar results for the Ext modules appear also in \[Mu\] and \[Ya\]. Our second example is that of the modules giving the sheaves $`\mathrm{\Omega }_X^i`$ on a smooth toric variety. Using this result and the relation between the local cohomology of a module and the Zariski cohomology of the corresponding sheaf, we deduce the various vanishing theorems.
In the first section of the paper we summarize the construction in \[Cox\] for the homogeneous coordinate ring suitably generalized to be applicable also to toric varieties defined by a degenerate fan. All the results can be easily extended to this context. We prove that every quasicoherent sheaf on a toric variety comes from a graded module over the homogeneous coordinate ring, generalizing the result in \[Cox\] for the simplicial case. We describe the relation between the local cohomology of modules and the cohomology of the associated sheaves. This is used in the second section to prove the vanishing results described above.
In order to apply these results, we need numerical characterizations for ampleness and numerical effectiveness for the toric case and in the third section we provide these results. In the simplicial case, a toric Nakai criterion is given in \[Oda\]. We show that the result holds for an arbitrary complete toric variety. We also prove that $`L\mathrm{Pic}(X)`$ is globally generated if and only if $`(LC)0`$, for every integral invariant curve $`CX`$. In particular, we see that $`L`$ is globally generated if and only if it is numerically effective. These results have been recently obtained also by Mavlyutov in \[Ma\]. We mention a generalization in a different direction due to Di Rocco \[DR\] who proved that on a smooth toric variety, $`L\mathrm{Pic}(X)`$ is $`k`$-ample if and only if $`(LC)k`$ for every invariant curve $`CX`$.
As a consequence of the above results, we deduce that $`L`$ is big and nef if and only if there is a map $`\varphi :XX^{}`$ induced by a fan refinement (therefore $`\varphi `$ is proper and birational) and $`L^{}\mathrm{Pic}(X^{})`$ ample such that $`L\varphi ^{}(L^{})`$. This easily implies the version of Kawamata-Viehweg vanishing theorem for nef and big line bundles.
The fourth section is devoted to the above generalization (in this context) of Fujita’s Conjecture and some related results. The proof goes by induction on the dimension of $`X`$, by taking the restriction to the invariant prime divisors. The result which allows the induction says that for every $`l1`$, if $`L`$ is a line bundle such that $`(LC)l`$ for every invariant integral curve $`CX`$, then for every invariant prime divisor $`D`$ and every $`CX`$ aforementioned, $`(L(D)C)l1`$. From the case $`l=1`$ we see that if $`L`$ is ample, then $`L(D)`$ is globally generated. We conclude this section by proving a related result, which characterizes the situation in which $`L`$ is ample and $`D`$ is a prime invariant divisor, but $`L(D)`$ is not ample.
A well-known ampleness criterion (see, for example, \[Fu\]) can be interpreted as saying that on a complete toric variety $`X`$, $`L\mathrm{Pic}(X)`$ is ample if and only if it is globally generated and the map induced by restrictions
$$H^0(L)\underset{xX_0}{}H^0(L|_{\{x\}})$$
is an epimorphism, where $`X_0`$ is the set of fixed points of $`X`$.
In the last section we generalize this property of ample line bundles under the assumption that $`X`$ is smooth. We prove that the analogous map is still an epimorphism if we replace $`X_0`$ with any set of pairwise disjoint invariant subvarieties. In this case, the blowing-up $`\stackrel{~}{X}`$ of $`X`$ along the union of these subvarieties is still a toric variety and we obtain the required surjectivity by applying to $`\stackrel{~}{X}`$ the results in the second section.
Acknowledgements. This work started from a joint project with David Eisenbud and Mike Stillman to understand the cohomology of sheaves on toric varieties. It is a pleasure to thank them for encouragement and generous support. We are also very grateful to William Fulton, Robert Lazarsfeld and Sorin Popescu for useful discussions and to Markus Perling for his comments on an earlier version of this paper. Last but not least, we acknowledge the referee’s numerous comments and suggestions which greatly improved the quality of our presentation.
## 1. The homogeneous coordinate ring of a toric variety
Let $`k`$ be a fixed algebraically closed field of arbitrary characteristic. We will use freely the definitions and results on toric varieties from \[Fu\]. We first review the notation we are going to use.
Let $`N^n`$ be a lattice and $`M=\mathrm{Hom}(N,)`$ the dual lattice. For a rational fan $`\mathrm{\Delta }`$ in $`V=N`$, we have the associated toric variety $`X=X(\mathrm{\Delta })`$. For every $`in`$, the set of cones in $`\mathrm{\Delta }`$ of dimension $`i`$ is denoted by $`\mathrm{\Delta }_i`$. The torus $`N_{}k^{}`$ acts on $`X`$, and by an invariant subvariety of $`X`$ we mean a subvariety which is invariant under this action.
The closed invariant subvarieties of $`X`$ of dimension $`i`$ are in bijection with the set $`\mathrm{\Delta }_{ni}`$. For each cone $`\tau \mathrm{\Delta }`$ we denote by $`V(\tau )`$ the corresponding subvariety. Recall that $`V(\tau )`$ is again a toric variety and $`\tau _1\tau _2`$ if and only if $`V(\tau _2)V(\tau _1)`$. In particular, the prime invariant Weil divisors $`D_1,\mathrm{},D_d`$ on $`X`$ correspond to the one dimensional cones in $`\mathrm{\Delta }`$. If $`X`$ is smooth, then so is each $`V(\tau )`$.
Let $`V^{}`$ be the vector space spanned by $`\mathrm{\Delta }`$, $`N^{}=NV^{}`$ and $`M^{}=\mathrm{Hom}(N^{},)`$ its dual lattice. We have an exact sequence:
$$0M^{}\mathrm{Div}_T(X)\mathrm{Cl}(X)0,$$
where $`\mathrm{Div}_T(X)=_{i=1}^dD_i^d`$ is the group of invariant Weil divisors and $`\mathrm{Cl}(X)`$ is the class group of $`X`$.
We fix a decomposition $`MM^{}\times ^e`$, where $`e`$ is the codimension of $`V^{}`$ in $`V`$. We correspondingly have a decomposition $`XX^{}\times (k^{})^e`$, where $`X^{}`$ is the toric variety defined by $`\mathrm{\Delta }`$ in $`N^{}`$.
The homogeneous coordinate ring of $`X`$ was introduced by Cox in \[Cox\] in the case when the fan $`\mathrm{\Delta }`$ is not degenerate, i.e., is not contained in a hyperplane. We slightly generalize this notion in order to allow an arbitrary toric variety, following the suggestion in \[Cox\]. We first review some of the definitions and the results in that paper, all of which can be easily generalized to this context.
For each $`i`$ with $`1id`$ we introduce an indeterminate $`Y_i`$, corresponding to the divisor $`D_i`$. We introduce also the indeterminates $`Y_j`$ with $`d+1jd+e`$, and the homogeneous coordinate ring of $`X`$ is the ring $`S=k[Y_1,\mathrm{},Y_d,Y_{d+1}^{\pm 1}\mathrm{},Y_{d+e}^{\pm 1}]`$. Note that the decomposition $`MM\times ^e`$ corresponds to a decomposition $`k[M]k[M^{}]k[Y_{d+1}^{\pm 1},\mathrm{},Y_{d+e}^{\pm 1}]`$.
For every effective divisor $`D=_{i=1}^da_iD_i`$, we write $`Y^D`$ for the corresponding monomial $`_{i=1}^dY_i^{a_i}S`$. On the ring $`S`$ we have a fine grading, the usual $`^{d+e}`$-grading by monomials. However, in this section we will consider exclusively a coarse $`\mathrm{Cl}(X)`$-grading defined by
$$\mathrm{deg}(\underset{i=1}{\overset{d+e}{}}Y_i^{a_i})=[\underset{i=1}{\overset{d}{}}a_iD_i]\mathrm{Cl}(X).$$
In the ring $`S`$ there is a distinguished ideal which is a reduced monomial ideal. For each cone $`\sigma \mathrm{\Delta }`$ we put $`D_{\widehat{\sigma }}=_{i;\tau _i\sigma }D_i`$, the sum being taken over the divisors corresponding to one dimensional cones outside $`\sigma `$ and $`Y^{\widehat{\sigma }}=Y^{D_{\widehat{\sigma }}}`$. If $`\mathrm{\Delta }_{\mathrm{max}}`$ is the set of maximal cones in $`\mathrm{\Delta }`$, then $`B=(Y^{\widehat{\sigma }}|\sigma \mathrm{\Delta }_{\mathrm{max}})`$.
As in the case of projective space, a graded $`S`$-module $`P`$ gives a quasicoherent sheaf on $`X`$ by the following procedure. $`X`$ is covered by the affine toric varieties $`U_\sigma =\mathrm{Spec}k[\sigma ^{}M]`$, for $`\sigma \mathrm{\Delta }`$. Using the above decomposition of $`k[M]`$ and the argument in \[Cox\], we obtain canonical isomorphisms $`k[\sigma ^{}M](S_{Y^{\widehat{\sigma }}})_0`$ for every $`\sigma \mathrm{\Delta }`$, which are pairwise compatible. Therefore if $`P`$ is a graded $`S`$-module, on the affine piece $`U_\sigma `$ we can consider the quasicoherent sheaf defined by $`(P_{Y^{\widehat{\sigma }}})_0`$. These sheaves glue together to give a quasicoherent sheaf $`\stackrel{~}{P}`$ on $`X`$. In this way we get an exact functor $`P\stackrel{~}{P}`$ from graded $`S`$-modules to quasicoherent sheaves. If $`P`$ is finitely generated, then $`\stackrel{~}{P}`$ is coherent.
In particular, if $`\alpha \mathrm{Cl}(X)`$, $`𝒪(\alpha )`$ is defined to be $`\stackrel{~}{S(\alpha )}`$. As in \[Cox\], if $`\alpha =[D]`$, then there is a natural isomorphism $`𝒪(\alpha )𝒪(D)`$. Moreover, we have an isomorphism of graded rings
$$S\underset{\alpha \mathrm{Cl}(X)}{}H^0(X,𝒪(\alpha )).$$
For a quasicoherent sheaf $``$, we put $`(\alpha ):=𝒪(\alpha )`$.
Remark. In general, if $`P`$ is a graded $`S`$-module, the natural morphism $`\stackrel{~}{P}𝒪(\alpha )\stackrel{~}{P(\alpha )}`$ is not an isomorphism. However, it is an isomorphism if $`\alpha \mathrm{Pic}(X)`$. Indeed, by taking a graded free presentation of $`P`$, we can reduce ourselves to the case when $`P=S(\beta )`$ for some $`\beta =[E]`$. Since $`\alpha =[D]`$ with $`D`$ locally invertible, $`𝒪(\alpha )`$ is invertible and the fact that the morphism $`𝒪(D)𝒪(E)𝒪(D+E)`$ is an isomorphism follows now directly from the definition.
We prove now that every quasicoherent sheaf is isomorphic to $`\stackrel{~}{P}`$ for some graded $`S`$-module $`P`$. This was proved in \[Cox\] under the assumption that $`X`$ is simplicial. With a slightly different definition for the homogeneous coordinate ring it was proved more generally for toric varieties with enough effective invariant divisors by Kajiwara in \[Ka\].
###### Theorem 1.1.
For every toric variety $`X`$ and every quasicoherent sheaf $``$ on $`X`$, there is a graded $`S`$-module $`P`$ such that $`\stackrel{~}{P}`$.
###### Proof.
We take $`P=_{\alpha \mathrm{Cl}(X)}H^0(X,(\alpha ))`$, which is clearly a graded $`S`$-module. For simplicity, we will use the notation $`P_\sigma =P_{Y^{\widehat{\sigma }}}`$.
For each $`\sigma \mathrm{\Delta }`$, there are canonical maps
$$\varphi _\sigma :(P_\sigma )_0H^0(U_\sigma ,),$$
defined as follows. If $`s/Y^D(P_\sigma )_0`$ such that $`sH^0(X,(\alpha ))`$ and $`D`$ is an effective divisor with $`[D]=\alpha `$ and $`\mathrm{Supp}DU_\sigma =\mathrm{}`$, then $`1/Y^D`$ defines a section in $`H^0(U_\sigma ,𝒪(\alpha ))`$ and $`\varphi _\sigma (s/Y^D)=(1/Y^D)s`$ is the image of $`(1/Y^D,s)`$ by the canonical pairing
$$H^0(U_\sigma ,𝒪(\alpha ))\times H^0(X,(\alpha ))H^0(U_\sigma ,).$$
These morphisms glue together to give $`\varphi :\stackrel{~}{P}`$ (note that $``$ is assumed to be quasicoherent). We will prove that $`\varphi `$ is an isomorphism by proving that $`\varphi _\sigma `$ is an isomorphism for each $`\sigma \mathrm{\Delta }`$.
We first show that $`\varphi _\sigma `$ is a monomorphism. Suppose that $`\varphi _\sigma (s/Y^D)=0`$ for some $`sH^0(X,(\alpha ))`$ and $`D`$ effective, $`[D]=\alpha `$.
We may assume that $`\mathrm{Supp}D=_{\tau _i\sigma }V(\tau _i)`$, and in this case we will prove that there is an integer $`N1`$ such that $`Y^{ND}s=0`$ in $`H^0(X,(\alpha +N\alpha ))`$. In fact, we will find for each $`\tau \mathrm{\Delta }`$ an integer $`N_\tau `$ such that $`Y^{N_\tau D}s|_{U_\tau }=0`$. Then it is clear that $`N=_\tau N_\tau `$ satisfies the requirement.
From now on, we fix also $`\tau \mathrm{\Delta }`$. Since $`\sigma \tau `$ is a face of $`\tau `$, we can write $`\sigma \tau =\tau u^{}`$ for some $`u\tau ^{}M`$. If for each $`vM`$, the corresponding element of $`k[M]`$ is denoted by $`\chi ^v`$, we consider the principal divisor $`D_0=\mathrm{div}(\chi ^u)`$. It is effective on $`U_\tau `$, where its support coresponds to the one-dimensional cones $`\tau _i\tau \sigma `$.
We consider the restrictions of all the sections from above to $`U_\tau `$: $`s|_{U_\tau }H^0(U_\tau ,(\alpha ))`$, $`Y^D|_{U_\tau }H^0(U_\tau ,𝒪(\alpha ))`$ and $`(1/Y^D)|_{U_\sigma U_\tau }H^0(U_{\sigma \tau },𝒪(\alpha ))`$.
Since $`\varphi _\sigma (s/Y^D)=0`$ in $`H^0(U_\sigma ,)`$, we have that $`s|_{U_\sigma }=0H^0(U_\sigma ,(\alpha ))`$, as the image of $`(Y^D,\varphi _\sigma (s/Y^D))`$ by the canonical pairing
$$H^0(U_\sigma ,𝒪(\alpha ))H^0(U_\sigma ,)H^0(U_\sigma ,(\alpha )).$$
In particular, we have $`s|_{U_{\sigma \tau }}=0`$. But $`U_\sigma U_\tau =U_{\sigma \tau }U_\tau `$ is a principal affine subset defined by $`Y^{D_0}H^0(U_\tau ,𝒪_X)`$. Therefore, we get an integer $`t1`$ such that $`Y^{tD_0}s=0`$ in $`H^0(U_\tau ,(\alpha ))`$.
If $`a_\tau ^{}`$ and $`a_\tau ^{}^0`$ are the coefficients of $`V(\tau ^{})`$ in $`D`$ and $`D_0`$, respectively, and $`N_\tau `$ is such that $`N_\tau a_\tau ^{}ta_\tau ^{}^0`$ for every one-dimensional face $`\tau ^{}\tau `$ (by the form of $`D`$ and $`D_0`$, we can choose such an $`N_\tau `$), then $`Y^{N_\tau D}s=0`$ in $`H^0(U_\tau ,(\alpha +N_\tau \alpha ))`$. This follows from the fact that if $`\tau ^{\prime \prime }`$ is an one-dimensional cone with $`\tau ^{\prime \prime }\tau `$, then $`𝒪(V(\tau ^{\prime \prime }))|_{U_\tau }`$ is invertible and $`Y^{V(\tau ^{\prime \prime })}`$ is an invertible section in it. This completes the proof of the fact that $`\varphi _\sigma `$ is a monomorphism.
We prove now that $`\varphi _\sigma `$ is an epimorphism. Let $`tH^0(U_\sigma ,)`$, and let $`D=_{\tau _i\sigma }D_i`$ and $`\alpha =[D]`$.
Using an analogous argument, we see that for each $`\tau \mathrm{\Delta }`$, there is an integer $`N_\tau `$ such that $`Y^{N_\tau D}t|_{U_{\sigma \tau }}H^0(U_{\sigma \tau },(N_\tau \alpha ))`$ can be extended to a section in $`H^0(U_\tau ,(N_\tau \alpha ))`$. Indeed, with the notation and arguments we used before, we first find $`N_\tau ^{}`$ such that $`Y^{N_\tau ^{}D_0}t`$ can be extended to $`U_\tau `$ and then find $`N_\tau `$, as claimed.
If we apply this to two cones $`\tau _1`$, $`\tau _2\mathrm{\Delta }`$ and take $`NN_{\tau _1}`$, $`N_{\tau _2}`$, we see that $`Y^{ND}t`$ can be extended to both $`U_{\tau _1}`$ and $`U_{\tau _2}`$, giving sections $`t_1`$ and $`t_2`$, respectively. Since $`(t_1t_2)|_{U_{\sigma \tau _1\tau _2}}=0`$, by applying to $`\tau _1\tau _2`$ the argument we used to show that $`\varphi _\sigma `$ is a monomorphism, we find $`N_{\tau _{12}}`$ such that $`Y^{N_{\tau _{12}}D}t_1=Y^{N_{\tau _{12}}D}t_2`$ on $`U_{\tau _1}U_{\tau _2}`$.
This shows that for large enough $`N`$, we can extend $`Y^{ND}t|_{U_{\sigma \tau }}`$ to $`t_\tau H^0(U_\tau ,(N\alpha ))`$ for every $`\tau \mathrm{\Delta }`$ such that $`t_{\tau _1}|_{U_{\tau \tau _2}}=t_{\tau _2}|_{U_{\tau _1\tau _2}}`$ for every $`\tau _1`$, $`\tau _2\mathrm{\Delta }`$. Therefore $`t`$ is in the image of $`\varphi _\sigma `$, which completes the proof. ∎
Using the same argument as in \[Cox\], we deduce the following corollary.
###### Corollary 1.2.
For every toric variety $`X`$ and every coherent sheaf $``$ on $`X`$, there is a finitely generated $`S`$-module $`P^{}`$ such that $`\stackrel{~}{P^{}}`$.
###### Proof.
With the notation in the proof of Theorem 1.1, we have seen that
$$\varphi _\sigma :(P_\sigma )_0H^0(U_\sigma ,)$$
is an isomorphism for every $`\sigma \mathrm{\Delta }`$.
Since $``$ is coherent, this implies that we can find a finitely generated graded submodule $`P^{}P`$ such that $`(P_\sigma ^{})_0=(P_\sigma )_0`$ for every $`\sigma \mathrm{\Delta }`$. It is clear that this $`P^{}`$ satisfies the assertion of the corollary. ∎
As in the case of projective space, the cohomology of the sheaf $`\stackrel{~}{P}`$ can be expressed as the local cohomology of the module $`P`$ at the irrelevant ideal $`B`$.
###### Proposition 1.3.
Let $`P`$ be a graded $`S`$-module. Then there exist an isomorphism of graded modules
$$H_B^{i+1}(P)\underset{\alpha \mathrm{Cl}(X)}{}H^i(X,\stackrel{~}{P(\alpha )}),$$
for every $`i1`$ and an exact sequence
$$0H_B^0(P)P\underset{\alpha \mathrm{Cl}(X)}{}H^0(X,\stackrel{~}{P(\alpha )})H_B^1(P)0.$$
###### Proof.
$`X`$ is covered by the affine open subsets $`U_\sigma `$, $`\sigma \mathrm{\Delta }_{\mathrm{max}}`$, and we compute the cohomology of $`\stackrel{~}{P}`$ as Cech cohomology with respect to this cover.
On the other hand, we can compute the local cohomology of $`P`$ at $`B`$ using the direct limit of Koszul complexes on the powers of the generators of $`B=(Y_{\widehat{\sigma }}|\sigma \mathrm{\Delta }_{\mathrm{max}})`$ (see \[Ei\], Appendix 4.1).
Since for $`\sigma _1,\mathrm{},\sigma _t\mathrm{\Delta }_{\mathrm{max}}`$, $`_{i=1}^tU_{\sigma _i}=U_\sigma `$, where $`\sigma =_{i=1}^t\sigma _i`$ and
$$H^0(U_\sigma ,\stackrel{~}{P(\alpha )})=(P(\alpha )_{Y^{\widehat{\sigma }}})_0=(P_{Y_{\widehat{\sigma }_1},\mathrm{},Y_{\widehat{\sigma }_t}})_\alpha ,$$
we conclude as in the case of the projective space (see \[Ei\], Appendix 4.1). ∎
Note. In the situation in Proposition 1.3, suppose that $`P`$ is in fact a $`^{d+e}`$-graded $`S`$-module, so that the corresponding sheaf $`\stackrel{~}{P}`$ is equivariant with respect to the torus action. In this case the local cohomology module $`H_B^{i+1}(P)`$ is $`^{d+e}`$-graded, too, and under the isomorphism in Proposition 1.3 this finer decomposition of $`H_B^{i+1}(P)`$ corresponds to the eigenspace decomposition of the Zariski cohomology of the different twists $`\stackrel{~}{P(\alpha )}`$.
## 2. Vanishing theorems
We keep the notation from the previous section. However, from now on we consider on $`S`$ the fine $`^{d+e}`$ grading by monomials and all $`S`$-modules are assumed to be $`^{d+e}`$-graded. Note that this implies that the associated sheaf is equivariant with respect to the torus action. The canonical basis of $`^{d+e}`$ will be denoted by $`f_1,\mathrm{},f_{d+e}`$.
For every subset $`I`$ and every graded $`S`$-module $`P`$, we will say that $`P`$ is $`I`$-homogeneous if for every $`\alpha =(\alpha _j)^{d+e}`$ with $`\alpha _jI`$, the multiplication by $`Y_j`$:
$$\nu _{Y_j}:P_\alpha P_{\alpha +f_j}$$
is an isomorphism. Our main example is $`S`$, which is obviously $`\{1\}`$-homogeneous.
###### Proposition 2.1.
If $`P`$ is an $`I`$-homogeneous $`S`$-module, then $`H_B^i(P)`$ is $`I`$-homogeneous.
###### Proof.
We compute the local cohomology module as the cohomology of a Cech-type complex (see, for example, \[Ei\], Appendix 4.1). Let us temporarily denote the generators of $`B`$ by $`m_1,\mathrm{},m_t`$. For a subset $`L\{1,\mathrm{},t\}`$, let $`m_L`$ be the least common multiple of $`\{m_l|lL\}`$. Since it is enough to prove the assertion at the level of complexes, we have to check that for every $`\alpha ^{d+e}`$ with $`\alpha _jI`$, the multiplication by $`Y_j`$:
$$\mu _{Y_j}:(P_{m_L})_\alpha (P_{m_L})_{\alpha +f_j}$$
is an isomorphism.
This is obvious if $`Y_j|m_L`$. Suppose now that $`Y_jm_L`$. Then the assertion is clear once we notice that in this case, if $`m/m_L^sP_{m_L}`$, then $`\mathrm{deg}(m/m_L^s)_j=\mathrm{deg}(m)_j`$, so that we can apply the fact that $`P`$ in $`I`$-homogeneous. ∎
We consider now an example of $`\{1,0\}`$-homogeneous $`S`$-modules. These are the modules which define the exterior powers $`\mathrm{\Omega }_X^i`$ of the cotangent sheaf. For simplicity, in this case we will assume that $`X`$ is smooth.
It is shown by Batyrev and Cox in \[BC\] that if the fan defining $`X`$ is nondegenerate, then the cotangent bundle on $`X`$ appears in an Euler sequence:
$$0\mathrm{\Omega }_X^1_{j=1}^d𝒪_X(D_j)𝒪_X^{dn}0.$$
In general, we have $`XX^{}\times (k^{})^e`$ with $`X^{}`$ as above and $`\mathrm{\Omega }_X^1p_1^{}(\mathrm{\Omega }_X^{}^1)𝒪_X^e`$. Therefore we can include $`\mathrm{\Omega }_X^{}^1`$ in an exact sequence:
$$0\mathrm{\Omega }_X^1(_{j=1}^d𝒪_X(D_j))𝒪_X^e𝒪_X^{dn+e}0.$$
We consider the graded morphism inducing the epimorphism in the second exact sequence:
$$E=(_{j=1}^dS(f_j))S^eF=S^{dn+e}.$$
For each $`i1`$, let $`M_i`$ be the kernel of the induced map $`^iE^{i1}EF`$.
###### Lemma 2.2.
With the above notation, we have
(i) $`\stackrel{~}{M}_i\mathrm{\Omega }_X^i`$.
(ii) $`M_i`$ is $`\{1,0\}`$-homogeneous.
###### Proof.
(i) The assertion follows easily from the above mentioned result of Batyrev and Cox and the fact that in the Euler sequence all the sheaves are locally free.
(ii) Since $`M_i`$ is a submodule of $`^iE`$, which is free, the multiplication by $`Y_j`$ on $`M_i`$ is injective.
Let $`\alpha =(\alpha _j)^{d+e}`$, $`\alpha _j\{1,0\}`$. Since $`^{i1}EF`$ is free, the surjectivity of the map $`\nu _{Y_j}:(M_i)_\alpha (M_i)_{\alpha +f_j}`$ follows from the surjectivity of the analogous map for $`^iE`$. The latter is surjective since $`^iE`$ is a direct sum of modules of the form $`S(f_{j_1}\mathrm{}f_{j_r})`$ with $`ri`$ and $`j_1<\mathrm{}<j_r`$. ∎
###### Proposition 2.3.
Let $`X`$ be an arbitrary toric variety.
(i) If $`P`$ is a $`\{1,0\}`$-homogeneous $`S`$-module and $`L\mathrm{Pic}(X)`$ is such that $`H^i(\stackrel{~}{P}L^m)=0`$ for some $`i0`$ and $`m1`$, then $`H^i(\stackrel{~}{P}L)=0`$. In particular, if $`X`$ is projective and $`L\mathrm{Pic}(X)`$ is ample, then $`H^i(\stackrel{~}{P}L)=0`$ for all $`i1`$.
(ii) Let $`P`$ be a $`\{1\}`$-homogeneous $`S`$-module such that for every $`\alpha \mathrm{Cl}(X)`$, $`\stackrel{~}{P(\alpha )}\stackrel{~}{P}𝒪(\alpha )`$. Suppose that $`D\mathrm{Div}_T(X)`$ and that there is $`E=_{j=1}^da_jD_j`$ with $`a_j`$ and $`0a_j1`$ such that $`m(D+E)`$ is integral and Cartier for some integer $`m1`$. If $`H^i(\stackrel{~}{P}𝒪_X(D+m(D+E)))=0`$ for some $`i0`$, then $`H^i(\stackrel{~}{P}𝒪_X(D))=0`$. In particular, if $`X`$ is projective and we have $`E`$ aforementioned such that $`D+E`$ is $``$-ample, then $`H^i(\stackrel{~}{P}𝒪_X(D))=0`$ for all $`i1`$.
###### Proof.
(i) If $`L=𝒪(\alpha )`$, $`H^i(\stackrel{~}{P}L)=H^i(\stackrel{~}{P(\alpha )})`$ (see the remark in the first section). We will restrict ourselves to the case $`i1`$ in order to apply the isomorphism in Proposition 1.3. When $`i=0`$, one can give a similar argument using the exact sequence in that proposition.
As already mentioned, we have
$$H^i(\stackrel{~}{P}L)\underset{\underset{¯}{b}}{}H_B^{i+1}(P)_{\underset{¯}{b}},$$
where the direct sum is taken over those $`\underset{¯}{b}=(b_1,\mathrm{},b_{d+e})^{d+e}`$ such that $`[_{i=1}^db_iD_i]=\alpha `$. Since by hypothesis $`H^i(\stackrel{~}{P}L^m)=0`$, for every $`\underset{¯}{b}`$ with $`[_{i=1}^db_iD_i]=\alpha `$ we have $`H_B^i(P)_{m\underset{¯}{b}}=0`$. Proposition 2.1 implies that
$$H_B^i(P)_{\underset{¯}{b}}H_B^i(P)_{m\underset{¯}{b}},$$
which proves the first assertion.
In the case of an ample line bundle $`L`$ on a projective toric variety, we have $`H^i(\stackrel{~}{P}L^m)=0`$ for $`i1`$ and $`m0`$, so that we are in the previous situation.
(ii) We proceed similarly. Using our hypothesis on $`P`$ and Proposition 1.3, for every $`i1`$ we have
$$H^i(\stackrel{~}{P}𝒪(\alpha ))H^i(\stackrel{~}{P(D)})\underset{\underset{¯}{b}}{}H_B^{i+1}(P)_{\underset{¯}{b}},$$
where the direct sum is taken over those $`\underset{¯}{b}=(b_1,\mathrm{},b_{d+e})^{d+e}`$ such that $`[_{i=1}^db_iD_i]=[D]`$.
Using again the hypothesis on $`P`$ and the fact that $`m(D+E)`$ is Cartier (see the remark in the first section), we get
$$H^i((P(D+m(D+E)))\stackrel{~}{})=0.$$
We fix some $`\underset{¯}{b}^{d+e}`$ with $`[_{i=1}^db_iD_i]=[D]`$. We have to prove that $`H_B^{i+1}(P)_{\underset{¯}{b}}=0`$. If $`\underset{¯}{b}^{}=\underset{¯}{b}+m(\underset{¯}{b}+\underset{¯}{a})`$, where $`\underset{¯}{a}=(a_1,\mathrm{},a_d,0\mathrm{},0)`$, then $`[_{i=1}^db_i^{}D_i]=[D+m(D+E)]`$, and therefore $`H_B^{i+1}(P)_{\underset{¯}{b}^{}}=0`$.
Proposition 2.1 implies that in order to complete the proof, it is enough to show that $`b_j0`$ if and only if $`(m+1)b_j+ma_j0`$. This follows easily from the fact that $`0a_j1`$. ∎
We apply Proposition 2.3 in conjunction with Lemma 2.2 for $`P=M_i`$ and for $`P=S`$.
###### Theorem 2.4.
(i) (Bott-Steenbrink-Danilov) If $`X`$ is a smooth toric variety and $`L\mathrm{Pic}(X)`$ is such that $`H^i(\mathrm{\Omega }_X^jL^m)=0`$ for some $`m1`$ and $`i0`$, then $`H^i(\mathrm{\Omega }_X^jL)=0`$. In particular, if $`X`$ is projective and $`L\mathrm{Pic}(X)`$ ample, then $`H^i(\mathrm{\Omega }_X^jL)=0`$ for every $`i1`$.
(ii) Let $`X`$ be an arbitrary toric variety, $`D\mathrm{Div}_T(X)`$ and $`E=_{j=1}^da_jD_j`$, with $`a_j`$ and $`0a_j1`$ such that for some integer $`m1`$ we have $`m(D+E)`$ integral and Cartier. If $`H^i(𝒪_X(D+m(D+E)))=0`$, then $`H^i(𝒪_X(D))=0`$. In particular, if $`X`$ is projective and there is $`E`$ aforementioned such that $`D+E`$ is $``$-ample, then $`H^i(𝒪_X(D))=0`$ for all $`i1`$.
Remark. As pointed out by the referee, in the case $`P=S`$ the assertion in Proposition 2.1 can be proved also via the combinatorial description of the cohomology of a sheaf of fractional ideals (see for example \[KKMS\], pg. 42). More precisely, the graded components $`H_B^{i+1}(S)_\alpha `$ and $`H_B^{i+1}(S)_{\alpha +f_j}`$ (or, equivalently, the corresponding eigenspaces of $`H^i(X,𝒪(\alpha ))`$ and $`H^i(X,𝒪(\alpha +f_j))`$) can be described as simplicial cohomology groups of certain subsets of $`^n`$. The assertion can be proved by showing that these spaces are homotopically equivalent. Note that the case $`P=S`$ is enough to give the statement of Proposition 2.4 (ii).
If $`D=_{j=1}^db_jD_j`$ is a $``$-divisor, we define
$$D:=\underset{j=1}{\overset{d}{}}b_iD_j,$$
where for any real number $`x`$, $`x`$ is the integer defined by $`xx<x+1`$. Similarily, we define
$$D:=\underset{j=1}{\overset{d}{}}b_jD_j,$$
where for every $`x`$, $`x`$ is the integer defined by $`x1<xx`$. $`K_X`$ denotes the canonical divisor $`_{j=1}^dD_j`$ so that $`\omega _X=𝒪(K_X)`$.
###### Corollary 2.5.
Let $`X`$ be a projective toric variety.
(i) (Kawamata-Viehweg) If $`D=_{j=1}^db_jD_j`$ is a $``$-Cartier ample $``$-divisor, then $`H^i(𝒪_X(K_X+D))=0`$ for every $`i1`$.
(ii) If $`D`$ is as above, then $`H^i(O_X(D))=0`$ for every $`i1`$.
(iii) Let $`L\mathrm{Pic}(X)`$ be an ample bundle. If $`D_{j_1},\mathrm{},D_{j_r}`$ are distinct prime invariant divisors, then $`H^i(L(D_{j_1}\mathrm{}D_{j_r}))=0`$ for every $`i1`$.
###### Proof.
All these are particular cases of Theorem 2.4 (ii). ∎
Remark. In the proof of Fujita’s Conjecture we will use the assertion in Corollary 2.5 for smooth varieties. As the referee pointed out, when $`X`$ is smooth it is possible to prove this assertion directly, by induction on dimension and descending induction on $`r`$, as for $`r=d`$ this is just Kodaira’s vanishing theorem.
As we mentioned in the Introduction, some particular cases of the above results can be proved by reducing the problem to a toric variety $`X`$ over a field of positive characteristic $`p`$ and prove that such a variety is Frobenius split. This means that if $`F`$ is the Frobenius morphism, then the canonical morphism $`𝒪_XF_{}𝒪_X`$ has a left inverse. With the description for the cohomology we used above this can be seen as follows.
First of all, by embedding $`X`$ as an open subvariety of a complete toric variety, we may suppose that $`X`$ is complete. Next, by taking a toric resolution of singularities, we may suppose that $`X`$ is also smooth (see \[MR\]). Moreover, an argument in that paper shows that in this case, if $`dim(X)=n`$, then $`X`$ is Frobenius split if and only if the morphism
$$f:H^n(\omega _X)H^n(\omega _X^p)$$
induced by the Frobenius morphism is not trivial. But $`H^n(\omega _X)H_B^{n+1}(S)_{(1,\mathrm{},1)}k`$, all the other components being zero. On the other hand,
$$H^n(\omega _X^p)\underset{[{\scriptscriptstyle (a_j+p)D_j}]=0}{}H_B^{n+1}(S)_{\underset{¯}{a}}$$
has by Proposition 2.1 the component $`H_B^{n+1}(S)_{(p,\mathrm{},p)}`$ canonically isomorphic with $`H_B^{n+1}(S)_{(1,\mathrm{},1)}`$ and therefore with $`k`$. It is easy to see that via these identifications, the corresponding component of $`f`$ is just the Frobenius morphism of $`k`$, and therefore $`f`$ is nonzero.
For a different approach to Frobenius splitting in the toric context and other applications we refer to Buch, Thomsen, Lauritzen and Mehta \[BTLM\].
## 3. Ample and numerically effective line bundles
Our main goal in this section is to give the condition for a line bundle to be ample or nef (i.e., numerically effective) in terms of the intersection with the invariant curves. For ampleness, this is the toric Nakai criterion which is proved in \[Oda\] for the smooth case and is stated also for the simplicial case. We obtain also a similar condition for the nef property, both the results holding for arbitrary complete toric varieties. In particular, we will see that on such a variety, a line bundle is nef if and only if it is globally generated. With a different proof, these results have been obtained also by Mavlyutov in \[Ma\]. We use the ideas in \[Oda\] together with the description for the intersection with divisors in the non-smooth case from \[Fu\].
We will apply these results to show that a line bundle $`L`$ on $`X`$ which is big and nef is a pull-back of an ample line bundle on $`X^{}`$, for a proper birational equivariant map of toric varieties $`\varphi :XX^{}`$. Recall that a line bundle $`L`$ on $`X`$ is called nef if for every curve $`CX`$, $`(LC)0`$.
###### Theorem 3.1.
If $`X`$ is a complete toric variety and $`L\mathrm{Pic}(X)`$, the following are equivalent:
(i) $`L`$ is globally generated.
(ii) $`L`$ is nef.
(iii) For every invariant integral curve $`CX`$, $`(LC)0`$.
###### Proof.
(i)$``$(ii) is true in general and (ii)$``$(iii) follows from the definition.
We now prove the implication (iii)$``$(i). Let $`D`$ be an invariant Cartier divisor such that $`L𝒪(D)`$. Recall that there is a function $`\psi =\psi _D:N`$ associated with $`D`$ which is linear on each cone $`\sigma \mathrm{\Delta }`$. It is defined in the following way: if $`D|_{U_\sigma }=\mathrm{div}(\chi ^{u_\sigma })|_{U_\sigma }`$, then $`\psi |_\sigma =u_\sigma |_\sigma `$ (the notation is that used in the first section).
A well-known result(see \[Fu\], Section 3.3) says that $`L`$ is globally generated if and only if $`\psi `$ is convex. Recall that $`dim(X)=n`$. To prove that $`\psi `$ is convex, it is enough to prove that for every $`\sigma _1`$, $`\sigma _2\mathrm{\Delta }_n`$ with $`dim(\sigma _1\sigma _2)=n1`$, $`\psi |_{\sigma _1\sigma _2}`$ is convex, i.e., for every $`x\sigma _1`$, $`y\sigma _2`$ and $`t[0,1]`$ such that $`tx+(1t)y\sigma _1\sigma _2`$, we have $`\psi (tx+(1t)y)t\psi (x)+(1t)\psi (y)`$.
It is clear, therefore, from the definition of $`\psi `$ that it is enough to prove that for each $`\sigma _1`$, $`\sigma _2`$ as above and each $`D_i=V(\tau _i)`$, with $`\tau _i\sigma _2\sigma _1`$ a one-dimensional cone,
$$u_{\sigma _2}(v_i)u_{\sigma _1}(v_i),$$
where $`v_i`$ is the primitive vector of $`\tau _i`$.
Let $`D=_{j=1}^da_jD_j`$. Note that by definition, if $`D_j=V(\tau _j)`$, $`\tau _j\sigma `$, then $`u_\sigma (v_j)=a_j`$. For $`\sigma _1`$ and $`\sigma _2`$ as above, let $`\tau =\sigma _1\sigma _2`$. Our hypothesis gives $`(DV(\tau ))0`$. By definition, $`(D+\mathrm{div}(\chi ^{u_{\sigma _1}}))|_{U_{\sigma _1}}=0`$. Therefore
$$D+\mathrm{div}(\chi ^{u_{\sigma _1}})=\underset{\tau _i\sigma _2\sigma _1}{}b_iD_i+\mathrm{},$$
where we wrote down only the divisors corresponding to cones in $`\sigma _1\sigma _2`$. Since $`a_i=u_{\sigma _2(v_i)}`$ for $`\tau _i\sigma _2`$, we get
$$b_i=u_{\sigma _1}(v_i)u_{\sigma _2}(v_i),$$
if $`\tau _i\sigma _2\sigma _1`$.
On the other hand, let us denote by $`\overline{e}`$ the generator of the one-dimensional lattice $`N/N_\tau `$ such that the classes of the primitive vectors of $`\tau _i`$ for $`\tau _i\sigma _2\sigma _1`$ are positive multiples of $`\overline{e}`$. Here $`N_\tau `$ denotes the subgroup of $`N`$ generated by $`N\tau `$. If for every $`\tau _i`$ aforementioned we write $`\overline{v}_i=c_i\overline{e}`$, then the intersection formula in \[Fu\], Section 5.1 shows that
$$(D+\mathrm{div}(\chi ^{u_{\sigma _1}})V(\tau ))=b_i/c_i,$$
for every $`\tau _i\sigma _1\sigma _2`$. Since $`0(DV(\tau ))=b_i/c_i`$ and $`c_i>0`$, we deduce that $`b_i0`$ for every $`\tau _i`$ aforementioned. From the formula for $`b_i`$ we see that the proof is complete. ∎
Remark. The equivalence between (i) and (ii) above can be deduced also from the result of Reid from \[Re\], which says that every effective one dimensional cycle on $`X`$ is rationally equivalent to an effective sum of invariant curves.
###### Theorem 3.2.
(Toric Nakai criterion) If $`X`$ is a complete toric variety, with $`dim(X)=n`$ and $`L\mathrm{Pic}(X)`$, then the following are equivalent:
(i) $`L`$ is ample.
(ii) For every invariant integral curve $`CX`$, $`(LC)>0`$.
###### Proof.
The proof of the relevant implication (ii)$``$(i) is the same as the above proof for the implication (iii)$``$(i). We have just to use the fact that $`L=𝒪(D)`$ is ample if and only if $`\psi _D`$ is strictly convex and to replace all the inequalities by strict inequalities. ∎
Recall that a line bundle $`L\mathrm{Pic}(X)`$ is called big if for a certain multiple $`L^m`$, the rational map it defines: $`\varphi _{L^m}:X\text{P}^N`$ has the image of maximal dimension $`n=dim(X)`$.
###### Proposition 3.3.
(i) If $`X`$ is a complete toric variety of dimension $`n`$ and $`L\mathrm{Pic}(X)`$ is a line bundle which is globally generated and big, then $`dim\varphi _L(X)=n`$.
(ii) $`L\mathrm{Pic}(X)`$ is globally generated and big if and only if there is a fan $`\mathrm{\Delta }^{}`$ such that $`\mathrm{\Delta }`$ is a refinement of $`\mathrm{\Delta }^{}`$ and $`L^{}\mathrm{Pic}(X^{})`$ ample, where $`X^{}=X(\mathrm{\Delta }^{})`$, and that if $`f:XX^{}`$ is the map induced by the refinement, $`f^{}(L^{})L`$.
###### Proof.
Let us fix an invariant Cartier divisor $`D`$ such that $`L𝒪(D)`$. If $`\psi _D`$ is the function which appeared in the proof of Theorem 3.1, it defines an associated convex polytope
$$P_D=\{uM|u\psi _D\mathrm{on}N\}.$$
If $`L`$ is globally generated, then $`dim\varphi _L(X)=dimP_D`$ (see \[Fu\], Section 3.4). But $`P_{mD}=mP_D`$, so that $`dim\varphi _L(X)=dim\varphi _{L^m}(X)`$, which completes the proof of (i).
Since a map as in (ii) is birational, the “if” part in (ii) is trivial. Let us suppose now that $`L`$ is globally generated and big. By the above argument, $`P=P_D`$ is an $`n`$-dimensional convex polytope. Such a polytope defines a complete fan $`\mathrm{\Delta }^{}`$ and an ample Cartier divisor $`D^{}`$ on $`X^{}=X(\mathrm{\Delta }^{})`$. The cones in $`\mathrm{\Delta }^{}`$ are in a one-to-one correspondence, reversing inclusions, with the faces of $`P`$: for a face $`Q`$ of $`P`$ we have the cone
$$C_Q=\{vN|u,vu^{},v\mathrm{for}\mathrm{all}uQ,u^{}P\}.$$
For every $`\sigma \mathrm{\Delta }_n`$, $`u_\sigma `$ is a vertex of $`P`$. Indeed, it is the intersection of $`P`$ with
$$\{uMu,v_i\psi _D(v_i)\mathrm{for}v_i\sigma \}.$$
In fact, every vertex of $`P`$ is of this form. Indeed, if $`u_0`$ is a vertex of $`P`$, then there is $`vN`$ such that $`u_0,v<u,v`$, for all $`uP\{u_0\}`$. In particular, we have $`\psi _D(v)=u_0,v`$. If $`\sigma \mathrm{\Delta }_n`$ is such that $`v\sigma `$, then $`u_0,v=u_\sigma ,v`$, so that $`u_0=u_\sigma `$.
Now it is easy to check that
$$C_{u_\sigma }=\underset{\tau \mathrm{\Delta }_n,u_\tau =u_\sigma }{}\tau .$$
Therefore $`\mathrm{\Delta }`$ is a refinement of $`\mathrm{\Delta }^{}`$. Moreover, the ample divisor $`D^{}`$ on $`X^{}`$ is defined by
$$\psi _D^{}(v)=\mathrm{min}_{\sigma \mathrm{\Delta }_n}u_\sigma ,v=\psi _D(v).$$
It follows that if $`f:XX^{}`$ is the map induced by the refinement, $`f^{}(D^{})=D`$, which completes the proof. ∎
It is easy to see that using the results of this section, we can extend the form of the Kawamata-Viehweg theorem we obtained in the previous section to the case of a divisor which is big and nef. For the proof, however, we have to assume that the divisor is Cartier.
###### Theorem 3.4.
(Kawamata-Viehweg) If $`X`$ is a projective toric variety and $`L\mathrm{Pic}(X)`$ is a line bundle which is nef and big, then $`H^i(\omega _XL)=0`$ for every $`i1`$.
###### Proof.
Since $`L`$ is a line bundle, the duality theorem gives
$$H^i(X,\omega _XL)H^{ni}(X,L^1),$$
where $`n=dim(X)`$ (see \[Fu\], Section 4.4).
Using Theorem 3.1 and Proposition 3.3, we get a morphism $`f:XX^{}`$, induced by a fan refinement, and $`L^{}\mathrm{Pic}(X^{})`$ ample such that $`f^{}(L^{})L`$. But then
$$H^{ni}(X,L^1)H^{ni}(X^{},L^1)H^i(X^{},\omega _X^{}L^{})=0,$$
by Corollary 2.5. ∎
###### Corollary 3.5.
Let $`X`$ be a complete toric variety and $`L`$ a line bundle on $`X`$. If the base locus of $`L`$ is nonempty, then it contains an integral invariant curve $`CX`$.
###### Proof.
This is an immediate consequence of Theorem 3.1, since for an integral curve $`CX`$, if $`C`$ is not contained in the base locus of $`L`$, then $`(LC)0`$. ∎
## 4. Fujita’s conjecture on toric varieties
The main result of this section is the following strong form of Fujita’s Conjecture in the toric case.
###### Theorem 4.1.
Let $`X`$ be an $`n`$-dimensional projective smooth toric variety, $`L\mathrm{Pic}(X)`$ a line bundle and $`D_1,\mathrm{},D_m`$ distinct prime invariant divisors.
(i) If $`(LC)n`$ for every invariant intergral curve $`CX`$, then $`L(D_1\mathrm{}D_m)`$ is globally generated, unless $`X\text{P}^n`$, $`L𝒪(n)`$ and $`m=n+1`$.
(ii) If $`(LC)n+1`$ for every invariant integral curve $`CX`$, then $`L(D_1\mathrm{}D_m)`$ is very ample, unless $`X\text{P}^n`$, $`L𝒪(n+1)`$ and $`m=n+1`$.
In particular, we have the following corollary.
###### Corollary 4.2.
Let $`X`$ be an $`n`$-dimensional projective smooth toric variety and $`L\mathrm{Pic}(X)`$.
(i) If $`(LC)n`$ for every invariant integral curve $`CX`$, then $`\omega _XL`$ is globally generated, unless $`(X,L)(\text{P}^n,𝒪(n))`$.
(ii) If $`(LC)n+1`$ for every invariant integral curve $`CX`$, then $`\omega _XL`$ is very ample, unless $`(X,L)(\text{P}^n,𝒪(n+1))`$.
We prove Theorem 4.1, using the numerical conditions for $`L`$ to be globally generated or ample, as well as the vanishing result in Corollary 2.5 (iii). The proof goes by induction on the dimension of $`X`$, based on the following proposition.
###### Proposition 4.3.
Let $`X`$ be a projective smooth toric variety with $`dim(X)=n`$, $`L\mathrm{Pic}(X)`$ and $`l1`$ an integer. If $`(LC)l`$ for every invariant integral curve $`CX`$, then for every prime invariant divisor $`D`$ and every $`C`$ aforementioned, $`(L(D)C)l1`$.
We first deal with the case $`l=1`$ of this proposition in the lemma below.
###### Lemma 4.4.
Let $`X`$ be a projective smooth toric variety, $`dim(X)=n`$. If $`L\mathrm{Pic}(X)`$ is ample and $`D`$ is an invariant prime divisor, then $`L(D)`$ is globally generated.
###### Proof of Lemma 4.4.
We prove the lemma by induction on $`n`$. For $`n=1`$, $`X=\text{P}^1`$ and the assertion is clear. If $`n2`$ and $`L(D)`$ is not globally generated, since the base locus of $`L(D)`$ is invariant, we can choose a fixed point $`x`$ in this locus.
Let $`D^{}`$ be a prime divisor distinct from $`D`$ and containing $`x`$. By Corollary 2.5 (iii) , the restriction map
$$H^0(L(D))H^0(L(D)|_D^{})$$
is surjective. On the other hand, $`D^{}`$ is a smooth toric variety of dimension $`n1`$ and $`DD^{}`$ is either empty or a prime invariant divisor on $`D^{}`$. Therefore the restriction map
$$H^0(L(D)|_D^{})H^0(L(D)|_x)$$
is also surjective.
Since the composition of the above maps is surjective, we get a contradiction to the assumption that $`x`$ is in the base locus of $`L(D)`$. ∎
We now give the proof of the proposition for an arbitrary $`l1`$.
###### Proof of Proposition 4.3.
We make induction on $`n`$, the case $`n=1`$ being trivial. Note that since $`l1`$, $`L`$ is ample
Let us assume now that $`n=2`$. Clearly, it is enough to prove that $`(L(D)D)l1`$. Since $`(L(D)D)=(LD)(D^2)`$, we may restrict ourselves to the case $`(D^2)2`$. From the description of the selfintersection numbers in terms of the fan $`\mathrm{\Delta }`$, it follows easily that if $`D^{}`$ and $`D^{\prime \prime }`$ are the divisors whose rays are adjacent to the ray corresponding to $`D`$, then $`(D^2)0`$ or $`(D^{\prime \prime 2})0`$.
But if, for example, $`(D^2)0`$, then $`L((l1)D^{})`$ is ample, so that Lemma 4.4 implies that $`L((l1)D^{}D)`$ is globally generated and therefore
$$0(L((l1)D^{}D)D)=(L(D)D)(l1),$$
which completes the case $`n=2`$.
Suppose now that $`n3`$ and let $`\tau \mathrm{\Delta }_{n1}`$ be such that $`C=V(\tau )`$. We can choose a prime invariant divisor $`D^{}`$ such that $`D^{}D`$ and $`CD^{}`$. Therefore $`(L(D)C)=(L(D)|_D^{}C)`$, and we may clearly restrict to the case when $`DD^{}\mathrm{}`$, so that it is a prime invariant divisor on $`D^{}`$. We apply the induction hypothesis for $`L|_D^{}`$; note that for every integral invariant curve $`C^{}D^{}`$,
$$(L|_D^{}C^{})=(LC^{})l.$$
This concludes the proof. ∎
We can now prove the strong form of Fujita’s conjecture for the toric case.
###### Proof of Theorem 4.1.
(i) It is clear that we may assume $`n2`$ and $`X\simeq ̸\text{P}^n`$. We make induction on $`n`$. If $`L(D_1\mathrm{}D_m)`$ is not globally generated, then
$$(L(D_1\mathrm{}D_m)V(\tau ))<0$$
for some $`\tau \mathrm{\Delta }_{n1}`$. We will show that this asumption implies $`X\text{P}^n`$, a contradiction.
We can immediately restrict ourselves to the following situation: $`2mn+1`$, $`D_1`$ and $`D_2`$ are the divisors corresponding to the rays spanning together with $`\tau `$ maximal cones and $`D_3,\mathrm{},D_{n+1}`$ are the divisors containing $`V(\tau )`$.
Claim. We have $`m=n+1`$, $`D_i\text{P}^{n1}`$ for every $`i`$, $`1in+1`$, and $`D_iD_j\mathrm{}`$ for every $`ij`$.
Fix $`i`$ such that $`im`$. Since $`n2`$, our hypothesis and Proposition 4.3 imply that $`L(D_i)`$ is ample. Hence Corollary 2.5 (iii) shows that the restriction map
$$H^0(L(D_1\mathrm{}D_m))H^0(L(D_1\mathrm{}D_m)|_{D_i})$$
is surjective. Since $`V(\tau )\mathrm{Bs}L(D_1\mathrm{}D_m)`$, it follows that $`L(D_1\mathrm{}D_m)|_{D_i}`$ is not globally generated.
Another application of Proposition 4.3 gives $`(L(D_i)C)n1`$ for every integral invariant curve $`CX`$. In particular, $`(L(D_i)|_{D_i}C^{})n1`$ for every integral invariant curve $`C^{}D_i`$. From the induction hypothesis we get $`D_i\text{P}^{n1}`$, $`m=n+1`$ and $`D_iD_j\mathrm{}`$ for $`ji`$.
It is now easy to see that $`X\text{P}^n`$. The claim implies that if $`D_i=V(\tau _i)`$,$`1in+1`$, and if $`\tau _0`$ is any other one-dimensional cone in $`\mathrm{\Delta }`$, then $`\tau _0`$ and $`\tau _i`$ do not span a cone in $`\mathrm{\Delta }`$ for any $`i`$. From this it follows that the only one-dimensional cones in $`\mathrm{\Delta }`$ are $`\tau _1,\mathrm{},\tau _{n+1}`$. Since $`X`$ is smooth, it follows that $`X\text{P}^n`$.
(ii) Since an ample line bundle on a complete smooth toric variety is very ample (see \[De\]), it is enough to prove that if $`L(D_1\mathrm{}D_m)`$ is not ample, then $`X\text{P}^n`$. Again we may assume $`n2`$.
If $`L(D_1\mathrm{}D_m)`$ is not ample, then there exists an invariant integral curve $`CX`$ such that
$$(L(D_1\mathrm{}D_m)C)0.$$
As above, we may assume that $`D_1`$ and $`C`$ correspond to cones in $`\mathrm{\Delta }`$ spanning together a maximal cone.
By Proposition 4.3, we may apply (i) to $`L(D_1)`$ and conclude that if $`X\simeq ̸\text{P}^n`$, then $`L(2D_1D_2\mathrm{}D_m)`$ is globally generated. In particular,
$$(L(2D_1D_2\mathrm{}D_m)C)0,$$
so that
$$(L(D_1\mathrm{}D_m)C)1,$$
a contradiction. ∎
We conclude this section by giving two results with the same flavour as those proved above. By Lemma 4.4, if $`L`$ is ample, then $`L(D)`$ is globally generated for every integral invariant divisor. The case $`X=\text{P}^n`$, $`L=𝒪(1)`$ shows that this is optimal. The next proposition gives the condition under which for $`L`$ ample we get $`L(D_1D_2)`$ globally generated for distinct divisors $`D_1`$ and $`D_2`$ as above.
###### Proposition 4.5.
Let $`X`$ be a projective smooth toric variety with $`dim(X)=n`$, $`L\mathrm{Pic}(X)`$ ample and $`D_1`$, $`D_2`$ distinct prime invariant divisors. Then $`L(D_1D_2)`$ is not globally generated if and only if there is an $`(n1)`$-dimensional cone $`\tau \mathrm{\Delta }`$ such that if $`\tau _1`$, $`\tau _2`$ are the one dimensional cones correponding to $`D_1`$ and $`D_2`$, then $`(\tau ,\tau _1)`$ and $`(\tau ,\tau _2)`$ span cones in $`\mathrm{\Delta }_n`$ and $`(LC)=1`$, where $`C=V(\tau )`$.
###### Proof.
The “if” part is clear, since in this case we have $`(L(D_1D_2)V(\tau ))=1`$.
Suppose now that $`L(D_1D_2)`$ is not globally generated. We prove the proposition by induction on $`n`$. The case $`n=1`$ is trivial, and therefore we may assume $`n2`$. Let $`xX`$ be a fixed point in the base locus of $`L(D_1D_2)`$.
Suppose first that there is an invariant prime divisor $`DD_1`$, $`D_2`$ such that $`xD`$. We apply the induction hypothesis for the smooth toric variety $`D`$, the line bundle $`L|_D`$ and the prime invariant divisors $`DD_1`$ and $`DD_2`$. By Corollary 2.5 (iii) , the restriction map
$$H^0(L(D_1D_2))H^0(L(D_1D_2)|_D)$$
is surjective, so that our hypothesis on $`x`$ and $`D`$ implies that $`x`$ is in the base locus of $`L(D_1D_2)|_D`$. Lemma 4.4 implies that $`DD_1`$ and $`DD_2`$ are nonempty. If $`D=V(\tau _0)`$, then by induction we find a cone $`\tau ^{}`$ in the fan $`\mathrm{Star}(\tau _0)`$ of $`D`$. This corresponds to a cone $`\tau \mathrm{\Delta }`$ which satisfies the requirements of the proposition.
Therefore it remains to consider the case when, for every fixed point $`x`$ in the base locus of $`L(D_1D_2)`$ and every divisor $`D`$ containing $`x`$, we have $`D=D_1`$ or $`D=D_2`$. Clearly this implies $`n=2`$ and the fact that the base locus consists of a point, the corresponding cone being generated by the rays defining $`D_1`$ and $`D_2`$. But this contradicts Corollary 3.5 and the proof is complete. ∎
As a consequence of Proposition 4.3 we get that if $`(LC)2`$ for every integral invariant curve on $`X`$, then $`L(D)`$ is ample for every prime invariant divisor $`D`$. The next result makes this more precise by giving the condition for an ample line bundle $`L`$ and a prime invariant divisor $`D`$ to have $`L(D)`$ not ample.
###### Proposition 4.6.
Let $`X`$ be a complete smooth toric variety with $`dim(X)=n`$, $`L\mathrm{Pic}(X)`$ an ample line bundle and $`D=V(\tau _0)`$ a prime invariant divisor. Then $`L(D)`$ is not ample if and only if there is $`\tau \mathrm{\Delta }_{n1}`$ such that $`\tau ,\tau _0\mathrm{\Delta }_n`$ and $`(LV(\tau ))=1`$.
###### Proof.
It is clear that if there exists $`\tau `$ as above, then $`(L(D)V(\tau ))=0`$, so that $`L(D)`$ is not ample.
Suppose now that $`L(D)`$ is not ample and therefore there exists $`\tau ^{}\mathrm{\Delta }_{n1}`$ such that $`(L(D)V(\tau ^{}))0`$. Since $`L(D)`$ is globally generated by Lemma 4.4, we must have $`(L(D)V(\tau ^{}))=0`$.
We must have $`(DV(\tau ^{}))0`$, and therefore we deduce that either $`\tau _0,\tau ^{}\mathrm{\Delta }_n`$ or $`V(\tau ^{})D`$. In the first case, we have $`(LV(\tau ^{}))=1`$ and may take $`\tau =\tau ^{}`$.
If $`V(\tau ^{})D`$, we choose a divisor $`D_1=V(\tau _1)`$ such that $`\tau _1,\tau ^{}\mathrm{\Delta }_n`$. Then $`(L(DD_1)V(\tau ^{}))<0`$, and Proposition 4.5 implies that there is $`\tau \mathrm{\Delta }_{n1}`$ such that $`\tau _0,\tau \mathrm{\Delta }_n`$ and $`(LV(\tau ))=1`$. ∎
## 5. Sections of ample line bundles
In this section we fix a globally generated line bundle $`L`$ on a complete toric variety $`X`$ and an invariant divisor $`D`$ such that $`L𝒪(D)`$. Since $`L`$ is globally generated, for each maximal cone $`\sigma `$ there is a unique $`u_\sigma M`$ such that $`\mathrm{div}(\chi ^{u_\sigma })+D`$ is effective and zero on $`U_\sigma `$. Equivalently, for each maximal cone $`\sigma `$, there is a nonzero section $`s_\sigma H^0(X,L)`$, unique up to scalars, which is an eigenvector with respect to the torus action and whose restriction to $`U_\sigma `$ is everywhere nonzero.
A well-known ampleness criterion (see \[Fu\], Section 3.4) says that $`L`$ is ample if and only if $`u_\sigma u_\tau `$ (or, equivalently, $`ks_\sigma ks_\tau `$) for $`\sigma \tau `$. From the unicity of the sections $`s_\sigma `$, this is equivalent to the fact that if $`\sigma \tau `$, then $`s_\sigma |_{U_\tau }`$ vanishes at some point. But in that case, it must vanish at the unique fixed point $`x_\tau `$ of $`U_\tau `$.
We consider the following map whose components are given by the restriction maps:
$$\varphi :H^0(L)\underset{\sigma \mathrm{\Delta }_{\mathrm{max}}}{}H^0(L|_{\{x_\sigma \}}).$$
Since $`\varphi `$ is an equivariant map under the torus action, the discussion above shows that $`L`$ is ample if and only if $`\varphi `$ is surjective.
Our goal in this section is to extend this property of ample line bundles in the case when $`X`$ is smooth to a set of higher dimensional subvarieties which are pairwise disjoint. More precisely, we have the following
###### Theorem 5.1.
Let $`X`$ be a projective smooth toric variety and $`L\mathrm{Pic}(X)`$ an ample line bundle. If $`Y_1,\mathrm{},Y_rX`$ are integral invariant subvarieties such that $`Y_iY_j=\mathrm{}`$ for $`ij`$ and
$$\psi :H^0(L)\underset{i=1}{\overset{r}{}}H^0(L|_{Y_i})$$
is induced by restrictions, then $`\psi `$ is surjective.
###### Proof.
Let $`Y=_{i=1}^rY_i`$. In order to prove that
$$\psi :H^0(L)H^0(L|_Y)$$
is surjective, it is enough to prove that $`H^1(L_{Y/X})=0`$.
Let $`\pi :\stackrel{~}{X}X`$ be the blowing-up of $`X`$ along $`Y`$ and $`E`$ the exceptional divisor. Then
$$H^1(X,L_{Y/X})H^1(\stackrel{~}{X},\pi ^{}L𝒪(E)).$$
Since $`X`$ is smooth, the blowing-up of $`X`$ along an integral invariant subvariety is still a smooth toric variety (\[Ew\]). Since $`Y_iY_j=\mathrm{}`$ for $`ij`$, $`\pi `$ is a composition of such transformations, and therefore $`\stackrel{~}{X}`$ is a toric variety. Moreover, from the description in Ewald it follows that if $`E_i=\pi ^1(Y_i)`$, then $`E_i`$ is an invariant prime divisor on $`\stackrel{~}{X}`$ and $`E=_{i=1}^rE_i`$.
Since $`L`$ is ample, Proposition 7.10 in \[Ha\] implies that there is an integer $`s1`$ such that $`\pi ^{}(L^s)𝒪(E)`$ is ample on $`\stackrel{~}{X}`$. We choose an invariant divisor $`D`$ on $`\stackrel{~}{X}`$ such that $`\pi ^{}L𝒪(D)`$. Then $`D(1/s)E=D(1/s)_{i=1}^rE_i`$ is $``$-ample and $`D(1/s)E=DE`$. Now Corollary 2.5 gives $`H^1(X,\pi ^{}L𝒪(E))=0`$. ∎ |
warning/0001/hep-lat0001016.html | ar5iv | text | # REFERENCES
Axial Anomaly and Ginsparg-Wilson fermions in the Lattice Dirac Sea picture
Srinath Cheluvaraja
Dept. of Physics and Astronomy, Louisiana State University, Baton Rouge, LA, 70808
N.D. Hari Dass
Institute of Mathematical Sciences, Chennai, 600113
ABSTRACT
The axial anomaly equation in 1+1 dimensional QED is obtained on the lattice for fermions obeying the Ginsparg-Wilson relation. We make use of the properties of the Lattice Dirac sea to investigate the connection between the anomaly and the Ginsparg-Wilson operator in the Hamiltonian picture. The correct anomaly is reproduced for gauge fields whose characteristic time is much larger than the lattice spacing, which is the regime where the adiabatic approximation applies. A non-zero Wilson $`r`$ parameter is necessary to get the correct anomaly. The anomaly is shown to be independent of $`r`$ for $`r>0.5`$. The generalization to 3+1 dimensions is also discussed.
PACS numbers:12.38Gc,11.15Ha,05.70Fh,02.70g
The lattice regularization is one of the few non-perturbative methods available for defining quantum field theories. Lattice gauge theories have revealed many interesting features of gauge theories that are not easily visible in the usual perturbative approach. Nevertheless, the lattice regulator has proved problematic if fermions have to be incorporated into the theory. A naive discretization of the fermionic theory suffers from the replication of fermion modes due to the ”doublers”. The doublers are degenerate in energy with the originally introduced fermions and though they have lattice momenta of the order of the cut-off ($`1/a`$ in lattice theories, $`a`$ is the lattice spacing), they mimic ordinary low energy fermions. These doubler modes cannot be ignored as they participate in physical processes, for instance they can be pair created, and can affect the value of physical quantities– such as the free energy. The first method to handle these doublers was given in and it uses an additional term in the action –the Wilson term– to lift the degeneracy of the fermions, thereby decoupling the doublers in the continuum limit. However, this method has the disadvantage of explicitly breaking chiral symmetry, and hinders the study of dynamical questions related to chiral symmetry breaking. A cure for the doubling problem that explicitly breaks chiral symmetry also makes the lattice regularization of chiral gauge theories, such as the standard model, much more difficult. The Nielsen-Ninomiya no-go theorem decrees that any chirally symmetric lattice Hamiltonian satisfying general properties like locality and hermiticity must result in a replication of fermion species. This theorem seems to suggest the impossibility of defining undoubled fermions on the lattice without breaking chiral symmetry. Recently, however, alternative methods for tackling these problems have emerged. One of them uses the so called Ginsparg-Wilson relation for Dirac fermions.The Ginsparg-Wilson (G-W) operator is obtained by the application of block-spin transformations to a chirally invariant Dirac operator, and which therefore suffers from fermion doubling, using a chirally non-invariant blocking kernel. Although the G-W operator is not chirally invariant, it contains the information of chiral symmetry because it has been obtained after blocking a chirally invariant lattice action. Its construction, by a renormalization group transformation of a chirally invariant action, is bound to leave the low energy properties related to the chiral symmetry unchanged. This approach of formulating lattice fermions has led to many recent developments , such as, lattice formulations of chiral symmetry, the search for chiral gauge theories on the lattice, methods of defining a lattice topological charge, and formulation of lattice index theorems etc.
The G-W operator has to satisy the following relation
$$D\gamma _5+\gamma _5D=aD\gamma _5D.$$
(1)
(There are different versions of the G-W relation depending on the precise form of the blocking kernel used. The above form is one of the simpler ones and is sufficient for the ensuing discussion. Here $`a`$ is the lattice spacing. ) The G-W operator clearly does not satisfy chiral symmetry (because $`\{D,\gamma _5\}0`$). Even though the G-W operator seems to share the properties of a chirally noninvariant mass term, it is a milder way to break the chiral symmetry on the lattice. This is because it is obtained by blocking a chirally symmetric action, although using a chirally non-invariant kernel. The low energy properties of the G-W operator on the lattice are the same as those of the chirally symmetric action.
Another approach to problems of chirality on the lattice is the overlap approach introduced in . It captures many essential elements of domain wall fermions as well as the one in which requires an infinite number of auxiliary fields. The chiral determinant is expressed as an overlap of the ground states of two many body Hamiltonians and a construction of chiral gauge theories involves regularizing the overlap . Though the original Ginsparg-Wilson approach and the overlap approach appear to have nothing in common, the overlap operator (which appears in the Hamiltonian) has been shown to satisfy the Ginsparg-Wilson relation . The overlap is also a real time approach since it involves the quantum mechanical scalar product of the ground states of two Hamiltonians. This approach has led to many further studies of chirality on the lattice .
Another useful way of looking at the fermion doubling problem on the lattice is to look at the chiral anomaly structure of the lattice theory and to see what it yields in the continuum limit. A symmetry is said to be anomalous if it is no longer present in the quantum theory although it is present in the classical theory. The anomaly manifests itself by a non-conservation of a classically conserved charge. Anomalies are an inescapable part of some quantum field theories and have many important physical consequences. Their origin is related to the problem of regularizing amplitudes in quantum field theories while maintaining their invariances. In the path integral formulation of quantum field theory they arise because of the non-invariance of the measure of the path integral . As is well known 1+1 dimensional QED has a chiral anomaly when massless fermions are present. The anomaly arises because it is not possible to find a regularization of the gauge theory which maintains both the gauge invariance and chiral symmetry. A simple way of demonstrating the anomaly in 1+1 dimensions is by using a gauge invariant point-split definition of the axial vector current which can be seen not to be conserved. There are no anomalies in the naive latticisation because it is a gauge invariant regulator which also maintains chiral symmetry, but it is impossible to put only fields of one chirality on the lattice. This is consistent with the fact that no regularisation exists which simultaneously preserves gauge invariance and chiral symmetry for arbitrary matter content. The continuum limit of the lattice theory, on the other hand, must be able to reproduce the correct anomaly structure of a given theory. The naive fermionic lattice action coupled to gauge fields gives an anomaly free theory in the continuum limit because the anomalies are cancelled between the naive and doubled modes . It was shown in that the Wilson term reproduces the correct anomaly on the lattice provided the symmetry breaking parameter $`r0`$, the anomaly in the continuum limit being given by the co-efficient of the Wilson term. This is quite a surprising result because the Wilson term explicitly breaks the chiral symmetry but yet reproduces the anomaly which is essentially a quantum mechanical breakdown of the classical chiral symmetry.
As stressed by Nielsen and Ninomiya , and Peskin, the Hamiltonian formulation provides a much clearer physical picture of the anomaly in terms of the energy level shifting of the filled Dirac sea ( for a very clear exposition see also ). In this picture the anomaly arises because pairs of net chirality are pumped out of the infinitely filled Dirac sea. If one tries to transcribe this picture on the lattice, as was done by Ambjorn et al , one finds that the lattice Dirac sea is always finite and the anomaly always gets cancelled by the doubler modes in the absence of the Wilson term. The physical picture of the anomaly presented in can be applied on the lattice with a Wilson mass term. The role of the Wilson mass term is to suppress the contributions to the chiral charge coming from the doubler modes resulting in a non-zero anomaly on the lattice .
Our aim is to carry out a similar analysis for fermions satisfying the G-W relation and to see how a non-zero anomaly comes about on the lattice. This should complement the derivation of the anomaly from the Ginsparg-Wilson action in the Euclideanised formalism where the anomaly is showed to arise out of the measure . We mention here that the axial anomaly in 1+1 dimensional QED is also reproduced in the overlap formulation . Our derivation, apart from being quite different from the methods employed in , also highlights the role played by the Wilson term in giving the correct anomaly. The discussion will be in the Hamiltonian framework and we will derive the anomaly equation for the abelian theory in the 1+1 dimensions. We will then comment on the extension of this picture to 3+1 dimensions.
The Lagrangian density for a Dirac fermion in Minkowski space is given by
$$L(\psi ,\overline{\psi })=\overline{\psi }(i\gamma _\mu _\mu m)\psi ;$$
(2)
the gamma matrices satisfy $`\gamma _0^{}=\gamma _0,\gamma _i^{}=\gamma _i`$, and obey the relation $`\{\gamma _\mu ,\gamma _\nu \}=2g_{\mu \nu }`$. The metric $`g_{\mu \nu }`$ is $`diag(1,1,1,1)`$. $`\overline{\psi }(x)`$ denotes the relativistic adjoint $`\psi ^{}(x)\gamma _0`$. We shall use the Weyl representation for the gamma matrices. The Hamiltonian density is given by
$$H=\overline{\psi }(x)(i\gamma _i_i+m)\psi (x)$$
(3)
Before we discuss the GW fermions in the Hamiltonian picture, it is instructive to briefly review how the unwanted doublers are handled in the Euclidean formalism with the help of the Wilson mass term.In Euclidean space the Lagrangian density becomes
$$L(\psi ,\overline{\psi })=\overline{\psi }(i\gamma _\mu _\mu m)\psi ;$$
(4)
the Euclidean gamma matrices satisfy $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$, and $`\gamma _\mu ^{}=\gamma _\mu `$. The naive lattice discretization of the Dirac action for a massless fermion in Euclidean space is given by
$$S=\underset{xi}{}\frac{i}{2a}\overline{\psi }(x)\gamma _i(\psi (x+i)\psi (xi)).$$
(5)
In momentum space (in $`d`$ dimensions) this becomes
$$S=_{BZ}d^dk\overline{\psi }(k)(\underset{i=1}{\overset{d}{}}(\frac{\mathrm{sin}(k_ia)}{a})\gamma _i)\psi (k).$$
(6)
$`BZ`$ denotes the range of integration to be the $`d`$ dimensional Brillouin zone of the lattice. As is well known, the above discretization suffers from the presence of additional fermions at the corners of the Brillouin zone ($`2^d`$ in $`d`$ Euclidean dimensions) leading to $`2^d`$ fermions in the $`a0`$ limit. A method for eliminating the unwanted fermions is to give them very high masses in the continuum limit. The oldest way of achieving this is by adding a Wilson term to the massless action
$$S_w=\frac{r}{2a}\underset{xi}{}\overline{\psi }(x)(\psi (x+i)+\psi (xi)2\psi (x)).$$
(7)
(The Wilson term mimicks a mass term although in a more subtle way; the mass terms are ”momentum dependent”) This leads to the modified propagator (in momentum space)
$$D_w(k)=\underset{i=1}{\overset{d}{}}\gamma _i(\frac{\mathrm{sin}k_ia}{a})+m+\frac{r}{a}\underset{i=1}{\overset{d}{}}(1\mathrm{cos}(k_ia)).$$
(8)
In the above expression we have also introduced a bare mass term $`m`$. The momentum dependent mass terms ensure that the modes at the corners of the Brillouin zone have masses of the order of $`1/a`$ and decouple from the low energy effects (in the limit $`a0`$). The price paid for eliminating these doublers is the lack of chiral symmetry in the fermion action (the Wilson term explicitly breaks chiral symmetry ).
In order to define the real time evolution of Ginsparg-Wilson fermions in the Hamiltonian formulation, we have to first construct a Hamiltonian operator starting from the Euclidean functional integral. In the transfer matrix formalism this is done by choosing a particular axis as the time direction (with a lattice spacing $`\tau `$) and then taking the so called $`\tau `$ continuum limit. The $`\tau `$ continuum limit ($`\tau 0`$) is taken on an anisotropic lattice with different spacings in the space and the time directions. In the Euclidean formulation, the Ginsparg-Wilson operator can be formally understood to have arisen out of a chirally non-invariant blocking transformation. One could have used a $`\tau `$-continuum procedure to obtain the Hamiltonian version of the GW prescription. However, since we are eventually interested only in the Hamiltonian, it is sufficient to do the block spinning only in the spatial directions which already yields an (anisotropic) lattice with a blocked action. This blocked action has the same partition function as the original action, though spatial correlation lengths are halved on this lattice. We can then proceed to construct a Hamiltonian operator via the transfer matrix in the usual way for fermions . This way we are able to study the real time evolution of Ginsparg-Wilson fermions. As is well known, the passage from the transfer-matrix formalism to the Hamiltonian formalism can be carried out in more than one way, depending on the simplicity or complexity one desires. Our Hamiltonian is invariant under the symmetry discussed by Luscher in and is a valid starting point for carrying out calculations of wave functions, ground states, excited states, etc. The chirally invariant action on the lattice is
$$S=\underset{m,n}{}\overline{\varphi }_mh_{mn}\varphi _n,$$
(9)
where $`h`$ satisfies $`\{h,\gamma _5\}=0`$. We can block only the spatial degrees of freedom by using the kernel (defined in )
$$K(\psi ,\varphi )=\mathrm{exp}\xi \underset{mn}{}(\overline{\psi }(m)\overline{\chi }(m))(\psi (m)\chi (m)).$$
(10)
The only difference is that the blocked field $`\chi (m)`$ is defined as
$$\chi (m,t)=\eta \underset{n}{}\varphi (n,t),$$
(11)
where the above summation is only over fields defined over a hypercube at the same instant. This is the condition which ensures that only the spatial degrees of freedom are blocked. This leads to the following effective action for the fermions on the blocked lattice
$$\mathrm{exp}(\stackrel{~}{A}[\overline{\psi }\psi ])=𝑑\overline{\varphi }𝑑\varphi \mathrm{exp}[A(\overline{\varphi },\varphi )]K(\psi ,\varphi ).$$
(12)
The action on the blocked lattice is given by
$$\stackrel{~}{A}[\overline{\psi },\psi ]=\underset{m,n}{}\overline{\psi }(m)\stackrel{~}{h}_{mn}^{}\psi (n),$$
(13)
where the fields $`\psi `$ are defined on the blocked lattice which has twice the spacing (in the spatial direction) of the original lattice. The propagator of the blocked lattice satisfies the following relations
$`\{\stackrel{~}{h}_t,\gamma _5\}=0`$
$`\{\stackrel{~}{h}_s,\gamma _5\}=a\stackrel{~}{h}_s\gamma _5\stackrel{~}{h}_s.`$
After the usual passage to the Hamiltonian via a transfer matrix formalism the Hamiltonian one obtains is simply
$$H=\underset{x,y}{}\psi ^{}(x)\gamma _0\stackrel{~}{h}_s(x,y)\psi (y).$$
(14)
Henceforth, the operator $`\stackrel{~}{h}_s`$ will be called $`D`$ and it satisfies the G-W relation. Any $`\stackrel{~}{h}_s`$ satisfying the Ginsparg-Wilson relation can be used for defining our Hamiltonian. An explicit choice for $`D`$ satisfying the Ginsparg-Wilson relation is the overlap operator given by Neuberger‘s construction
$$D=\frac{1}{a}(1\frac{A}{\sqrt{A^{}A}});$$
(15)
A is defined in terms of the Wilson operator $`D_w`$ as
$$A=1aD_w.$$
(16)
It can be easily checked that the operator $`D`$ satisfies the G-W relation. It is worth adding here that any $`A`$ satisfying the following properties: $`A^{}=\gamma _5A\gamma _5`$ and $`A^{}A`$ commutes with $`A`$ will give a $`D`$ satisfying the GW relation. This may be useful in a more general context.
Though any $`D`$ satisfying the Ginsparg-Wilson relation can be used to construct our Hamiltonian, we have used the above explicit form proposed by Neuberger et al for our calculations. It should be stressed that apart from this choice, the considerations of this paper are independent of the overlap formalism. The above operator relations can be translated into momentum space and it is in momentum space that we will make most of our manipulations. It should be emphasised that in general, with external fields, momentum space description is not very economical.But in the $`1+1`$ dimensional abelian case with only an electric field, and in the $`3+1`$ dimensional case with uniform electric and magnetic fields, momentum space description is still useful.
The Hamiltonian of the lattice field theory is
$$H=\underset{x,y}{}\overline{\psi }(x)D(x,y)\psi (y).$$
(17)
$`\overline{\psi }`$ and $`\psi `$ can be interpreted as field operators in the usual sense. The wave equation for the fermion fields is
$$i\frac{\psi (x,t)}{t}=\underset{y}{}\gamma _0D(x,y)\psi (y).$$
(18)
Using the properties of $`D`$ and $`\gamma _0`$ it is easy to show that the Hamiltonian is hermitian, and therefore the evolution is unitary. We are basically interested in how the anomaly arises in this model. To get the anomaly we must of course couple the fermions to an external gauge field and then look for non-conservation of the chiral charge. It is well known that the anomaly can be extracted by treating the gauge fields as a classical variable and quantizing only the fermions. In order to be able to study the problem of fermions in an external field we will have to make some approximations which will be described shortly. The analysis will be presented in 1+1 dimensions.
In 1+1 dimensions the only effect of the external gauge field is to shift the momentum variable $`k`$ to $`kga(t)`$ where $`a(t)`$ is the time dependent component of the vector potential ($`A_1(x,t)=a(t)andA_0=0`$). This is true only in the Hamiltonian picture, even in the $`1+1`$ case, because in the Euclidean case, $`D^{}D=(kea(t))^2+\gamma _5`$, where $``$ is the electric field. We define the chiral charge operator on the lattice in the usual manner as
$$Q_n=a\underset{x}{}\psi ^{}(x,t)\gamma _5\psi (x,t).$$
(19)
All the operators are defined in the Heisenberg representation. The time dependence of the chiral charge is given by
$$\dot{Q}_n=\frac{Q_n}{t}+i[H,Q_n].$$
(20)
Since only the second term on the righthand side contributes to $`\dot{Q}_n`$ we have
$$\dot{Q}_n=i[H,Q_n].$$
(21)
Evaluating the commutator this becomes
$$\dot{Q}_n=i\underset{x,y}{}\overline{\psi }(x,t)\{D,\gamma _5\}\psi (x,t).$$
(22)
Using the Ginsparg-Wilson relation this becomes
$$\dot{Q}_n=ia\underset{x,y}{}\overline{\psi }(x,t)D\gamma _5D\psi (y,t).$$
(23)
The spatial indices of $`D`$ and all Dirac indices have been suppressed for ease in reading. Since $`D\gamma _5D`$ is non-zero we see that the chiral charge is not in general conserved, as expected. It remains to be seen if this non-conservation of the chiral charge reproduces the correct anomaly. For this the Dirac equation in Eq. 18 has to be solved in the presence of an external potential and the solutions have to be examined. We consider a time dependent potential which rises from zero to a constant value $`A_\tau `$ in a time $`\tau `$. $`\tau `$ is a time scale in the problem and two cases can be easily analyzed, the sudden limit $`\tau 0`$ and the adiabatic limit $`\tau \mathrm{}`$, as was also done in . First one writes the free field $`\psi (x,t)`$ as a superposition of positive and negative energy spinors
$$\psi (x,t)=_{BZ}\frac{1}{2\pi }[b(k)u(k)\mathrm{exp}(ikx)+d^{}(k)v(k)\mathrm{exp}(ikx)].$$
(24)
$`kx`$ is short for $`k_0x_0k_1x_1`$. $`E=k_0`$. As usual $`u(k)`$ and $`v(k)`$ represent positive and negative energy spinors. The operator $`b(k)`$ destroys an electron of momentum $`k_1`$ and the operator $`d^{}(k)`$ creates a positron of momentum $`k_1`$. Putting these spinors in the equation for the axial charge we have
$$|\dot{Q}_n|=_{bz}𝑑k\frac{1}{2\pi }\overline{v}(k,t)D(k,t)\gamma _5D(k,t)v(k,t)d(k)d^{}(k).$$
(25)
The angular brackets denote expectation values in the vacuum which is the state with zero electrons and positrons. Using the definition of $`D`$ and the properties of the gamma matrices it is easy to show that
$$D(k)\gamma _5D(k)=\gamma _5D^{}(k)D(k).$$
(26)
Now , $`D^{}(k)D(k)`$ is a c-number and acts trivially on the Dirac spinors. ($`D^{}D`$ is not a c-number in general.For the specific class where $`D`$ is of the form $`A_i\gamma _i+B`$ with $`A_i,B`$ commuting, this is true.Already in the $`3+1`$ dimensional case this is no longer true even for a uniform magnetic field.) In order to evaluate $`\overline{v}(k,t)\gamma _5v(k,t)`$ we have to determine the evolution of the negative energy spinor in the external field $`a(t)`$. Before we calculate this quantity it is instructive to calculate the same without any field. In the absence of an external field the positive and negative energy states evolve as
$`u(k,t)=\mathrm{exp}(iE(k)t)\stackrel{~}{u}(k)`$
$`v(k,t)=\mathrm{exp}(iE(k)t)\stackrel{~}{v}(k).`$
The spinors $`\stackrel{~}{u}(k)`$ and $`\stackrel{~}{v}(k)`$ satisfy the time independent Schroedinger equations with positive and negative energies.
$`\gamma _0D(k)\stackrel{~}{u}(k)=E(k)\stackrel{~}{u}(k)`$
$`\gamma _0D(k)\stackrel{~}{v}(k)=E(k)\stackrel{~}{v}(k).`$
The Weyl representation for the $`2`$ dimensional $`\gamma `$ matrices is
$$\gamma _1=i\gamma _5=i\sigma _3\gamma _0=\sigma _1\gamma _5=\sigma _2.$$
(27)
$`\sigma _1`$ and $`\sigma _3`$ are the Pauli matrices. The eigen values of the spinors are given by $`E(k)^2=D^{}(k)D(k)`$. Using the definition of $`D(k)`$ we get $`D^{}(k)D(k)`$ to be a c-number. $`D(k)`$ can be written as
$$D(k)=\frac{1}{a}(1\frac{(1am(k))}{\sqrt{(}A^2(k))})+\frac{1}{a}(\frac{\mathrm{sin}(k_1a)}{\sqrt{(}A^2(k))})\gamma _1.$$
(28)
If we now write $`D(k)`$ as
$$D(k)=g(k,t)+\gamma _1f(k,t),$$
(29)
then $`E^2(k)=D^{}(k)D(k)=f^2(k,t)+g^2(k,t)`$. $`f`$ and $`g`$ are functions given by
$`g(k)={\displaystyle \frac{1}{a}}(1{\displaystyle \frac{(1am(k))}{\sqrt{(}A^2(k))}})`$
$`f(k)={\displaystyle \frac{1}{a}}({\displaystyle \frac{\mathrm{sin}(k_1a)}{\sqrt{(}A^2(k))}})`$
$`.`$
The function $`f`$ is an odd function of $`k`$ whereas the function $`g`$ is an even function of $`k`$. The time independent spinors can be normalized to satisfy
$`\overline{\stackrel{~}{u}(k)}\gamma _5\stackrel{~}{u}(k)=0`$
$`\overline{\stackrel{~}{v}(k)}\gamma _5\stackrel{~}{v}(k)=0`$
$`\overline{\stackrel{~}{u}(k)}\gamma _5\stackrel{~}{v}(k)=1.`$
This means that $`\dot{Q}_n=0`$ in the absence of an external field. This is as expected, there is no anomaly in zero external field even though the Hamiltonian is not chirally invariant.
In the presence of an external field the only change in the operator $`D`$ is a replacement of $`k`$ by $`kga(t)`$. Since the structure of $`D`$ is not affected by an external field the Ginsparg-Wilson relation is still satisfied. Although the evolution of Dirac spinors in an arbitrary external field can only be analyzed numerically, two limiting cases admit a simpler analysis. These are the adiabatic limit and the sudden limit, and we shall examine these two cases separately. When an external field is turned on slowly (compared to the time scales in the system) we can use the adiabatic approximation. The sudden approximation is useful when the field is turned on faster than the fastest time scale in the system. The rate at which the field is turned on can be controlled by introducing a parameter $`\tau `$ defined as follows
$$A(t)=0t0$$
(30)
$$A(t)=At>\tau .$$
(31)
The precise form of $`A(t)`$ for $`0<t<\tau `$ is not very important. The sudden limit corresponds to $`\tau 0`$ and the adiabatic limit corresponds to $`\tau \mathrm{}`$. In the adiabatic approximation the form of the positive and negative energy spinors for times $`0<t<\tau `$ is given by
$$u(k,t)=a(k,t)\xi ^{}(t)u^0(k,t)+b(k,t)\xi (t)v^0(k,t)$$
(32)
$$v(k,t)=c(k,t)\xi ^{}(t)u^0(k,t)+d(k,t)\xi (t)v^0(k,t).$$
(33)
The above equation is written for the individual fourier components of the positive and negative energy spinors in an external field, $`u^0(k,t)`$ and $`v^0(k,t)`$ are the positive and negative energy spinors satisfied by the instanteneous Schrodinger equation at the instant $`t`$. We closely follow the notation of and we have also corrected some of the misprints which occur therein. $`\xi (t)^{}`$ and $`\xi (t)`$ are the phase factors for the positive and negative energy spinors. The phase factor $`\xi (t)`$ is given by
$$\xi (t)=\mathrm{exp}(i_0^tE(k,t^{})dt^{}.$$
(34)
$`a,b,c,andd`$ are called Bogoulobov coefficients. The Bogoulobov co-efficients are chosen to satisfy the following boundary conditions
$$a(0)=d(0)=1b(0)=c(0)=0.$$
(35)
These boundary conditions ensure that we are looking at the evolution of the positive and the negative energy states before the external field is switched on. As mentioned before, $`u^0(k,t)`$ and $`v^0(k,t)`$ are positive and negative energy spinors having momentum $`k`$ and $`k`$ respectively, and they satisfy the instanteneous Schrodinger equations given by
$`\gamma _0D(k)u^0(k,t)=E(k,t)u^0(k,t)`$
$`\gamma _0D(k)v^0(k,t)=E(k,t)v^0(k,t).`$
Substituting the expressions in Eq. 32 and Eq. 33 in the wave equation for the fermions and using the previously mentioned boundary conditions we get
$$c(k,t)=_0^t\alpha (k,t^{})d(k,t^{})\xi ^2(t^{})𝑑t^{}$$
(36)
$$d(k,t)=1_0^t\alpha (k,t^{})c(k,t^{})\xi _{}^{}{}_{}{}^{2}(t^{})𝑑t^{}$$
(37)
$$b(k,t)=c^{}(k,t)$$
(38)
$$a(k,t)=d^{}(k,t).$$
(39)
The quantity $`\alpha (k,t)`$ is defined by
$$\dot{u}^0(k,t)=\alpha (k,t)v^0(k,t),$$
(40)
and is
$$\alpha (k,t)=\frac{(g\dot{f}f\dot{g})}{2E^2(k,t)}.$$
(41)
After using the stated normalizations of the spinors, $`\overline{v}(k,t)\gamma _5v(k,t)`$ is given
$$\overline{v}(k,t)\gamma _5v(k,t)=c^{}d\xi ^2c.c.$$
(42)
The co-efficients $`c(k,t)`$ and $`d(k,t)`$ can be approximated by (for small values of $`\alpha (k,t)`$)
$$c(k,t)=i\frac{\alpha (k,t)}{2E(k,t)}\xi ^2(t)d(k,t)=1+O(\alpha ^2).$$
(43)
In the adiabatic approximation the quantities inside the integrand on the right hand side of Eq. 36 and Eq. 37 are evaluated at $`t=0`$. Substituting for the values of $`f(k,t)`$ and $`g(k,t)`$ and using the relation
$$D(k)\gamma _5D(k)=\gamma _5D^{}(k)D(k)=\gamma _5E^2(k)$$
(44)
the r.h.s of Eq. 25 becomes
$$_{BZ}\frac{1}{2\pi }𝑑kC(k)(aag\dot{a}(t))$$
(45)
where $`C(k)`$, a complicated expression, is given in the appendix along with the expressions for $`f,g,\dot{f},\dot{g}`$. The function $`C(k)`$ can be plotted and the integral of $`C(k)`$ over the Brillouin zone can be estimated numerically. The function $`C(k)`$ depends on $`r,k,a,m`$. We plot $`C(k)`$ as a function of $`k`$ in Fig. 1 to Fig. 5. We first plot it for zero mass and then for a non-zero value ($`m=5`$). The eigenvalues of the spinors (in a zero external field) are also plotted in Fig. 6 to Fig. 10. The first thing we observe is that when $`r=0`$ the eigenvalue spectrum does not distinguish between the modes at $`k=0`$ and $`k=\pi /a`$, and the integral of $`C(k)`$ over the Brillouin zone is zero. For $`r0`$ the modes at $`k=0`$ and $`k=\pi /a`$ have different energies and an asymmetry develops in the function $`C(k)`$. For $`r=1`$, the integral of $`C(k)`$ over the Brillouin Zone has the value $`2`$. Substituting this in Eq. 25 we get the anomaly equation
$$\dot{Q}_n=\frac{g}{2\pi }d^2xϵ_{\mu \nu }F_{\mu \nu }$$
(46)
in the continuum limit.
We have studied the function $`C(k)`$ for different values of $`r`$ and we find that the anomaly is independent of $`r`$ for large $`r`$ but vanishes for a smaller and, in particular, a zero value of $`r`$. The value $`r=0.5`$ seems to separate the region with and without the anomaly. A non-zero Wilson $`r`$ parameter is necessary to get the anomaly in the continuum limit. This means that a Neuberger like operator for $`D`$ where the naive Dirac operator is used in place of $`D_w`$ in Eq. 16 will not reproduce the correct anomaly in the continuum limit inspite of satisfying the Ginsparg-Wilson relation (the operator $`D`$ with $`D_w`$ replaced by $`D_{naive}`$ also satisfies the Ginsparg-Wilson relation). The case $`m0`$ can also be analyzed along the same lines and it turns out that the integral of $`C(k)`$ is very small, consistent with zero. It appears that in this case we have an exact cancellation of the bare mass term with the anomaly term to give a zero rate of change of chiral charge. Nevertheless, the anomaly term is still present and so is the mass term, but the two appear with opposite signs. The adiabatic approximation is justified when the switching time $`\tau `$ is much greater than the characteristic time periods of the system. In our example $`2\pi /E(k)`$ is the characteristic time period of the system and the adiabatic approximation is justified when $`\tau >>a`$. To summarize, a zero bare mass term with a non-zero $`r`$ parameter gives an anomaly independent of $`r`$ for $`r>0.5`$. When a bare mass term is included we have a cancellation of the anomaly term with the mass term though both terms are still present.
The rate of change of the chiral charge can also be calculated in the sudden approximation in which an external field is turned on infinitely fast. This approximation corresponds to the limit $`\tau 0`$. In this limit the spinors $`u^0(k,t),v^0(k,t)`$ are unchanged immediately after the field is turned on and only their evolution is governed by the new Hamiltonian (with a constant field). The normalization of the spinors in a constant external field (with value $`A_\tau `$) can be made just as in the zero field case and there is no anomaly in this limit. This limit corresponds to the case $`\tau <<a`$.
The time $`\tau `$ is a characteristic time associated with the gauge fields on the lattice and our calculation clearly shows that in order to get the correct anomaly we have to ensure that we are not in the regime of the sudden approximation. If we are in the intermediate region we will see a crossover from one limit to another limit. So far we have only studied fermions and (abelian) gauge fields in 1+1 dimensions. It was pointed out in that the anomaly in 3+1 dimensions factorizes into a 1+1 dimensional part and an extra factor coming from the additional dimensions. We briefly review the argument in . To get the anomaly in 3+1 dimensions we first turn on a magnetic field in the, say z, direction. This leads to the usual Landau levels for the fermions which are labelled by integers. We then turn on an electric field paralell to the magnetic field. The important point is that the fermions in the lowest Landau level in the presence of this electric field behave like fermions in the 1+1 dimensional case that we have just analyzed. Hence the same 1+1 dimensional anomaly is present but with an additional degeneracy factor coming from the Landau levels. As shown in the argument goes through for the lattice Dirac sea case for the case of an uniform magnetic field. The degeneracy factor is a geometrical quantity that is in general dependent on the details of the lattice Hamiltonian which is more complicated for Ginsparg-Wilson fermions. However, when the Ginsparg-Wilson operator is constructed as $`aD_{GW}=(1\frac{A}{\sqrt{A^{}A}})`$ with $`A=1aD_W`$, the degeneracies of $`D_{GW}`$ and $`D_W`$ are the same, and in the limit of zero lattice spacing the degeneracy is just
$$L_1L_2gH/(2\pi ),$$
(47)
the number of states in the square $`L_1L_2`$ perpendicular to the magnetic field ($`H`$). The above factor simply multiplies the $`1+1`$ dimensional anomaly and gives the correct anomaly in 3+1 dimensions.
The main aim of this note was to show that a Hamiltonian analysis of Ginsparg-Wilson fermions leads to a non-zero rate of chiral charge and gives the anomaly equation in the continuum limit. The doubler modes are suppressed by the Wilson parameter, infact the Wilson parameter $`r`$ plays a crucial role in yielding the correct anomaly. A quantum mechanical analysis supplemented by an adiabatic approximation was necessary to get the anomaly. It is noteworthy that if we are not in the adiabatic regime we will get other contributions ($`\dot{a}(t)`$ and higher time derivatives) to the chiral charge and this will not reproduce the anomaly equation. It may be useful to compare our derivation with that of the overlap method. In the overlap method the anomaly is extracted by looking at the scalar product of the ground states of two different many body Hamiltonians, whereas in our approach we study the dynamical picture behind the anomaly by using the properties of the Dirac sea in an external electric field in the adiabatic limit.
Appendix
In this appendix we collect together some expressions which are necessary to get the function $`C(k)`$.
$$f(k)=\frac{1\frac{1a\left(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a}\right)}{\sqrt{12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2}}}{a}$$
(48)
$$g(k)=\frac{\mathrm{sin}(ak)}{a\sqrt{12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2}}$$
(49)
$`\dot{f}(k)={\displaystyle \frac{\left(1m\frac{r\left(1\mathrm{cos}(ak)\right)}{a}\right)\left(2r\mathrm{sin}(ak)+2ar\left(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a}\right)\mathrm{sin}(ak)+2\mathrm{cos}(ak)\mathrm{sin}(ak)\right)}{2a(12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2)^{\frac{3}{2}}}}+`$
$`{\displaystyle \frac{r\mathrm{sin}(ak)}{a(12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2)}}`$
$`\dot{g}(k)={\displaystyle \frac{\mathrm{sin}(ak)}{2a(12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+\mathrm{sin}(ak)^2)^{\frac{3}{2}}}}`$
$`+{\displaystyle \frac{\mathrm{cos}(ak)}{a(12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+\mathrm{sin}(ak)^2)}}`$
$`E^2(k)={\displaystyle \frac{\mathrm{sin}(ak)^2}{a^2(12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2)}}+`$
$`{\displaystyle \frac{\left(1\frac{1a\left(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a}\right)}{\sqrt{12a(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})+a^2(m+\frac{r\left(1\mathrm{cos}(ak)\right)}{a})^2+\mathrm{sin}(ak)^2}}\right)^2}{a^2}}`$
$$C(k)=(1/(2E(k))(g(k)\dot{f}(k)f(k)\dot{g}(k))$$
(50) |
warning/0001/nucl-th0001037.html | ar5iv | text | # From bound states to resonances: analytic continuation of the wave function
## Abstract
Single–particle resonance parameters and wave functions in spherical and deformed nuclei are determined through analytic continuation in the potential strength. In this method, the analyticity of the eigenvalues and eigenfunctions of the Schrödinger equation with respect to the coupling strength is exploited to analytically continue the bound–state solutions into the positive–energy region by means of Padé approximants of the second kind. The method is here applied to single–particle wave functions of the $`{}_{}{}^{154}\mathrm{Sm}`$ and $`{}_{}{}^{131}\mathrm{Eu}`$ nuclei. A comparison of the results with the direct solution of the Schrödinger equation shows that the method can be confidently applied also in coupled–channel situations requiring high numerical accuracy.
In recent years there has been an increasing experimental activity on nuclei far from the stability line. These nuclear systems are usually unbound or weakly bound, and often exhibit resonances with a pronounced single–particle character, since the Fermi level is close to or even immersed in the continuum. This phenomenology has prompted people to look for more and more efficient methods for the description of unbound states, which could compete with the techniques presently available for the discrete part of nuclear spectra. A possible approach, which proved to be rather successful, exploits the fact that resonances can be described by wave functions with purely outgoing behavior and complex eigenvalues $`E_Ri\mathrm{\Gamma }/2`$ (Gamow states) . One is then confronted with the solution of a complex eigenvalue problem for a non–Hermitian Hamiltonian. Even if Gamow states are characterized by the at first sight unpleasant feature of an exponentially growing oscillatory behavior at large distances, they can be however normalized by a suitable generalization of the quantum–mechanical inner product. This can be achieved in several, essentially equivalent ways, either introducing convergence factors in the integrals expressing their norm , or by analytic continuation from the upper half of the $`k`$–plane to the resonance poles . These recipes can be implemented by solving the Schrödinger equation along a deformed contour in the complex $`r`$–plane , which amounts to an analytic continuation $`𝐫exp(i\theta )𝐫`$ in configuration space. This is very convenient under a computational point of view, resonance states being transformed into square–integrable wave functions, while leaving untouched the corresponding poles of the S–matrix in the energy plane. If this is done for a value of $`r`$ so large that the effective nucleon–core interaction can be neglected (exterior complex scaling), one has the advantage that the potential itself remains unchanged in the contour deformation. Complex–coordinate rotation methods have found interesting applications both in nuclear and in atomic physics.
Recently, unbound states in exotic nuclei have been studied by means of Analytic Continuation in the Coupling Constant (ACCC) . This method, proposed already several years ago by Kukulin and co–workers , starts from the intuitive expectation that, for an attractive potential, a resonance state will become a bound state as the coupling strength is increased. Under a mathematical point of view, one can prove that the wave number $`k=\sqrt{2\mu E}/\mathrm{}`$ is an analytic function of the strength $`\lambda `$, with the same restrictions on the potential which guarantee the analyticity of the Jost function . Here and in the following $`\mu `$ denotes the reduced mass of the nucleon–core system. Near the value $`\lambda _0`$ for which $`k(\lambda _0)=0`$ (i.e. near the scattering threshold), one has
$$k(\lambda )i\sqrt{\lambda \lambda _0},k(\lambda )i(\lambda \lambda _0),$$
(1)
for $`l0`$, and for $`l=0`$, respectively. These properties suggest the analytic continuation of $`k`$ in the complex $`\lambda `$–plane from the bound–state region into the resonance region through the employment of Padé approximants of the second kind
$$kk^{(N,M)}(x)=i\frac{c_0+c_1x+c_2x^2+\mathrm{}+c_Mx^M}{1+d_1x+d_2x^2+\mathrm{}+d_Nx^N},$$
(2)
where $`x\sqrt{\lambda \lambda _0}`$. In practice, the following procedure can be followed to find the resonance parameters for an interaction $`V`$, when $`l0`$. One endows $`V`$ with a strength parameter $`\lambda `$, $`V\lambda V`$, and solves the bound–state problem for $`\lambda V`$ in correspondence to $`N+M+1`$ different values $`\lambda _i`$ of the coupling strength. Given the threshold value $`\lambda _0`$, the $`N+M+1`$ coefficients in the Padé approximant (2) can be determined by equating $`k^{(N,M)}(x_i)`$ to the actual values $`k_i`$ of the wave number. The approximant can then be used to estimate the resonance wave number $`k_r`$, and hence the resonance position and width, in correspondence to the “physical” value $`\lambda =1`$ of the potential strength. This procedure can be easily modified if the potential $`V`$ supports an $`s`$–wave resonance, in which case the corresponding pole leaves the negative imaginary axis at $`k(\overline{\lambda }_0)=i\overline{\chi }_0`$ . The method has been applied with some success to unbound states in <sup>5</sup>He and <sup>5</sup>Li, as well as in three–cluster nuclei . An extrapolation procedure similar the one described above can be used in the complex $`k`$–plane to analytically continue the bound–state wave function $`\psi _l^{(B)}(kr)`$ into the scattering region for any value of the radial variable $`r`$.
In this paper we study the application of this technique to single–particle resonances in deformed nuclei, a situation quite common in the drip–line region. The numerical evaluation of resonance (Gamow) states for nonspherical nuclei is a major challenge, which has been only recently solved . Bound–state wave functions, on the other hand, can be calculated nowadays through very efficient and quick algorithms . Here, we shall enquire whether the solution of the coupled–channel bound–state problem, obtained for a set of values of the coupling strength $`\lambda `$, can be analytically continued into the unbound region. We shall consider the Padé extrapolation both for the resonance parameters and for the wave function. The latter case is particularly interesting for proton decay, since one has to deal with lifetimes $`\tau 1\mu s`$, which implies resonance widths $`\mathrm{\Gamma }=\mathrm{}/\tau `$ smaller than $`10^{16}`$ MeV. Such small widths are difficult to obtain with enough accuracy starting directly from the energy eigenvalue, and can be best estimated from the wave–function behavior .
To test the validity of the ACCC method in the present case, we have solved exactly the problem both for bound–state energies and in the continuum, and we have compared the results with the outcome of the Padé extrapolation obtained starting from the bound–state region, for different, decreasing values of the coupling strength. The radial Schrödinger equation with outgoing–wave boundary conditions has been solved in the standard way, namely starting from the origin and from the outer region, and matching the logarithmic derivative of these functions at some radius $`R`$. For real, negative energies one gets normalized wave functions, with the proper exponentially decreasing tail, in correspondence to the energy eigenvalues of the bound system; in the scattering case, on the other hand, one gets purely outgoing states for complex eigenvalues $`E_Ri\mathrm{\Gamma }/2`$ . The normalization of our unbound wave functions agrees with the Zel’dovich or Gyarmati–Vertse prescriptions for Gamow states.
As a first step, let us consider single–particle resonances in spherical nuclei. In Fig. 1 we report the results of our calculations for the $`f_{5/2}`$ neutron state in the nucleus $`{}_{}{}^{154}\mathrm{Sm}`$. The neutron is assumed to move in a spherical Saxon–Woods potential, supporting a bound state for $`\lambda =1`$. Fig. 1(a) exhibits the real and imaginary part of the energy $`E`$, as a function of the decreasing strength parameter $`\lambda `$. The full dots represent the outcome of the numerical solution of the Schrödinger equation, whereas the crosses correspond to the bound–state energies used for the Padé extrapolation into the scattering region, the extrapolated results being given by the full lines. As it can be seen, the neutron is bound for $`\lambda `$ decreasing from 1 down to $`\lambda _00.84`$, where the neutron state becomes unbound and the energy acquires a non–vanishing imaginary part. With a (7,7) approximant of the form (2), the agreement between the exact and the extrapolated results is quite good both for $`Re(E)`$ and $`Im(E)`$ in the whole considered region of the strength–parameter values. In Figs. 1(b) and 1(c) we show the calculated (full dots) and extrapolated values of the real and imaginary part of $`r\psi _{lj}(r)`$ at $`r=7`$ fm and $`r=15`$ fm, respectively. The results are plotted as functions of the real part $`Re(E)`$ of the corresponding eigenvalue. As in Fig. 1(a), the crosses denote the bound–state points used as input to the extrapolation procedure. Note that at $`r7`$ fm the neutron is feeling the strongest effect from the nuclear potential, and the wave function is attaining its largest value; for $`r15`$ fm, on the other hand, one is far away enough from the nuclear core, to have an indication of the quality of the extrapolation in reproducing the tail of the wave function. This is crucial in order to obtain an accurate evaluation of the resonance width. Indeed, in the spherical case the partial width for the decay to the channel $`lj`$ can be related to the value of the wave function $`\psi _{lj}(r)`$ at $`r=R`$ by
$$\mathrm{\Gamma }_{lj}=\frac{\mathrm{}^2k}{\mu }\frac{R^2|\psi _{lj}(R)|^2}{|G_l(R)+iF_l(R)|^2},$$
(3)
where $`F_l`$ and $`G_l`$ represent the free regular and irregular radial wave functions, respectively. If $`R`$ is large enough to be outside of the range of the potential, Eq. 3 provides a width which is actually independent from the value of $`R`$, as it should be . These considerations can be extended to the coupled–channel case . As Figs. 1(b) and 1(c) clearly exhibit, the results of the (7,7) extrapolation agree very well with the numerical calculation both in the inner and in the tail region. In particular, the non–trivial behavior of the real part of the inner wave function in the threshold region is accurately reproduced by the extrapolation, an indication that the Padé analytic continuation performs well in transferring the information from the bound–state into the scattering region. Similarly, the steep increase of the wave–function tail when going to positive energies is successfully reproduced by the Padé approximation. We have also considered the behavior of the Padé extrapolation for lower–rank approximants, and verified that the analytic continuation still compares well with the exact results when a (4,4) Padé approximant is employed.
Similar calculations have been done for a proton state in the $`{}_{}{}^{154}\mathrm{Sm}`$ nucleus, with the proton moving in a $`h_{9/2}`$ state. In this case, when solving the Schrödinger equation, the inner wave function has to be matched to the radial Coulomb wave function $`r\psi _{lj}^{out}(r)=N_{lj}[G_l^{(C)}(r)+iF_l^{(C)}(r)]`$, where $`G_l^{(C)}`$ and $`F_l^{(C)}`$ are the usual regular and irregular Coulomb functions and $`N_{lj}`$ is a normalization factor. The outcome of the calculations is reported in Fig. 2. As for the neutron case, Fig. 2(a) displays the real part of the energy eigenvalue for decreasing $`\lambda `$, whereas Figs. 2(b) and 2(c) refer to the real part of $`r\psi _{lj}(r)`$ at $`r=5`$ fm and $`r=15`$ fm, respectively, the symbols having the same meaning as in Fig. 1. The imaginary parts of both the energy and the wave function remain vanishingly small in passing from the bound–state into the scattering region, so that they are not even reported in the figures. The proton state becomes unbound at $`\lambda 1`$, with $`Re(E)`$ scaling linearly with the strength parameter. Similarly, the wave function exhibits a smooth behavior when the scattering threshold is crossed. It is not surprising, therefore, that a low–rank (4,4) Padé approximant can provide an extremely good reproduction of the exact results. Note that, because of the Coulomb barrier, the proton wave function in the exterior region is much smaller than in the neutron case.
The ACCC method has been finally tested for proton resonances in deformed nuclei. As is well–known, such a calculation is much more challenging than in the spherical case, since it requires the solution of a coupled–channel problem to determine the intrinsic, single–particle wave functions . In Fig. 3 (left panel) we report the real part of the energy of the $`3/2+`$ state in $`{}_{}{}^{131}\mathrm{Eu}`$ as a function of $`\lambda `$. Improved experimental data have been recently given for this highly deformed proton emitter, leading to the identification of a fine structure splitting in the radioactive decay from the ground state . The process can be successfully described in the framework of a model where the proton is emitted from a deformed single–particle Nilsson level . The energy of the $`3/2+`$ state turns out to be a linear function of $`\lambda `$, much as in the spherical situation, the imaginary part remaining always very close to zero. The relevant wave–function components at $`r=15`$ fm are given in the right panel of Fig. 3. For the proton decay of the $`{}_{}{}^{131}\mathrm{Eu}`$ nucleus from the $`K=3/2+`$ state they are the $`d3/2`$ and $`d5/2`$ components of the $`3/2+`$ Nilsson orbital (higher–$`l`$ partial waves giving a negligible contribution to the decay width ). The former determines the $`{}_{}{}^{131}\mathrm{Eu}`$ proton decay to the $`{}_{}{}^{130}\mathrm{Sm}`$ ground state, while the latter is by far the dominant term in the decay to the $`2^+`$ excited state of the daughter nucleus . The bound–state solution has been analytically continued into the scattering region by means of a simple (3,3) Padé approximant. One can see that the extrapolation reproduces the outcome of the numerical solution extremely well both for the energy and for the wave–function components. Note that, in passing from the bound–state into the positive–energy region the imaginary part of the wave functions is again so small, that it is not given in the figure.
In summary, we have applied a Padé extrapolation to determine the resonance parameters and wave function for single–particle resonances in nuclei, starting from bound–state calculations. This has been obtained by varying the coupling strength so as to analytically continue the bound–state energy eigenvalue and wave function into the positive–energy region. The method has been applied to the single–particle decay of both spherical and deformed nuclei, which entails the solution of a single– or a coupled–channel problem, respectively. A comparison of the extrapolated results with the outcome of the direct solution of the Schrödinger equation with the proper boundary conditions shows that the ACCC method can be confidently applied to these situations, where high numerical accuracy is required in order to have a meaningful comparison with the experimental data. |
warning/0001/gr-qc0001059.html | ar5iv | text | # Untitled Document
Unification of Weak and Gravitational Interactions Stemming from Expansive Nondecelerative Universe Model
Miroslav Sukenik<sup>a</sup>, Jozef Sima<sup>a</sup> and Julius Vanko<sup>b</sup>
<sup>a</sup>Slovak Technical University, Dep. Inorg. Chem., Radlinskeho 9, 812 37 Bratislava, Slovakia
<sup>b</sup>Comenius University, Dep. Nucl. Physics, Mlynska dolina F1, 842 48 Bratislava, Slovakia
Abstract. There is a deep interrelationship of the General Theory of Relativity and weak interactions in the model of Expansive Nondecelerative Universe. This fact allows an independent determination of the mass of vector bosons Z and W, as well as the time of separation of electromagnetic and weak interactions.
In the early stage of the Universe creation, i.e. in the lepton era an equilibrium of protons and neutrons formation existed at the temperature about 10<sup>9</sup> \- 10<sup>10</sup> K which corresponds to the energy of 1 MeV (a lower side of the range of weak interaction energies). The amount of neutrons was stabilized by processes including antineutrinos (eq.1) or neutrinos (eq.2) such as
$`\nu +p^+n+e^+`$ (1)
$`e^{}+p^+n+\nu `$ (2)
The cross section $`\sigma `$related to the above processes can be expressed as
$`\sigma \frac{g_F^2.E_w^2}{(\mathrm{}.c)^4}`$ (3)
where $`g_F`$ is the Fermi constant (10<sup>-62</sup> J.m<sup>3</sup>), $`E_w`$is the energy of weak interactions that, based on (3), can be formulated by relation
$`E_w\frac{r.\mathrm{}^2.c^2}{g_F}`$ (4)
where $`r`$ represents the effective range of weak interactions. Stemming from relation (4) it holds that in limiting case when
$`r=\frac{\mathrm{}}{m_{ZW}.c}`$ (5)
the maximum energy of weak interaction is given by
$`E_wm_{ZW}.c^2`$ (6)
Relations (5) and (6) represent the Compton wavelength of the vector bosons Z and W, and their energy, respectively. Equations (4), (5) and (6) lead to expression for the mass of the bosons Z and W
$`m_{ZW}^2\frac{\mathrm{}^3}{g_F.c}\left|100GeV\right|^2`$ (7)
providing the value that is in good agreement with the known actual value.
For the density of gravitational energy $`\epsilon _g`$ it follows from the ENU model \[3 - 5\]
$`ϵ_g=\frac{R.c^4}{8\pi .G}=\frac{3m.c^2}{4\pi .a.r^2}`$ (8)
where $`\epsilon _g`$ is the gravitational energy density of a body with the mass $`m`$ in the distance $`r,R`$ is the vector curvature, $`a`$ is the gauge factor that reaches at present
$`a10^{26}m`$ (9)
As a starting point for unifying the gravitational and weak interactions, the conditions in which the weak interaction energy $`E_w`$ and the gravitational energy $`E_g`$of a hypothetic black hole are of identical value
$`E_w=\left|E_g\right|`$ (10)
can be chosen. Based on relations (4), (8) and (10) in such a case it holds
$`\frac{r.\mathrm{}^2.c^2}{g_F}=\left|ϵ_g𝑑V\right|\frac{m_{BH}.c^2.r}{a}`$ (11)
where $`m_{BH}`$ is the mass of a black hole and $`r`$ is the range of weak interaction. It follows from (11) that
$`m_{BH}\frac{a.\mathrm{}^2}{g_F}`$ (12)
The above relation manifests that the mass of a black hole depends on the gauge factor, i.e. it is increasing with time. On the other hand, the black hole mass may not be lower than the Planck mass $`m_{Pc}`$
$`m_{BH}m_{Pc}`$ (13)
that approximates
$`m_{Pc}=\left(\frac{\mathrm{}.c}{G}\right)^{1/2}10^{19}GeV`$ (14)
The gravitational radius $`l_{Pc}`$ of a black hole having the minimum mass $`m_{Pc}`$ is
$`l_{Pc}=\left(\frac{G.\mathrm{}}{c^3}\right)^{1/2}10^{35}m`$ (15)
If there is a mutual relationship of the gravitational and weak interactions, there had to be a time $`t_x`$ corresponding to a gauge factor $`a_x`$ when $`m_{BH}`$ and $`m_{Pc}`$ were of identical value. It was the time of weak interactions formation. In such a case it stems from (12) and (13) that
$`a_x\frac{m_{Pc}.g_F}{\mathrm{}^2}10^2m`$ (16)
and
$`t_x10^{10}s`$ (17)
This is actually the time when, in accordance with the current knowledge, electromagnetic and weak interactions were separated (it might be worth mentioning that its value represents also typical duration of weak interaction processes). In the time $`t_x`$ it had to hold
$`\frac{m_{Pc}}{m_{ZW}}=\left(\frac{a_x}{l_{Pc}}\right)^{1/2}`$ (18)
Substitution of (16) into (18) leads to (7) which means that the mass of the vector bosons Z and W as well as the time of separation of the electromagnetic and weak interactions are directly obtained, based on the ENU model, in an independent way.
Conclusions:
1.The Vaidya metrics based ENU model allowing to localize the gravitational energy exhibits its capability to manifest some common features of the gravitational and weak interactions.
2.The paper presents an independent mode of determination of the mass of vector bosons Z and W, as well as the time of separation of the electromagnetic and weak interactions.
3.The paper follows up our previous contributions showing the unity of the fundamental physical interactions.
References
1. I. L. Rozentahl, Adv. Math. Phys. Astr. 31 (1986) 241
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3.V. Skalsky, M. Sukenik, Astrophys. Space Sci., 236 (1991) 169
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6.P.C. Vaidya: Proc. Indian Acad. Sci., A33 (1951) 264 |
warning/0001/math0001135.html | ar5iv | text | # ANALYTIC TOPOLOGY of groups, actions, strings and varietes
## Introduction
This paper is devoted to an application of Analysis to Topology. The latter is very broadly understood and includes geometric theory of finitely generated groups, group cohomology, Kazhdan groups, actions of groups on manifolds, superrigidity, fundamental groups of Kähler and quaternionic Kähler manifolds and conformal field theory. The motivation and philosophy which has led to the present research will be reflected upon in \[Reznikov 7\] and here we will merely say that we believe Analysis to be a major tool in studying finitely generated groups. An alternative look is provided by arithmetic method, notably by passing to a pro-p completion and using Galois cohomology. This will be described in \[Reznikov 8\].
Each of six chapters which constitute this paper opens with a short overview; a global picture is as follows. Chapter I and III treat analytic aspects of geometry of finitely generated groups. Given an immersion $`MN`$ of negatively curved manifolds ($`M`$ compact) there is a boundary map $`\stackrel{~}{M}\stackrel{~}{N}`$, and it has remarkable regularity properties. Invoking the Thurston theory, we show that the actions $`A`$ of pseudoAnosovs on $`W_p^{1/p}(S^1)/const`$ have striking properties from the viewpoint of functional analysis, namely,
$$\underset{n}{}A^nv^p<\mathrm{}$$
for some $`v0`$. We apply this to a classical problem: when a surface fibration is negatively curved and derive a strong necessary condition.
We then develop a theory of quantization for the mapping class group. A classical work on $`𝒟iff^{\mathrm{}}(S^1)`$ suggests a two-step quantization: first, obtaining a symplectic representation in $`Sp(W_2^{1/2}(S^1)/const)`$ with image in the restricted symplectic group \[Pressley-Segal 1\] and then using the Shale-Weil representation. The first step meets obstacles and the second step breaks down completely for the mapping class group: first, because $`ap_{g,1}`$ does not act smoothly on $`S^1`$, so it’s unclear why it can be represented in $`Sp(W_2^{1/2}(S^1)/const)`$, second, if even it can (this happens to be the case), there is no way to show that the image lies in the restricted symplectic group (it almost certainly does not). The solution comes at the price of abandoning the classical scheme and developing a theory of a new object which we call bicohomology spaces $`_{p,g}`$. The mapping class group $`ap_g`$ act in $`_{p,g}`$ and the latter shows remarkable properties, like duality and existence of vacuum. The last property is translated into the fact that $`H^1(ap_{g,1},W_p^{1/p}(S^1)/const)`$ is not zero. Finally, we find $`ap_g`$-equivariant maps of the space of all discrete representations of the surface group into $`PSL_2()`$ to our spaces $`_{p,g}`$.
Chapter II uses Analysis to study groups, acting on the circle (we need $`𝒟iff^{1,\alpha }`$ regularity, so $`ap_{g,1}`$ is not included). Our first main theorem says roughly that Kazhdan groups do not act on the circle. Very special cases of this result, for lattices in Lie groups, were recently found (see the references in Chapter II). The Hilbert transform, which played a major role in Chapter I, is crucial for the proof of this result as well. We then develop a theory of higher characteristic classes for subgroups of $`𝒟iff^{1,\alpha }`$ (the first being classically known as an integrated Godbillon-Vey class). All this classes vanish on $`𝒟iff^{\mathrm{}}(S^1)`$. It is safe to say that the less is smoothness, the more interesting is the geometry ”of the circle”.
Chapter III brings us back to asymptotic geometry of finitely generated groups. We propose, for a non-Kazhdan group, to study the asymptotic behaviour of unitary cocycles. We prove a general convexity result which shows that an embedding of $`G`$ in the Hilbert space, given by a unitary cocycle, is ”uniform”. We then prove a growth estimate for unitary cocycles of a surface group, using very heavy machinery from complex analysis, adjusted for our situation. Similar result for cocycles in $`H^1(G,l^p(G))`$ has already been given in Chapter I.
Chapter IV studies symplectomorphism groups. There is a misterious similarity between groups acting on the circle and groups acting symplectically on a compact symplectic manifold. In parallel with the above mentioned result in Chapter II we show, roughly, that transformations of a Kazhdan group acting on a symplectic manifold must satisfy a partial differential equation. An example is $`Sp(2n,)`$ acting linearly on $`T^{2n}`$ and, very probably, $`ap_g`$ acting on the space of stable bundles over a Riemann surface. (I don’t know for sure if $`ap_g`$ is Kazhdan). In dimension 2 the result is very easy and was known before. We also introduce new charateristic classes for symplectomorphism groups, in addition to the two series of classes defined in our previous papers, and use them to express a volume of a negatively curved manifold through the Busemann function on the universal cover.
Chapter V studies volume-preserving actions. We introduce a new technique into the subject, that of (infinite-dimensional, non-positively curved) spaces of metrics. We define a invariant of an action which is an infinum of a displacement in the space of metrics and show that for an action of a Kazhdan group which does not fix a $`\mathrm{log}L^2`$-metric, this invariant is positive (a weak version of this result for the special case of lattices was known before). We then turn to a major open problem, that of non-linear superrigidity and prove what seems to be first serious breakthrough after many years of effort.
Chapter VI deals with fundamental groups of Kähler and quaternionic Kähler manifolds. The situation is exactly the opposite to the studied in Chapters II and IV, namely, these groups tend to be Kazhdan. We first extend our rationality theorem for secondary classes of flat bundles over projective varietes to the case of quasiprojective varietes, answering a question posed to us by P.Deligne. We then prove that a fundamental group of a compact quaternionic Kähler manifold is Kazhdan, therefore providing a very strong restriction on its topology. We also discuss polynomial growth of the group cohomology classes for Kähler groups, proved nontrivial in a previous paper.
The paper uses many different analytic techniques. Within each chapter, there is a certain coherence in the point of view adopted for study.
I started this project on a chilly evening of November, 1998 in an African café in Leipzig and finished it on a hot afternoon of July, 1999 in Jerusalem. The manuscript has been written up by August, 11, 1999; I would appreciate any mentioning of a possible overlap with any paper/preprint which appeared before this date. During the long time when the paper was being typed and then polished, I found a proof of several statements which had been conjectured in the paper, in particular a construction of a cocycle for the group of quasisymmetric homeomorphisms valued in $`W_p^{1/p}(S^1)/const`$, which was conjectured in Chapter I. The proofs will appear in a sequel to this paper.
## Chapter 1 Analytic topology of negatively curved manifolds, quantum strings and mapping class groups
Chapter I opens with simple observations concerning the cohomology $`H^1(G,l^{\mathrm{}}(G))`$ for a finitely-generated group. If $`G`$ is amenable we produce plenty of polynomial cohomology classes in $`H^{}(G,)`$ given by an explicit formula (Theorem 1.2.1). Then we prove a convexity theorem 1.2.2 saying that if there are Euclidean-type quasigeodesics in the Cayley graph of $`G`$, then $`G/[G,G]`$ is infinite.
We then review some standard facts on $`l^p`$-cohomology in sections 2–4. One defines an asymptotic invariant of a finitely generated group $`G`$, called a constant of coarse structure $`\alpha (G)`$, as an infinum of $`p`$, $`1p\mathrm{}`$ such that $`H^1(G,l^p(G))0`$. For all noncusp discrete groups of motions of complete manifolds of pinched negative curvature, $`\alpha (G)<\mathrm{}`$. For discrete subgroups $`G`$ of $`SO^+(1,n)`$, $`\alpha (G)\delta (G)`$, where $`\delta `$ is the exponent of the group. In section 4 we review function spaces. A classical result in weighted Sobolev spaces may be reformulated as an identification of the $`l^p`$-cohomology of cocompact real hyperbolic lattices: $`H^1(G,l^p(G))=W_p^{(n1)/p}(S^{n1})/const`$ . It follows that $`\alpha (G)=n1`$.
In section 5 we prove a first result within a program to classify groups according to the cocycle growth. We show for surface groups, that if $`L_gl^p(G)`$ for all $`gG`$, then $`|(g)|const[length(g)]^{1/p^{}}`$. Here $`:G`$ is any function (Theorem 5.6). This result with no doubt generalizes to higher-dimensional cocompact lattices in simple Lie groups of rank one.
In section 6 we present a new theory for boundary maps of negatively curved spaces, associated with immersions of closed manifolds. The most striking is a partial regularity result (Theorem 6.1, part 4).
As is well known, the group of quasisymmetric ($`n=2`$) or quasiconformal ($`n3`$) homeomorphisms of $`S^{n1}`$ act on $`W_p^{\frac{n1}{p}}(S^{n1})`$ for $`p>n1`$. The action of $`𝒢_1`$ on $`W_2^{1/2}(S^1)/const`$ is in fact symplectic. We give application to the regularity of quasisymmetric homeomorphisms (Theorem 8.2). In Corollary 9.2 we prove that the unitary representation of a subgroup $`G`$ of $`SO(1,2)`$ in $`W_2^{1/2}(S^1)/const`$ is an invariant of a component of the Teichmüller space $`𝐓(G)`$.
In Theorem 9.3 we show striking properties of invertible operators $`A`$ in Banach spaces $`W_p^{1/p}(S^1)/const,p>2`$, induced by quasiAnosov maps in $`ap_{g,1}`$, namely
$$\underset{k}{}A^kv^p<\mathrm{}$$
for some $`0v`$. In Theorem 9.5 we find a new inequality in topological Arakelov theory, based on the work of \[Matsumoto-Morita 1\]. In Theorem 9.6 we find very strong restrictions on a subgroup $`Gap_g`$, such that an induced group extension $`\stackrel{~}{G}`$:
$$1\pi _1(\mathrm{\Sigma }_g)\stackrel{~}{G}G1$$
is a fundamental group of a negatively curved compact manifold (this is a classical problem). In section 10 we extend the theory to the limit case $`p=1`$, introducing an $`L^1`$-analogue of Zigmund spaces, which we call $`_{k,\alpha }`$.
In section 11 we start a new theory of secondary quantization of Teichmüller spaces. We introduce the bicohomology spaces $`_{g,p}`$ and show that $`ap_g`$ acts on these spaces. We show (difficult!) that $`_{2,p}`$ is an infinite-dimensional Hilbert space and there is a symmetric bilinear nondegenerate form of signature $`(\mathrm{},m)`$ which is $`ap_g`$-invariant. What is the value of $`m`$, we don’t know at the time of writing of this introduction (August,1999). So does the secondary quantization lead to ghosts? We provide a holomorphic realization in the space of $`L^2`$-holomorphic 2-forms on $`^2\times ^2/G`$ and $`^2\times \overline{^2}/G`$ (Theorem 11.12). In section 12 we interpret $`_{p,g}`$ as operator spaces (proposition 12.2), and prove the existence of vacuum (Theorem 12.5). We prove that $`H^1(ap_{g,1},W_p^{1/p}(S^1)/const)0`$ for $`p2`$. It still may be true that $`ap_{g,1}`$ is Kazhdan, because the action is not orthogonal.
In section 13 we construct $`ap_g`$-equivariant maps of the space of discrete representations of the surface group in $`SO^+(1,3)=PSL_2()`$ to our spaces $`_{p,g}`$ (Theorem 13.1). In Theorem 13.2 we summarize our knowledge of the functional-analytic structure coming from hyperbolic 3-manifolds which fiber over the circle.
### 1.1 Metric cohomology
1.1.1. Let $`G`$ be a finitely generated group. Let $`𝕂=,`$. Let $`V`$ be a locally convex topological $`𝕂`$-vector space which is a $`G`$-module, that is, there is a homomorphism $`GAut(V)`$. If $`\{g_i\},i=1,\mathrm{},n`$ is a finite set of generators of $`G`$, then the evaluation map $`f\{f(g_i)\}`$ establishes an injective homomorphism $`Z^1(G,V)\mathrm{\Pi }_{i=1}^nV`$ of the space of 1-cocycles of $`G`$ in $`V`$. One calls the induced topology in $`Z^1(G,V)`$ the cocycle topology; it does not depend on the choice of generators. A coboundary map $`VZ^1(G,V)`$ may have an image $`\overline{B^1(G,V)}`$ which is not closed in $`Z^1(G,V)`$; the quotient $`Z^1(G,V)/\overline{B^1(G,V)}`$ is called reduced first cohomology space. One way to produce nontrivial cohomology classes is to consider limits of coboundaries, that is, elements of $`\overline{B^1(G,V)}/B^1(G,V)`$. That amounts to considering nets $`\{v_\alpha V\}`$ such that $`g_iv_\alpha v_\alpha l(g_i)`$ for $`i=1,\mathrm{},n`$. If $`V`$ is a Banach space and $`G`$ acts isometrically without invariant vectors, then $`B^1(G,V)`$ is closed in $`Z^1(G,V)`$ if and only if there are no almost invariant vectors, that is, sequences $`v_j,v_j=1`$, such that $`g_iv_jv_j0`$ for all $`i=1,\mathrm{},n`$. This statement is an immediate consequence of the Banach theorem and is called Guichardet’s lemma \[Guichardet 1\]. So if there are almost invaiant vectors, then $`H^1(G,V)0`$, though the reduced cohomology $`H_{red}^1(G,V)`$ may be zero.
If $`V`$ is Banach and $`G`$ acts isometrically, let $`lZ^1(G,V)`$ be a cocycle. Then
$$l(g)\underset{i=1}{\overset{n}{\mathrm{max}}}l(g_i)length(g)$$
, where $`length(g)`$ is the length of the element $`g`$ in the word metric, induced by $`\{g_i\}`$. The proof is immediate by induction, using the cocycle equation $`l(gh)=gl(h)+l(g)`$.
Now let $`V_j,j=1,\mathrm{},m`$ be a collection of Banach spaces on which $`G`$ acts isometrically and let $`\phi :_{j=1}^mV𝕂`$ be a map continuous in a sense that $`\phi (_{j=1}^mv_j)const\mathrm{\Pi }_{j=1}^mv_j`$. Let $`l_jZ^1(G,V_j)`$ and let $`lZ^m(G,𝕂)`$ be the cup product $`l(g_1,\mathrm{},g_m)=\phi (_{j=1}^ml_j(g_j))`$.
Lemma 1.1.$`lZ^m(G,𝕂)`$ is of polynomial growth, more precisely
$$|l(g_1,\mathrm{},g_m)|const\mathrm{\Pi }_{i=1}^mlength(g_i).$$
Proof.— is immediate from the remarks made above.
A general definition of polynomial cohomology is to be found in \[Connes-Moscovici 1\]. As we will see, Lemma 1.1 is a very powerful tool for constructing cocycles of polynomial growth in concrete situations.
Proposition 1.1.— Let $`G`$ be an infinite finitely generated group. Consider a left action of $`G`$ on $`l^{\mathrm{}}(G)`$. Then $`H^1(G,l^{\mathrm{}}(G))0`$. Moreover, $`H^1(G,l_0^{\mathrm{}}(G))0`$.
Proof.— Let $`\{g_i\}`$ be a finite set of generators of $`G`$, and let $`length(g)`$ be a word length of an element $`g`$. Define a right-invariant word metric by $`\rho (x,y)=length(xy^1)`$. Let $`x_0G`$ and let $`F(x)=\rho (x_0,x)`$. Obviously, $`F`$ is unbounded. Now let $`l(g)=L_gFF`$ where $`L_g`$ is a left action on functions, that is, $`l(g)(x)=F(g^1x)F(x)`$. We find
$$|l(g)(x)|=|\rho (x_0,g^1x)\rho (x_0,x)||\rho (g^1x,x)|=\rho (g^1,1).$$
So $`l`$ is a cocycle of $`G`$ in $`l^{\mathrm{}}(G)`$. If it were trivial, we would have a bounded function $`f`$ such that $`L_gFF=L_gff`$ that is, $`Ff`$ would be invariant, therefore constant, a contradiction. The second statement of the Proposition will be proved later in section 1.3.
1.1.2. Now let $`G`$ be amenable. In this case we have a continuous map $`\phi :\mathrm{\Pi }_{j=1}^ml^{\mathrm{}}(G)𝕂`$ given by $`(f_1,\mathrm{},f_m)_Gf_1\mathrm{}f_m`$. By an integral we mean a left-invariant normalized mean of bounded functions. We obtain
Theorem 1.2.1.— Let $`G`$ be a finitely generated amenable group, let $`\rho _j,j=1,\mathrm{},m`$ be a collection of right-invariant word metrics on $`G`$. A formula
$$l(g_1,\mathrm{},g_m)=_G\mathrm{\Pi }_{j=1}^m[\rho _j(x_0,g_j^1x)\rho _j(x_0,x)]$$
defines a real-valued $`m`$-cocycle on $`G`$ of polynomial growth:
$$|l(g_1,\mathrm{},g_m)|const\mathrm{\Pi }_{j=1}^mlength(g_j)$$
for any word length $`length()`$.
Examples.— Let $`G=`$. If we choose generators $`\{1,1\}`$, then $`length(g)=|g|,`$ and
$$\rho (x_0,g^1x)\rho (x_0,x)=|x_0x+g||x_0x|\pm |g|$$
as $`x\pm \mathrm{}`$ and
$$_{}(|x_0x+g||x_0x|)=0.$$
However, if we choose generators $`\{1,2\}`$, then
$$length\left(g\right)=\{\begin{array}{cc}\left|g\right|,& g0\\ \frac{g}{2},& g0\text{ and even }\\ \frac{g+1}{2},& g0\text{ and odd }\end{array}$$
Then
$$length(x_0+gx)length(x_0x)$$
for $`g>0`$ and even will have limits $`\frac{g}{2}`$ when $`x\mathrm{}`$ and $`g`$ when $`x\mathrm{}`$, so
$$_{}[length(x_0+gx)length(x_0x)]=\frac{g}{4}.$$
So we obtain a cocycle $`l:`$ given by $`g\frac{g}{4}`$. Now, if $`G=^k`$, $`k2`$, let $`\rho _j,j=1,\mathrm{},k`$ be a word metric defined by a set of generators
$$\{e_1^{\pm 1},e_2^{\pm 1},\mathrm{},e_j^1,e_j^2,e_{j+1}^{\pm 1},\mathrm{},e_k^{\pm 1}\}$$
where $`e_s`$ is a generator of the s-th factor. If $`1j_1<j_2<\mathrm{}<j_mk`$ is a set of indices, then Theorem 2 provides a cocycle
$$l(g_1,\mathrm{},g_m)=_^k\mathrm{\Pi }_{r=1}^m[\rho _{j_r}(x_0,g_j^1x)\rho _{j_r}(x_0,x)].$$
If $`\pi _i:^k`$ is a projection to i-th factor, then $`l(g_1,\mathrm{},g_m)=(\frac{1}{4})^m\mathrm{\Pi }_{r=1}^m\pi _{j_r}(g_r)`$. It follows that classes of cocycles, given by Theorem 1.2.1, generate the real cohomology space of $`^k`$.
Remark.— If $`G`$ is amenable, $`\rho `$ is a right-invariant word metrics and for some $`x_0,gG`$,
$$_G[\rho (x_0,g^1x)\rho (x_0,x)]0,$$
then $`H_1(G,)0`$ and in fact $`g[G,G]`$ for all $`s0`$. This is a direct corollary of Theorem 1.2.1. A more interesting structure theorem is given below.
Theorem 1.2.2.— Let $`G`$ be a finitely generated amenable group, $`\rho `$ a right-invariant word metric. Let $`gG`$, assume a following convexity condition: there is some $`C>0`$, such that for any $`xG`$ there exists $`N0`$ such that $`\rho (g^k,g^1x)\rho (g^k,x)C`$ for $`kN`$. Then $`H_1(G,)0`$ and moreover, $`g^s[G,G]`$ for all $`s0`$.
Corollary 1.2.3.— Let $`G`$ be a Heisenberg group $`\{x,y,z|[x,y]=z,[x,z]=[y,z]=1\}`$. Then for any right-invariant word metric $`\rho `$, there exists $`aG`$ such that $`lim\; inf_k\mathrm{}[\rho (z^k,z^1a)\rho (z^k,a)]0`$.
Proof of the Corollary.— Since $`z[G,G]`$, the result follows from Theorem 1.2.2. Indeed $`G`$ is nilpotent, therefore amenable.
Proof of the Theorem.— Consider a 1-cocycle $`l(\gamma )(x)=\rho (x_0,\gamma ^1x)\rho (x_0,x)`$, $`lZ^1(G,l^{\mathrm{}}(G)).`$ Set $`x_0=g^n`$, so
$$l_n(g)(x)=\rho (g^n,g^1x)\rho (g^n,x).$$
If for any $`x`$ and sufficiently big $`n`$, $`\rho (g^n,g^1x)\rho (g^n,x)>C`$ then a pointwise limit $`lim_n\mathrm{}l_n(g)(x)`$ exists and is $`C`$. Since $`|l_n(z)(x)|\rho (z^1,1)`$, there is a subsequence $`n_k`$ such that $`l_{n_k}(z)`$ converges pointwise for any $`z`$ to a bounded function $`l(z)`$. One sees immediately that $`l:Gl^{\mathrm{}}(G)`$ is a cocycle, so $`z_Gl(z)`$ is a homomorphism from $`G`$ to $``$. Since $`l(g)(x)C>0`$ for all $`x`$, $`_Gl(g)C>0`$, so $`H_1(G,)0`$ and $`g^s[G,G]`$, as desired.
1.1.3. Let $`\phi :R_+R_+`$ be a smooth function such that $`\phi (x)\mathrm{}`$ as $`x\mathrm{}`$ and $`\phi ^{}(x)0`$. Let $`G`$ be a finitely generated group and let $`\rho `$ be a right-invariant word metric. Consider $`F(x)=\phi (\rho (x_0,x))`$ where $`x_0G`$ is a fixed element. Since
$$\begin{array}{cc}\hfill \left|\left(L_gFF\right)\left(x\right)\right|& =\left|F\left(g^1x\right)F\left(x\right)\right|\hfill \\ & =\left|\phi \left(\rho (x_0,g^1x)\right)\phi \left(\rho (x_0,x)\right)\right|\hfill \\ & sup_{tI}\left|\phi ^{}\left(t\right)\right|\left|\rho (g^1x,x)\right|\hfill \\ & sup_{tI}\left|\phi ^{}\left(t\right)\right|\rho (g^1,1)\hfill \end{array}$$
where $`I=[\mathrm{min}(\rho (x_0,x),\rho (x_0,g^1x)),\mathrm{max}(\rho (x_0,x),\rho (x_0,g^1x))]`$, we see that $`L_gFFl_0^{\mathrm{}}`$. Therefore $`H^1(G,l_0^{\mathrm{}}(G))0`$, because the cocycle $`L_gFF`$ cannot be trivial as a cocycle valued in $`l_0^{\mathrm{}}`$ (by the same reasons as in the proof of the first statement of Proposition 1.1 ) . The proof of Proposition 1.1 is now complete.
Notice that, since $`\rho (u,v)=length(uv^1)`$,
$$\rho (x_0,x)length(g)\rho (x_0,g^1x)\rho (x_0,x)+length(g),$$
so that
$$|(L_gFF)(x)|\underset{|t\rho (x_0,x)|length(g)}{sup}|\phi ^{}(t)|\times \rho (g^1,1).$$
Remark.— Let $`S(N)=\{g|length(g)=N\}`$. If $`S(N)/S(N1)1`$ and $`_{k=1}^NS(k)/S(N)\mathrm{}`$ as $`N\mathrm{}`$, then for $`p>1`$ there is a radial function $`F(x)=\phi (\rho (x))`$ such that $`L_gFFl^p(G)`$ and the cocycle $`l:Gl^p(G)`$ defined by $`gL_gFF`$ is not a coboundary. Note that $`G`$ is automatically amenable. On the other hand, if $`S(N)e^{cN}`$, then such radial function does not exist. This follows at once from Hardy’s inequality. To produce classes in $`H^1(G,l^p(G))`$, one needs to use some more elaborate geometry than just distance function. In the next section we produce such classes for negatively curved groups/manifolds, using the visibility angles.
### 1.2 Constants of coarse structure for negatively curved groups
1.2.1. Throughout this section we assume that $`G`$ is a finitely generated, non-amenable group, therefore $`B^1(G,l^p(G))`$ is closed in $`Z^1(G,l^p(G))`$ for $`p1`$.
Definition 2.1.— A number $`\alpha (G)=inf_{1p\mathrm{}}\{p|H^1(G,l^p(G))0\}`$ is called a constant of coarse structure of $`G`$.
Remark.— The definition makes sense since by Proposition 1.1, $`H^1(G,l^{\mathrm{}}(G))0`$. We will need a proof of the following well-known fact (see, for example \[Pansu 1\]). The argument below is a slightly modified, from nonpositive curvature to negative curvature, version of a classical argument of \[Mishchenko 1,2\].
Proposition 2.1.— Let $`M^n`$ be a complete Riemannian manifold of negative curvature, not a cusp, satisfying $`K(M)1,Ric(M)(n1)K`$. Let $`G=\pi _1(M)`$. Then $`\alpha (G)(n1)\sqrt{K}`$.
Proof.— Let $`q_0\stackrel{~}{M}`$. Consider a map of $`G`$ onto an orbit $`𝒪`$ of $`q_0:ggq_0`$; it is equivariant with respect to the left action of $`G`$ on itself. Let $`q𝒪`$ and let $`v_q(s)`$ be an outward pointing vector from $`q`$ to $`s`$, that is, a unit vector in $`T_s\stackrel{~}{M}`$, tangent to geodesic segment joining $`q`$ and $`s`$. Consider for $`xG`$, $`F(x)=v_q(xq_0)`$. Notice that $`F(x)`$ takes values in $`T_{xq_0}\stackrel{~}{M}`$. We can consider the restriction of $`T\stackrel{~}{M}`$ on $`𝒪`$ as an equivariant vector bundle over $`𝒪`$. Pulling back to $`G`$, we obtain an left-equivariant vector bundle over $`G`$, equipped with an equivariant Euclidean structure. Then $`F`$ is a section of this bundle. Now consider $`(L_gFF)(x)`$. Since the action of $`G`$ on sections is given by $`L_gF(x)=g_{}F(g^1x)`$, where $`g_{}`$ is the derivative map ( $`g_{}:T_s\stackrel{~}{M}T_{gs}\stackrel{~}{M}`$), we get $`(L_gFF)(x)=g_{}F(g^1x)F(x)=g_{}v_q(g^1xq_0)v_q(xq_0)=v_{gq}(xq_0)v_q(xq_0)`$. So $`(L_gFF)(x)=|2\mathrm{sin}\frac{1}{2}\mathrm{}(gq,xq_0,q)|\mathrm{}(gq,xq_0,q)`$.
Let $`E|G`$ be the equivariant Euclidean vector bundle considered above (the pullback of $`G`$ of $`T\stackrel{~}{M}|𝒪`$). Let $`L^p(E)`$ be a Banach space of $`L^p`$-sections of $`E`$. We claim $`L_gFFL^p(E)`$ for $`p>(n1)\sqrt{K}`$. Let $`r(x)=dist_{\stackrel{~}{M}}(q_0,xq_0)`$. For $`g,q_0,q`$ fixed we have $`\mathrm{}(gq,xq_0,q)const_1e^{r(x)}`$ by a standard comparison theorem, since $`K(M)1`$. On the other hand, for fixed $`\delta >0`$, $`\mathrm{\#}(x|r\delta r(x)r+\delta )const_2e^{(n1)\sqrt{K}r}`$ by the Bishop’s theorem. Therefore $`L_gFFL^p(E)`$ for $`p>(n1)\sqrt{K}`$. Note we only need that $`G`$ acts discretely in $`\stackrel{~}{M}`$.
A map $`l:GL^p(E)`$ defined by $`l(g)=L_gFF`$ is obviously a cocycle. If it were trivial, we would have an $`L^p`$-section $`sL^p(E)`$, such that $`Fs`$ is invariant. That means $`g_{}(v_q(g^1xq_0)s(g^1x))=v_q(xq_0)s(x)`$, or $`v_{gq}(xq_0)g_{}s(g^1x)=v_q(xq_0)s(x)`$. Notice that since $`F(x)=1`$, $`Fs`$ is invariant and $`s(g)0`$ as $`length(g)\mathrm{}`$, $`(Fs)(x)=1`$ for all $`x`$. In particular, $`w=v_q(xq_0)s(x)`$ has norm one. Fix $`x`$ and let $`g`$ vary. We get $`v_{gq}(xq_0)w=g_{}s(g^1x)0`$ as $`length(g)\mathrm{}`$. Let $`P_+,P_{}`$ be an attractive and repelling fixed points of $`g`$ on the sphere at infinity of $`\stackrel{~}{M}`$.
Let $`w_+,w_{}`$ be unit vectors in $`T_{xq_0}(\stackrel{~}{M})`$, tangent to geodesics, joining $`xq_0`$ with $`P_+,P_{}`$. Then $`v_{g^nq}(xq_0)w_\pm 0`$ if $`n\pm \mathrm{}`$. It follows that $`w_\pm =w`$. Therefore all elements of $`G`$ are parabolic and have a common fixed point at infinity. So $`M`$ is a cusp, a contradiction. So $`H^1(G,l^p(E))0`$. However, $`l^p(E)`$ is equivariantly isometric to $`l^P(G)T_{p_0}(\stackrel{~}{M})`$. So $`H^1(G,l^p(E))H^1(G,l^p(G))T_{p_0}(\stackrel{~}{M})`$. We deduce that $`H^1(G,l^p(G))0`$.
The estimate of the Proposition is sharp. We will see later that if $`G`$ is a cocompact lattice in $`SO^+(1,n)`$, i.e. $`K(M)=1`$, then $`\alpha (G)`$ is exactly $`(n1)`$. Let now $`G`$ be a discrete nonamenable subgroup of $`SO^+(1,n)`$, or, equivalently, $`K(M)=1`$. Recall that the exponent $`\delta (G)`$ is defined by $`\delta (G)=inf\{\lambda |_{gG}e^{\lambda r(g)}<\mathrm{}\}`$ where $`r(g)=dist_{\stackrel{~}{M}}(p_0,gp_0)`$ for some fixed $`p_0\stackrel{~}{M}`$. If $`G`$ is geometrically finite, then by a well-known theorem \[Nicholls 1\] $`\delta (G)`$ is equal to the Hausdorff dimension of the limit set $`dim(\mathrm{\Lambda }(G))S^{n1}`$. Note that if $`\mathrm{\Lambda }(G)S^{n1}`$, then $`dim\mathrm{\Lambda }(G)<n1`$ by \[Sullivan 1\] and \[Tukia 1\].We now have
Proposition 2.2.— Let $`G`$ be a discrete subgroup of $`SO^+(1,n)`$, not a cusp group. Then $`\alpha (G)\delta (G)`$.
Proof.— The Proposition follows from the proof of the Proposition 2.1. Indeed, we only need that $`_{gG}e^{pr(g)}<\mathrm{}`$ to conclude that one has a cocycle $`l:Gl^p(G)`$. It has been proven already that this cocycle is not a coboundary.
Remark.— The relation of the constant of coarse structure to “conformal dimension at infinity” is discussed in \[Pansu 2\].
### 1.3 Function spaces: an overview
For $`s0`$ an integer and fractional part of $`s`$ are denoted $`[s]`$ and$`\{s\}`$ respectively. A Sobolev-Slobodec̆ky space $`W_p^s(^n)`$,$`(p>1)`$ consists of measurable locally integrable functions $`f`$ on $`^n`$ such that $`D^\alpha fL^p(R^n)`$ for $`|\alpha |[s]`$ and
$$\underset{|\alpha |=[s]}{}\frac{|D^\alpha f(x)D^\alpha f(y)|^p}{|xy|^{n+\{s\}p}}𝑑x𝑑y<\mathrm{}$$
A space of Bessel potentials $`H_p^s`$ consists of functions $`f`$ for which a Liouville-type operator
$$𝒟^sf=((1+|\xi |^2)^{s/2}\widehat{f}(\xi ))^{}$$
satisfies $`𝒟^sfL^p`$. Warning: $`H_p^sW_p^s`$ if $`s`$ is not an integer. For $`p=2`$ the condition is equivalent to
$$(1+\mathrm{})^{s/2}\widehat{f}L^2(^n).$$
Here $`f(x)\widehat{f}(\xi )`$ is the Fourier transform and $`\mathrm{}=\frac{^2}{x_i^2}`$.
A space of BMO functions BMO$`(^n)`$ is defined as a space of functions $`f`$ for which
$$\underset{Q}{sup}\frac{1}{|Q|}_Q|f(x)f_Q|𝑑x<\mathrm{},$$
where $`Q`$ runs over all cubes in $`^n`$ and
$$f_Q=\frac{1}{|Q|}_Qf(x)𝑑x,$$
$`|Q|=_Q1𝑑x`$. One has $`W_p^{n/p}BMO`$ for all $`1<p<\mathrm{}`$, and moreover $`H_p^{n/p}H_{p_1}^{n/p_1}`$ for $`1<p<p_1<\mathrm{}`$ (this follows from Theorem 2.7.1 of \[Triebel 1\]. In some sense $`BMO`$ is a limit of $`H_p^{n/p}`$ as $`p\mathrm{}`$.
If $`fW_p^1`$ the restriction of $`f`$ on hyperplanes $`\{x_n=ϵ\}R^n`$ (where $`(x_1,\mathrm{},x_n)`$ are Euclidean coordinates ) have both $`L^p`$ and nontangential limits a.e. on $`^{n1}=\{x|x_n=0\}`$ and the limit function $`f|_{^{n1}}`$, called trace of $`f`$, satisfies $`f|_{^{n1}}W_p^{11/p}`$. By a nontangential limit we mean the following. Let $`y^{n1}`$ and let $`C_\delta `$ be a Stolz angle centered at $`y`$, that is, a set $`\{z,x_n|x_n\delta |z|\}`$ for $`\delta >0`$. Then a function $`f`$ defined in $`_+^n=\{x_n>0\}`$ has a nontangential limit $`f(y)`$ at $`y`$ if
$$f(x)\underset{\underset{xC_\delta }{xy}}{}f(y)$$
for all $`\delta `$. Note that the points in $`C_\delta `$ are within a bounded distance from any geodesic of a hyperbolic metric
$$\frac{_{i=1}^ndx_i^2}{x_n^2},$$
which has $`y`$ as a point at infinity. The trace theorem mentioned above may be found in \[Triebel 1\], section 2.7.2. Notice that functions in $`W_2^1(^2)`$ have traces in $`W_2^{1/2}(^1)`$.
Now let $`\mathrm{\Omega }^n`$ be a bounded domain with a smooth boundary. We define $`W_p^s(\mathrm{\Omega })`$ as a space of locally integrable functions with $`D^\alpha fL^p`$ for $`|\alpha |s`$ and such that
$$\underset{|\alpha |=[s]}{}\frac{|D^\alpha f(x)D^\alpha f(y)|^p}{|xy|^{n+\{s\}p}}<\mathrm{}.$$
Equivalently, $`W_p^s(\mathrm{\Omega })`$ is a space of restrictions of function from $`W_p^s(^n)`$ on $`\mathrm{\Omega }`$. \[Triebel 1, Chapter 3\]. One also defines $`H_p^s(\mathrm{\Omega })`$ as a space of restrictions of $`H_p^s(^n)`$ on $`\mathrm{\Omega }`$. For a compact smooth manifold $`M`$ without boundary (in particular, for the boundary $`\mathrm{\Omega }`$) one easily defines the spaces $`W_p^s(M)`$ and $`H_p^s(M)`$ \[Triebel 1, Chapter 3\] ($`H_p^s`$ is $`F_{p,2}^s`$ in Triebel’s notations ).
If $`M`$ is compact and $`g`$ a Riemannian metric on $`M`$, let $`\mathrm{}_g`$ be a corresponding Laplace-Beltrami operator. One can construct a space of Bessel potentials $`(1+\mathrm{})^{s/2}(L_p(M))`$. It is known \[Rempel-Schulze 1, Theorem 1, section 2.3.2.5\], \[Hörmander 1\], that this space coincides with $`W_p^s`$ (and not $`H_p^s`$). Warning: our $`W_p^s`$ is called $`H^{p,s}`$ in \[Rempel-Schulze 1\] and in many other sources. In particular, $`W_2^s(S^1)`$ consists of functions $`f=_na_ne^{in\theta }`$ , such that $`|n|^{2s}|a_n|^2<\mathrm{}`$. We will see that $`W_2^{1/2}(S^1)`$ is especially important in topology.
If $`fW_p^1(\mathrm{\Omega })`$ then $`f`$ has an $`L^p`$ and nontangential limit a.e. on $`\mathrm{\Omega }`$ and $`f|_\mathrm{\Omega }W_p^{11/p}(\mathrm{\Omega })`$. In particular, for a unit disc $`D^2`$, and a function $`fW_2^1(D)`$, $`f|_{S^1}W_2^{1/2}(S^1)`$.
We will need trace theorems for weighted Sobolev-Lorentz spaces \[Kudryavcev 1,2\], \[Vasharin 1\], \[Lions 1\], \[Lizorkin 1,2\], \[Uspenski 1\]. Let $`\mathrm{\Omega }`$ be as above and let $`\rho (x)=dist(x,\mathrm{\Omega })`$. Consider $`L_p^1(\mathrm{\Omega },\rho ^\alpha )`$ as a space of functions $`f`$ such that $`_\mathrm{\Omega }|f|^p\rho ^\alpha 𝑑x<\mathrm{}`$. Then $`f`$ has a nontangential limit a.e. on $`\mathrm{\Omega }`$ and
$`1)f|_\mathrm{\Omega }=0\text{ if }\alpha 1`$
$`2)f|_\mathrm{\Omega }W_p^{\frac{p1\alpha }{p}}(\mathrm{\Omega }),\alpha >1`$.
Moreover,
$$f_{W_p^{\frac{p1\alpha }{p}}}const_\mathrm{\Omega }|f|^p\rho ^\alpha 𝑑x$$
and for any $`fW_p^{\frac{p1\alpha }{p}}(\mathrm{\Omega })`$ and harmonic $`h`$, $`h|_\mathrm{\Omega }=f`$, one has
$$_\mathrm{\Omega }|h|^p\rho ^\alpha 𝑑xconstf_{W_p^{\frac{p1\alpha }{p}}}.$$
### 1.4 $`l^p`$-cohomology of cocompact real hyperbolic lattices
The following result is an immediate corollary of the Poincaré inequality in hyperbolic space, which is equivalent to Hardy inequality, and the classical results on traces of functions in weighted Sobolev spaces, reviewed in the previous section. It first appeared in print, with a different proof, in \[Pansu 1\]. We include a proof here, as many parts of it will be used in the theory later.
Theorem 4.1, part 1.— Let $`GSO^+(1,n)`$ be a cocompact (uniform) lattice. Then there is a $`G`$-equivariant isomorphism of Banach spaces
$$H^1(G,l^p(G))W_p^{\frac{n1}{p}}(S^{n1})/const$$
for $`p>n1`$. For $`1<pn1`$, $`H^1(G,l^p(G))=0`$.
Corollary 4.1.— The constant of fine structure $`\alpha (G)`$ is equal $`n1.`$
Remarks.
1)Theorem 4.1 is a first step in the program of linearization of 3-dimensional topology, which we will develop below in this chapter. A crucial fact is that $`W_2^{1/2}(S^1)`$ admits a natural action of the extended mapping class group $`ap_{g,1}`$. This will be proved in section 7 below.
2)Let $`^n`$ be a hyperbolic space. Since $`G=\pi _1(^n/G)`$, by the work of \[Golds̆tein-Kuzminov-Shvedov 1 \] we know that $`H^1(G,l^p(G))`$ equals $`L^p`$-cohomology of $`^n`$. So Theorem 4.1 computes the $`L^p`$-cohomology of the hyperbolic space.
Recall that any class $`l`$ in $`H^1(G,l^p(G))`$ has a primitive function $`:G`$ defined up to a constant, such that $`l(g)=L_gFF`$. This follows from the fact that a module of all functions $`^G`$ is coinduced from the trivial subgroup and therefore cohomologically trivial \[Brown 1\].
Theorem 4.1, part 2.— Let $`G`$ be a cocompact lattice in $`SO^+(1,n)`$ and let $`lH^1(G,l^p(G))`$, let $`:G`$ be a primitive function for $`l`$ (unique up to a constant). Let $`GS^{n1}`$ be the boundary of $`G`$ as a word-hyperbolic group. Then for almost all points $`xG`$, $`(g)`$ has nontangential limits as $`gx`$, and the limit function $`|_{S^{n1}}W_p^{\frac{n1}{p}}(S^{n1}).`$
Corollary 4.2.— If $`1<p<p_1<\mathrm{}`$, then a natural map $`H^1(G,l^p(G))H^1(G,l^{p_1}(G))`$ is injective. In fact, for $`n1<p<p_1<\mathrm{}`$ one has a commutative diagram
$$\begin{array}{ccc}\hfill H^1(G,l^p\left(G\right))& \stackrel{}{}& W_p^{\frac{n1}{p}}\left(S^{n1}\right)/const\hfill \\ \hfill & & \hfill \\ \hfill H^1(G,l^{p_1}\left(G\right))& \stackrel{}{}& W_{p_1}^{\frac{n1}{p_1}}\left(S^{n1}\right)/const\hfill \end{array}$$
where the right vertical arrows exists by an embedding theorem of Sobolev-Slobodec̆ki space \[Triebel 1, 2.7.1 \].
Proof of the Corollary 4.2.— The commutative diagram is implied by the proof of the Theorem 4.1. The injectivity follows immediately.
Proof of the Theorem.— Though a shorter proof of part 1 of the Theorem can be given by using \[Goldshtein-Kuzminov-Shvedov 1\], in order to prove part 2 we need to make an isomorphism $`H^1(G,l^p(G))L^pH^1(^n)`$ explicit. Here $`L^pH^1(V)`$ is the $`L^p`$-cohomology of a complete Riemannian manifold $`V`$.
Let $`l`$ be a cocycle in $`Z^1(G,l^p(G))`$. We have then an affine isometric action $`g\stackrel{𝜋}{}(vL_gv+l(g))`$ of $`G`$ on $`l^p(G)`$. To it corresponds a smooth locally trivial affine Banach bundle over $`M=^n/G:E=[\stackrel{~}{M}\times l^p(G)]/\text{diagonal action}`$. By local triviality, smooth partition of unity and affine structure on fibers one constructs a smooth section $`s`$ of this bundle. It can be interpreted as an equivalent smooth map $`s:\stackrel{~}{M}l^p(G)`$, that is, $`s(g^1x)=L_gs(x)+l(g)`$. We note that there is some sonstant $`C>0`$ such that $`s(x)<C`$ for all $`x\stackrel{~}{M}^n`$ ($`sT_x^{}\stackrel{~}{M}l^p(G)`$ ). This is because $`M`$ is compact. Now let $`R^G`$ be a primitive for $`l`$, i.e. $`l(g)=L_g`$. Put $`\sigma (x)=s(x)+`$: this is a function $`\sigma :\stackrel{~}{M}^G`$ with the same derivative as $`s`$ in the sense that for all $`gG`$, $`\sigma _g=s_g`$ where $`\sigma _g`$, $`s_g`$ means $`g`$-th coordinate. Next, we claim that $`\sigma `$ is invariant, i.e. $`\sigma (g^1x)=L_g\sigma (x)`$. In fact, $`l(g)=L_g`$, so $`s(g^1x)=L_gs(x)+L_g`$, so $`\sigma (g^1x)=L_g\sigma (x)`$. So for $`x\stackrel{~}{M}`$ and $`g,hG`$ we have $`\sigma (g^1x)(h)=\sigma (x)(g^1h)`$. Let $`f(x)=\sigma (x)(1)`$, then $`\sigma (x)(g)=f(gx)`$. Since $`\sigma (x)=s(x)l^p`$ and is bounded in norm, we have for all $`x\stackrel{~}{M}`$ that $`_{gG}|f(gx)|^p<C`$. In particular,
$$_{\stackrel{~}{M}}|f|^p=_{\stackrel{~}{M}/G}\underset{gG}{}|f(gx)|^p<CVol(M).$$
In other words, $`|f|L^p(^n)`$. Now, we can use a Poincaré model for the hyperbolic space, that is, the unit ball $`B^n^n`$ with a hyperbolic metric
$$g_h=\frac{g_e}{(1r^2)^2}.$$
If $`\mu _e,\mu _h`$ denote a Euclidean and a hyperbolic measure respectively, $`|f|_e,|f|_h`$ denote a norm of a gradient of a function in the Euclidean and hyperbolic metric respectively, $`\rho (z)=1r(z)`$ denote a Euclidean distance to the boundary $`B^nS^{n1}`$, then
$$const_2\rho ^{pn}|f|_e^p\mu _e|f|_h^p\mu _hconst_1\rho ^{pn}|f|_e^p\mu _e,$$
so we have $`_{B^n}\rho ^{pn}|f|_e^p\mu _e<\mathrm{}`$.
By a theorem of Kudryavcev-Vasharin-Lizorkin-Uspenski-Lions mentioned above, we find that $`f|_{(1ϵ)S^{n1}}`$ has an $`L^p`$-limit $`f|_{S^{n1}}`$ to which it converges nontangentially a.e. , and moreover $`f|_{S^{n1}}W_p^{\frac{n1}{p}}(S^{n1})`$ if $`p>n1`$ and $`f|_{S^{n1}}=const`$ if $`pn1`$. We claim that a map $`lf|_{S^{n1}}`$ is a well-defined bounded linear operator from $`H^1(G,l^p(G))`$ to $`W_p^{\frac{n1}{p}},p>n1`$. First, we notice that since $`s:\stackrel{~}{M}l^p(G),\sigma (x)=s(x)+`$ and $`\sigma (x)(g)=f(gx)`$, we have for almost all $`x\stackrel{~}{M}`$, $`f(gx)(g)l^p(G)`$ (as a function of $`g`$). In particular, $`f(gx)(g)0`$ as $`length(g)\mathrm{}`$. This proves that, identifying $`G`$ with an orbit of $`x`$, $`(g)`$ has a nontangential limit a.e. on the boundary $`GS^{n1}`$ and $`|_{S^{n1}}=f|_{S^{n1}}W_p^{\frac{n1}{p}}`$. In particular, $`f|_{S^{n1}}\text{ mod constants}`$ does not depend on the choice of a section $`s`$. Since changing $`l`$ by a coboundary leads to an isomorphc affine $`l^p(G)`$-bundle, $`f|_{S^{n1}}\text{mod constants}`$ depends only on the class $`[l]H^1(G,l^p(G))`$. So we get a well-defined operator $`H^1(G,l^p(G))W_p^{\frac{n1}{p}}/const`$. We claim it is bounded. An affine flat bundle $`E`$ has been defined as $`\stackrel{~}{M}\underset{𝐺}{\times }l^p(G)`$, where $`G`$ acts on $`l^p(G)`$ by $`vL_gv+l(g)`$. It is enough to show, that there is a constant $`C`$, depending only on $`G`$ but not on $`l`$, such that $`E`$ possesses a Lipschitz section $`s`$ with $`s<Cl`$, where $`l=sup_il(g_i)`$ for a choice of generators $`g_i,i=1,\mathrm{},m`$. We note that $`l`$ effectively controls the monodromy of the flat connection in $`E`$. A construction of $`s`$ mentioned above, that is, a choice of an open covering $`U_\alpha =M`$, flat sections $`s_\alpha `$ over $`U_\alpha `$, a partition of unity $`f_\alpha =1`$ with $`suppf_\alpha U_\alpha `$, so that $`s=f_\alpha s_\alpha `$, gives a bound of $`|s|`$ in terms of monodromy, as desired.
We note that by \[Golds̆tein-Kuzminov-Shvedov 1\], $`H^1(G,l^p(G))=L^pH^1(^n)`$, so to any class in $`H^1(G,l^p(G))`$ we have associated a function $`f`$ such that $`df`$ is in $`L^p`$, or, equivalently, $`_^n|f|_h^p\mu _h<\mathrm{}`$. What we in fact did above was an explicit construction of this correspondence between $`l^p`$\- and $`L^p`$-cohomology.
So far we have constructed a bounded operator $`H^1(G,l^p(G))W_p^{\frac{n1}{p}}(S^{n1})`$, $`p>n1`$. We wish to show that this operator is in fact an isomorphism of Banach spaces. To this end, we will need a Poincaré inequality in hyperbolic space.
Proposition 4.5 (Poincaré inequality in $`^n`$).— Let $`f`$ be a locally integrable measurable function with $`_^n|f|_h^p𝑑\mu _h<\mathrm{}`$. Then
1)If $`pn1`$, then $`_^n|fc|^p𝑑\mu _h<\mathrm{}`$ for some constant $`c`$;
2)If $`p>n1`$ and $`f|_{S^{n1}}`$ as an element of $`W_p^{\frac{n1}{p}}(S^{n1})`$ is zero, then $`_^n|f|^p𝑑\mu _h<\mathrm{}`$.
Proof.— A much more general theorem is contained in \[Strichartz 1\].
We now claim that $`H^1(G,l^p(G))=0`$ for $`pn1`$. This in fact follows immediately from $`H^1(G,l^p(G))=L^pH^1(^n)`$ \[Golds̆tein-Kuzminov-Shvedov 1 \] and Proposition 4.5. Now, if $`p>n1`$, then we claim that the operator $`H^1(G,l^p(G))W_p^{\frac{n1}{p}}(S^{n1})/const`$ constructed above is injective. In fact, if $`f|_{S^{n1}}=0`$, then by Proposition 4.5, $`fL^p(^n,\mu _h)`$, so $`_M\mathrm{\Sigma }_g|f(gx)|^p𝑑\mu _h<\mathrm{}`$, so for almost all $`x\stackrel{~}{M}`$, $`_g|f(gx)|^p<\mathrm{}`$. But $`f(gx)(g)l^p(G)`$, so $`l^p(G)`$ and $`[l]=0`$. Now, if $`hW_p^{\frac{n1}{p}}(S^{n1})`$, we denote by $`H`$ its harmonic extension into $`B^n`$. Then \[Uspenski 1\], \[Lizorkin 2\], $`_{B^n}\rho ^{pn}|H|^p𝑑\mu _e<h_{W_p^{\frac{n1}{p}}(S^{n1})}`$, so $`dH`$ is an $`L^p`$ 1-form on $`^n`$. This shows that the injective operator $`H^1(G,l^p(G))=L^pH^1(^n)W_p^{\frac{n1}{p}}(S^{n1})/const`$ has a bounded right inverse, so it is an isomorphism by Banach theorem. This proves Theorem 4.1.
Corollary 4.5.— Let $`G`$ be a cocompact lattice in $`SO^+(1,n)`$ and let $`:G`$ be such that $`L_gl^p(G)`$, for all $`gG(p>n1)`$. Then the limit function $`u=|_{S^{n1}}`$ belongs to $`L^q(S^{n1})`$ for all $`q>1`$. In fact,
$$\underset{1<q<\mathrm{}}{sup}\left(\frac{n1}{q}\right)^{1/q^{}}u_{L^q(S^{n1})}<\mathrm{}$$
Moreover, $`u`$ is in the linear hull of all functions $`f`$ satisfying
$$_{S^{n1}}exp(|f|^p^{})<\mathrm{}.$$
Proof. is an immediate corollary of Theorem 4.1 and the properties of the Orlicz space $`L_{\mathrm{}}(\mathrm{log}L)_a`$ and the fact that $`W_p^{(n1)/p}(S^{n1})L_{\mathrm{}}(\mathrm{log}L)_a(S^{n1})`$ for $`a1/p^{}`$, (see \[Edmunds-Triebel 1\]).
We will use this corollary in a sequel to this paper \[Reznikov 10\] in analyzing the local behaviour of the Cannon-Thurston Peano curves, corresponding to fibers of the hyperbolic 3-manifolds, fibered over the circle.
### 1.5 Growth of primitives for $`l^p`$-cocycles on the surface group
Theorem 5.1.— Let $`G`$ be a cocompact lattice in $`SO(2,1)`$ and let $`:G`$ be such that for all $`gG`$, $`L_gl^p(G)`$, $`p>1`$. Then for any word length on $`G`$,
$$|(g)|const[length(g)]^{1/p^{}}.$$
Proof follows from Theorem 4.1 and a following lemma.
Lemma 5.2.— Let $`u`$ be a harmonic function in the unit disc such that $`u|_{S^1}W_p^{1/p}(S^1)`$. Then
$$|u(z)|const[\mathrm{log}(1|z|)]^{1/p^{}}.$$
Proof of the lemma.— Here we only treat the case $`p=2`$. The full proof will be given in Section 11. Let $`u(e^{i\theta })=_na_ne^{in\theta }`$. Since $`(1+\mathrm{})^{1/4}uL^2`$, we have $`\{|n|^{1/2}a_n\}l^2()`$, therefore for $`|z|<1(b_n=|a_n|+|a_n|)`$.
$$\begin{array}{cc}\hfill u\left(z\right)a_0& _{n>0}\left(\left|a_n\right|+\left|a_n\right|\right)\left|z\right|^n\hfill \\ & =\left|n\right|^{1/2}b_n\frac{1}{\left|n\right|^{\frac{1}{2}}}\left|z\right|^n\hfill \\ & \left(\left|n\right|b_n^2\right)^{1/2}\left(\frac{1}{\left|n\right|}\left|z\right|^{2n}\right)^{1/2}\hfill \\ & const\left[\mathrm{log}\left(1\left|z\right|\right)\right]^{1/2}.\hfill \end{array}$$
### 1.6 Embedding of negatively curved manifolds and the boundaries of their universal covers
A problem of fundamental importance in topology is the following: let $`M^m\stackrel{𝜑}{}N^n`$ be a smooth $`\pi _1`$-injective embedding of manifolds of nonpositive curvature. Let $`\stackrel{~}{\phi }:\stackrel{~}{M}\stackrel{~}{N}`$ be a lift of $`\stackrel{~}{\phi }`$. Is there a limit map $`S^{m1}\stackrel{~}{M}\stackrel{\stackrel{~}{\phi }}{}\stackrel{~}{N}S^{n1}`$ and if there is, how smooth it is? For instance, let $`N^3`$ be a compact hyperbolic 3-manifold, and $`M^2`$ be an incompressible embedded surface in $`N^3`$. Then there always exists a limit continuous map $`S^1\stackrel{\stackrel{~}{\phi }}{}S^2`$. Moreover, if $`M`$ is not a virtual fiber of a fibration over the circle, then $`\stackrel{~}{\phi }(S^1)`$ is a quasifuchsian Jordan curve. If $`M`$ is a virtual fiber, then $`\stackrel{~}{\phi }:S^1S^2`$ is a Peano curve in a sense that its image fills $`S^2`$ \[Cannon-Thurston 1\]. This deep dichotomy follows from the result of \[Bonahon 1\]. We have a following very general theorem 9the embedding condition is superfluous but makes the proof more transparent):
Theorem 6.1.— Let $`M^m\stackrel{𝜑}{}N^n`$ be a smooth $`\pi _1`$-injective embedding of complete Riemannian manifolds, of pinched negative curvature. Suppose $`M`$ is compact. Let $`\stackrel{~}{\phi }:\stackrel{~}{M}\stackrel{~}{N}`$ be a lift of $`\stackrel{~}{\phi }`$. Let $`p_0\stackrel{~}{N}`$ and $`\pi :\stackrel{~}{N}\backslash \{p_0\}S^{n1}(T_{p_0}\stackrel{~}{N})`$ be a radial geodesic projection of $`\stackrel{~}{N}\backslash \{p_0\}`$ onto the unit tangent sphere. Identify $`T_{p_0}\stackrel{~}{N}`$ with $`^n`$. Let $`q_0\stackrel{~}{M}`$. Then:
1) For almost all unit tangent vectors $`vT_{q_0}(\stackrel{~}{M})`$, the restriction of $`\pi \stackrel{~}{\phi }`$ on a geodesic $`\gamma (q_0,v)`$ starting at $`q_0`$ with a tangent vector $`v`$ has an $`L^1`$-derivative as a map $`\stackrel{~}{\phi }|_{\gamma (q_0,v)}:_+^n.`$
2) For almost all $`v`$ there exists a limit $`lim_t\mathrm{}\pi \stackrel{~}{\phi }[\gamma (q_0,v)(t)]`$.
3) The resulting measurable map $`\stackrel{~}{M}S^{m1}\stackrel{\stackrel{~}{\phi }}{}S^{n1}\stackrel{~}{N}`$ does not depend on the choice of $`p_0,q_0`$.
4) If both $`M`$,$`N`$ are (real) hyperbolic, then for any $`p>n1`$, $`\stackrel{~}{\phi }`$ induces a bounded linear operator
$$\stackrel{~}{\phi }_{}:W_p^{\frac{n1}{p}}(S^{n1})W_p^{\frac{m1}{p}}(S^{m1}).$$
5) If $`M`$ is hyperbolic and $`KK(N)1`$, then for $`p>(n1)\sqrt{K}`$, $`\stackrel{~}{\phi }`$ induces a bounded linear operator
$$\stackrel{~}{\phi }_{}:C^{\mathrm{}}(S^{n1})W_p^{\frac{m1}{p}}(S^{m1})$$
for $`p>(n1)\sqrt{K}`$.
Theorem 6.2.— Let $`N^3`$ be a compact oriented hyperbolic three-manifold, let $`M^2\stackrel{𝜑}{}N^3`$ be an incompressible immersed surface, and let $`x_1,x_2,x_3`$ be Euclidean coordinates on $`S^2N^3`$. Then
1) If $`\stackrel{~}{\phi }`$ is quasifuchsian, then $`x_i\stackrel{~}{\phi }:S^1`$ are in $`W_p^{1/p}`$ for $`p2.`$
2) If $`M^2`$ is a virtual fiber then $`x_i\stackrel{~}{\phi }:S^1`$ are in $`W_p^{1/p}`$ for $`p>2`$ (but probably not in $`W_2^{1/2}`$ ).
Proof of the Theorem 6.1.— We will assume $`kK(M)1,KK(N)1`$. For $`x\stackrel{~}{N}`$ let $`r(x)=\rho (p_0,x)`$.
Lemma 6.3.— For $`r_0>0`$ and $`r(x)>r_0`$, $`|\pi (x)|const(r_0)e^{r(x)}`$, where we view $`\pi `$ as a map $`N\backslash \{p_0\}^n`$.
Proof is an immediate application of the comparison theorem, mentioned above in the proof of Proposition 2.1.
Lemma 6.4.
$$_{\stackrel{~}{N}\backslash B(p_0,r_0)}|\pi (x)|^p<\mathrm{}\text{ for }p>(n1)\sqrt{K}.$$
Proof repeats the argument in the proof of Proposition 2.1.
Now consider a tubular neighbourhood of $`M`$ in $`N`$. There exists an embedding of $`M\times IN`$ where $`I=[1,1]`$. Moreover, the restriction of the metric $`g_N`$ of $`N`$ onto $`M\times I`$ is equivalent to the product metric $`g_M+dx^2`$ (we say two metrics are equivalent if each one is bounded above by another one times a constant). It follows that there is an embedding
$$\stackrel{~}{M}\times I\stackrel{\Phi }{}\stackrel{~}{N}$$
such that $`g_{\stackrel{~}{N}}|\stackrel{~}{M}\times I`$ is equivalent to $`g_{\stackrel{~}{M}}+dx^2`$. Since $`\phi `$ is $`\pi _1`$-injective, for any $`r_0>0`$ there is $`r_1>0`$ such that if $`\rho _M(q_0,z)>r_1`$, then $`\rho _N(p_0,\mathrm{\Phi }(z,t))>r_0`$ for $`t[1,1]`$. It follows that
$$_{\stackrel{~}{M}\backslash B(q_0,r_1)\times I}|\pi \mathrm{\Phi }|^p𝑑Vol(\stackrel{~}{M})𝑑t<\mathrm{}$$
Therefore for almost all $`t_0I`$,
$$_{\stackrel{~}{M}\backslash B(q_0,r_1)}|(\pi \mathrm{\Phi }(z,t_0))|^p𝑑Vol(\stackrel{~}{M})<\mathrm{}$$
Fix such $`t_0`$ and let $`f=\pi \mathrm{\Phi }(z,t_0):\stackrel{~}{M}\backslash B(q_0,r_1)^n`$. We know that
$$_{\stackrel{~}{M}\backslash B(q_0,r_1)}|f|^p𝑑Vol(\stackrel{~}{M})<\mathrm{}.$$
Expressing the integral in polar coordinates and taking into account that $`K(M)1`$ we have
$$_{S^{m1}(T_{q_0}\stackrel{~}{M})}𝑑v_{r_1}^{\mathrm{}}e^{(m1)t}|f|^p𝑑t<\mathrm{}.$$
In particular, for almost all $`vS^{m1}(T_{q_0}\stackrel{~}{M})`$,
$$_{r_1}^{\mathrm{}}e^{(m1)t}|\frac{f}{t}|^p𝑑t<\mathrm{}.$$
In other words, for such $`v`$, $`|\frac{f}{t}|e^{\frac{(m1)}{p}t}L^p[r_1,\mathrm{}]`$, therefore $`|\frac{f}{t}|L^1[r_1,\mathrm{}]`$, since $`e^{\frac{m1}{p}t}L^p^{}[r_1,\mathrm{}]`$. This proves 1). The statements 2) and 3) follow directly.
Now suppose $`K(M)=K(N)=1`$. Let $`uW_p^{\frac{n1}{p}}(S^{n1})`$, $`p>n1`$. Then a harmonic extension $`g`$ of $`u`$ satisfies
$$_{\stackrel{~}{N}}|g|^p<\mathrm{}$$
as we know from \[Lizorkin 1\], \[Uspenski 1\] and the proof of Theorem 5.1. By the argument above, there is a $`t_0I`$, such that the composite function $`g\mathrm{\Phi }(z,t_0)`$ satisfies
$$_{\stackrel{~}{M}}|(g\mathrm{\Phi }(z,t_0)|^p<\mathrm{}$$
But then the trace $`g\mathrm{\Phi }(z,t_0)|_{\stackrel{~}{M}}`$ lies in $`W_p^{\frac{m1}{p}}(S^{m1})`$. This proves part 4) of Theorem 6.1. A proof of part 5) is identical. The Theorems 6.1 and 6.2,2) are proved. To prove Theorem 6.2, 1) we notice that a restriction of any function $`uW_2^1(S^2)`$ on a quasicircle belongs to the class $`W_2^{\frac{1}{2}}(S^1)`$. This follows immediately from the invariance of $`W_2^1(S^2)`$ under quasiconformal homeomorphisms, and a fact that functions from $`W_2^1(B^2)`$ have traces in $`W_2^{\frac{1}{2}}(S^1)`$(notice that the Dirichlet energy of a function of two variables is an invariant of the conformal class of a metric).
As the reader has noticed, we could assume $`\pi _1(M)=\pi _1(N)`$, so that $`\pi _1(M)`$ acts discretely in $`\stackrel{~}{N}`$ and $`N=\stackrel{~}{N}/\pi _1(M)`$. On the other hand, the proof does not use the fact that $`M`$ is embedded, so the Theorem 6.1 stays true for (finite-to-one)immersions in $`N`$.
We will outline now, having in mind the applications in the sequel to this paper, how to study the limit maps from the point of view of ergodic theory. The results thus obtained are weaker then those proved above, but apply to non-discrete representations. Our treatment can be seen as a development of a vague remark of \[Thurston 1, 6.4.4\]. Let $`M^m`$ be a compact hyperbolic mainfold, $`\stackrel{~}{N}=^n`$ and $`\rho :\pi _1(M)Iso(\stackrel{~}{N})`$ a discrete faithful representation. Let $`N=\stackrel{~}{N}/\rho (\pi _1(M))`$. We would like to study a boundary map $`\stackrel{~}{\phi }:\stackrel{~}{M}\stackrel{~}{N}`$ where $`\phi `$ is a smooth map $`MN`$, inducing $`\rho `$.
Lemma 6.5.— There exists a $`\pi _1(M)`$ equivariant measurable map $`\psi `$ from $`\stackrel{~}{M}=S^{m1}`$ to the space of probability measures on $`\stackrel{~}{N}=S^{n1}`$.
Proof.— For any compact Riemannian manifold $`M`$, any compact metric space $`X`$ and any representation $`\rho :\pi _1(M)Homeo(X)`$, there is a $`\pi _1(M)`$-equivariant harmonic function from $`\stackrel{~}{M}`$ to the affine space of charges on $`X`$, taking values in probability measures. This simple fact is various degrees of generality has been shown in \[Furstenberg 1\], \[L.Garnett 1\], \[Kaimanovich-Vershik 1\]. If $`M`$ is hyperbolic, then the Poisson boundary of $`\stackrel{~}{M}`$ is $`\stackrel{~}{M}`$, and the result follows.
Now let $`\psi _0+\psi _c`$ be the decomposition of $`\psi `$ into atomic and non-atomic parts. Obviously, $`\psi _c`$ is also $`\pi _1(M)`$-equivariant. We claim $`\psi _c=0`$. First, $`\psi _c`$ is a $`\pi _1(M)`$-invariant function on $`\stackrel{~}{M}=S^{n1}`$, whence a constant, since $`\pi _1(M)`$ acts on $`S^{n1}`$ ergodically. So if $`\psi _c0`$ we may assume $`\psi _c`$ is a probability measure. Second, let $`G`$ be a center of gravity map from the nonatomic measures on $`\stackrel{~}{N}`$ to $`N`$ \[Furstenberg 2\]. Then $`G\varphi _c`$ is a $`\pi _1(M)`$-equivariant map from $`\stackrel{~}{M}`$ to $`N`$. In particular, $`\rho (G\psi _c(x),G\psi _c(y))`$ is a $`\pi _1(M)`$-invariant function on $`\stackrel{~}{M}\times \stackrel{~}{M}`$ whence a constant by \[Hopf 1\] and \[Sullivan 3\]. It follows easily that $`G\psi _c=const`$ which is impossible since $`\rho `$ is discrete. So $`\psi _c=0`$.
We deduce that $`\psi `$ is atomic, $`\psi (z)=_{i=1}^{\mathrm{}}m_i\delta (\psi _i(z)),m_1m_2\mathrm{}`$. Though $`\psi _i(z):\stackrel{~}{M}\stackrel{~}{N}`$ are not uniquely defined, $`m_i:\stackrel{~}{M}`$ are. It follows that $`m_i`$ are $`\pi _1(M)`$-invariant, whence constant. Since $`m_i=1`$, there is some $`i`$ such that $`m_{i+1}<m_1`$. Choose first such $`i`$. Then $`m_1=\mathrm{}=m_i`$ and we get a measurable equivariant map
$$\stackrel{~}{M}=S^{m1}\underset{i}{\underset{}{S^{n1}\times \mathrm{}\times S^{n1}}}/S_i,$$
where $`S_i`$ is the symmetric group in $`i`$ letters.
So far we did not use the fact that $`\rho `$ is discrete, but only that $`\rho (\pi _1(M))`$ does not have fixed points in $`\stackrel{~}{N}=^m`$. So:
Propostion 6.6.— Let $`M^m`$ be a compact hyperbolic manifold and let $`\rho :\pi _1(M)SO^+(1,n)`$ be such that $`\rho (\pi _1(M))`$ does not have fixed points in $`^n`$. Then there exists a $`\pi _1(M)`$-equivariant measurable map
$$S^{m1}=\stackrel{~}{M}\stackrel{𝜓}{}(\text{subsets of cardinality }i\text{ of }S^{n1}=^n)$$
for some $`i1`$.
Using cross-ratios and the ergodicity of the action of $`\pi _1(M)`$ on $`\stackrel{~}{M}\times \stackrel{~}{M}`$ , one can easily show $`i=1`$. Now to any $`x\stackrel{~}{M}`$ one associates a Poisson measure $`\mu _x`$ on $`S^{m1}`$. Its pushforward $`\psi _{}\mu _x`$ is a probability measure on $`S^{n1}`$. The pushforward of a measure by a multivalued map is defined by
$$_{S^{n1}}fd[\psi _{}\mu ]=_{S^{m1}}\underset{y\psi (x)}{}f(y)d\mu ,$$
where $`fC(S^{n1})`$.
Now under some natural conditions $`\psi _{}\mu _x`$ does not have atoms and using the baricenter map $`G`$ in $`^n`$ one can define $`s(x)=G(\psi _{}\mu _x)`$. This can easily be shown to be continuous equivariant map $`\stackrel{~}{M}\stackrel{𝑠}{}^n`$, again under some natural assumption on $`\rho `$. The multivalued map $`\psi `$ should be regarded as a weak radial limit of $`s`$, but we will not pursue this point any further.
### 1.7 The action of quasisymmetric and quasiconformal homeomorphisms on $`W_p^{(n1)/p}`$
A well known result \[Reimann 1\] characterizes quasiconformal maps between domains $`D_1,D_2`$ in $`^n,n>2`$ as those which induce an isomorphism of Banach spaces $`BMO(D_1)`$ and $`BMO(D_2)`$. We will see now that this result in case $`D_1=D_2=^n`$ is a limit as $`p\mathrm{}`$ of the following result which establishes a quasiconformal invariance of fractional Sobolev spaces $`W_p^{n/p}`$. Of special importance is the fact that the result holds for $`n=1`$ and quasisymmetric homeomorphisms of $`S^1`$. The proof of the following lemma is “almost” contained in remarks made in \[Pansu 1–3\].
Lemma 7.1.— Let $`𝒢_{n1},n2`$ be a group of quasisymmetric ($`n=2`$) or quasiconformal ($`n3`$) homeomorphisms of $`S^{n1}`$. Then for any $`p>1(n=2)`$ or $`pn1(n3)`$, $`𝒢_{n1}`$ leaves invariant a Sobolev-Slobodec̆ki space $`W_p^{n1/p}(S^{n1})`$. For any $`\mathrm{\Phi }𝒢_{n1}`$, the corresponding map
$$\mathrm{\Phi }_{}:W_p^{n1/p}(S^n)W_p^{n1/p}(S^{n1})$$
is an automorphism of the Banach space $`W_p^{n1/p}(S^{n1})`$.
Theorem 7.2.— There exists for any $`n2`$ a bounded antisymmetric polylinear map
$$\underset{n}{\underset{}{W_n^{\frac{n1}{n}}(S^{n1})/const\times \mathrm{}\times W_n^{\frac{n1}{n}}(S^{n1})/const}},$$
defined on the smooth functions by $`f_1,\mathrm{},f_n_{S^n}f_1𝑑f_2\mathrm{}𝑑f_n`$, which is invariant under $`𝒢_{n1}.`$
In particular, we have
Corollary 7.3.— There exists a representation
$$𝒢_1Sp(W_2^{1/2}(S^1)/const),$$
defined by $`\mathrm{\Phi }(f)=f\mathrm{\Phi }^1`$.
Proof of the Lemma 7.1.— We will need a result, proved for $`n=2`$ in \[Ahlfors-Beurling 1\] for $`n=3`$ in \[Carleson 1\] and for $`n4`$ in \[Tukia-Väisälä 1\]:
Theorem.— Let $`\varphi :S^{n1}S^{n1}`$ be quasisymmetric ($`n=2`$) or quasiconformal ($`n3`$). Then there exists an extension $`\stackrel{~}{\varphi }`$ of $`\varphi `$ as a homeomorphism of $`B^n`$, which is a quasiisometry of the hyperbolic metric:
$$const_2g_h\stackrel{~}{\varphi }_{}g_hconst_1g_h.$$
Now let $`fW_p^{\frac{n1}{p}}(p>n1)`$. Let $`u`$ be a harmonic function in $`B^n`$, extending $`f`$. We know that
$$|u|_h^p𝑑\mu _hconst_3f_{W_p^{\frac{n1}{p}}(S^{n1})}.$$
It follows that
$$|(u\stackrel{~}{\varphi })|_h^p𝑑\mu _hconst_4f_{W_p^{\frac{n1}{p}}}<\mathrm{},$$
and by the trace theorem,
$$u\stackrel{~}{\varphi }_{W_p^{\frac{n1}{p}}}const_5f_{W_p^{\frac{n1}{p}}},$$
which proves the theorem for $`p>n1`$. For $`p=n1,n3`$, the result is standard.
Proof of the Theorem 7.2.— Let $`f_1,\mathrm{},f_nW_n^{\frac{n1}{n}}(S^{n1})`$. Let $`u_i`$ be a harmonic extension of $`f_i`$. The result follows at once from the formula
$$_{S^{n1}}f_1𝑑f_2\mathrm{}𝑑f_n=_{B^n}𝑑u_1𝑑u_2\mathrm{}𝑑u_n$$
. Since $`|u_i|_h^n𝑑u_h<\mathrm{}`$, the integral $`_{B^n}𝑑u_1\mathrm{}𝑑u_n`$ is finite by Hölder inequality. The invariance is obvious.
Proof of the Corollary 7.3.— A formula $`<f_1,f_2>=_{S^1}f_1𝑑f_2`$ gives $`W_2^{1/2}/const`$ a structure of a symplectic Hilbert space. This means that a map
$$W_2^{1/2}/const(W_2^{1/2}/const)^{}$$
given by $`f<f,>`$ is an isomorphism (not isometry) of Hilbert spaces. By $`Sp(W_2^{1/2}/const)`$ we mean a group of invertible bounded operators which leaves this symplectic form invariant. The result now follows from Lemma 7.1 and Theorem 7.2 .
### 1.8 Boundary values of quasiconformal maps and regularity of quasisymmetric homeomorphisms
Proposition 8.1.— Let $`\varphi `$ be a quasiconformal map, defined in a neighborhood of the unit ball $`B^n`$. Then $`\varphi |_{S^n}`$ as a map $`S^n^n`$ belongs to a class $`W_n^{\frac{n1}{n}+\delta }`$ for some $`\delta >0`$. In particular if $`n=2`$ and $`\varphi (e^{i\theta })=_na_ne^{in\theta }`$ then $`|n|^{1+\delta }|a_n|^2<\mathrm{}`$. If $`\varphi `$ is just defined in $`B^n`$ then for almost all $`\alpha S^n`$ there exists a limit $`lim_{r1}\varphi (rx)`$ and $`\varphi |_{S^{n1}}W^{\frac{n1}{n}}`$. In particular, for $`n=2`$ and $`\varphi (e^{i\theta })=_na_ne^{in\theta }`$, $`|n||a_n|^2<\mathrm{}`$.
Remark.— The last statement for conformal maps is the ”Flachensatz”.
Proof.— Since $`\varphi `$ as a map $`B^n^n`$ belongs to $`W_n^1`$, the last statement follows immediatedly from the trace theorem. To prove the first, recall that $`\varphi `$ is locally in $`W_{n+\delta ^{}}^1,\delta ^{}>0`$ \[Bojarski 1\], \[Gehring 2\]. Therefore $`\varphi |_{S^{n1}}W_n^{\frac{n1}{n}+\delta }`$, again by the trace theorem.
Theorem 8.2.— Let $`\phi :S^1S^1`$ be a quasisymmetric homeomorphism. Then as a map $`S^1^2,\phi W_p^{1/p+\delta (p)},\delta (p)>0,`$ for all $`p>1`$. If $`\phi (e^{i\theta })=_na_ne^{in\theta }`$, then $`_n|n|^{p^{}/p+\delta }|a_n|^p^{}<\mathrm{}`$ for all $`1<p2`$.
Proof.— Let $`\mathrm{\Phi }:D^2D^2`$ be a quasiisometry of the hyperbolic plane, extending $`\phi `$. We know that $`\mathrm{\Phi },\mathrm{\Phi }^1`$ are Hölder in Euclidean metric. Let $`f`$ be a smooth function defined in a neighbourhood of $`D^2`$. Then for $`p>1`$
$$_{D^2}|f|_h^p\rho _e^ϵ(x,D^2)𝑑\mu _h<\mathrm{}$$
for $`ϵ>0`$ small enough (one needs $`ϵ<p1`$).
Since $`\mathrm{\Phi }`$ is a quasiisometry for the hyperbolic metric and biHölder for the Euclidean metric, we have for $`g=f\mathrm{\Phi }`$:
$$_{D^2}|g|_h^p\rho _e^\beta (y,D^2)𝑑\mu _h<\mathrm{}$$
for some $`\beta >0`$. Rewriting in Euclidean terms, we have
$$_{D^2}|g|_e^p[\rho (y,D^2)]^{p\beta 2}<\mathrm{},$$
therefore $`g|_{S^1}W_p^{\frac{1}{p}+\delta }`$ by the trace theorem for weighted Sobolev spaces. Letting $`f`$ be an Euclidean coordinate function, we get $`\phi W_p^{\frac{1}{p}+\delta }`$. The last statement follows from Young-Hausdorff theorem.
Remark 8.3.— It had been a famous problem in fifties if $`\phi `$ is absolutely continuous ( that is, in $`W_1^1`$). Though the answer is well-known to be negative, we see that $`\phi `$ is as close to be absolutely continuous as one wishes. We will use Theorem 8.2 in a sequel to this paper to prove the existence of the vacuum vector for quantized moduli space for $`p>1`$. We also notice that the argument above together with the proof of Theorem 6.1 shows the following: if $`\phi :MN`$ is an $`\pi _1`$-injective immersion of hyperbolic manifolds, $`M`$ compact, such that for $`g\pi _1(M)`$ and some fixed $`z_0\stackrel{~}{N}`$
$$\rho (z_0,\phi _{}(g)z_0)constlength(g),$$
then $`\stackrel{~}{\phi }`$ is of class $`W_p^{(m1)/p+\delta }`$ and therefore Hölder continuous. It is not enough, though, to prove a continuity if the Cannon-Thurston curve. See \[Reznikov 10\] for futher study.
### 1.9 Teichmüller spaces and quantization of the mapping class group, I
We denote $`ap_g`$ the mapping class group of genus $`g`$ and $`ap_{g,1}`$ the extended mapping class group. If $`\mathrm{\Sigma }^g`$ is a closed oriented surface of genus $`g`$, $`\mathrm{\Gamma }_g=\pi _1(\mathrm{\Sigma }^g)`$ then $`ap_{g,1}=Aut(\mathrm{\Gamma }_g)`$ and one has an exact sequence
$$1\mathrm{\Gamma }_gap_{g,1}ap_g1.$$
Proposition 9.1(Quantization of the moduli space)— For any $`p>1`$ there exists a representation
$$ap_{g,1}\stackrel{\pi _p}{}Aut(W_p^{1/p}(S^1)/const)$$
given by the formula
$$\pi _p(\phi )(f)=f\mathrm{\Phi }^1,$$
where $`\mathrm{\Phi }`$ is a quasisymmetric homeomorphism of $`S^1`$, induced by $`\phi `$ and a choice of a hyperbolic structure in $`\mathrm{\Sigma }^g`$. For $`p=2`$ the representation
$$\pi _2:ap_{g,1}Aut(W_2^{1/2}(S^1)/const)$$
is symplectic, that is, $`\pi _2(ap_{g,1})Sp(W_2^{1/2}(S^1)/const).`$
Proof.— Fix a hyperbolic structure on $`\mathrm{\Sigma }`$. Then by a classical theorem of Nielsen, one gets a representation $`ap_{g,1}𝒢_1`$. The theorem now follows from Theorem 7.1.
Now let $`G\stackrel{\pi _0}{}PSL_2()`$ be a Fuchsian group, possibly infinitely generated. We recall that a Teichmüller space $`𝐓(G)`$ is defined as follows : points of $`𝐓(G)`$ are discrete representation $`G\stackrel{𝜋}{}PSL_2()`$, defined up to conjugation by an element of $`PSL_2()`$, which are quasiconformally conjugate to $`\pi _0`$, that is, there is a quasisymmetric homeomorphism $`\mathrm{\Phi }`$ of $`S^1`$ such that $`\pi =\mathrm{\Phi }\pi _0\mathrm{\Phi }^1`$. Notice that this definition is equivalent to the standard one by a result of \[Douady-Earle 1\].
Corollary 9.2.— Let $`\pi _0,\pi `$ be two discrete representation of a group $`G`$. Then if $`\pi `$ lies in the Teichmüller space of $`\pi _0`$, then the unitary representations
$$G\stackrel{\pi _0}{}PSL_2()\stackrel{𝛽}{}U(W_2^{1/2}(S^1)/const)$$
and
$$G\stackrel{𝜋}{}PSL_2()\stackrel{𝛽}{}U(W_2^{1/2}(S^1)/const)$$
are unitarily equivalent.
Remark 1.— The fact that $`PSL_2()`$ acts in $`W_2^{1/2}(S^1)/const`$ by unitary operators (with respect to the complex structure given by the Hilbert transform) is well-known \[Nag 1\]. In fact, this unitary representation belongs to the discrete series and may be realized in $`L^2`$ holomorphic 1-forms in $`B^2`$.
Proof.— Since $`\pi =\mathrm{\Phi }\pi _0\mathrm{\Phi }^1`$ and $`𝒢_1`$ act in $`W_2^{1/2}(S^1)/const`$, we get an invertible operator $`A`$ such that $`\beta \pi =A\beta \pi _0A^1`$. By polar decomposition $`A=UP`$ where $`P`$ is positive self-adjoint, $`U`$ is unitary, $`P`$ commutes with $`\beta \pi _0`$ and $`U`$ intertwines $`\beta \pi _0`$ and $`\beta \pi `$, as desired.
The following special case is very important. Let $`\pi _0:GPSL_2()`$ be a Fuchsian group corresponding to a Riemann surface of finite type (that is, a torsion-free lattice in $`PSL_2()`$). Let $`\mathrm{\Sigma }=^2/G`$ and let $`\phi ap(\mathrm{\Sigma },x_0),x_0\mathrm{\Sigma }`$. Let $`\mathrm{\Phi }`$ be a quasisymmetric homeomorphism of $`S^1`$ which is the boundary value of a quasiconformal homeomorphism $`\mathrm{\Psi }`$ of $`(\mathrm{\Sigma },x_0)`$, representing $`\phi `$. Then
$$\pi _0\phi ^1=\mathrm{\Phi }\pi _0\mathrm{\Phi }^1.$$
Let $`A_\phi :W_2^{1/2}/constW_2^{1/2}/const`$ be an invertible operator, representing $`\phi `$. Let $`P_\phi ^2=A_\phi ^{}A_\phi `$.Then $`P_\phi `$ commutes with $`\beta \pi _0`$. We obtained the following
Theorem 9.2.— Let $`\pi _0:GPSL_2()`$ be a torsion-free lattice. Let $`\mathrm{\Sigma }=^2/G,x_0\mathrm{\Sigma },\phi ap(\mathrm{\Sigma },x_0)`$, $`\mathrm{\Psi }`$ a quasiconformal homeomorphism inducing $`\phi `$, $`\mathrm{\Phi }`$ the trace of its lift to $`^2`$ on $`S^1`$, $`A_\phi `$ an invertible opreator in $`W_2^{1/2}(S^1)/const`$ given by $`A_\phi (f)=f\mathrm{\Phi }^1`$. Then a self-adjoint bounded operator
$$P_\phi ^2=A_\phi ^{}A_\phi $$
commutes with $`\beta \pi _0`$. If $`P_\phi ^2=\lambda 𝑑E_\lambda `$ is the spectral decomposition then $`E_\lambda `$ commute with $`\beta \pi _0`$.
Remarks.
1) If $`G`$ is cocompact, then we know that $`W_2^{1/2}/constH^1(G,l^2(G))`$, so $`W_2^{1/2}/const`$ is a Hilbert module over the type II factor defined by $`G`$ of dimension $`dim_GW_2^{1/2}/const=L^2b_1(G)=2g2`$.
2) In practice, finding $`A_\phi `$ is difficult. The reason is that $`\mathrm{\Phi }`$ is not a diffeomorphism, so the explicit formulae of Chapter 2 do not make sense. Moreover, $`\mathrm{\Phi }`$ is given in a very implicit way as a boundary value of a quasiconformal map, defined by a quadratic differential on $`^2`$ which is $`G`$-invariant!
We will now show that for $`p>2`$ the operator $`A_\phi `$ shows very unusual properties, from the point of view of functional analysis.
Theorem 9.3.— Let $`G`$ be a fundamental group of a closed hyperbolic surface $`\mathrm{\Sigma }^g`$. Let $`\phi ap(\mathrm{\Sigma },x_0)`$ be such that its image in $`ap(\mathrm{\Sigma })`$ is pseudo-Anosov. Let $`A`$ be the operator, representing $`\phi `$ in $`W_p^{1/p}(S^1),p>2`$. Then there is an element $`0vW_p^{1/p}(S^1)`$ such that
$$\underset{k}{}A_\phi ^k(v)^p<\mathrm{}.$$
Proof.— Let $`M`$ be a mapping torus of $`\mathrm{\Psi }`$, that is , $`\times \mathrm{\Sigma }/`$ where $`1`$ acts by $`(t,x)(t+1,\mathrm{\Psi }(x))`$. Then $`M`$ is hyperbolic \[Thurston 2\]. We will view $`M`$ as a fibration over a circle $`/`$ with coordinate $`t,0t<1`$; the fiber over $`t`$ will be called $`\mathrm{\Sigma }_t`$. We can trivialize $`M\stackrel{𝜓}{}/`$ over $`I=[0,1/2]`$ so that $`(t,x_0),0t1/2`$ will be a horizontal curve. Let $`g`$ be the hyperbolic metric on $`M`$ and $`g_0`$ be some hyperbolic metric on $`\mathrm{\Sigma }`$, then $`g`$ and $`g_0+dt^2`$ are equivalent on $`\mathrm{\Sigma }\times [0,1/2]\psi ^1([0,1/2])M`$. Lifting to $`\stackrel{~}{M}=^3`$, we get a fibration $`^3\stackrel{\stackrel{~}{\psi }}{}`$ with $`\stackrel{~}{\psi }^1(t)=\stackrel{~}{\mathrm{\Sigma }}_t`$. Let $`GPSL_2()`$ be the mondromy element, corresponding to $`\phi `$. Let $`f:^3`$ be such that
$$_^3|f|^p𝑑\mu _h<\mathrm{}.$$
We then have
$$\underset{k}{}_{G^k(\stackrel{~}{\psi }^1[0,1/2])}|f|^p𝑑\mu _h_^3|f|^p𝑑\mu _h<\mathrm{};$$
on the other hand the left hand side is
$$\underset{k}{}_{\stackrel{~}{\psi }^1[0,1/2]}|(fG^k)|^p𝑑\mu _hconst_0^{\frac{1}{2}}𝑑t_{\stackrel{~}{{\scriptscriptstyle }}_t}\underset{k}{}|(fG^k)|^pdVol(g_0).$$
It follows that for some $`t_0`$,
$$_{\stackrel{~}{\mathrm{\Sigma }}_{t_0}}\underset{k}{}|(fG^k)|^pdVol(g_0)<\mathrm{}.$$
Since $`g_0`$ is a hyperbolic metric, for any function $`F`$ on $`\stackrel{~}{\mathrm{\Sigma }}`$
$$_{\stackrel{~}{\mathrm{\Sigma }}}|F|^p𝑑Vol(g_0)=constF|\stackrel{~}{\mathrm{\Sigma }}_{W_p^{1/p}(S^1)/const}^p,$$
actually, we may let the LHS be a definition of the norm in $`W_p^{1/p}(S^1)/const`$, making the constant equal one. So
$$\underset{k}{}fG^k|\stackrel{~}{\mathrm{\Sigma }}_{t_0}_{W_p^{1/p}(S^1)/const}^p<\mathrm{}.$$
We will now identify $`fG^k|\stackrel{~}{\mathrm{\Sigma }}_{t_0}`$. We have a boundary map
$$\stackrel{~}{\mathrm{\Sigma }}_{t_0}=S^1\stackrel{𝛼}{}S^2=^3.$$
We know that $`G^k\alpha =\alpha \phi ^k`$, so $`fG^k=A_\phi ^kf`$ and finally
$$\underset{k}{}A_\phi ^k(f|\stackrel{~}{\mathrm{\Sigma }}_t)_{W_p^{1/p}(S^1)/const}^p<\mathrm{}.$$
Now, for any $`uW_p^{2/p}(S^2)`$ we can take $`f`$ its harmonic extension. In particular, any smooth function $`u`$ will do. Since $`\alpha :S^1S^2`$ is continuous and nonconstant, we can take $`u`$ such that $`V=u\alpha `$ is nonconstant. Then
$$\underset{k}{}A_\phi ^kv_{W_p^{1/p}(S^1)/const}^p<\mathrm{},$$
as desired.
We remark that $`_kA_\phi ^kv^p<\mathrm{}`$ will hold for all $`v`$ which are in the image of the bounded operator
$$W_p^{2/p}(S^2)W_p^{1/p}(S^1),$$
induced by $`\stackrel{~}{\mathrm{\Sigma }}^3`$.
Corollary 9.4.— Suppose that the space of fixed vectors of $`A_\phi `$ acting in $`W_p^{1/p}/const`$ possesses a complementary invariant subspace $`W`$. Then the spectre of $`A_\phi `$ in $`W`$ satisfies
$$\sigma (A_\phi |W)S^1\varphi .$$
Proof.— Suppose the opposite, then $`W=W_+W_{}`$ such that $`A_\phi ^k|W_+`$ and $`A_\phi ^k|W_{}`$ are strict contractions for some $`k>0`$. But then
$$\underset{k}{}A_\phi ^kv^p=\mathrm{}$$
for all $`vW_p^{1/p}/const`$.
We now turn to a generalization. Let $`\stackrel{~}{G}ap_{g,1}`$ be a subgroup, which containes $`\pi _1(\mathrm{\Sigma }_g)`$, so that we have an extension
$$1\pi _1(\mathrm{\Sigma }_g)\stackrel{~}{G}G1.$$
Notice that $`Gap_g`$. A well-known problem in hyperbolic topology is : when there exists a fibration
$$\begin{array}{ccc}\hfill \mathrm{\Sigma }& & Q\hfill \\ & & \hfill \\ & & T\hfill \end{array}$$
with $`\pi _1(Q)=\stackrel{~}{G}`$ such that $`Q`$ is a compact manifold of negative curvature. In case $`T`$ is a closed surface, a corollary F.3 to the Theorem F.1 of \[Reznikov 9\] provided some necessary condition. This condition is unfortunately void, as we will show now.
Theorem 9.5.— Let $`\mathrm{\Sigma }^{g_1}Q\mathrm{\Sigma }^{g_2}`$ be a surface fibration over a surface ($`\mathrm{\Sigma }^{g_i}`$ are hyperbolic and oriented ). Let $`\mathrm{\Sigma }`$ be a section of this fibration. Then
$$|\mathrm{\Sigma }\mathrm{\Sigma }|2g_22.$$
Proof.— Let $`\xi `$ be a vertical tangent bundle for $`\mathrm{\Sigma }`$, $`e(\xi )`$ its Euler class, then $`\mathrm{\Sigma }\mathrm{\Sigma }=(e(\xi ),[\mathrm{\Sigma }])`$. We have a natural homomorphism $`\pi _1(Q)ap_{g_1,1}`$ and a composite homomorphism
$$\pi _1(\mathrm{\Sigma })\pi _1(Q)ap_{g,1},$$
which we call $`\phi `$. An inclusion $`ap_{g_1,1}𝒢_1`$ induces an Euler class $`ϵ`$ in $`H^2(ap_{g,1})`$ coming from the action of $`𝒢_1`$ on $`S^1`$. By \[Matsumoto-Morita 1\], \[Morita 2\], $`\phi ^1ϵ=e(\xi )`$. Moreover, as is well known (and obvious ) $`ϵ`$ is a bounded class, in fact, for any homomorphism $`\pi _1(\mathrm{\Sigma }^g)\stackrel{𝜑}{}Homeo(S^1),|(\phi ^{}ϵ,[\mathrm{\Sigma }])|2g2`$. This proves the theorem.
Remarks.
1) If the fibration $`Q\mathrm{\Sigma }^{g_2}`$ is holomorphic and the action of $`\pi _1(\mathrm{\Sigma }^{g_2})`$ on $`H_1(\mathrm{\Sigma }^{g_1},)`$ is simple, then a famous inequality of Arakelov \[Arakelov 1\] reads $`\mathrm{\Sigma }\mathrm{\Sigma }<0`$ for all holomorphic sections. By Theorem 9.5,
$$(2g_22)<\mathrm{\Sigma }\mathrm{\Sigma }<0.$$
We now have a following result, which seems to be a very strong restriction on $`G`$.
Theorem 9.6.— Let $`1\pi _1(\mathrm{\Sigma }^g)\stackrel{~}{G}G`$ be an extension. Suppose $`\stackrel{~}{G}`$ is a fundamental group of a compact manifold of negative curvature
$$KK(Q^n)1.$$
Then for $`p>(n1)\sqrt{K}`$ there is a vector $`constvW_p^{1/p}(S^1)`$, such that
$$\underset{gG}{}A_gv_{W_p^{1/p}/const}^p<\mathrm{}$$
$`().`$
Proof.— Since the proof is essentially identical to the proof of Theorem 9.3, we will only indicate the differences. Let $`q_0\stackrel{~}{Q}`$ and let $`u:S^{n1}(T_{q_0}\stackrel{~}{Q})`$ be a smooth function. Composing with a geodesic projection $`\stackrel{~}{Q}\backslash \{0\}S^{n1}(T_{q_0}\stackrel{~}{Q})`$ we arrive to a function $`f:\stackrel{~}{Q}\backslash B(q_0,r)`$ with $`_{\stackrel{~}{Q}}|f|^p𝑑Vol<\mathrm{}`$ for $`p>(n1)\sqrt{K}`$. Since $`\mathrm{\Sigma }`$ is embedded in $`Q`$, one has a limit map $`\stackrel{~}{\mathrm{\Sigma }}=S^1S^{n1}=\stackrel{~}{Q}`$ be Theorem 6.1. Let $`v=u\alpha `$, where we identified $`\stackrel{~}{Q}`$ and $`S^{n1}(T_{q_0}\stackrel{~}{Q})`$. Then $`vW_p^{1/p}(S^1)`$ by Theorem 6.2. As in Theorem 9.3 we have the inequality $`()`$. Finally, if $`v=const`$, for any choice of $`u`$, then $`\alpha `$ is almost everywhere a constant map, say to $`zS^{n1}`$. Since $`\alpha `$ is equivariant, it follows that $`\pi _1(\mathrm{\Sigma }_g)`$ stabilizes $`z`$. This is obviously impossible.
### 1.10 Spaces $`_{k,\alpha }^{(n1)}`$ and cohomology with weights
In this section we will describe a limit form of Theorem 4.1 when $`p=1`$, and discuss $`l^{n1}`$-cohomology with weights of cocompact lattices in $`SO^+(1,n)`$.
Let $`G`$ be a finitely generated group, $`w:GR_+`$ a function such that $`w(g)\mathrm{}`$ as $`length(g)\mathrm{}`$. Consider a space $`l^p(G,w)`$ defined by $`fl^p(G,w)`$ iff $`_g|f(g)|^pw^1(g)<\mathrm{}`$. Suppose $`L_gw/w=O(1)`$ for all $`gG`$, and the same for $`R_gw/w`$. Then $`l^p(G,w)`$ becomes a $`G`$-bimodule.
Example 1.— If $`r(g)=length(g)`$ then consider $`w(g)=r^\alpha (g),\alpha >0`$ or $`w(g)=r(g)^\alpha \mathrm{log}r(g)\mathrm{log}\mathrm{log}r(g)\mathrm{}\underset{k}{\underset{}{\mathrm{log}\mathrm{log}\mathrm{}\mathrm{log}}}r(g),\alpha >0`$.
2.— Consider $`w(g)=e^{\alpha r(g)},\alpha >0`$.
Now let $`G`$ be a cocompact lattice in $`SO^+(1,n)`$, We know by Theorem 5.1, that $`H^1(G,l^p(G))0`$ exactly for $`p>n1`$. In particular, $`H^1(G,l^{n1}(G))=0`$. However, by introducing of weights the situation is changed.
Theorem 10.1.— Let $`G`$ be a cocompact lattice in $`SO^+(1,n)`$, then for any $`k1`$ and $`\alpha >0`$,
$$H^1(G,l^p(G,w))0$$
for $`p=n1`$ and $`w=r(g)\mathrm{log}r(g)\mathrm{}(\underset{k}{\underset{}{\mathrm{log}\mathrm{log}\mathrm{}\mathrm{log}}}r(g))^\alpha ,\alpha >1,k1`$.
Proof.— essentially repeats the argument of Proposition 2.1. Let $`u:S^{n1}`$ be any smooth function and denote again by $`u`$ its harmonic extension in $`B^n`$. We have $`|u|_e<const`$, therefore
$$|u|_h(z)<const\rho _e(z,S^{n1})^1$$
Let $`(h)=u(h^1z_0)`$, then a direct computation shows that $`L_gl^{n1}(G,w)`$ and $`constl^{n1}(G,w)`$ so $`l(g)=L_g`$ is a nontrivial cocycle if $`u`$ is one of the coordinate functions on $`S^{n1}`$, as in Theorem 2.1.
We would like to compute $`H^1(G,l^{n1}(G,w))`$. A construction of Theorem 4.1 produces from any class in $`H^1(G,l^p(G,w))`$ a function in $`L_w^1(^n)`$, where the latter space is defined as a space of locally integrable function $`f`$ with distributional derivatives such that
$$_^n|f|^{n1}w^1(z)<\mathrm{}$$
$`()`$
where $`w(z)=\rho _h(z_0,z)\mathrm{log}\rho _h(z_0,z)\mathrm{}(\underset{k}{\underset{}{\mathrm{log}\mathrm{log}\mathrm{}\mathrm{log}}}\rho _h(z_0,z))^\alpha `$.
Definition.— A space $`_{k,\alpha }^{(n1)}`$ is defined as a Banach space of traces of $`L_w^1(^n)`$ on $`S^{n1}`$. A norm in $`_{k,\alpha }^{(n1)}`$ is defined as infinum of integrals $`()`$ taken over the set of all functions $`f`$ with a given trace.
Remark.— The norm just defined depends on $`z_0`$. Therefore a natural action of $`SO^+(1,n)`$ in $`_{k,\alpha }^{(n1)}`$ is not isometric.
We will describe $`_{k,\alpha }^1`$ as a Zygmund-type space. One can analogously describe $`_{k,\alpha }^{(n1)}`$ for $`n>2`$, of course, but we will not need it.
Theorem 10.2.$`_{k,\alpha }^1`$ consists of all function $`u:S^1`$ for which $`(a>0)`$
$$_0^a𝑑h_0^{2\pi }\frac{|u(x+h)u(x)|}{h^2\mathrm{log}h\mathrm{}\underset{k}{\underset{}{\mathrm{log}\mathrm{}\mathrm{log}}}^\alpha h}<\mathrm{}.$$
Proof is a word-to-word repetition of Uspenski’s argument in \[Uspenski 1\]. One does not need to use Hardy’s inequality, since $`p=1`$.
Theorem 10.3.$`𝒢_1`$ acts on $`_{k,\alpha }^1`$ by
$$A_\mathrm{\Phi }u(x)=u\mathrm{\Phi }^1.$$
Corollary 10.4.— If $`\mathrm{\Phi }:S^1S^1`$ is quasisymmetric, then as a function $`S^1^2`$, $`\varphi _{k,\alpha }^1`$.
We suggest the reader to compare this result to \[Carleson 2\] and \[Gardiner-Sullivan 1\]
Proof.— Let $`\psi :B^2B^2`$ be a quasiisometry of the hyperbolic metric, extending $`\mathrm{\Phi }`$. If $`u`$ satisfies $`()`$ then $`u\mathrm{\Phi }^1`$ satisfies $`()`$ as well, whence the result.
Embedding $`ap_{g,1}𝒢_1`$ we obtain a representation
$$ap_{g,1}Aut(_{k,\alpha }^1),$$
which is a limit case of Theorem 9.1.
### 1.11 Bicohomology and the secondary quantization of the moduli space
We will now introduce a very important notion of bicohomology spaces which to an extent linearize 3-dimensional topology.
Definition.— Let $`G`$ be a finitely generated group. For $`p>1`$ define
$$_p(G)=H^1(G_r,H^1(G_l,l^p(G)),$$
where $`r`$ and $`l`$ stand for the right and left action, respectively.
Proposition 11.1.— A group $`Out(G)`$ of outer automorphism of $`G`$ acts naturally in $`_p(G)`$.
Proof.— By definition, $`Out(G)=Aut(G)/(G/Z(G))`$. Obviously $`Aut(G)`$ acts on $`H^1(G_l,l^p(G))`$ extending the right action of $`G`$, so $`Aut(G)/(G/Z(G))`$ will act on $`H^1(G_r,(H^1(G_l,l^p(G)))`$.
For a surface group $`\pi _1(\mathrm{\Sigma }_g)`$ we write $`_{p,g}=_p(G)`$.
Theorem 11.2.— There exists a natural representation
$$ap_gAut(_{p,g}).$$
Moreover, for $`p>1`$, $`_{p,g}`$ is a nontrivial Banach space. For $`p=2`$, $`_{p,g}`$ is an infinite-dimensional Hilbert space. There is a pairing
$$_{p,g}\times _{p^{},g},$$
which is $`ap_g`$-invariant. For $`p=p^{}=2`$ this pairing is a nondegenerate symmetric bilinear form. One has therefore a representation
$$ap_g\stackrel{𝜓}{}O(\mathrm{},m),0m\mathrm{},$$
which we call a secondary quantization of the moduli space of Riemann surfaces.
The proof of the theorem will occupy the rest of this section.
For a compact oriented manifold $`M`$ let $`\mathrm{\Omega }^{1/p}`$ be a space of measurable $`1/p`$-powers of densities such that for $`\omega \mathrm{\Omega }^{1/p}`$
$$_M|\omega |^p<\mathrm{}.$$
Then $`\mathrm{\Omega }^{1/p}`$ is Banach, and for $`p=2`$, Hilbert. Let $`G`$ be a finitely generated group acting in $`M`$.
Lemma 11.3.— Suppose that any element $`gG`$ has finitely many repelling points, say $`x_1^{},\mathrm{},x_n^{}`$ and finitely many attractive points, say $`x_1^+,\mathrm{},x_m^+`$ such that for any set of neighbourhoods $`U_i^{},U_+^+`$ of $`x_i^\pm `$, there is $`N`$ such that for $`kN`$, $`g^k(M\backslash U_i^{})U_i^+`$. Suppose there are $`g_1,g_2,g_3,g_4G`$ such that $`U_{i,s}^{}U_{i,s}^+`$ are disjoint for different $`s=1,2,3,4`$. Then the action of $`G`$ in $`\mathrm{\Omega }^{1/p}`$ does not have almost-invariant unit vectors.
Proof.— Suppose the opposite, then there is a sequence $`\omega _j\mathrm{\Omega }^{1/p}`$, $`\omega _j=1`$ and $`g_s^k\omega _j\omega _j\underset{j\mathrm{}}{}0`$ for all $`s,k`$. Choose $`k_s,U_{i,s}^\pm `$ such that
$$g_s^{k_s}(M\backslash U_{i,s}^{})U_{i,s}^+$$
and $`U_{i,s}^{}`$(respectively $`U_{i,s}^+`$ ) don’t intersect for different $`i`$. Let $`\omega `$ be such that $`\omega =1`$ and
$$g_s^{k_s}(\omega )\omega <(2/3)^{1/p}(1/3)^{1/p}.$$
For $`EM`$ define $`C(E,\omega )=_E|\omega |^p`$. We claim that
$$C(M\backslash U_{s,i}^{}\backslash U_{s,i}^+,\omega )<2/3.$$
Suppose the opposite, then by the invariance of the density $`|\omega |^p`$,
$$\begin{array}{c}C(M\backslash U_{s,i}^{}\backslash U_{s,i}^+,\omega g_s^{k_s})\\ =C(g_s^{k_s}(M\backslash U_{s,i}^{}\backslash U_{s,i}^+),\omega )\\ C(g_s^{k_s}(M\backslash U_{s,i}^{}),\omega )1/3.\end{array}$$
It follows that
$$\begin{array}{c}\left[_{M\backslash U_{s,i}^{}\backslash U_{s,i}^+}\left|\omega \omega g_s^{k_s}\right|^p\right]^{1/p}\\ \left|\left[_{M\backslash U_{s,i}^{}\backslash U_{s,i}^+}\left|\omega ^p\right|\right]^{1/p}\left[_{M\backslash U_{s,i}^{}\backslash U_{s,i}^+}\left|\omega g_s^{k_s}\right|^p\right]^{1/p}\right|\left(2/3\right)^{1/p}\left(1/3\right)^{1/p},\end{array}$$
a contradiction.
So $`C(U_{s,i}^{},\omega )+C(U_{s,i}^+,\omega )1/3`$. Since $`U_{s,i}^\pm `$ are disjoint for different $`s`$, we get
$$1\underset{s=1}{\overset{4}{}}C(U_{s,i}^{},\omega )+C(U_{s,i}^+,\omega )4/3,$$
a contradiction. This proves the lemma.
Corollary 11.4.— Let $`GSO^+(1,n)`$ be a cocompact lattice. Then the natural isometric action of $`G`$ in $`W_p^{(n1)/p}(S^{n1})`$ does not have almost-invariant vectors. In particular, $`H^1(G,W_p^{(n1)/p}(S^{n1}))`$ is Banach for $`p>(n1)`$.
Proof.— For $`uW_p^{(n1)/p}(S^{n1})/const`$ , let $`f`$ be a harmonic extension of $`u`$ so that
$$u=_^n|f|^p.$$
Since the energy density $`|f|^pd\mu _h`$ is invariant under isometries of $`^n`$, the proof of the Lemma 11.3 applies directly.
Corollary 11.5.$`_{p,g}`$ is Banach (Hilbert for $`p=2`$ ).
Proof.$`H^1(G_l,l^p(G))=W_p^{1/p}(S^1)/const`$.
We now describe the pairing
$$_{p,g}\times _{p^{},g}.$$
This is given by the cup-product in cohomology
$$H^1(G_r,H^1(G,l^p(G)))\times H^1(G_r,H^1(G_l,l^p^{}(G)))H^2(G_r,H^2(G_l,l^p(G)l^p^{}(G)))$$
followed by the duality $`l^p(G)\times l^p^{}(G)`$ and evaluating twice on the fundamental cycle in $`H_2(G,)`$. We have also an analytic description, namely a pairing
$$W_p^{1/p}(S^1)/const\times W_p^{}^{1/p^{}}(S^1)/const$$
is given on smooth function by $`f,g_{S^1}f𝑑g`$ and then extended as a bounded bilinear form. This induces a pairing
$$H^1(G,W_p^{1/p}/const)\times H^1(G,W_p^{}^{1/p^{}}/const).$$
Lemma 11.6.\[Korevaar-Schoen 1\]— Let $`G`$ be a finitely presented group which is realized as a fundamental group of a compact Riemannian manifold $`M`$. Let $`\rho :GO()`$ be an orthogonal representation, which does not have almost-invariant vectors. Let $`[l]H^1(G,)`$. Let $`E`$ be a flat vector bundle with fiber $``$ over $`M`$, corresponding to $`\rho `$. Then there is a harmonic 1-form $`\omega \mathrm{\Omega }^1(M,E)`$, corresponding to $`[l]`$.
Proof.—This is a reformulation of \[Korevaar-Schoen 1\].
Corollary 11.7.— Let $`M`$ be Kähler. Then if $`\rho `$ is as in the previous lemma, then
1) There is a natural complex structure in $`H^1(G,)`$, making it a complex Hilbert space ;
2) A pairing
$$H^1(G,)\times H^1(G,),$$
given by $`[l_1],[l_2]([\omega ]^{n1}([l_1],[l_2]),[M])`$ where $`[\omega ]`$ is a Kähler class, $`[M]`$ is the fundamental class and $`([l_1],[l_2])H^2(G,)`$ is a cup-product composed with the scalar product $`\times `$, is a non-degenerate symplectic structure in $`H^1(G,)`$.
Proof is the same as for finite-dimensional $``$, once we have the Hodge theory of the previous lemma.
We now ready to prove that the symmetric pairing
$$_{2,g}\times _{2,g}$$
is nondegenerate. Realize $`G`$ as a lattice in $`SO(1,2)`$. Then $`H^1(G,l^p(G))=W_p^{1/p}(S^1)/const`$. Let $`H`$ denote the Hilbert transform. It is a bounded operator
$$H:L^p(S^1)/constL^p(S^1)/const(p>1)$$
defined as follows: for $`uL^p(S^1)`$ let $`f`$ be its harmonic extension and $`g`$ a conjugate harmonic function, then $`Hu=g|S^1`$. Since
$$_^2|f|^p=_^2|g|^p.$$
$`H`$ restricts to $`W_p^{1/p}(S^1)`$ as an isometry.
Now, the symplectic duality $`f𝑑g`$ in $`W_2^{1/2}(S^1)/const`$ is simply equal to $`(Hf,g)`$. Moreover, $`H`$ is $`SO(1,2)`$-invariant. Then the Corollary 11.7 implies that the pairing of Theorem 11.2 is also nondegenerate.
We still have to prove that $`_{g,p}0`$ and for $`p=2`$ is infinite-dimensional. We first describe an element of $`_{g,p}`$ associated to a given realization $`GSO(1,2)`$ as a cocompact lattice,which we will call a principal state.
Recall that if $`M`$ is a smooth compact oriented manifold, $`𝒟iff^1(M)`$ a group of orientation-preserving diffeomorphism of class $`C^1`$, then one has a cocycle $`l`$ in $`Z^1(𝒟iff^1(M),C^0(M))`$ defined as \[Bott 1\]
$$l(\varphi )=\mathrm{log}\frac{\varphi _{}\mu }{\mu },$$
where $`\mu `$ is any smooth density on $`M`$, and $`\varphi _{}\mu `$ a left action. The class $`[l]H^1(𝒟iff^1(M),C^0(M))`$ does not depend on $`\mu `$. For $`r1`$ one similarly gets a class in $`H^1(𝒟iff^r(M),C^{r1}(M))`$. Now, let $`M=S^{n1}`$ and consider a standard conformal action of $`SO^+(1,n)`$ on $`S^{n1}`$. We get a class
$$[l]_pH^1(SO^+(1,n),W_p^{(n1)/p}(S^{n1})/const)$$
for all $`p>1`$ simply because $`C^{\mathrm{}}(S^{n1})W_p^{(n1)/p}(S^{n1})`$. We claim $`[l]_p0`$ for $`n=2`$ and $`p>n1`$. Since the action is isometric, it follows from the following lemma (we prove and use it only for $`n=2`$).
Lemma 11.8.— Fix $`z_0B^n`$ and let $`r(g)=\rho _h(z_0,g^1z_0)`$. Then for any fixed $`\mu `$, $`l(g)_{W_p^{(n1)/p}(S^{n1})/const}\mathrm{}`$ as $`r(g)\mathrm{}`$.
Proof.(Only for $`n=2`$)— We choose for $`\mu `$ the harmonic (Poisson) measure $`\mu _0`$, associated with $`z_0`$. Then $`l(g)=\mathrm{log}\frac{g_{}\mu _0}{\mu _0}`$. For $`\beta S^{n1}`$, $`l(g)(\beta )=B_\beta (z_0,gz_0)`$ where $`B_\beta (z_0,)`$ is a Busemann function of $`B^n`$ corresponding to $`\beta B^n`$ and normalized at $`z_0`$, that is, $`B_\beta (z_0,z_0)=0`$ (see, for example, \[Besson-Courtois-Gallot 1\]).
We will make the computation only for $`n=2`$. Let $`z_0=0`$, $`gz_0=w`$, then
$$B_\beta (0,w)=\mathrm{log}\frac{1|w|^2}{|w\beta |^2}.$$
Notice that $`\mathrm{log}|\frac{\beta w}{1\overline{w}\beta }|=0`$, since $`|\beta |=1`$, so
$$B_\beta (0,w)=\mathrm{log}(1|w|^2)2\mathrm{log}|1\overline{w}\beta |=2\mathrm{log}|1\overline{w}\beta |(\text{mod const}).$$
Notice that $`\mathrm{log}|1\overline{w}z|`$ is defined and is harmonic in $`|z|1`$, so
$$\begin{array}{c}B_\beta (0,w)_{W_p^{1/p}\left(S^1\right)/const}^p=2^p_{B^2}[\left(\mathrm{log}\right|1\overline{w}z)|)]_h^pd\mu _h=\\ =2^p_{B^2}\frac{\left|w\right|^p}{\left|1\overline{w}z\right|^p}\frac{1}{\left(1\left|z\right|^2\right)^{2p}}𝑑z𝑑\overline{z}\end{array}$$
$`().`$
Sublemma.— An integral $`()`$ grows as $`\mathrm{log}(1|w|)`$ as $`|w|1`$.
Proof.— Computing in polar coordinates, we have
$$_0^1𝑑r\frac{1}{(1r^2)^{2p}}_0^{2\pi }\frac{d\theta }{|1r|w|e^{i\theta }|^p}.$$
It is elementary to see that the inner integral grows as $`\frac{1}{(1r|w|)^{p1}}`$, so we arrive at
$$_0^1𝑑r\frac{1}{(1r)^{2p}}\frac{1}{(1r|w|)^{p1}}_0^a\frac{ds}{s^{2p}(A+s)^{p1}}$$
where $`a>0`$ is fixed and $`A=1|w|`$. Further we have ($`s=At`$)
$$_0^{a/A}\frac{dt}{t^{2p}(1+t)^{p1}}_0^{a/A}\frac{dt}{t}\mathrm{log}|A|,$$
which proves the Sublemma.
Finally,
$$B_\beta (0,w)_{W_p^{1/p}(S^1)/const}[\mathrm{log}(1|w|)]^{1/p},$$
where $``$ means that the ratio converges to a constant.
The proof for $`n>2`$ will be given elsewhere.
Notice that for $`p=2`$ we have (for $`n=2`$)
$$l(g)_{W_2^{1/2}(S^1)/const}g^{1/2}$$
where $`g`$ is a hyperbolic length of a (pointed) geodesic loop, representing $`g`$. This exponent in the RHS is the maximal possible. We will later prove a general theorem (Theorem III.3.1) showing that for any orthogonal or unitary representation of $`G=\pi _1(\mathrm{\Sigma })`$ in a Hilbert space $``$ and any cocycle $`lZ^1(G,)`$,
$$l(g)const\mathrm{𝑙𝑒𝑛𝑔𝑡ℎ}(g)^{1/2}\mathrm{log}\mathrm{log}\mathrm{𝑙𝑒𝑛𝑔𝑡ℎ}(g)$$
as $`g`$ converges nontangentially to almost all $`\theta S^1=G`$.
Coming back to principal states $`[l]_pH^1(G,W_p^{(n1)/p}(S^{n1})/const)`$, let $`E`$ be a flat affine bundle over $`M=^n/G`$ with fiber $`W_p^{(n1)/p}(S^{n1})`$, associated with an affine action
$$gR_g+l(g).$$
Notice that
$$s:zlog\frac{\mu (z)}{\mu (z_0)}$$
is an $`G`$-equivariant section of the lift of $`E`$ on $`\stackrel{~}{M}=^n`$, or, equivalently, defines a section of $`E`$. We claim that this section is harmonic. This immediately reduces to a statement that $`B_\beta (z_0,z)`$ is harmonic mod const as a function of $`z`$. In the upper half-plane model it simply means that $`(x,y)\mathrm{log}y`$ is harmonic mod const. The harmonic section just defined does not lift to a harmonic section of the flat affine bundle with fiber $`W_p^{(n1)/p}(S^{n1})`$. For $`n=2`$ we can say more. Let
$$H:W_p^{1/p}(S^1)/constW_p^{1/p}(S^1)/const$$
be the Hilbert transform. It makes $`W_p^{1/p}(S^1)/const`$ into a complex Banach space. Then a direct inspection shows that the section of $`E`$ defined above is (anti)holomorphic (depending on the choice of a sign of $`H`$). This will be used later. Equivalently, $`ds`$ is an (anti)-holomorphic one-form on $`^2/G`$, valued in $`E`$. Again, this holomorphic form does not lift to a $`d`$ and $`\delta `$-closed form of a flat bundle with fiber $`W_p^{1/p}(S^1)`$ even for $`p=2`$. This latter bundle is a flat bundle with fiber a Hilbert space, but whose monodromy is not orthogonal. The Hodge theory of \[Korevaar-Schoen 1\] and \[Jost 1\] does not apply and in fact not every cohomology class is represented by a $`d`$ and $`\delta `$-closed form. We will discuss these subtle obstructions to the Hodge theory in a sequel to this paper \[
Reznikov 10\].
As an application of the computation made above, we will complete the proof of Lemma 5.6 for $`p>1`$. Let $`uW_p^{1/p}(S^1)/const`$ and let $`f:B^2`$ be a harmonic extension of $`u`$. We claim that
$$|f(w)|c[\mathrm{log}(1|w|)]^{1/p^{}}.$$
Since the Hilbert transform is invertible in $`W_p^{1/p}(S^1)/const`$, we can assume that the Fourier coefficients $`\widehat{u}(n)=0`$ for $`n<0`$, so that $`f`$ is holomorphic:
$$\begin{array}{cc}& \left|f\left(w\right)\right|=\left|\frac{1}{2\pi i}_{S^1}\frac{u\left(\xi \right)d\xi }{\xi w}\right|=\frac{1}{2\pi }\left|_0^{2\pi }\frac{u\left(e^{i\theta }\right)e^{i\theta }}{e^{i\theta }w}𝑑\theta \right|=\hfill \\ \hfill =& \frac{1}{2\pi }\left|_0^{2\pi }\frac{u\left(e^{i\theta }\right)d\theta }{1we^{i\theta }}\right|=\frac{1}{2\pi }\left|_0^{2\pi }\frac{u\left(e^{i\theta }\right)d\theta }{1we^{i\theta }}\right|=\hfill \\ \hfill =& \left|\frac{1}{2\pi i}_0^{2\pi }\frac{\left[u\left(e^{i\theta }\right)e^{i\theta }\right]ie^{i\theta }d\theta }{1we^{i\theta }}\right|=\left|\frac{1}{2\pi iw}_0^{2\pi }\left[u\left(e^{i\theta }\right)e^{i\theta }\right]\left[\mathrm{log}\left(1we^{i\theta }\right)\right]^{}𝑑\theta \right|=\hfill \\ \hfill =& |\frac{1}{2\pi iw}<u\left(e^{i\theta }\right)e^{i\theta },\mathrm{log}(1we^{i\theta })>|\hfill \\ \hfill & \frac{1}{2\pi \left|w\right|}u\left(e^{i\theta }\right)e^{i\theta }_{W_p^{1/p}\left(S^1\right)/const}\mathrm{log}\left(1we^{i\theta }\right)_{W_p^{}^{1/p^{}}/const}\hfill \\ \hfill & cu_{W_p^{1/p}\left(S^1\right)/const}\left|\mathrm{log}\left(1\left|w\right|\right)\right|^{1/p^{}}.\hfill \end{array}$$
It is very plausible that the result is true, for $`uW_p^{\frac{n1}{p}}(S^{n1})/const`$ for $`n3`$. Our proof obviously does not work in this case.
We now start to prove that $`_{2,g}`$ is infinite-dimensional. Let $`M_0,M_0^{}`$ be factors generated by the left (respectively, right) actions of $`G`$ in $`l^2(G)`$ \[Murray-von Neumann 1\]. Notice that $`H^1(G_l,l^2(G))`$ can be viewed as a cohomology of a complex
$$l^2(G)\stackrel{d_0}{}\underset{i=1}{\overset{2g}{}}l^2(G)\stackrel{d_1}{}l^2(G)$$
$`(),`$
computed from the standard CW-decomposition of $`\mathrm{\Sigma }^g`$ with one zero-dimensional cell, $`2g`$ one-dimensional cells and one two-dimensional cell. Notice that $`d_0,d_1`$ are given by matrices with entries in $`[G]`$, acting on $`l^2(G)`$ from the left. Letting $`\mathrm{\Delta }_l=d_0d_0^{}+d_1^{}d_1`$ we can view $`H^1(G,l^2(G))`$ as $`Ker\mathrm{\Delta }_l`$. Notice that $`\mathrm{\Delta }_lM_0`$. It follows that $`H^1(G,l^2(G))`$ is a module over $`M_0^{}`$. Now, since $`M_0`$ is type II, there is a decomposition
$$W_2^{1/2}(S^1)/const=H^1(G_l,l^2(G))=Ker\mathrm{\Delta }_l=\underset{j=1}{\overset{m}{}}H_m,$$
for any $`m1`$ where $`H_m`$ are isomorphic right $`G`$-modules. Since we know already that $`H^1(G_r,W_2^{1/2}(S^1)/const)0`$, and $`H_j`$ are all isomorphic, it follows that $`H(G_r,H_j)0`$ for all $`j`$, therefore $`dimH^1(G,W_2^{1/2}(S^1))/constm`$. This finally proves Theorem 11.2.
There are natural invariant von Neumann algebras acting in $`_{2,g}`$. Indeed, let $`M_1^{}`$ be a double commutant of $`M_0^{}`$ in $`H^1(G,l^2(G))=Ker\mathrm{\Delta }_l`$ and $`M_1`$ be a commutant of $`M_0^{}`$. We could define $`M_1^{}`$ as a von Neumann algebra , generated by the right action of $`G`$ in $`H^1(G,l^2(G))`$ and $`M_1`$ as a commutant of $`M_1^{}`$. It follows that $`M_1,M_1^{}`$ do not depend on the choice of the complex $`()`$ and therefore $`ap_{g,1}=Aut(G)`$ acts in $`H^1(G,l^2(G))`$ leaving $`M_1^{},M_1`$ invariant. Now consider $`_{2,g}=H^1(G_r,H^1(G_l,l^2(G))`$. Then $`_{2,g}=Ker\mathrm{\Delta }_r:_{i=1}^{2g}H^1(G_l,l^2(G))_{i=1}^{2g}H^1(G_l,l^2(G))`$ where a right Laplacian is defined exactly as the left one. It follows that $`_{2,g}`$ is a module over $`M_1`$. Let $`M_2`$ be a double commutant of $`M_1`$ and $`M_2^{}`$ be its commutant. We have proved a following theorem, except for the last statement.
Theorem 11.9.— There are infinite-dimensional von Neumann algebras $`M_2,M_2^{}`$ acting in $`_{2,g}`$, which are invariant under the action of $`ap_g`$. Moreover, there is an involution $`\tau `$ of $`_{2,g}`$ which commutes with the $`ap_g`$-action and permutes $`M_2,M_2^{}`$.
Proof.— Everything is already proved except for the last statement. Notice that there is an involution $`\tau :l^2(G)l^2(G)`$ defined by $`\tau f(g)=f(g^1)`$. A Lyndon-Serre-Hochschild spectral sequence of the extension $`1GG\times GG1`$ shows that $`_{2,g}=H^2(G\times G,l^2(G))`$. Let $`\sigma `$ be an involution of $`G\times G`$ defined by $`\sigma (g,h)=(h,g)`$. Then one has $`\tau [(g,h)v]=(\sigma (g,h))\tau (v)`$ where $`g,hG`$ and $`vl^2(G)`$. It follows that $`\tau `$ induces an involution, which we also call $`\tau `$, in $`_{2,g}`$, which obviously commutes with $`ap_g`$-action and permutes $`M_2`$ and $`M_2^{}`$. This completes the proof of Theorem 11.9.
Note that since the unitary representation of $`G`$ in $`H^1(G_l,l^2(G))=W_2^{1/2}(S^1)/const`$ extends to an irreducible representation of $`PSL_2()`$, the commutator $`M_1`$ of $`G`$ in $`W_2^{1/2}(S^1)/const`$ possesses a faithful trace defined by
$$tr(a)Id=_{PSL_2()/G}gag^1𝑑g.$$
Proposition 11.10.— Let $`\stackrel{~}{}_{2,g}`$ be a completion of $`M_1`$ under the norm $`trxx^{}`$. Then $`\stackrel{~}{}_{2,g}`$ is a Hilbert space and there is a representation
$$\stackrel{~}{\rho }:ap_gAut(\stackrel{~}{}_{2,g}),$$
leaving invariant a nondegenerate form $`xtrx^2`$.
I don’t know at the time of writing if $`\stackrel{~}{}_{2,g}`$ is isomorphic to $`_{2,g}`$ as $`ap_g`$-module.
We now turn to the holomorphic realization of $`_{2,g}`$. Fix a realization of $`G`$ as a cocompact lattice in $`SO^+(1,2)`$, then $`_{2,g}=H^1(G,W_2^{1/2}(S^1)/const)`$. Recall that $`G`$ commutes with the Hilbert transform in $`W_2^{1/2}(S^1)/const`$. Let $`S=^2/G`$, then $`S`$ is a hyperbolic Riemann surface, homeomorphic to $`\mathrm{\Sigma }^g`$. For any element $`wH^1(G,W_2^{1/2}(S^1)/const)`$ we have by Lemma 11.3 and Lemma 11.6 a unique harmonic form in a flat Hilbert bundle $`E`$ with fiber $`W_2^{1/2}(S^1)/const`$, associated with the action of $`G`$.
Uniqueness should be explained. We have a following general fact.
Lemma 11.11.— Let $`M`$ be a compact Riemannian manifold. $`\rho :\pi _1(M)O(H)`$ an orthogonal representation in a real Hilbert space, without fixed vectors, $`\omega H^1(\pi _1(M),H)`$. Then there at most one harmonic form, $`\omega \mathrm{\Omega }^1(M,E)`$, representing $`\omega `$.
Proof.— If $`\omega _1,\omega _2`$ are two such forms, then $`\omega _1\omega _2`$ is a derivative of a harmonic section of $`M`$. But standard Bochner vanishing theorem shows that such section should be self-parallel, so $`\rho `$ has a fixed vector, a contradiction.
Notice that $`H`$ makes $`W_2^{1/2}(S^1)/const`$ into a complex Hilbert space. Then $`\frac{1}{2}(\omega H(\omega J))`$, where $`J`$ is a complex structure on $`S`$, will be a holomorphic 1-form in $`E`$, whereas $`\frac{1}{2}(\omega +H(\omega J))`$ will be an anti-holomorphic 1-form. Let $`_{2,g}^\pm `$ be the spaces of holomorphic(respectively, anti-holomorphic) 1-forms in $`E`$, then $`_{2,g}=_{2,g}^+_{2,g}^{}`$. Now, $`W_2^{1/2}(S^1)/const`$ is identified with exact $`L^2`$-harmonic 1-forms in the hyperbolic plane $`^2`$, which is isomorphic as a complex Hilbert space ( with a complex structure, defined by the Hodge star operator) to the space of exact $`L^2`$-holomorphic 1-form in $`^2`$. So any element in $`_{2,g}^+`$ defines a holomorphic 1-form on $`S`$ valued in a bundle with fibers $`L^2`$-holomorphic 1-forms on $`^2`$. In other words, let $`G`$ act diagonally in $`^2\times ^2`$ and
$$Q=^2\times ^2/G,$$
then we have an $`L^2`$ holomorphic 2-form on $`Q`$. The space $`_{2,g}^+`$ therefore is identified with the space of $`L^2`$ holomorphic 2-forms on $`Q`$. Similarly, $`_{2,g}^{}`$ is identified with the space of $`L^2`$ holomorphic 1-form on
$$Q^{}=^2\times \overline{^2}/G,$$
where $`\overline{^2}`$ is obtained from $`^2`$ by reversing the complex structure (i.e. $`\overline{J}=J`$). Notice that as complex surfaces, $`Q`$ and $`Q^{}`$ are not biholomorphic: $`Q`$ contains a compact curve (the quotient of the diagonal) whereas $`Q^{}`$ does not. We have proved the following:
Theorem 11.12.(Holomorphic realization of quantum moduli space)— Fix an embedding $`GSO^+(1,2)`$ as a cocompact surface, then $`_{2,g}`$ splits as $`_{2,g}^+_{2,g}^{}`$ where $`_{2,g}^+`$ (respectively, $`_{2,g}^{}`$) is identified with a space of $`L^2`$ holomorphic 2-forms on $`Q=^2\times ^2/G`$ (respectively, $`Q^{}=^2\times \overline{^2}/G`$). Moreover, the splitting is orthogonal with respect to the canonical symmetric scalar product in $`_{2,g}`$ and the restriction of this scalar product on $`_{2,g}^\pm `$ is positive (respectively, negative).
Example.— The principal state $`[l]_2`$ lies in $`_{2,g}^{}`$. We do not know at the time of writing if $`_{2,g}^+=0.`$
### 1.12 $`_{p,g}`$ as operator spaces and the vacuum vector
In this section we will develop an algebraic and an analytic theory of $`_{p,g}`$ as spaces of operators between $`W_q^{1/q}(S^1)/const`$, which commute with the action of $`G`$. We use rather rough estimates of matrix elements, so the ranges of indices for which the action is established is certainly not the best possible. We start with a lemma.
Lemma 12.1.— Let $`uW_p^{1/p}(S^1)`$ and $`al^q(G)`$. Then $`a(g)R_guW_r^{1/r}(S^1)`$ where $`\frac{1}{p}+\frac{1}{q}1=\frac{1}{r}`$.
Remark.$`R_g`$ means the action of $`G`$ in $`W_p^{1/p}(S^1)`$ (a reminiscent of the actions from the right in $`l^2(G)`$).
Proof.— Let $`f`$ be a harmonic extension of $`u`$, so that
$$_^2|f|^p𝑑\mu _h<\mathrm{}.$$
By Young-A.Weil inequality \[Hewitt-Ross 1\], $`l^pl^ql^r`$, so if $`h=a_gR_gf`$, we have
$$\begin{array}{ccc}\hfill _^2\left|h\right|^r𝑑\mu _h& =& _{^2/G}𝑑\mu _h\left(z\right)_g\left|h\left(g^1z\right)\right|^r\hfill \\ & & _{^2/G}𝑑\mu _h\left(z\right)h\left(g^1z\right)_{l^r}\hfill \\ & & c_{^2/G}𝑑\mu _h\left(z\right)a\left(g\right)_{l^q\left(G\right)}f\left(g^1z\right)_{l^p}\hfill \\ & & ca\left(g\right)_{l^q\left(G\right)}_^2\left|f\right|^p𝑑\mu _h.\hfill \end{array}$$
The result follows with an estimate
$$a(g)R_gu_{W_r^{1/r}(S^1)/const}ca(g)_{l^q(G)}u_{W_p^{1/p}(S^1)/const}.$$
Now recall that we have a canonical pairing
$$B:W_r^{1/r}(S^1)/const\times W_r^{}^{1/r^{}}(S^1)/const,$$
so that a formula
$$(u,v)(aB(a(g)R_gu,v))$$
defines a map
$$W_p^{1/p}(S^1)/const\times W_r^{}^{1/r^{}}/constl^q^{}(G).$$
Now an element of $`W_r^{}^{1/r^{}}(S^1)/const`$ defines an element of
$$Hom_G(W_p^{1/p}(S^1)/const,l^q^{}(G))$$
and an induced map
$$\begin{array}{c}Hom(H^1(G,W_p^{1/p}\left(S^1\right)/const,H^1(G,l^q^{}\left(G\right))=\\ =Hom(H^1(G,W_p^{1/p}\left(S^1\right)/const,W_q^{}^{1/q^{}}\left(S^1\right)/const).\end{array}$$
In other words, we have a map
$$H^1(G,W_p^{1/p}(S^1)/const)Hom_{}(W_r^{}^{1/r^{}}(S^1)/constW_q^{}^{1/q^{}}(S^1)/const),$$
and it is immediate to check that the image lies in $`Hom_G`$. So we have a
Proposition 12.2.— The construction above defines a map
$$_{p,g}Hom_G(W_r^{}^{1/r^{}}(S^1)/const,W_q^{}^{1/q^{}}(S^1)/const)$$
for $`p,q^{},r^{}`$ satisfying $`\frac{1}{p}1+\frac{1}{q^{}}\frac{1}{r^{}}`$, which is $`ap_g`$-equivariant.
An induced map in $`H^1(G,)`$ produces a bounded $`ap_g`$-equivariant product
$$_{p,g}\times _{r^{},g}_{q^{},g}.$$
We stress again that the range of indices for which this product is defined should be improved. We will see that viewing $`_{p,g}`$ as an operator space helps to understand $`ap_g`$-action. We turn now to an analytic description of the above. Let $`lZ^1(G,W_p^{1/p}(S^1)/const)`$. A construction of Theorem 4.1 produces a smooth map
$$F:^2W_p^{1/p}(S^1)/const,$$
satisfying $`F(g^1z)=R_gF(z)+l(g),gG`$, In particular, $`g^{}(F)(g^1z)=R_g(F)(z)`$. Now let $`vW_r^{}^{1/r^{}}(S^1)/const`$, where $`rp`$. Then we have a scalar function
$$<F,v>:^2,$$
where $`<,>`$ is a pairing $`W_r^{1/r}(S^1)/const\times W_r^{}^{1/r^{}}(S^1)/const`$ defined in Theorem 7.2. Since $`rp`$, $`W_p^{1/p}(S^1)W_r^{1/r}(S^1)`$, so $`<F,v>`$ is defined. Without futher assumption one can only say that
$$|<F,v>|const,$$
but if we assume $`r>p`$, say $`\frac{1}{p}=1+\frac{1}{q^{}}\frac{1}{r^{}}`$, then $`<F,v>`$ will satisfy
$$\underset{g}{}|<F,v>(gz)|^q^{}<const$$
for all $`z^2`$. Integrating over $`^2/G`$, we get
$$_^2|<F,v>|^q^{}<\mathrm{},$$
so there exists $`<F,v>|S^1W_q^{}^{1/q^{}}(S^1)`$. This defines a desired map
$$H^1(G,W_p^{1/p}(S^1)/const)Hom_G(W_r^{}^{1/r^{}}(S^1)/const,W_q^{}^{1/q^{}}(S^1)/const).$$
We will use this description now to compute the operator, associated with the principal state
$$[l]_pH^1(G,W_p^{1/p}(S^1)/const).$$
Proposition 12.3.— For $`p,r^{},q^{}>1,\frac{1}{p}1+\frac{1}{q^{}}\frac{1}{r^{}}`$, an operator in
$$Hom_G(W_r^{}^{1/r^{}}(S^1)/const,W_q^{}^{1/q^{}}(S^1)/const),$$
associated to the principal state $`[l]_p`$ is proportional to the Hilbert transform
$$H:W_r^{}^{1/r^{}}(S^1)/constW_r^{}^{1/r^{}}(S^1)/const,$$
followed by the embedding $`W_r^{}^{1/r^{}}(S^1)W_q^{}^{1/q^{}}(S^1)`$.
Proof.— First, we notice that the Hilbert transform acts as an isometric operator in $`W_p^{1/p}(S^1)/const`$ for all $`p>1`$. This follows at once from the definition of the norm as
$$u=_^2|f|^p𝑑\mu _h,$$
where $`\mathrm{\Delta }f=0`$ and $`f|S^1=u(\text{mod const})`$. We will prove the proposition by a direct unimaginative computation. Let
$$g(z)=\frac{z+z_0}{1+\overline{z}_0z},|z_0|<1,|z|<1,$$
so that $`g(0)=z_0`$. Then the Jacobian of $`g`$ on the unit circle is
$$\frac{1|z_0|^2}{|zz_0|^2},$$
so
$$l(g)=\mathrm{log}(1|z_0|^2)\mathrm{log}|zz_0|^2.$$
Let $`\phi :S^1`$ be smooth. Then
$$<\phi ,l(g)>=_{S^1}\phi ^{}(\theta )[\mathrm{log}(1|z_0|^2)\mathrm{log}|e^{i\theta }re^{i\phi }|^2]𝑑\theta ,$$
where $`z_0=re^{i\phi }`$. Obviously,
$$_{S^1}\phi ^{}(\theta )\mathrm{log}(1|z_0|^2)=0,$$
so
$$\begin{array}{ccc}\hfill <\phi ,l\left(g\right)>& =& _{S^1}\phi \left(\theta \right)\left[\mathrm{log}\left|e^{i\theta }re^{i\phi }\right|^2\right]^{}\hfill \\ & =& \phi \left(\theta \right)\frac{2r\mathrm{sin}\left(\theta \phi \right)}{1+r^22r\mathrm{cos}\left(\theta \phi \right)}.\hfill \end{array}$$
As $`z_0=re^{i\phi }\underset{r1}{}e^{i\phi }`$, this converges to
$$v.p.\phi (\theta )\frac{2\mathrm{sin}(\theta \phi )}{22\mathrm{cos}(\theta \phi )}=v.p.\phi (\theta )\frac{1}{\mathrm{𝑡𝑔}\frac{\theta \phi }{2}}=\pi H\phi (\theta )$$
almost everywhere on $`S^1`$. The Proposition is proved, since smooth functions are dense in $`W_r^{}^{1/r^{}}(S^1)`$.
Notice that since $`H`$ commutes with the action of $`SO^+(1,2)`$, for any cocycle $`mZ^1(G,W_p^{1/p}(S^1)/const)`$, $`Hm`$ is also an cocycle. In particular, $`H[l]_pH^1(G,W_p^{1/p}(S^1)/const)`$. We wish to compute a corresponding operator in $`Hom(W_r^{}^{1/r^{}}(S^1)/constW_q^{}^{1/q^{}}(S^1)/const)`$. Let $`F`$, as above, be a smooth map
$$F:^2W_p^{1/p}(S^1)/const,$$
satisfying $`F(g^1z)=R_gF(z)+l_p(g)`$. For $`vW_r^{}^{1/r^{}}(S^1)`$ we need to find a limit on the boundary of $`<HF,v>`$. But $`H`$ repects the pairing $`<,>`$ and $`H^2=1`$, so $`<HF,v>=<F,Hv>`$, whose limit on $`S^1`$ is $`\pi H(Hv)=\pi v`$. We have proved the following lemma.
Lemma 12.4.— For $`p,r^{},q^{}>1,\frac{1}{p}1+\frac{1}{q^{}}\frac{1}{r^{}}`$, an operator in
$$Hom_G(W_r^{}^{1/r^{}}(S^1)/const,W_q^{}^{1/q^{}}(S^1)/const),$$
associated with $`\frac{1}{\pi }H[l]_p`$, is the identity.
Theorem 12.5.A.— An element $`v=H[l]_2_{2,g}`$ does not depend on the choice of the lattice $`GSO^+(1,2)`$
B.— The action of $`ap_g`$ in $`_{2,g}`$ fixes $`v`$.
Remark.— The Theorem is beyond doubt true for all $`p>1`$ and not only $`p2`$, however I can’t prove this at the moment of writing this paper (July,1999). (Added January, 2000). This is in fact true. The proof will appear in \[Reznikov 10\]).
The vector $`v`$ is called a vacuum vector.
Proof.— Consider two embeddings $`i_1,i_2:GSO^+(1,2)`$ as cocompact lattices and let $`v_1,v_2`$ be corresponding elements. We view $`v_1,v_2`$ as elements of $`H^1(G_r,H^1(G_l,l^2(G))`$. Let $`A_1,A_2`$ be associated operators
$$A_1,A_2:H^1(G_l,l^r^{}(G))H^1(G_l,l^q^{}(G)).$$
We know that $`A_1=A_2=id`$. It follows that an operator, associated with $`v_1v_2`$ is zero. We are going to show that $`v_1v_2`$ is zero. Since
$$v_1v_2H^1(G_r,V),$$
where $`V`$ stands for $`H^1(G_l,l^2(G))W_2^{1/2}(S^1)/const`$, by a result of \[Korevaar-Schoen 1\] cited above (Lemma 11.6) there exists a harmonic section $`F`$ of the affine Hilbert bundle over $`M=^2/G`$ with fiber $`V`$ and nonodromy
$$gR_g()+m(g),$$
where $`m(g)`$ is any cocycle, representing $`v_1v_2`$. Let $`vW_r^{}^{1/r^{}}(S^1)/const`$, then, denoting by $`F`$ again the lift of this section on $`\stackrel{~}{M}=^2`$, we see that $`<F,v>`$ is a harmonic function such that
$$_^2|(<F,v>)|^q^{}d\mu _h<\mathrm{},$$
and the trace of $`<F,v>`$ on $`S^1`$ is constant. It follows that $`<F,v>`$ is constant itself, therefore ($`v`$ is arbitrary!) $`F=w=const`$ and
$$m(g)=R_gww,$$
so $`v_1v_2=0`$. This proves A. Now, if $`\varphi ap_{g,1}=Aut(\pi _1(\mathrm{\Sigma }^g))`$, simply apply A to $`i_1`$ and $`i_1\varphi `$.
We wish to compute $`v`$. $`[l]_2`$ is given by a cocycle
$$g2\mathrm{log}|\beta w|,$$
$`\beta S^1,w=g(0)`$. This is equal to $`2\mathrm{log}|1\overline{w}\beta |`$. The latter function is a real part of $`2\mathrm{log}(1\overline{w}z)`$ which is holomorphic in $`|z|1`$, so the Hilbert transform is $`2Arg(1\overline{w}\beta )`$. This means that a cocycle
$$m(g)(\beta )=2Arg(W\beta )(\text{ mod const})$$
where $`W=1/\overline{w},w=g(0)`$, represents $`v`$.
Theorem 12.6.$`H^1(ap_{g,1},H^1(G_l,l^p(G)))0`$ for $`p2`$.
Proof.— We embed $`G`$ as a lattice in $`SO^+(1,2)`$ and identify $`H^1(G_l,l^p(G))`$ and $`W_p^{1/p}(S^1)/const`$. We know that
$$H^0(ap_g,H^1(G_r,W_p^{1/p}(S^1)/const)v0.$$
Notice that $`H^0(G_r,W_p^{1/p}(S^1)/const)=0`$ since any $`G`$-invariant harmonic 1-form in $`^2`$ has infinte $`p`$-energy. So in the spectral sequence
$$E_{i,j}^2:H^i(ap_g,H^j(G_r,W_p^{1/p}(S^1)/const))H^{i+j}(ap_{g,1},W_p^{1/p}(S^1)/const)$$
the second differential
$$d_2:H^0(ap_g,H^1(G_r,W_p^{1/p}(S^1)/const)H^2(ap_g,H^0(G_r,W_p^{1/p}(S^1)/const)$$
must be zero. Therefore the vacuum vector $`v`$ survives in $`E^{\mathrm{}}`$.
It is plausible that, in fact,
$$H^1(𝒢_1,W_p^{1/p}(S^1)/const)0(p>1)$$
for the group $`𝒢_1`$ of quasisymmetric homeomorphisms. (Added January, 2000). This is in fact true. A formula
$$\mathrm{\Phi }Arg\mathrm{\Phi }^1(\beta )Arg(\beta )\text{mod const}$$
defines a cocycle of $`𝒢_1`$ in $`W_p^{1/p}(S^1)/const`$ for any $`p>1`$. The proof will in \[Reznikov 10\]).
### 1.13 Equivariant mapping of the Teichmüller Space, a space of quasifuchsian representations and a space of all discrete representations into $`_{p,g}`$
Theorem 13.1.A.— A map which associates to a discrete cocompact representation
$$GSO^+(1,2)$$
its principal state
$$[l]_pH^1(G_r,H^1(G_l,l^p(G))$$
is an $`ap_g`$-equivariant map of the Teichmüller space $`𝐓_{6g6}`$ to $`_{p,g}`$ for all $`p>1`$.
B.— Let $`\phi :GSO^+(1,3)`$ be a discrete representation. Let $`\alpha _\phi :S^1S^2`$ be the limit map of the boundaries
$$S^1=\stackrel{~}{\mathrm{\Sigma }}^3=S^2,$$
defined in section 6, associated to $`\phi `$. For $`p>2`$ let
$$[l]_pH^1(SO^+(1,3),W_p^{2/p}(S^2)/const)$$
be the principle state. A map
$$\phi A_\phi \phi ^{}[l]_pH^1(G_r,H^1(G_l,l^p(G))),$$
defined by first pulling back $`[l]_p`$ to $`\phi ^{}[l]_pH^1(G,W_p^{2/p}(S^2)/const)`$ and then applying the operator
$$A_\phi :W_p^{2/p}(S^2)/constW_p^{1/p}(S^1)/const,$$
induced by $`\alpha _\phi `$ and defined in section 6, is an $`ap_g`$ equivariant map
$$Hom_{discrete}(G,SO^+(1,3))/SO^+(1,3)_{p,g}$$
for all $`p>2`$.
C.— A restriction of the map, defined in $`B`$ to
$$Hom_{quasifuchsian}(G,SO^+(1,3))/SO^+(1,3)$$
is contained in $`_{2,g}`$.
Proof.— is already contained in section 6–12. We notice that from the operator viewpoint the map of A sends any realization of $`G`$ as a lattice in $`SO^+(1,2)`$ to a Hilbert transform of $`W_p^{1/p}(S^1)/const`$, followed by an identification
$$H^1(G_l,l^p(G))W_p^{1/p}(S^1)/const,$$
which depends on the lattice. In other words, fix one lattice embedding
$$\beta _0:GSO^+(1,2).$$
Then any other lattice embedding
$$\beta :GSO^+(1,2)$$
can be written as
$$\beta (g)=\mathrm{\Phi }_{\beta _0,\beta }\beta _0(g)\mathrm{\Phi }_{\beta _0,\beta }^1,$$
where $`\mathrm{\Phi }_{\beta _0,\beta }𝒢_1`$ is a quasisymmetric map. Then an operator, associated with $`\beta `$ is
$$\mathrm{\Phi }_{\beta _0,\beta }H\mathrm{\Phi }_{\beta _0,\beta }^1Aut(W_p^{1/p}(S^1)/const).$$
This gives an $`ap_g`$-equivariant map
$$𝐓_{6g6}Aut_G(W_p^{1/p}(S^1)/const)$$
For $`p=2`$ one gets a map
$$𝐓_{6g6}Sp_G(W_2^{1/2}(S^1)/const)$$
because the Hilbert transform and $`𝒢_1`$-action are symplectic(section 7), which can be described as follows. First, one embeds $`𝐓_{6g6}`$ in the universal Teichmüller space
$$𝐓=𝒢_1/SO^+(1,2).$$
Then using the representation
$$𝒢_1Sp(W_2^{1/2}(S^1)/const)$$
defined in section 7, one defines an embedding to $`Sp/U`$:
$$𝐓Sp(W_2^{1/2}(S^1)/const)/U$$
where $`U`$ is a group of operator in $`Sp`$ which commutes with $`H`$ seen as a complex structure in $`W_2^{1/2}(S^1)/const`$. Finally, one uses the Cartan embedding
$$Sp/USp.$$
Theorem 13.2.(Linearization of pseudoAnosov automorphisms )— Let $`\varphi ap_{g,1}=Aut(\pi _1(\mathrm{\Sigma }^g))`$ is a pseudoAnosov automorphism. Then for any $`p>1`$ there exists a nontrivial element $`S_p_{p,g}`$ with the following properties:
1) for $`p_1<p_2`$, $`S_{p_2}`$ is an image of $`S_{p_1}`$, under the natural map $`_{p_1,g}_{p_2,g}`$;
2) $`S_p`$ is invariant under $`\overline{\varphi }ap_g`$;
3) there is a cocycle $`\stackrel{~}{l}_pZ^1(G,W_p^{1/p}(S^1)/const)`$ , representing $`S_p`$, such that for any $`gG`$
$$\underset{n}{}\stackrel{~}{l}_p(g)\mathrm{\Phi }^n_{W_p^{1/p}(S^1)/const}<\mathrm{}$$
where $`\mathrm{\Phi }:S^1S^1`$ is a quasisymmetric homeomorphism, associated with $`\varphi `$ (or, in other words,
$$\underset{n}{}A_\phi ^m\stackrel{~}{l}_p(g)<\mathrm{}$$
where $`A_\phi Aut(H^1(G_l,l^p(G))`$ is induced by $`\varphi `$)
Proof.— is an immediate corollary of \[Thurston 2\] (see also an exposition in \[Otal 1\] ), which shows that the mapping torus of any homeomorphism $`\mathrm{\Psi }:\mathrm{\Sigma }\mathrm{\Sigma }`$, representing $`\phi `$ is a hyperbolic 3-manifold, Theorem 13.1, Theorem 9.1 and Theorem 9.3.
It is plausible that such $`S_p`$ is unique up to a multiplier. Knowing $`S_p`$ is essentially equivalent to knowing the hyperbolic volumes of all ideal simplices with vertices on the limit curve $`S^1S^2`$.
## Chapter 2 A theory of groups acting on the circle
Our first main result in this Chapter is Theorem 1.7 which says, roughly, that Kazhdan group cannot act on the circle. This general theorem draws a line after many years of study and various special results concerning the actions of lattices in Lie groups, see \[Witte 1\], \[Farb-Shalen 1\], \[Ghys 1\]. One can see here a historic parallel with a similar, but easier, general theorem of \[Alperin 1\] and \[Watatani 1\] concerning Kazhdan groups acting on trees, which also followed a study of the actions of lattices. Our technique is absolutely different from the cited papers and uses a fundamental cocycle, introduced and studied in section 1. We also use standard facts from Kazhdan groups theory \[de la Harpe-Valette 1\].
In Sections 2,3 we quantize equivariant maps between boundaries of universal covers, studied in Chapter I, Section 6. Our main tool is a harmonic map theory into infinite-dimensional spaces, as developed in \[Korevaar-Schoen 1\], see also \[Jost 1\]. In Section 4 we review some facts about Banach-Lie groups and regulators. In Section 5 we describe a series of higher characteristic classes of subgroups of $`𝒟iff^{1,\alpha }(S^1)`$. There are two construction given. One uses an extension to a restricted linear group of a Hilbert space of classes originally defined in \[Feigin-Tsygan 1\] for infinite Jacobian matrices. Another construction uses the action of a restricted symplectic group $`Sp(W_2^{1/2}(S^1)/const)`$ on the infinite-dimensional Siegel half-plane. In both construction we use an embedding of $`𝒟iff^{1,\alpha }`$ into a restricted linear group, by the unitary and symplectic representation of $`𝒟iff`$, respectively. Using the geometry of the Siegel half-plane, we prove that our classes have polynomial growth.
There is a striking similarity between the theory of this Chapter and a theory of symplectomorphism group, see Chapter IV, \[Reznikov 2\] and \[Reznikov 4\]. Note that the extended mapping class group action is not $`C^{1,\alpha }`$ smooth, so the results of this Chapter do not apply to this group. On the other hand, $`ap_g`$ does act symplectically on a smooth compact symplectic manifold.
### 2.1 Fundamental cocycle
By $`𝒟iff^{1,\alpha }(S^1)`$ we denote a group of orientation-preserving diffeomorphisms with derivative in the Hölder space $`C^\alpha (S^1)`$, which consists of functions $`f`$ such that
$$|f(x)f(y)|<c|xy|^\alpha .$$
There is a series of unitary representations of $`𝒟iff^{1,\alpha }(S^1)`$ in $`L_{}^2(S^1,d\theta )`$ given by
$$(\pi (g)(f))(x)=f(g^1x)[(g^1)^{}(x)]^{\frac{1}{2}+i\beta },\beta .$$
We will mostly consider $`\beta =0`$, in which case one has an orthogonal representation in $`L_{}^2(S^1,d\theta )`$. An invariant meaning is, of course a representation in half-densities on $`S^1`$. Now consider a Hilbert transform $`H`$ as an operator in $`L_{}^2(S^1,d\theta )`$ given by a usual formula
$$Hf(\phi )=\frac{1}{\pi }v.p._{S^1}\frac{f(\theta )}{\mathrm{𝑡𝑔}\frac{\phi \theta }{2}}𝑑\theta .$$
We wish to consider $`[H,\pi (g^1)]`$. This is a bounded operator in $`L^2(S^1,d\theta )`$ given by an integral kernel which we are going to compute. Notice that
$$\frac{1}{\mathrm{𝑡𝑔}\frac{\phi \theta }{2}}=\frac{2}{\phi \theta }+\text{smooth kernel}.$$
A computation of \[Pressley-Segal 1\] shows that
$$H[\pi (g)f](\phi )=\frac{2}{\pi }v.p._{S^1}\frac{d\theta }{\phi \theta }f(g^1(\theta ))[(g^1(\theta ))^{}]^{1/2}+\text{smooth kernel}\pi (g),$$
so
$$\begin{array}{c}\left(\pi \left(g^1\right)H\pi \left(g\right)f\right)\left(\phi \right)=\left[g^{}\left(\phi \right)\right]^{1/2}\frac{2}{\pi }v.p._{S^1}\frac{d\theta f\left(g^1\left(\theta \right)\right)\left[\left(g^1\left(\theta \right)\right)^{}\right]^{1/2}}{g\left(\phi \right)\theta }\hfill \\ +\pi \left(g^1\right)\text{smooth kernel}\pi \left(g\right)\hfill \end{array}.$$
Letting $`\theta =g(\eta )`$ we have
$$\begin{array}{c}\left(\pi \left(g^1\right)H\pi \left(g\right)f\right)\left(\phi \right)=\left[g^{}\left(\phi \right)\right]^{1/2}\frac{2}{\pi }v.p._{S^1}\frac{f\left(\eta \right)\left[g^{}\left(\eta \right)\right]^{1/2}}{g\left(\phi \right)g\left(\eta \right)}𝑑\eta +\hfill \\ +\pi \left(g^1\right)\text{smooth kernel}\pi \left(g\right)=\hfill \\ =\frac{2}{\pi }v.p._{S^1}\frac{\left[g^{}\left(\phi \right)g^{}\left(\eta \right)\right]^{1/2}}{g\left(\phi \right)g\left(\eta \right)}f\left(\eta \right)𝑑\eta +\pi \left(g^1\right)\text{smooth kernel}\pi \left(g\right)\hfill \end{array}$$
Finally,
$$\begin{array}{c}\left[\left(\pi \left(g^1\right)H\pi \left(g\right)H\right)\right]\left(\phi \right)=\hfill \\ =\frac{1}{\pi }_{S^1}\frac{\left[g^{}\left(\phi \right)g^{}\left(\eta \right)\right]^{1/2}\left(\phi \eta \right)\left(g\left(\phi \right)g\left(\eta \right)\right)}{\left(g\left(\phi \right)g\left(\eta \right)\right)\left(\phi \eta \right)}f\left(\eta \right)𝑑\eta +\pi \left(g^1\right)\text{smooth kernel}\pi \left(g\right)+\hfill \\ +\text{smooth kernel}.\hfill \end{array}$$
$`\left(1.1\right)`$
For a Hilbert space $``$ and $`p1`$ we denote by $`J_p()`$ a Shatten class of operators such that a sum of the $`p`$-th powers of their singular numbers converges. By $`J_{p+}()`$ we mean the intersection of all $`J_q()`$ with $`q>p`$.
Now recall that $`g𝒟iff^{1,\alpha }(S^1)`$. A following proposition sharpens that of \[Pressley-Segal 1\] for $`𝒟iff^{\mathrm{}}(S^1)`$:
Propositin 1.1.
A. For $`\alpha >1/2,\pi (g^1)H\pi (g)HJ_2(L^2(S^1,d\theta ))`$.
B. For $`\alpha >0,\pi (g^1)H\pi (g)HJ_{1/\alpha +}(L^2(S^1),d\theta ).`$
Proof.— As $`\phi \eta 0`$,
$$\frac{[g^{}(\phi )g^{}(\eta )]^{1/2}(\phi \eta )(g(\phi )g(\eta ))}{(g(\phi )g(\eta ))(\phi \eta )}<const(\phi \eta )^{\alpha 1},$$
so the kernel in (1.1) is in $`L^2(S^1\times S^1,d\theta d\theta )`$ for $`\alpha >1/2`$. This proves A.
To prove B we notice that by \[Pietsch 1\], the estimate on the kernel implies that the operator lies in $`𝒥_{1/\alpha +}`$. Strictly speaking, the conditions of \[Pietsch 1\] require $`C^{\mathrm{}}`$ smoothness off the diagonal, whereas we have only the Hölder continuity, but the result stays true.
Now notice that $`GL(L^2(S^1,d\theta ))`$ acts in $`J_p`$ by conjugation. We deduce the following
Proposition 1.2.— A map
$$l:g\pi (g)H\pi (g^1)H$$
is a 1-cocycle of $`𝒟iff^{1,\alpha }(S^1)`$ in $`J_p(L^2(S^1,d\theta ))`$ for $`p>1/\alpha `$. In particular, $`l`$ is a 1-cocycle of $`𝒟iff^{1,\alpha }(S^1)`$ in $`J_2`$ for $`\alpha >1/2`$.
We will call $`l`$ a fundamental cocycle of $`𝒟iff^{1,\alpha }(S^1)`$.
Now let $`G`$ be a subgroup of $`𝒟iff^{1,\alpha }(S^1)`$. We obtain a class in $`H^1(G,J_p(L^2(S^1,d\theta ))`$ by restricting $`l`$ on $`G`$. We are going to show that this class is never zero, except for completely pathological actions of $`G`$ on $`S^1`$.
Proposition 1.3.— Let $`G`$ be a subgroup of $`𝒟iff^{1,\alpha }(S^1)`$, $`0<\alpha <1`$. Suppose $`p>1/\alpha `$. If $`[l]H^1(G,J_p)`$ zero, then the unitary action of $`G`$ in $`L_{}^2(S^1,d\theta )`$ is reducible. Moreover, if $`H^1(G,J_p)=0`$ then $`L_{}^2(S^1,d\theta )`$ a direct sum of countably many closed invariant subspaces.
Proof.— If $`[l]=0`$ then there is $`AJ_p`$ such that
$$\pi (g)H\pi (g^1)H=\pi (g)A\pi (g^1)A$$
so that $`[\pi (g),HA]=0`$. Since $`H`$ has two different eigenvalues with infinitely-dimensional eigenspaces, $`HAconstId`$, so the action of $`G`$ in $`L_{}^2(S^1,d\theta )`$ is reducible.
Next, consider an operator $`R`$ in $`L^2(S^1,d\theta )`$ with a kernel
$$K(\phi ,\eta )=\frac{1}{|\mathrm{𝑡𝑔}\frac{\phi \eta }{2}|}.$$
One sees immediately that $`R`$ is a self-adjoint unbounded operator. Repeating the computation above, we deduce that $`\pi (g)R\pi (g^1)RJ_p`$, so $`\stackrel{~}{l}(g)=\pi (g)R\pi (g^1)R`$ is another cocycle. If this cocycle is trivial, then we get an unbounded self-adjoint operator $`RA`$ which commutes with the action of $`G`$. An application of the spectral theorem shows that $`L^2(S^1,d\theta )`$ is a countable sum of invariant subspaces.
Corollary 1.4.— A restriction of $`l,\stackrel{~}{l}`$ on $`SO^+(1,2)`$ is not zero, for all $`\alpha >0`$.
Proof.$`SO^+(1,2)`$ act in $`L_{}^2(S^1,d\theta )`$ as a representation of principal series, which are irreducible.
We now specialize for $`\alpha =1/2`$ and $`p=2`$. Since $`[\stackrel{~}{l}]H^1(SO^+(1,2),J_2)`$ is nonzero, $`\stackrel{~}{l}(g)_{J_2}`$ is unbounded as a function of $`g`$ \[de la Harpe-Valette 1\]. In fact, one has the following
Proposition 1.5.— Let $`\pi :SO^+(1,2)U(H)`$ be a unitary representation and let $`l:SO^+(1,2)H`$ be a continuous cocycle. Suppose $`[l]0`$. Then
A. For any cocompact lattice $`GSO^+(1,2)`$, $`[l]|G0`$.
B. $`l(g^n)`$ is unbounded as $`n\mathrm{}`$ for any hyperbolic $`g`$.
C. $`l(\gamma ^n)`$ is unbounded as $`n\mathrm{}`$ for any parabolic $`\gamma 1`$.
Proof.— Let $`VSO^+(1,2)`$ be compact and such that $`VG=SO^+(1,2)`$. For $`vV,gG`$ we have
$$l(vg)=\pi (v)l(g)+l(v),$$
so $`l(vg)l(g)+l(v)`$. If $`l|G`$ is bounded, then so is $`l`$. This proves A. Next, let $`P`$ be the image of $`SO^+(1,2)/K`$ under Cartan embedding, where $`K`$ is a maximal compact subgroup. By the same reason as above, $`l|P`$ is unbounded. Let $`S^1P`$ be a nontrivial orbit of $`K`$ in $`P^2`$. Notice that $`P`$ is closed under raising into an integral power and there is a compact $`VSO^+(1,2)`$ such that
$$P\underset{n1}{}(S^1)^nV$$
where $`(S^1)^n`$ is an image of $`S^1`$ under raising to $`n`$-th power. We deduce that $`l|_{n1}(S^1)^n`$ is unbounded. Let $`\gamma S^1`$. Then any element in $`(S^1)^n`$ is of the form $`k\gamma ^nk^1`$, $`kK`$, so
$$l(k\gamma ^nk^1)l(k)+l(k^1)+l(\gamma ^n)$$
So $`l(\gamma ^n)`$ is unbounded. Since $`\gamma `$ can be any hyperbolic element, B follows. Notice that we proved that $`l(g_k)`$ is unbounded for any sequence $`g_kP`$, which escapes all compact sets. Now let $`gSO^+(1,2)`$ be parabolic $`1`$, and let $`\tau `$ be the involution fixing $`K`$. Then $`\tau (g^n)g^nP`$ and escapes all compact sets, so $`l[(\tau g^n)g^n]`$ is unbounded. It follows that either $`l(\tau g^n)`$ or $`l(g^n)`$ is unbounded. But all parabolics are conjugate in $`SO^+(1,2)`$, so $`C`$ follows.
Proposition 1.6.— Let $`G𝒟iff^{1,\alpha }(S^1),\alpha >1/2`$. Suppose that $`G`$ contains an element $`g`$ which is conjugate in $`𝒟iff^{1,\alpha }(S^1)`$ to a hyperbolic or a nontrivial parabolic fractional-linear transformation. Then $`[l]|G0`$ in $`H^1(G,J_2)`$.
Proof.— Any such $`g`$ is conjugate in $`𝒟iff^{1,\alpha }(S^1)`$ to an element $`g^{}SO^+(1,2)`$ for which $`l(g_{}^{}{}_{}{}^{n})`$ is unbounded, so $`l(g^n)`$ is unbounded as well.
We are ready to formulate the main result of this section.
Theorem 1.7.— Let $`G𝒟iff^{1,\alpha }(S^1),\alpha >1/2`$. Suppose that either
1) a natural unitary action ($`\beta =0`$) of $`G`$ in $`L^2(S^1,d\theta )`$ given by
$$\pi (g)(f)(\phi )=f(g^1(\phi ))[(g^1(\phi ))^{}]^{1/2},$$
is irreducible or is a direct sum of finitely many irreducible factors, or
2) $`G`$ contains an element, conjugate in $`𝒟iff^{1,\alpha }(S^1)`$ to a hyperbolic fractional-linear transformation, or
3) $`G`$ contains an element, conjugate in $`𝒟iff^{1,\alpha }(S^1)`$ to a parabolic ($`1`$) fractional-linear transformation, or
4)
$$\underset{gG}{sup}_{S^1}\left[\frac{\sqrt{g^{}(\phi )g^{}(\eta )}(\phi \eta )(g(\phi )g(\eta ))}{(g(\phi )g(\eta ))(\phi \eta )}\right]^2𝑑\phi 𝑑\eta =\mathrm{}$$
Then $`G`$ is not Kazhdan.
Proof follows from the formula (II.1.1), Proposition 1.3, Proposition 1.5 and Proposition 1.6.
### 2.2 Construction of $`N=2`$ quantum fields with lattice symmetry
It is possible that the physical time-space is discrete. Correspondingly, in the axiomatic quantum field theory it is possible that the fields must yield invariance not under the whole Poincaré group, but only under a lattice in it. See \[Michailov 1\], \[Belavin 1\] in this respect. We are going to construct mathematical objects, which yield such invariance on one hand, and quantize the equivariant measurable maps considered in I.6.3, on the other.
Theorem 2.1.— Let $`G`$ be a cocompact lattice in $`SO^+(1,2)`$. Let $`=L_{}^2(S^1,d\theta )`$ with the orthogonal action $`\pi `$, corresponding to $`\beta `$. Then there exists a measurable map to the space of bounded operators
$$S^1\stackrel{𝜌}{}(H)$$
with the following properties.
1) Equivariance: for $`sS^1`$ and $`gG`$
$$\rho (gs)=\pi (g)\rho (s)\pi (g^1)$$
almost everywhere on $`S^1`$.
2)One has
$$_{S^1}(\rho (s)H)\psi (s)𝑑sJ_2$$
for $`\psi C^{\mathrm{}}(S^1)`$.
3)There exists $`JJ_2()`$ such that $`\rho (s)`$ is a weak nontangential limit
$$\rho (s)=\underset{gs}{lim}\pi (g)(H+J)\pi (g^1)$$
as $`gG`$ converges nontangentially to $`sS^1=G`$ a.e. on $`S^1.`$
Proof.— As a Hilbert space with orthogonal $`G`$-action, $`J_2=L^2(S^1\times S^1,d\theta d\theta )`$. By the proof of Lemma I.11.3, $`G`$ does not have almost invariant vectors in $`J_2`$. Let $`\mathrm{\Sigma }=^2/G`$ and let $`E`$ be a flat affine vector bundle over $`\mathrm{\Sigma }`$ with a fiber $`J_2`$ and monodromy
$$gAd\pi (g)+l(g).$$
Then by a result of \[Korevaar-Schoen 1\], and \[Jost 1\] (lemma I.11.6), there exists a harmonic map
$$\stackrel{~}{f}:^2J_2$$
satisfying
$$\stackrel{~}{f}(gx)=\pi (g)\stackrel{~}{f}(x)\pi (g^1)+l(g)$$
Consider $`f(x)=\stackrel{~}{f}(x)+H`$. Then
$$f(gx)=\pi (g)f(x)\pi (g^1),$$
in particular, $`f(x)`$ is bounded in operator norm. An operator version of Fatou theorem \[Naboko 1 and references therein \] shows that $`f`$ has nontangential limit values a.e. on $`S^1`$, say $`\rho (s)`$. Obviously, $`\rho `$ is $`G`$-invariant. On the other hand, $`\stackrel{~}{f}`$ is a Bloch harmonic $`J_2`$-valued function, that is,
$$\underset{x^2}{sup}\stackrel{~}{f}_{J_2}<\mathrm{}.$$
It follows that $`\stackrel{~}{f}(w)_{J_2}<c\mathrm{log}(1|w|)`$, $`wB^2=^2`$. This implies by a standard argument (see e.g. \[Gorbačuk 1\] that $`\stackrel{~}{f}`$ has a limit on $`S^1`$ as an element of $`𝒟^{}(S^1,J_2)`$. So for $`\psi C^{\mathrm{}}(S^1)`$,
$$_{S^1}(\rho H)\psi J_2,$$
which proves the Theorem.
Remarks.— 1)As was mentioned above, the invariant meaning of the representation $`\pi `$ is that $`L^2(S^1,d\theta )`$ should be regarded as a space of half-densities. Correspondingly, an integral operator is defined by a kernel which is a half-density on $`S^1\times S^1`$ of the type $`K(\phi ,\eta )(d\phi d\eta )^{1/2}`$. If $`K(\phi ,\eta )`$ is smooth and has a zero of second order on the diagonal $`\mathrm{\Delta }S^1\times S^1`$, then one has an invariant definition of its residue or second derivative, which is a quadratic differential. A direct computation which we leave to the reader shows that for $`g𝒟iff^{\mathrm{}}(S^1)`$
: 1) $`l(g)=\pi (g)H\pi (g^1)H`$ is given by a kernel which has a zero of second order on $`\mathrm{\Delta }`$;
2) a corresponding residue $`S(g)`$ is the Schwartzian of $`g`$.
This shows that $`l(g)`$ is a quantization of the Schwartzian cocycle. The operator field $`\rho (s)`$ of Theorem 1.8 seems therefore to be related to objects axiomatized, but not constructed, in \[Belavin-Polyakov-Zamolodchikov 1\]. 2)The Theorem and the proof stay valid for any representation
$$\phi :G𝒟iff^{1,\alpha }(S^1),$$
$`\alpha >1/2`$, such that the action on $`S^1\times S^1`$ satisfies the very mild conditions of Lemma I.11.3.
### 2.3 Construction of $`N=3`$ quantum fields with lattice symmetry
A theory developed have for $`𝒟iff(S^1)`$ does not generalize to $`𝒟iff(S^n),n2`$. The reason is that the action of $`𝒟iff(S^1)`$ on $`S^1`$ is conformal. There are two ways to generalize various aspects of the theory to higher dimensions, by either considering $`SO^+(1,n)`$ acting on $`S^{n1}`$ or, very surprisingly, a group of symplectomorphisms of a compact symplectic manifold $`M`$ (see Chapter IV). Here we consider the action of $`SO^+(1,3)PSL_2()`$ on $`S^2`$. We set $`d(x,y)`$ to be a spherical distance in $`S^2`$. Let $`d\theta `$ denote the spherical measure and let $`=L^2(S^2,d\theta )`$. For $`gSO^+(1,3)`$ let $`\mu _g(x)`$ denote a conformal factor, that is $`\mu _g^2(x)`$ is a Jacobian of $`g`$ with respect to $`d\theta `$. A formula
$$\pi (g)f(x)=f(g^1(x))\mu _{g^1}^{1+i\beta }(x),\beta ,$$
defines a unitary representation of $`SO^+(1,3)`$ in $``$. Now we introduce an operator $`H`$ with the kernel
$$K(\phi ,\theta )=\frac{1}{d^2(\phi ,\eta )}.$$
This operator is self-adjoint and unbounded. Our goal is to compute
$$\pi (g)H\pi (g^1)H=l(g).$$
Proposition 3.1.$`l(g)J_2`$ for all $`gSO^+(1,3)`$ and $`\beta =0`$.
Proof.— A direct computation. One needs to show that as $`d(x,y)0`$,
$$d^2(g(x),g(y))\mu _g(x)\mu _g(y)d^2(x,y)$$
is of order $`d^4(x,y)`$. In other words, for a fractional-linear trangformation $`g`$ of $``$ one needs to show that as $`xy`$, $`Imx,Imy>0,g(x)=x`$,
$$\left|\frac{g(x)g(y)}{g(x)\overline{g(y)}}\right|^2|g^{}(x)||g^{}(y)|\frac{Imy}{Img(y)}\left|\frac{xy}{x\overline{y}}\right|^2$$
is of order $`|xy|^4`$. This verifies the result for hyperbolic metric instead of spherical metric, which is of course equivalent. One computes directly using Taylor series for holomorphic function $`g`$.
Now arguing as in section 2 we arrive at the following result.
Theorem 3.2.— Let $`G`$ be a cocompact lattice in $`SO^+(1,3)`$. Let $`=L_{}^2(S^2,d\theta )`$ with orthogonal action of $`G`$ corresponding to $`\beta =0`$. Then there exists a harmonic map
$$^3\stackrel{𝜓}{}J_2()$$
with the property that $`z\psi (z)+H`$ is equivariant:
$$\psi (gz)+H=\pi (g)(\psi (z)+H)\pi (g^1)$$
for all $`gG`$ and $`z^3`$.
Since $`H`$ is unbounded, the boundary value of $`\psi (z)+H`$ does not exist as a measurable map to the space of bounded operators. It is possible that there is a more clever choice of a conformally natural singular integral operator which is bounded, but I don’t know how to do it. Note in this respect that there is a very different realization of an orthogonal representation of $`SO^+(1,3)`$ in the space of functions on $`S^2`$ , discovered in \[Reznikov 1\]. Namely, look at the natural action of $`SO^+(1,3)`$ on smooth half co-densities, that is, sections of $`\sqrt{\mathrm{\Lambda }^2TS^2}`$. Using the spherical metric, we can identify this space with $`C^{\mathrm{}}(S^2)`$. Then the above-mentioned action leaves invariant a nonnegative quadratic form
$$Q(f)=_{S^2}((\mathrm{\Delta }f)^22|f|^2)𝑑area$$
whose kernel consists of constants and linear functions. It is possible that there are $`G`$-equivariant quantum fields valued in operators acting in the associated Hilbert space.
### 2.4 Banach-Lie groups and regulators: an overview
A Banach-Lie group is a Banach manifold with a compatible group structure. Usual Lie theory largely extends to this case. In particular, if $`𝒢`$ is a Banach-Lie group and $`𝔤`$ its Banach-Lie algebra, then a continuous n-cocycle on $`𝔤`$ defines a left-invariant closed form on $`𝒢`$, so that one has a homomorphism
$$H_{cont}^n(𝔤,𝕂)H_{top}^n(𝒢,𝕂)$$
where $`H_{top}^{}`$ is a cohomology of a topological space. In \[Reznikov 2\] we defined $`𝕂`$-homotopy groups of a Lie algebra, so that there is a map
$$\pi _i(𝒢)𝕂\pi _i(𝔤)$$
which in the case $`𝒢=SL_n(C^{\mathrm{}}(M))`$, $`M`$ a compact manifold, $`n>>1`$, reduces to the Chern character
$$K_i^{top}(M)HC_i(C^{\mathrm{}}(M))=\mathrm{\Omega }^i(M)/d\mathrm{\Omega }^{i1}(M)H^{i2}(M,𝕂)\mathrm{}$$
($`𝒢`$ is not a Banach-Lie group but a Frechét-Lie group in this case). More interesting is a secondary class (=regulator) map. Define an algebraic $`K`$-theory of $`𝒢`$ as
$$K_i^{alg}(𝒢)=\pi _i((B𝒢^\delta )^+)$$
and the augmented $`K`$-theory as a kernel of the map $`K_i^{alg}K_i^{top}`$:
$$0\overline{K}_i^{alg}(𝒢)K_i^{alg}(𝒢)\pi _i(B𝒢)=\pi _{i1}(𝒢).$$
Then the regulator map is a homomorphism
$$r:\overline{K}_i^{alg}(𝒢)coker(\pi _i(𝒢)𝕂\pi _i(𝔤)).$$
Lifting this map to cohomology, that is , constructing a map
$$H_{cont}^{}(𝔤,𝕂)H^{}(𝒢^\delta ,𝕂)$$
meets obstructions described in the van Est spectral sequence. If $`𝒦𝒢`$ is a closed subgroup such that $`𝒢/𝒦`$ is contractible, then these obstructions vanish and one gets a map
$$H_{cont}^{}(𝔤,𝔨)H^{}(𝒢^\delta )$$
given explicitly by a Dupont-type construction \[Dupont 1\]. This is essentially the same as geometric construction of secondary classes of flat $`𝒢`$-bundles, described in \[Reznikov 3\]. In case $`𝒢=SL_n(C^{\mathrm{}}(M))`$ this gives a usual regulator map in algebraic $`K`$-theory. However, for various diffeomorphism groups one construct new interesting classes. For symplectomorphism groups two series of classes, mentioned in the Introduction to Chapter 4, were constructed in \[Reznikov 2\] and \[Reznikov 4\], and a new class associated to a Lagrangian submanifold, will be constructed in Chapter 4. The symmetric spaces for $`Sympl(M)`$, used in \[Reznikov 2\] are sort of continuous direct products of finite-dimensional Siegel upper half-planes. On the other hand, a symmetric space which we will use in this chapter to construct classes in $`H^{}(𝒟iff^{1,\alpha }(S^1))`$ is an infinite-dimensional Siegel half-plane. The trouble is, however, that, for a compact manifold $`Y`$, (say, $`S^1`$) a group of diffeomorphisms of finite smoothness, like $`𝒟iff^k(Y)`$, is not a Banach-Lie group: the multiplication from the right is not a diffeomorphism (the multiplication from the left is). This is neatly explained in \[Adams-Ratiu-Schmid 1\]. Luckily, to construct secondary classes we only use the fact that the multiplication from the left is a diffeomorphism.
### 2.5 Charateristic classes of foliated circle bundles
As is well known, the continuous cohomology of $`𝒟iff^{\mathrm{}}(S^1)`$ is generated by the Euler class and by the integrated Godbillon-Vey class \[Geldfand-Fuks 1\], \[Fuks 1 and references therein\]. Moreover, the square of the Euler class is zero. This already shows that the degree of smoothness is crucial. For if one considers the action of the extended mapping class group
$$ap_{g,1}𝒢_1Homeo(S^1),$$
then the pull-back of the Euler class has nonzero powers to a degree which goes to infinity with $`g`$ \[Miller 1\], \[Morita 1\], \[Mumford 1\]. It appears that the scarcity of the cohomology of $`𝒟iff(S^1)`$ is a consequence of an (artificial) restriction of excessive degree of smoothness. Notice that the proofs in \[Fuks 1\] depend hopelessly on $`C^{\mathrm{}}`$\- smoothness. We will give two constructions of a series of new classes in $`H^{}(𝒟iff^{1,\alpha }(S^1))`$, $`0<\alpha <1`$ using both the unitary representation in $`L^2(S^1,d\theta )`$ and the symplectic representation in $`Sp(W_2^{1/2}(S^1)/const)`$. As in the case of the powers of the Euler class, a nonvanishing of these classes is an obstruction to smoothability, i.e. to a conjugation to a subgroup of $`𝒟iff^{\mathrm{}}(S^1)`$. We will also prove that our classes are of polynomial growth if $`\alpha >1/2`$. A related result (but not the argument) for $`C^{\mathrm{}}`$ Gelfand-Fuks cohomology in all dimensions is to be found in \[Connes-Gromov-Moscovici 1\]. Both in spirit and technology, the construction of the classes in $`H^{}(𝒟iff^{1,\alpha }(S^1))`$ resembles our construction of a series of classes in $`H_{cont}^k(Sympl(M),)`$, $`k=2,6,10,\mathrm{}`$, where $`M`$ is a compact symplectic manifold and $`Sympl(M)`$ is its symplectomorphism group \[Reznikov 4\].
We start with the construction using the unitary representation. By Proposition 1.1, $`\pi (g)H\pi (g^1)HJ_p`$ where $`g𝒟iff^{1,\alpha }(S^1),p>1/\alpha `$, $`\pi `$ is a unitary action in $`L_{}^2(S^1,d\theta )`$, and $`H`$ is a complexification of the Hilbert transform. That is $`H(e^{in\theta })=sgn(n)e^{in\theta }`$. The group of $`\mathrm{\Phi }GL()`$, $`=L_{}^2(S^1,d\theta )`$ such that $`\mathrm{\Phi }H\mathrm{\Phi }^1HJ_p`$ will be denoted $`GL_{J_p}()`$, following \[Pressley-Segal 1\]. The unitary subgroup of $`GL_{J_p}()`$ is denoted $`U_{J_p}()`$. Let $`_+,_{}`$ be the eigenspaces of $`H`$ with eigenvalues $`+1`$ and $`1`$ respectively. By $`Gr_{J_p}()`$ we denote the restricted Grassmanian $`U_{J_p}/U(_+)\times U(_{})`$. Then $`Gr_{J_p}()`$ is a Banach manifold, modelled by the Banach space $`J_p`$. The Banach-Lie group $`GL_{J_p}()`$ acts smoothly on $`Gr_{J_p}()`$. On the other hand, though $`𝒟iff^{1,\alpha }(S^1)`$ is a group and a Banach manifold, it is not a Banach-Lie group \[Adams-Ratiu-Schmid 1\]. However, multiplication from the left $`L_g(h)=gh`$ is a diffeomorphism (but not a multiplication from the right). The embedding
$$𝒟iff^{1,\alpha }(S^1)U_{J_p}GL_{J_p}()$$
is not continuous. However, an induced action of $`𝒟iff^{1,\alpha }(S^1)`$ on $`Gr_{J_p}()`$ is smooth \[Pressley-Segal 1\].
We will introduce a series of $`U_{J_p}`$-invariant differential forms on $`Gr_{J_p}()`$. These forms induce cohomology classes in the Lie algebra cohomology $`H^{}(Lie(U_{J_p}))`$, extending the classes introduced in \[Feigin-Tsygan 1\] for the Lie algebra of Jacobian matrices. Notice that a tangent space to the origin of $`Gr_{J_p}()`$ can be identified with matrices of the form
$$C=\left(\begin{array}{cc}0& B\\ A& 0\end{array}\right)$$
where $`AJ_p(_+,_{})`$ and $`BJ_p(_{},_+)`$. Let $`C_1,\mathrm{},C_{2k},(k\text{ odd})`$ be a collection of such matrices. Define
$$\mu _k(C_1,\mathrm{},C_{2k})=\underset{\sigma S_{2k}}{}sgn(\sigma )P_k(\rho (C_{\sigma (1)},C_{\sigma (2)}),\mathrm{},\rho (C_{\sigma (2k1)},C_{\sigma (2k)})$$
where $`P_k`$ is the $`k`$-th invariant symmetric functions of $`k`$ matrices, which is a polarization of $`trA^k`$ (not an elementary symmetric polynomial, as in \[Fuks 1\] ). Now, $`\rho (C_1,C_2)`$ is defined as follows: let $`\pi (C)`$ is the left upper corner of $`C`$, i.e. an operator in $`(_+)`$. Then $`\rho (C_1,C_2)=\pi ([C_1,C_2])[\pi (C_1),\pi (C_2)]`$. The ”meaning” of $`\pi `$ is that of a connection of a principal bundle on something like the classifying space of the Lie algebra $`Lie(GL_{J_p})`$, and of $`\rho `$ is that of the curvature of this connection. Then $`\mu _k`$ becomes a characteristic class, somewhat analogous to the characteristic classes in the standard Chern-Weil theory. Notice that $`\mu _k`$ is defined for all $`k[1/\alpha ]+1`$. In \[Feigin-Tsygan 1\], $`\rho (C_1,C_2)𝔤𝔩(\mathrm{},𝕂)`$ and $`\mu _k`$ is defined for all $`k`$. The form $`\mu _2`$ defines the famous “Japanese cocycle”, \[Verdier 1\].
Lemma 5.1.$`\mu _k`$ is $`U_{J_p}`$-invariant and closed.
Proof.— The invariance is obvious. The proof of closedness is standard and left to the reader, see the remarks above and \[Feigin-Tsygan 1\].
Pulling back to $`𝒟iff^{1,\alpha }(S^1)`$ (this is possible by the remarks made above) we obtain a left-invariant closed differential form on $`𝒟iff^{1,\alpha }(S^1)`$. Pulling back to the universal cover $`\stackrel{~}{𝒟iff}^{1,\alpha }(S^1)`$, we obtain a left-invariant closed differential form $`\stackrel{~}{\mu }_k`$ on $`\stackrel{~}{𝒟iff}^{1,\alpha }(S^1)`$. A following theorem follows.
Theorem 5.2.— The secondary characteristic class, corresponding to $`\stackrel{~}{\mu }_k`$ is a well-defined class $`r(\stackrel{~}{\mu }_k)`$ in $`H^{2k}([\stackrel{~}{𝒟iff}^{1,\alpha }]^\delta ,)`$.
Proof.$`\stackrel{~}{𝒟iff}^{1,\alpha }(S^1)`$ is contractible.
Notice that for $`\alpha >1/2`$ the class $`\mu _1H^2(𝒟iff^{1,\alpha }(S^1))`$ is defined, which is just the integrated Godbillon-Vey class.
Our second construction uses the symplectic action. For simplicity, we only treat the case $`\alpha >1/2`$. Recall (Corollary I.7.3) that $`𝒢_1`$ acts symplectically in $`V=W_2^{1/2}(S^1)/const`$. Restricting on $`𝒟iff^{1,\alpha }(S^1)`$, we obtain a representation
$$𝒟iff^{1,\alpha }(S^1)\stackrel{𝜋}{}Sp(V).$$
Let $`H`$ be the Hilbert transform in $`V`$, normalized such that $`H^2=1`$. Denote by $`Sp_{J_p}`$ a subgroup of $`ASp(V)`$ such that $`[A,H]J_p`$. Denote $`U=U(V)`$ the unitary group of such $`A`$ that $`[A,H]=0`$. Denote
$$X=Sp_{J_p}/U$$
a restricted Siegel half-plane. This is a Banach contractible manifold \[Palais 1\]. For $`p=2`$ this is a Hilbert manifold with canonical $`Sp_{J_2}`$-invariant Riemannian metric of nonpositive curvature. The metric is defined as follows. The tangent space $`T_H(X)`$ is identified with operators $`A`$ such that $`ALie(Sp_{J_2})`$ and $`AH=HA`$. It follows that $`AJ_2`$, and $`A=A^{}`$. Then the metric is defined as $`trA^2`$. This definition is dimension-free and so the proof that the curvature is nonpositive follows from the explicit formulae, as in finitely-dimentional case.
Lemma 5.3.— For $`\alpha >1/2`$, $`\pi (𝒟iff^{1,\alpha }(S^1))Sp_{J_2}(V).`$
Proof.— We will use the computation of \[Segal 1\]. Let $`g𝒟iff^{1,\alpha }(S^1)`$. We need to show that
$$S=\underset{n,m>0}{}\frac{m}{n}\left|_{S^1}e^{i(ng(\theta )+m\theta )}𝑑\theta \right|^2<\mathrm{}.$$
As in \[Segal 1\] we have, using a trick of Kazhdan,
$$S=\underset{N=1}{\overset{\mathrm{}}{}}\underset{m=1}{\overset{N1}{}}\frac{m}{n}\left|_{S^1}e^{iN\phi }[g_\beta ^1]^{}(\phi )𝑑\phi \right|^2,$$
where $`\beta =\frac{n}{N},n=Nm,g_\beta (\theta )=\beta g(\theta )+(1\beta )\theta ,\theta S^1=/2\pi `$. For $`0\beta 1`$, $`g_\beta ^1`$ are uniformly in $`𝒟iff^{1,\alpha }(S^1)`$ with $`\alpha >1/2`$, so
$$_{S^1}e^{iN\phi }[g_\beta ^1]^{}(\phi )𝑑\phi constN^\alpha c_N$$
with $`_{N=1}^{\mathrm{}}c_N^2<\mathrm{}`$. Since $`_{m=1}^{N1}\frac{m}{n}\mathrm{log}N`$, we have
$$Sconst\underset{N=1}{\overset{\mathrm{}}{}}N\mathrm{log}NN^{2\alpha }c_N^2<\mathrm{}.$$
Now let $`k`$ be odd, $`A_1,\mathrm{},A_{2k}T_HX`$ and
$$\nu _k(A_1,A_2,\mathrm{},A_{2k})=$$
Lemma 5.4.$`\nu _k`$ is closed and $`Sp_{J_2}`$-invariant.
Proof.— is identical to the finite-dimensional case \[Borel 1\].
Theorem 5.5.— The secondary characteristic class, corresponding to $`\nu _k`$ defines an element $`r(\nu _k)`$ in $`H_{cont}^{2k}(Sp_{J_2}(V))`$ and in $`H^{2k}([𝒟iff^{1,\alpha }(S^1)]^\delta ,)`$, $`\alpha >1/2`$. All these classes are of polynomial growth.
Proof.— Only the last statement needs a proof. For $`x_0,\mathrm{},x_sX`$ denote a geodesic span $`\sigma (x_0,\mathrm{},x_s)`$ in the following inductive way: $`\sigma (x_0,x_1)`$ is a geodesic segment joining $`x_0`$ and $`x_1`$ and $`\sigma (x_0,\mathrm{},x_s)`$ is a union of geodesic segments joining $`x_0`$ and points of $`\sigma (x_1,\mathrm{},x_s)`$. By standard comparison theorems $`Vol_s(\sigma (x_0,\mathrm{},x_s))const[\mathrm{max}_{0ijs}\rho (x_i,x_j)]^s`$, where $`\rho (,)`$ is the distance function (this is where we use non-positive curvature). By \[Dupont 1\], $`r(\nu _k)`$ can be represented by a cocycle
$$g_1,\mathrm{},g_{2k}_{\sigma (x_0,g_1x_0,g_1g_2x_0,\mathrm{},g_1,g_2\mathrm{}g_{2k}x_0)}\nu _k$$
where $`g_iSp_{J_2}`$ and $`x_0X`$ is fixed. Since $`\nu _k`$ is uniformly bounded, the result follows.
We will give an independent proof of polynomial growth of $`\mu _2H^2(𝒟iff^{1,1}(S^1))`$. Let $`Var(S^1)`$ be a space of functions of bounded variation on $`S^1`$ mod constants. Then for $`f_1,f_2Var(S^1)`$,
$$_{S^1}f_1𝑑f_2cf_1_{Var}f_2_{Var}.$$
Now, $`Homeo(S^1)`$ acts isometrically in $`Var(S^1)`$ and there is a cocycle $`\psi H^1(𝒟iff^{1,1}(S^1),Var)`$ given by $`g\mathrm{log}(g^1)^{}`$. By an formula of Thurston, $`\mu _2`$ can be represented as
$$_{S_1}\psi (g_1)𝑑\psi (g_2).$$
The result now follows from Lemma I.1.1. For $`\mu _2`$ as a class in $`H^2(𝒟iff^{\mathrm{}}(S^1))`$ see also \[Connes-Gromov-Moscovici 1\].
### 2.6 Examples
A typical example of a group in $`𝒟iff^{1,\alpha }(S^1)`$ is a following one. Let $`K`$ be a subfield (i.e., a number field). By $`S^1(K)`$ denote $`K`$-rational points of $`S^1^2`$. Define $`G_K`$ as a group of $`C^1`$-diffeomorphism $`g`$ such that there are points $`x_0,\mathrm{},x_n=x_0S^1(K)`$ in this order such that $`g_k=g_{[x_k,x_{k+1}]}`$ is a restriction of an element of $`PSL_2(K)`$. The $`C^1`$-condition simply means that $`g_k^{}(x_{k+1})=g_{k+1}^{}(x_{k+1})`$. Then automatically $`G_K𝒟iff^{1,1}(S^1)`$. Groups of this type, or rather their obvious analogues which act by piecewise-affine transformations on $`S^1`$ viewed as $`/`$ appeared in \[Thompson 1\], \[Greenberg-Sergiesku 1,2\], \[Brown-Georghegan 1\], etc. where various properties were studied. The “proper” Thompson group can be smoothed , that is, embedded in $`𝒟iff(S^1)`$ \[Ghys-Sergiesku 1\] so that the Theorem 1.7 applies. However, it also acts on a tree so it is not Kazhdan already by a result of \[Alperin 1\], \[Watatani 1\]. Generally, subgroups of $`𝒟iff^{1,\alpha }(S^1)`$ like described above, do not have any obvious action on a tree and one needs our Theorem 1.7 to show that they are not Kazhdan. A parallel theorem for symplectomorphism groups will be given in Chapter IV. Notice also that the proof that our characteristic classes constructed in Section 5 are in polynomial cohomology agrees with a recent result on the growth of the Dehn function of the Thompson group \[Guba 1\].
## Chapter 3 Geometry of unitary cocycles
In this Chapter we return to the asymptotic geometry of finitely generated groups. If $`G`$ is not Kazhdan, then an orthogonal cocycle $`lZ^1(G,)`$ should be viewed as a way to linearize the geometry of $`G`$. Our first result is a convexity theorem 2.1 which says that the embedding of $`G`$ into the Hilbert space $``$ given by $`l`$ coarsely respects the geometry in a sense that inner points of big ”domains” in $`G`$ are mapped inside the convex hull of the image of boundary points.
We have seen in Chapter I that primitive functions $`:G`$ of cocycles in $`Z^1(G,l^p(G))`$ of a surface group satisfy
$$|(g)|<clength(g)^{1/p^{}}$$
Here, we start a general study of cocycle growth. We show in Theorem 3.1 that for any orthogonal cocycle $`l:G`$,
$$l(g)<c(\theta )[length(g)\mathrm{log}\mathrm{log}length(g)]^{1/2}$$
for almost all $`\theta GS^1`$ and $`g\theta `$ nontangentially. We use in proof an adjusted version of Makarov’s law of iterated logarithm. The result extends to all complex hyperbolic cocompact lattices of any dimension.
Using another deep result of Makarov, we show the following in Proposition 3.3. Let $`G`$ be a surface group, $`\beta :G`$ a surjective homomorphism and $`G_0=Ker\beta `$. Then the conical limit set of $`G_0`$ has Hausdorff dimension 1, in particular, the exponent $`\delta (G_0)=1`$. We do not know if this set has a full Lebesgue measure ( it is certainly a doable problem).
Notice that the proof of Lemma I.11.8 shows that the estimate on $`l(g)`$ is essentially sharp. It also shows that this estimate does not hold in other Banach spaces. However, imposing various restricitons on a Banach space, one still hopes to get an estimate, reflecting a fine structure of $`G`$.
### 3.1 Smooth and combinatorial harmonic sections
Let $`G`$ be a finitely generated group. $`\pi :GO()`$ an orthogonal representation without almost invariant vectors and $`l:G`$ a nontrivial cocycle. If $`M`$ is a compact Riemannian manifold with $`\pi _1(M)=G`$ (so that $`G`$ is finitely presented) then one forms a flat affine bundle $`E`$ over $`M`$ with fiber $``$ and monodromy
$$g(v\pi (g)v+l(g))$$
A result of \[Korevaar-Schoen 1\] and \[Jost 1\] (Lemma I.11.6) states that there is a harmonic section $`f`$ of $`E`$. If $`M`$ is Kähler then there is another cocycle $`m:G`$ so that a complex affine bundle with fiber $`E`$ and monodromy
$$g(v+iw\pi (g)v+i\pi (g)w+l(g)+im(g))$$
admits a holomorphic section. Our first result is a combinatorial version of this theorem.
Let $`\{\gamma _i\}`$ be a finite set of generators for $`G`$. Let $`V`$ be a space of ”sections”, that is, $`G`$-equivariant maps
$$f:G.$$
This simply means that $`f(g^1x)=\pi (g)f(x)+l(g)`$. Obviously, such map is determined by $`f(1)`$. Therefore, $`V`$. A combinatorial Laplacian is defined as
$$\mathrm{}f(x)=\underset{i}{}f(\gamma _ix)+f(\gamma _i^1x)2f(x).$$
Proposition 2.1.— There exists an equivariant $`f:G`$ with $`\mathrm{}f=0`$.
Proof.— Let $`v=f(1)`$, then $`f(x^1)=xv+l(x)`$. Therefore
$$\begin{array}{ccc}\hfill \mathrm{}f\left(x^1\right)& =& f\left(\gamma _ix^1\right)+f\left(\gamma _i^1x^1\right)2f\left(x^1\right)\hfill \\ & =& x\gamma _i^1v+l\left(x\gamma _i^1\right)+x\gamma _iv+l\left(x\gamma _i\right)2xv2l\left(x\right)\hfill \\ & =& x\left(\gamma _i^1v+\gamma _iv2v\right)+xl\left(\gamma _i^1\right)+l\left(x\right)+xl\left(\gamma _i\right)+l\left(x\right)2l\left(x\right)\hfill \\ & =& x\left(\gamma _i^1+\gamma _i2\right)v+x\left[l\left(\gamma _i^1\right)+l\left(\gamma _i\right)\right],\hfill \end{array}$$
so that we need only to solve an equation
$$(\gamma _i^1+\gamma _i2)v=[l(\gamma _i^1)+l(\gamma _i)].$$
Notice that $`\stackrel{~}{\mathrm{}}:`$ defined by $`v(\gamma _i^1+\gamma _i2)v`$ is selfadjoint. Moreover, since $`\stackrel{~}{\mathrm{}}=(\pi (\gamma _i)1)^{}(\pi (\gamma _i)1)`$, $`\stackrel{~}{\mathrm{}}`$ is nonpositive and if $`0spec(\stackrel{~}{\mathrm{}})`$, then $`\pi :GO()`$ has almost invariant vectors. Therefore, $`\stackrel{~}{\mathrm{}}`$ is invertible and the result follows.
### 3.2 A convexity theorem
We keep the notation of 3.1. Any cocycle $`l:G`$ can be seen as an embedding of $`G`$ in the Hilbert space. If $`l(g)\mathrm{}`$ as $`length(g)\mathrm{}`$, then this embedding is uniform in the sense that $`l(g)l(h)\mathrm{}`$ as $`\rho (g,h)\mathrm{}`$ for any word left-invariant metric on $`G`$. For instance, Proposition I.2.1 implies that any group $`G`$ acting discretely (but possibly not cocompactly) on an Hadamard manifold of pinched negative curvature, admits a uniform embedding into $`l^p(G),p>1`$. We are, however, interested in a finer geometry of the cocycle embeddings.
For a finite $`AG`$ and $`C>0`$, a $`C`$-interior $`int_C(A)`$ is defined as $`\{x|\rho (x,y)<CyA\}`$. A $`C`$-boundary $`_C(A)`$ is defined as $`A\backslash int_C(A)`$.
Theorem 2.2.— Let $`\pi :G𝒪()`$ be an orthogonal representation without almost-invariant vectors. Let $`l:G`$ be a cocycle for $`\pi `$. Then there are constants $`C_1,C_2(l)>0`$ such that for any finite $`AG`$ and any $`xA`$,
$$dist_{}(l(x)\overline{conv}(l(_{C_1}A)))C_2.$$
$`()`$
Proof.— Let $`f:G`$ be an equivariant harmonic map of Proposition 1.1. Since $`f(x^1)l(x)=f(1)=const`$, we can replace $`()`$ by a condition
$$dist_{}(f(x)\overline{conv}f(_{C_1}(A))C_2^{},$$
where however, one uses a right-invariant word metric on $`G`$ in definition of $`_C(A)`$. This result follows from the maximum principle of harmonic functions. Indeed, let $`xint_{C_1}(A)`$ be such that $`dist_{}(f(x)\overline{conv}f(_{C_1}(A)))`$ is maximal possible (and $`>C_2`$)( a choice of $`C_1,C_2`$ will be made later). Let $`v`$ be a unit vector, such that
$$(f(x)y,v)=dist_{}(f(x)\overline{conv}f(_{C_1}(A))$$
for some $`y\overline{conv}f(_{C_1}(A))`$. Let $`h(z)=(f(z)y,v)`$. Then $`h(x)>C_2`$ and $`h(_{C_1}(A))(\mathrm{},0]`$. Moreover, $`\stackrel{~}{\mathrm{}}h=0`$ and $`h(z)h(x)`$ for $`zint_{C_1}(A)`$. It follows that $`h(\gamma _ix)=h(x)`$ for all $`i`$. Replacing $`x`$ by $`\gamma _ix`$ and continuing until we hit $`_{C_1}A`$, we arrive to a contradiction with $`C_1=2`$, $`C_2=2f(1)+1`$.
### 3.3 Cocycle growth for a surface group
In this section we continue, for general representations, a subject started in I.5.2. Recall that, for any group $`G`$, any primitive function $`:G`$ of a class in $`H^1(G,l^p(G))`$ satisties
$$|(g)|constlength(g)$$
at least of $`G`$ is finitely presented. However, if $`G=\pi _1(\mathrm{\Sigma })`$, a surface group, then one has much finer estimate, established in Theorem I.5.2:
$$|(g)|constlength(g)^{1/p^{}}.$$
Theorem 3.1.— Let $`G=\pi _1(\mathrm{\Sigma })`$ be a surface group. Let $`\pi :GO()`$ be an orthogonal representation without almost-invariant vectors and let $`l:G`$ be a cocycle. Then for almost all $`\theta S^1G`$,
$$l(g)const(\theta )[length(g)\mathrm{log}\mathrm{log}length(g)]^{1/2}$$
$`()`$
as $`g`$ converges nontangentially to $`\theta `$. Here ”almost all” corresponds to a Lebesgue measure on $`G`$, identified with $`S^1`$ under some lattice embedding $`GSO^+(1,2)`$.
Remark.— Nontangential convergence of points of $`B^2`$ to $`\theta B^2`$ is an invariant of quasi-conformal homeomorphism \[???\]. Therefore $`()`$ is $`ap_{g,1}`$-invariant. Let $`AS^1`$ be an exceptional set where $`()`$ does not hold. It follows that the Lebesgue measure:
$$meas\phi (A)=0$$
for all $`\phi ap_{g,1}`$, considered as a quasisymmetric homeomorphism of $`S^1`$.
Proof.— Complexifying, we find a holomorphic section of an affine bundle $`E_{}`$ as in section 1. Lifting to $`^2`$, we obtain an equivariant holomorphic map (we replace $``$ by $``$)
$$\stackrel{~}{f}:^2.$$
Notice that $`\stackrel{~}{f}`$ is a Bloch function, that is, $`\stackrel{~}{f}const`$. The result now follows from a version of the Makarov law \[Makarov 1\] of iterated logarithm for Hilbert-space-valued Bloch functions.
Proposition.— Let $`\psi :B^2`$ be holomorphic and $`\psi _hconst`$. Then for almost all $`\theta S^1`$,
$$\underset{z\theta }{lim\; sup}\frac{\psi (z)}{\sqrt{\mathrm{log}(1|z|)\mathrm{log}\mathrm{log}\mathrm{log}(1|z|)}}<\mathrm{}.$$
Proof.— We will simply note which changes should be made in a proof for complex-valued functions \[Pommerenke 1\]. The Hardy identity \[Pommerenke 1, page 174\] holds in the following form. Let $`S`$ be a Riemannian surface, $`z_0S`$, $`g:S`$ a holomorphic function, $`(x,y)`$ normal coordinates in the neighbourhood of $`z_0`$. Let $`n`$ be a positive integer. Then
$$\frac{}{x}(g,g)^{n+1}=(n+1)(g,g)[(g_x^{},g)+(g,g_x^{})],$$
$$\begin{array}{ccc}\hfill \frac{^2}{x^2}(g,g)^{n+1}=n(n+1)(g,g)^{n1}[g_x^{},g)+(g,g_x^{})]^2+& & \\ \hfill \left(n+1\right)(g,g)^n\left[2(g_x^{},g_x^{})+(g_x^{\prime \prime },g)+(g,g_x^{\prime \prime })\right]& & \end{array}$$
and the same for $`\frac{^2}{y^2}`$. Summing up, we have
$$\begin{array}{ccc}\hfill \mathrm{}(g,g)^{n+1}& =& \left(\frac{^2}{x^2}+\frac{^2}{y^2}\right)(g,g)^{n+1}\hfill \\ & =& n\left(n+1\right)(g,g)^{n1}4\left|(g^{},g)\right|^2+\left(n+1\right)(g,g)^n2(g^{},g^{}),\hfill \end{array}$$
because $`\mathrm{}g=0`$ and $`g_y^{}=\sqrt{1}g_x^{}`$. If $`S`$ is a unit disc then in polar coordinates $`z=re^{it}`$
$$\mathrm{}=\frac{^2}{r^2}+\frac{1}{r}\frac{}{r}+\frac{1}{r^2}\frac{^2}{t^2}=\frac{1}{r}\frac{}{r}(r\frac{}{r})+\frac{1}{r^2}\frac{^2}{t^2}=\frac{1}{r^2}[(r\frac{}{r})^2+\frac{^2}{t^2}].$$
So
$$\frac{1}{r^2}((r\frac{}{r})^2+\frac{^2}{t^2})(g,g)^{n+1}=4n(n+1)(g,g)^{n1}|(g^{},g)|^2+2(n+1)(g,g)^n|g^{}|^2.$$
Integrating over $`0t2\pi `$ and using Cauchy-Schwartz inequality, we arrive at the inequality of \[Pommerenke 1, Theorem 8.9\]. The rest of the proof will go unchanged once we know the Hardy-Littlewood maximal theorem for $`(g,g)^n`$, which is used in \[Pommerenke 1, page 187\]. Let
$$g^{}(s,\xi )=\underset{0r1e^s}{\mathrm{max}}|g(r\xi )|,es<\mathrm{},\xi S^1.$$
Since $`g:B^2`$ is holomorphic, it is also harmonic and yields the Poisson formula. Then a proof of the Hardy-Littlewood maximal theorem given in \[Koosis 1\] applies, since it reduces it to the Hardy-Littlewood inequality for the maximal function of $`|g|`$.
Remark 3.2.— Theorem 3.1 holds for complex hyperbolic cocompact lattices. This is because Makarov’s law of iterated logarithm holds for the complex hyperbolic space, as we can see by passing to totally geodesic spaces of complex dimension 1. It is plausible that a version of Theorem 3.1 holds for real hyperbolic lattices (but not quaternionic and Cayley, as these are Kazhdan, see a new proof in Chapter VI). On the other hand, another deep result of \[Makarov 2\] saying that Bloch functions are nontangentially bounded for a limit set of Hausdorff dimension one, fails for Hilbert space valued functions. In fact, we have shown in Chapter I that there are unitary cocycles on a surface group such that $`l(g)\mathrm{}`$ as $`length(g)\mathrm{}`$.
If $`G`$ is any finitely generated group, and we are given an orthogonal representation $`\pi :GO()`$ and a cocycle $`lZ^1(G,)`$ with a control on $`l(g)`$ from below, then for any embedding of the surface group $`\pi _1(\mathrm{\Sigma })`$ into $`G`$ we immediately have a comparison inequality between the word lengthes of elements of $`\pi _1(\mathrm{\Sigma })`$ in $`\pi _1(\mathrm{\Sigma })`$ and $`G`$. To get a nontrivial result, we need a low bound on $`l(g)`$ better then $`[length(g)\mathrm{log}\mathrm{log}length(g)]^{1/2}`$. To find such groups and cocycles seems to be a very attractive problem.
We will now use similar ideas to estimate the Hausdorff dimension of limit sets of some infinite index subgroups of $`G=\pi _1(\mathrm{\Sigma })`$.
Theorem 3.2.— Let $`\beta :G`$ be a surjective homomorphism and let $`G_0=Ker\beta `$. Let $`A`$ be a conical limit set of $`G_0`$. Then $`dimA=1`$.
Proof.— Let $`[\beta ]H^1(G,)`$ be an induced class. Realize $`G`$ as a cocompact lattice in $`SO^+(1,2)`$ so that $`S=^2/G`$ is a hyperbolic surface. Let $`\omega `$ be a holomorphic 1-form on $`S`$ such that $`Re[\omega ]=[\beta ]`$ and let $`\stackrel{~}{\omega }`$ be a lift of $`\omega `$ on $`^2`$. Let $`f:^2`$ be holomorphic with $`df=\omega `$. Then $`f`$ is a Bloch function. By a result of \[Makarov 2\], there is a set $`BS^1`$ with $`dimB=1`$ such that
$$\underset{z\theta }{lim\; sup}|f(z)|<\mathrm{}$$
for any $`\theta B`$ and nontangential convergence. Notice that $`f(gz)=f(z)+([\omega ],[g])`$ where $`gG`$ and $`[g]`$ is an image of $`g`$ in $`H_1(G,)`$. Now it is clear that $`BA`$, so $`dimA=1`$.
Remark.— This result does not contradict a theorem of \[Sullivan 1\] and \[Tukia 1\] because $`G_0`$ is infinitely generated.
In the opposite direction we have the following . Let $`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2`$ be two closed surfaces and let $`\psi :\mathrm{\Sigma }_1\mathrm{\Sigma }_2`$ be a smooth ramified covering. Let $`G_i=\pi _1(\mathrm{\Sigma }_i)`$ and let $`G_0=Ker\psi _{}:G_1G_2`$. Let $`G_1SO^+(1,2)`$ be a realization of $`G_1`$ as a lattice. Then for any $`zB^2`$,
$$\underset{gG_0}{}|1gz|<\mathrm{}.$$
In other words, either $`\delta (G_0)<1`$ or $`\delta (G_0)=1`$ and $`G_0`$ is of convergence type. In the latter case, the Patterson measure of the conical limit set of $`G_0`$ is zero. To see this, notice that we can find hyperbolic structures on $`\mathrm{\Sigma }_i,i=1,2`$ so that $`\psi `$ is holomorphic. Let $`\stackrel{~}{\psi }`$ be a lift of $`\psi `$ as a map $`\stackrel{~}{\psi }:B^2B^2`$. Since $`\stackrel{~}{\psi }`$ is a bounded holomorphic function, $`\stackrel{~}{\psi }`$ has limit values almost everywhere on $`S^1`$. By I.6.5., $`|\stackrel{~}{\psi }|S^1|=1`$ almost everywhere. So $`\stackrel{~}{\psi }`$ is an inner function. Let $`CB^2`$ be a countable set of zeros of $`\stackrel{~}{\psi }`$. We claim that $`C`$ is a finite union of orbits of $`G_0`$. First, it is clear that $`C`$ is $`G_0`$-invariant. Let $`QB^2`$ be compact which contains a fundamental domain for $`G_1`$. Then $`\stackrel{~}{\psi }(Q)`$ is compact so there is a finite set $`RG_2`$ such that $`g(0)\stackrel{~}{\psi }(Q)`$ if $`gR`$. Let $`TG_1`$ be finite and such that $`\psi _{}(T)R`$. Let $`Q_1=_{gT^1}gQ`$ so that $`Q_1`$ is compact and therefore $`CQ`$ is finite. Let $`xC`$, then $`x=gy`$ with $`yQ`$. So $`0=\stackrel{~}{\psi }(x)=\psi _{}(g)\stackrel{~}{\psi }(y)`$, i.e., $`\psi _{}(g^1)(0)=\stackrel{~}{\psi }(y)\stackrel{~}{\psi }(Q)`$. This means $`\psi _{}(g^1)R`$ so $`g^1TG_0`$, and $`gG_0T^1`$, say $`g=g_0t^1`$, $`g_0G_0`$, $`tT`$. Then $`t^1yC`$ and $`t^1yQ_1`$, so there are finitely many options for $`t^1y`$.
We deduce that there are $`x_1,\mathrm{},x_n`$ such that $`C=_{i=1}^nG_0x_i`$. The decomposition formula for inner functions implies that
$$\stackrel{~}{\psi }(z)=c\underset{\stackrel{g_0G_0}{1in}}{}\frac{\overline{g_0x_i}}{g_0x_i}\frac{zg_0x_i}{1\overline{g_0x_i}z}$$
which gives an explicit formula for holomorphic maps between hyperbolic Riemann surfaces (one still needs to find $`x_i`$). By a well-known result on zeros of a bouded holomorphic function \[Koosis 1, IV: B, Theorem 1\],
$$\underset{g_0G_0}{}(1|g_0x_i|)<\mathrm{}.$$
The rest follows from \[Nicholls 1\].
## Chapter 4 A theory of groups of symplectomorphisms
We already have noticed an intriguing similarity between groups acting on the circle and groups acting symplectically on a compact sympletic manifold. The two leading topics studied in Chapter 2, namely, (non-) Kazhdan groups acting on $`S^1`$ and characteristic classes, have exact analogues for Sympl(M). In fact, a theory of characteristic classes parallel to II.5 , has already been presented in \[Reznikov 2\] and \[Reznikov 4\]. In the second cited paper, we noticed that the Kähler action of Sympl(M) on the twistor variety allows us to define a series of classes in $`H_{cont}^{2k}(Sympl(M),)`$, $`k`$ odd , which are highly non-trivial. In the first cited paper,we introduced bi-invariant forms on Sympl(M) and the classes in $`H_{top}^{odd}(Sympl(M))`$ and $`H^{odd}`$($`Sympl(M)^\delta ,/A)`$ (cohomology of a topological space and a discrete group) where $`A`$ is a group of periods of the above-mentioned forms. Here we present a fundamental class in $`H^1`$$`(Sympl(M),L^2(M))`$ whose nontriviality on a subgroup $`GSympl(M)`$ implies that $`G`$ is not Kazhdan, similarly to the situation in $`𝒟iff^{1,\alpha }(S^1)`$. ¿From the nature of our class it is clear that its vanishing imposes severe restriction on the symplectic action, roughly, the transformations of $`G`$ should satisfy a certain PDE . We give an explicit formula for our class in the case of a flat torus.
We then introduce a characteristic class in $`H^{n+1}(Sympl^\delta (M^{2n}),)`$ associated with a compact Lagrangian immersed submanifold. This class is a sympletic counterpart, and a generalization, of the Thurston-Bott class \[Bott 1\]. We use this class to give a formula for the volume of compact negatively curved manifold through Euclidean volumes of “Busemann bodies” (the images of the manifold under Busemann functions).
### 4.1 Deformation quantization : an overview
Let $`F`$ be a field and $`A|F`$ a (commutative) algebra . A deformation of A is an algebra structure over $`F[[\mathrm{}]]`$ of $`A[[\mathrm{}]]`$ extending that $`A`$, so that if $`x,yA`$,
$$xy=xy+b_1(x,y)\mathrm{}+b_2(x,y)\mathrm{}^2+\mathrm{}$$
where $`xy`$ is a multiplication in $`A`$ and $`xy`$ is a deformed multiplication .
If $`F=`$, $`A=C^{\mathrm{}}(M)`$, where $`M`$ is a symplectic manifold , then a deformation quantization is a deformation of $`A`$ with $`b_1(f,g)=\{f,g\}`$, a Poisson bracket . A deformation quantization always exists by a result of \[Moyal 1\], \[Vey 1\], \[Bayen-Flato-Fronsdal-Lichnerowicz-Sternheimer 1\] \[Fedosov 1\]. For any algebra $`A|F`$ one defines a Hochschild collomology $`HH^k(A)=Ext_{AA}^k(A,A)`$. There is a natural Lie superalgebra structure in $`HH^{}(A)`$ \[Gerstenhaber 1\]. There exists a simple explicit complex , computing $`HH^k(A)`$ with $`C^k(A)=Hom_F(_{i=1}^kA,A)`$. In particular , $`b_1`$ above is a cocycle (for any deformation ). If $`F=`$ and A is a topological algebra, one modifies the definitions to obtain topological Hochschild cohomology . If $`M`$ is a smooth manifold and $`A=C^{\mathrm{}}(M)`$ with a pointwise multiplication , then
$$HH^k(A)=\mathrm{\Gamma }(M,\mathrm{\Lambda }^kTM),$$
a space of poly-vector fields . The Lie superalgebra structure coincides with a classical bracket of poly-vector fields.
We will need an explicit form of the cocycle condition for a 2-cocycle $`b:AA`$ :
$$xb(y,z)b(xy,z)+b(x,yz)b(x,y)z=0.$$
### 4.2 A fundamental cocycle in $`H^1(Sympl(M),L^2(M))`$
Let $`(M^{2n},\omega )`$ be a compact symplectic manifold . Fix a deformation quantization
$$fg=fg+\{f,g\}\mathrm{}+\underset{i=2}{\overset{\mathrm{}}{}}c_i(f,g)\mathrm{}^i$$
Let $`\mathrm{\Phi }:MM`$ be symplectic and let
$$\begin{array}{ccc}\hfill f\stackrel{~}{}g& =& \left(f\mathrm{\Phi }^1g\mathrm{\Phi }^1\right)\mathrm{\Phi }\hfill \\ & =& fg+\{f,g\}\mathrm{}+_{i=2}^{\mathrm{}}c_i^{}(f,g)\mathrm{}^i\hfill \end{array}$$
Lemma 2.1.— Let $`A|F`$ be an algebra and let
$$fg=fg+c_1(f,g)\mathrm{}+\mathrm{}+c_{k1}(f,g)\mathrm{}^{k1}+\underset{i=k}{\overset{\mathrm{}}{}}c_i(f,g)\mathrm{}^i$$
and
$$f\stackrel{~}{}g=fg+c_1(f,g)\mathrm{}+\mathrm{}+c_{k1}(f,g)\mathrm{}^{k1}+\underset{i=k}{\overset{\mathrm{}}{}}c_i^{}(f,g)\mathrm{}^i$$
be two deformations, which coincide up to the order $`\mathrm{}^{k1}`$. Then
$$c_ic_i^{}:AAA$$
is a Hochschild cocycle .
Proof.
$$(fg)p(f\stackrel{~}{}g)\stackrel{~}{}p=c_k(f,g)p+c_k(fg,p)c_k^{}(f,g)pc_k^{}(fg,p)(\text{mod }\mathrm{}^{k+1})$$
Similarly,
$$f(gp)f\stackrel{~}{}(g\stackrel{~}{}p)=fc_k(g,p)+c_k(f,gp)fc_k^{}(g,p)c_k^{}(f,gp)(\text{mod }\mathrm{}^{k+1})$$
So for $`c=c_kc_k^{}`$,
$$fc(g,p)+c(f,gp)c(f,g)pc(fg,p)=0,$$
which means that $`c`$ is a 2-cocycle.
Lemma 2.2.— A formula
$$\mathrm{\Phi }\left[(f,g)c_2(f\mathrm{\Phi }^1,g\mathrm{\Phi }^1)\mathrm{\Phi }c_2(f,g)\right]$$
defines a smooth cocycle of Sympl(M) in the space $`Z^2(C^{\mathrm{}}(M)`$, $`C^{\mathrm{}}(M))`$ of Hochschild 2-cocycles for $`C^{\mathrm{}}(M)`$.
Proof. Follows from Lemma 2.1.
Passing to Hochschild cohomology, we obtain a 1-cocycle of $`Sympl(M)`$ in
$$HH^2(C^{\mathrm{}}(M))=\mathrm{\Gamma }(M,\mathrm{\Lambda }^2TM).$$
Using the symplectic structure ,we identify $`\mathrm{\Gamma }(M,\mathrm{\Lambda }^2TM)`$ with $`\mathrm{\Omega }^2(M)`$, a space of 2-forms on $`M`$. Multilying by $`\omega ^{n1}`$ we obtain a cocycle
$$\mu H^1(Sympl(M),C^{\mathrm{}}(M)).$$
### 4.3 Computation for a flat torus and the main theorem
If $`M`$ is a coadjoint orbit of a compact Lie group, one can find an explicit formula for the deformation quantization $`fg`$. A classical case $`M=T^{2n}`$, a flat torus, is due to H.Weyl.
Proposition 3.1.— One has a following deformatiom quantization on $`T^{2n}`$ :
$$fg=\underset{k=0}{\overset{\mathrm{}}{}}\frac{1}{k!}\left(\frac{i\mathrm{}}{2}\right)^k\sigma ^{i_1j_1}\mathrm{}\sigma ^{i_kj_k}\frac{^kf}{y_{i_1}\mathrm{}y_{i_k}}\frac{^kg}{y_{j_1}\mathrm{}y_{j_k}}$$
where $`\sigma ^{ij}`$ are elements of the matrix, inverse to the matrix $`(\sigma _{ij})`$ of a (constant) symplectic form , and the “repeated indices” summation agreement is applied .
Now, since our definition of a fundamental cocycle is completely explicit, one can derive an explicit formula for $`\mu `$ in this case. We give an answer for $`T^2`$ (the formula for $`T^{2n}`$ is completely analogous). The computation is tedious (takes several pages) but straightforward and is left to reader . Here is the formula for $`T^2`$ :
$$\mathrm{\Phi }\frac{^2\mathrm{\Phi }_2}{y_1^2}\frac{^2\mathrm{\Phi }_1}{y_2^2}\frac{^2\mathrm{\Phi }_1}{y_1^2}\frac{^2\mathrm{\Phi }_2}{y_2^2}$$
where $`\mathrm{\Phi }=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)`$ a symplectomorphism of the form $`T^2`$. Summing up, we have
Theorem 3.2.— Let $`M^{2n}`$ be a compact symplectic manifold , let $`Sympl(M)`$ its symplectomorphism group , acting orthogonally on a Hilbert space $`L^2(M)`$. There exists a cocycle
$$\mu Z^1(Sympl(M),L^2(M)),$$
defined canonically by a given deformation quantization of $`C^{\mathrm{}}(M)`$ with the following properties :
A. Let
$$\begin{array}{ccc}\hfill f\stackrel{~}{}g& =& (f\mathrm{\Phi }^1g\mathrm{\Phi }^1)\mathrm{\Phi }=fg+\{f,g\}\mathrm{}+c_2^{}(f,g)\mathrm{}^2+\mathrm{},\hfill \\ \hfill fg& =& fg+\{f,g\}\mathrm{}+c_2(f,g)\mathrm{}^2+\mathrm{}\hfill \end{array}$$
and let us indentify the class of the Hochschild cocycle $`c_2^{}c_2`$ with a section $`\nu `$ of $`\mathrm{\Lambda }^2TM`$. Let $`\widehat{\nu }`$ be a 2-form obtained from $`\nu `$ by lifting the indices using the symplectic form . Then
$$\mu (\mathrm{\Phi })\omega ^n=\widehat{\nu }.\omega ^{n1}.$$
B. $`\mu (\mathrm{\Phi })`$ depends only on the second jet of $`\mathrm{\Phi }`$.
C. For $`M=T^2`$ and the Weyl deformation quantization , $`\mathrm{\Phi }=(\mathrm{\Phi }_1,\mathrm{\Phi }_2)`$,
$$\mu (\mathrm{\Phi })=\frac{^2\mathrm{\Phi }_2}{y_1^2}\frac{^2\mathrm{\Phi }_1}{y_2^2}\frac{^2\mathrm{\Phi }_1}{y_1^2}\frac{^2\mathrm{\Phi }_2}{y_2^2}.$$
D. If $`G`$ is a Kazhdan subgroup of $`Sympl(M)`$, then
$$\mu (\mathrm{\Phi })_{L^2}<const(\mathrm{\Phi }G).$$
Examples.
1) $`M=T^{2n}`$, $`G=Sp(2n,Z)`$ (Kazhdan for $`n2`$). Then $`\mu `$ is identically zero.
2)Let $`\mathrm{\Gamma }`$ be a surface group, and let $`M`$ be a component of
$$Hom(\mathrm{\Gamma },SO(3))/SO(3),$$
consisting of representations with nontrivial Stiefel-Whitney class. Then $`M`$ is a compact symplectic manifold and $`Map_g`$ acts symplectically on $`M`$. We do not know if part D of Theorem 3.2 holds in this case and if $`Map_g`$ is Kazhdan or not. There is a “Teichmüller structure” on $`M`$ defined by a holomorphic map of the Teichmüller space into the twistor variety of $`M`$, described in \[Reznikov 4\], see also Chapter 5.
Remark.— The case of two-dimensional $`M^2`$ is much easier, simply because $`SL_2()`$ is not Kazhdan . If $`Sympl(M,x_0)`$ is a subgroup fixing $`x_0M`$ then one gets a nontrivial unitary cocycle on $`Sympl(M,x_0)`$ by pulling back from $`SL_2()`$ under the tangent representation. Using the measurable transfer (= Shapiro’s lemma) one constructs a cocycle of $`Sympl(M)`$. See \[Zimmer 1\] for details .
### 4.4 Invariant forms on the space of Lagrangian immersions and new regulators for symplectomorphism groups
In this section we will “symplectify” the Thurston-Bott class in the cohomology of diffeomorphism groups. Let $`M`$ be any (possibly noncompact) symplectic manifold, and let $`L_0M`$ be a Lagrangian immersion of a compact oriented manifold $`L_0`$. Let $`Lag(L_0,M)`$ be a space of Lagrangian immersions of $`L`$ into $`M`$ which can be jointed to $`L_0`$ by an exact Lagrangian homotopy. This means the following. If $`f_tM`$ is a smooth family of Lagrangian immersions, than $`{\displaystyle \frac{d}{dt}}f_t|_{t=0}`$ is a vector field along $`L_0`$. Projecting to the normal bundle $`NL_0`$ and accounting that $`NL_0`$ is canonically isomorphic to $`T^{}L_0`$, we get a 1-form on $`L_0`$ which is immediately seen to be closed. A Lagrangian homotopy $`f_t`$ is exact, if this form is exact for all t. There is therefore a well-defined function $`F`$ (mod const ) on $`L`$ which can be seen as a tangent vector of the deformation .
Definition.— A canonical (n+1)-form $`\nu `$ on $`Lag(L,M)`$ is defined by
$$\nu (F_0\mathrm{}F_n)=_LF_0𝑑F_1\mathrm{}F_n=Vol_{n+1}(\stackrel{~}{Q})(),$$
where $`\stackrel{~}{Q}`$ is any chain in $`^{n+1}`$ spanning $`(F_0,\mathrm{},F_n)(L)`$.
Proposition 4.1$`\nu `$ is closed.
Proof is an exercise for reader.
Let $`Sympl_0(M)`$ be a group of Poissonian transformations of $`M`$. Then $`Lag(L,M)`$ is invariant under $`Sympl_0(M)`$.
Proposition 4.2.$`\nu `$ is $`Sympl_0(M)`$-invariant.
Proof is obvious.
A standard theory of regulators \[Reznikov 3\], \[Reznikov 2\] implies that, first, one has an induced class in $`H^{n+1}(𝔤,)`$. where $`𝔤`$ = $`Lie(Sympl_0(M))`$= $`C^{\mathrm{}}(M)/const`$ given by $`(),`$ where now $`F_iC^{\mathrm{}}(M)`$ and second, a class in
$$Hom(\pi _{n+1}(BSympl_0^\delta (M)^+,/A))(n+15),$$
where $`A`$ is a group of periods of $`\nu `$ on maps $`\mathrm{\Sigma }^{n+1}Sympl_0(M)`$ of homology spheres to $`Sympl_0(M)`$. This class often lifts to a class in $`H^{n+1}(Sympl_0^\delta (M),)`$ under suitable conditions on topology of $`Sympl_0(M)`$ (see the discussion in the papers cited above ).
As an example, let $`Q`$ be a compact oriented simply connected manifold, $`M=T^{}Q`$ and $`L_0=Q`$, a zero section. Then we obtain a class $`[\nu ]`$ in $`H^{n+1}(Sympl_0(T^{}Q),)`$. Notice that the restriction of this class on $`𝒟iff(Q)Sympl_0(T^{}Q)`$ is zero, as $`𝒟iff(Q)`$ fixes the zero section. However , our class is an extension of Thurston-Bott class \[Bott 1\] in $`𝒟iff(Q)`$ by means of the following construction. Let $`GSympl_0(T^{}Q)`$ be a subgroup of symplectomorphisms of the form
$$p_x\varphi ^{}p_x+df(x),$$
where $`fC^{\mathrm{}}(Q)`$, $`\varphi `$ $`𝒟iff(Q)`$, $`xQ`$, $`p_xT_x^{}Q`$. Then $`G`$ is an extension
$$0C^{\mathrm{}}(Q)/constG𝒟iff(Q)1.$$
Any 1-cocycle $`\psi Z^1(𝒟iff(Q),C^{\mathrm{}}(Q)/const)`$ induces a spliting of this exact sequence:
$$S_\psi :𝒟iff(Q)G.$$
Now let $`\mu `$ be a smooth density on $`Q`$ then $`\psi =\frac{\varphi _{}\mu }{\mu }`$ is a 1-cocycle, so it defines such a splitting. A pull-back $`S_\psi ^{}([\nu ]|G)`$ of our class on $`𝒟iff(Q)`$ is precisely the Thurston-Bott class.
We sum up :
Theorem 4.3.
A. A formula
$$\nu (F_0,\mathrm{},F_n)=_LF_0𝑑F_1\mathrm{}𝑑F_n=Vol_{n+1}(\stackrel{~}{Q})$$
defines an $`Sympl_0(M)`$-invariant closed $`(n+1)`$-form in $`Lag(L,M)`$. It induces a class $`[\nu ]H^{n+1}(Lie(Sympl_0(M),)`$ and a regulator
$$[\nu ]:\pi _{n+1}(BSympl_0^+(M)),n+15,$$
which lifts to a class
$$[\nu ]H^{n+1}(Sympl_0^\delta (M),)$$
if $`\stackrel{~}{H}_i(Lag(L,M),))=0,0in+1`$.
B. In particular , if $`Q`$ is a smooth oriented simply-connected closed manifold, then
$$[\nu ]H^{n+1}(Sympl^\delta (T^{}Q),)$$
pulls back to the Thurston-Bott class under any splittting
$$𝒟iff(Q)C^{\mathrm{}}(Q)/const𝒟iff(Q),$$
coming from a smooth density on $`Q`$.
### 4.5 A volume formula for negatively-curved manifolds
This section is ideologically influenced by \[Hamenstädt 1\] and discussions with G.Besson (Grenoble, 1996). Let $`N^n`$ be an Hadamard manifold. Let $`CN`$ be the space of oriented geodesic of $`N`$, which is a symplectic manifold of dimension $`2n2`$. Any point $`xN`$ defines a Lagrangian sphere $`S_xCN`$ of geodesics passing through x.
Lemma 5.1.— A pull-back $`S^{}\nu `$ of the form $`\nu \mathrm{\Omega }^n(CN)`$ to $`N`$ is the Riemannian volume form on $`N`$ times a constant .
Proof.— An exercise in Jacobi fields.
Now if $`G`$ acts discretely and cocompactly on $`N`$, we have
$$[S^{}\nu ,\text{fundamental class of}N/G]=cVol(N/G).$$
Corollary 5.2.$`[\nu ]0`$ in $`H^n(Sympl^\delta (N),).`$
Now we assume that the curvature of $`N`$ is strictly negative and moreover, the induced action of $`G`$ on the sphere at infinity $`S_{\mathrm{}}`$ is of class $`C^{1,\frac{n1}{n}}`$. For $`n=2`$ this is always the case \[Hurder-Katok 1\], whereas for $`n3`$ seems to require a pinching of the curvature. Notice that the map
$$s_+:CNS_{\mathrm{}},$$
sending any geodesic $`\gamma (t)`$ to $`\gamma (\mathrm{})`$, is a Lagrangian fibration. Therefore if we fix a Lagrangian section of $`s_+`$, we will have a symplectomorphism $`CNT^{}(S_{\mathrm{}})`$. Fix $`p_0N`$, then $`S_{p_0}`$ is such a section. Notice that an induced homomorphism $`GSympl(T^{}S_{\mathrm{}})`$ is given by ,
$$g(z\pi (g)z+dF(p_0,g^1p_0,\theta )),$$
where $`gG,zT_\theta ^{}S_{\mathrm{}},\pi :G𝒟iff^{1,\frac{n1}{n}}(S_{\mathrm{}})Sympl(T^{}S_{\mathrm{}})`$ is induced by the action of $`G`$ on $`S_{\mathrm{}}`$ and $`B(p_0,g^1p_0,\theta ))`$ is the Buseman function. Our assumption imply that $`B(p_0,p_1,)C^{\frac{n1}{n}}(S_{\mathrm{}})W_{n}^{}{}_{}{}^{\frac{n1}{n}}(S_{\mathrm{}})`$. Recall that for $`F_1,\mathrm{},F_nW_{n}^{}{}_{}{}^{\frac{n1}{n}}`$ we have an n-form
$$_S_{\mathrm{}}F_1𝑑F_2\mathrm{}𝑑F_n=_{B^n}𝑑u_1\mathrm{}𝑑u_n,$$
where $`u_i`$ is a harmonic extension of $`F_i`$.
We derive a
Corollary 5.3.— Let $`N^n/G`$ be a compact negatively curved manifold such that the induced action of $`G`$ on $`S_{\mathrm{}}`$ is of class $`C^{1,\frac{n1}{n}}`$. If the fundamental class of $`G`$ is
$$\underset{i}{}[g_{1}^{}{}_{}{}^{(i)}\mathrm{}g_{n}^{}{}_{}{}^{(i)}],$$
then the following volume formula holds :
$$Vol(N/G)=c(n)\underset{i}{}_S_{\mathrm{}}F_{1}^{}{}_{}{}^{(i)}𝑑F_{2}^{}{}_{}{}^{(i)}\mathrm{}𝑑F_{n}^{}{}_{}{}^{(i)},$$
where $`F_{k}^{}{}_{}{}^{(i)}(\theta )=B(p_0,\left(g_{k}^{}{}_{}{}^{(i)}\right)^1p_0,\theta ).`$
One can say that a volume of a negatively curved manifold is a sum of Euclidean volumes of Busemann bodies in $`^n`$ bounded by $`(F_1,\mathrm{},F_n)(S_{\mathrm{}})`$.
Replacing the Busemann cocycle by a Jacobian cocycle $`\frac{g\mu }{\mu }`$, where $`\mu `$ is a smooth density on $`S_{\mathrm{}}`$, we arrive to a similar formula for Godbillon-Vey-Thurston-Bott invariant of $`N/G`$,under the same regularity assumptions. This seems to have been also accomplished in a preprint \[Hurder 1\] cited in \[Hurder-Katok 1\], though I was unable to obtain this paper from its author. The case $`n=2`$ is ,however, covered in \[Hurder-Katok 1\].
## Chapter 5 A theory of groups of volume-preserving diffeomorphisms and the nonlinear superrigidity alternative
In this chapter, we shift the focus from linear functional analytic techniques to nonlinear PDE, notably harmonic maps into nonlocally compact spaces, a theory recently developed in \[Korevaar-Schoen 1\] and \[Jost 1\]. The main idea is to use twistor varietes, which were in a center of the characteristic classes construction of \[Reznikov 4\], for a deeper study of volume-preserving actions of groups. We introduce an invariant of a volume-preserving action, which we call $`\mathrm{\Lambda }`$, which is a sort of a $`\mathrm{log}L^2`$-version of a $`sup`$-displacement studied in \[Zimmer 2\]. Our first main result, Theorem 2.3, states that if $`G`$ is a Kazhdan group acting on a compact manifold $`M`$ preserving volume, then either $`\mathrm{\Lambda }>0`$ or $`G`$ fixes a $`\mathrm{log}L^2`$-metric. A much weaker analogue of this result for the special case of lattices in Lie groups and $`sup`$-displacement was known before \[Zimmer 2, Theorem 4.8\].
We then apply our technique to a major open problem in the field, that of the nonlinear superrigidity of volume-preserving actions of lattices in Lie groups. From a nonlinear version of Margulis theorem given in \[Zimmer 3\] one knows that a volume preserving action of a lattice in a semisimple Lie group of rank $`2`$ on a low dimensional (with respect to the group) manifold fixes a measurable Riemannian metric. Since measurable metrics do not define a geometry on a manifold, one wishes ,of course, to prove a much stronger result: that the action preserves a smooth metric. Zimmer noticed \[Zimmer 2 and references therein\] that such strong result would follow if one is able to find an invariant metric whose dilations with respect to any smooth metric are in the class $`L_{loc}^2`$. The central question of how to find such a “bounded” invariant metric was left completely open. We present a completely new approach to the problem which leads to Theorem 3.1. It states that if a cocompact lattice acts on $`M`$ preserving volume, then either it nearly preserves a $`\mathrm{log}L^2`$-metric, or a sort of $`G`$-structure. This theorem, though constituting a clear progress in solution of the main problem is still less than what one wants in two respects: first, we deal with $`\mathrm{log}L^2`$-metrics, not $`L^2`$-metrics, second, we leave open a very delicate situation when an action nearly preserves a $`\mathrm{log}L^2`$-metric, but does not exactly preserve such a metric. This situation is purely infinite-dimensional (if an action on a finitely dimensional space of nonpositive curvature nearly preserves a point, it actually preserves a point at infinity). As already said, we use a heavy machinery : harmonic maps into twistor varietes and vanishing results of \[Mok-Siu-Yeung 1\] and \[Corlette 1\]. These results will also be applied in the next Chapter to study quaternionic Kähler groups.
As is well-known, an original Kolmogorov’s definition of entropy used extremum over all partitions and only became computable after it had been realized by Kolmogorov and Sinai that certain partitions realize entropy. In a way of a pleasant similarity, we show how to compute our invariant $`\mathrm{\Lambda }`$ for $`G=`$ in case $`G`$ leaves a geodesic in the twistor space invariant, like a hyperbolic element of $`SL(n,)`$ acting on $`T^n`$. This clearly shows an advantage of $`\mathrm{log}L^2`$-displacement over sup-displacement.
### 5.1 $`\mathrm{log}L^2`$-twistor spaces
Twistor varietes $`(C^{\mathrm{}})`$ were used in \[Reznikov 4\] to define secondary characteristic classes for volume-preserving and symplectic actions. More specifically, we have defined, for a compact oriented manifold $`M`$ equipped with a volume form $`\nu `$, a series of classes in $`H_{cont}^{}(𝒟iff_\nu (M))`$ of dimension $`5,9,13,\mathrm{}`$ (where $`𝒟iff_\nu (M))`$ is a group of volume-preserving diffeomorphisms). Likewise, for a compact symplectic manifold $`M`$ we have defined classes in $`H_{cont}^{}(Sympl(M))`$ of dimensions of $`2,6,10,\mathrm{}`$. For purposes of the present paper, we will need to work with a $`\mathrm{log}L^2`$-version of twistor varietes, defined below.
Remark 1.1.— I would like to use an opportunity to note that by some strange reason I have overlooked an integrated Euler class in $`H_{cont}^n(𝒟iff_\nu (M^n))`$. The definition is exactly like that in \[Reznikov 4\] for classes in dimensions 5, 9, $`\mathrm{}`$ if one realizes that there exists an $`n`$-form on the twistor variety for $`M`$, which is $`𝒟iff_\nu `$-invariant. Alternatively, if $`𝒟iff_\nu (M,p_0)`$ is a subgroup fixing a point $`p_0`$, then one pulls back the Euler class of $`SL_n()`$ under the tangent representation
$$𝒟iff_\nu (M,p_0)SL_n(),$$
and then applies a measurable transfer (see the above cited paper). The just defined class viewed as a class in $`H^n(𝒟iff_\nu ^s(M))`$ is bounded. This follows from the fact that the Euler class is bounded \[Sullivan 2\] exactly in the same manner as in \[Reznikov 4\].
We now define the $`\mathrm{log}L^2`$-twistor variety $`X`$ for $`(M,\nu )`$. First, one defines a bundle $`𝒫`$ of metrics with volume form $`\nu `$ as an $`SL(n)/SO(n)`$-bundle, associated with a principal $`SL(n)`$-bundle, defined by $`\nu `$. Fix a smooth section (=a Riemannian metric with volume form $`\nu `$) $`g_0`$ of this bundle. For any other measurable section $`g`$ of $`𝒫`$ define
$$\rho ^2(g_0,g)=_M\rho _x^2(g_0,g)𝑑\nu ,$$
$`()`$
where $`\rho _x`$ is a distance in $`𝒫_x`$ induced by (fixed once forever) $`SL(n)`$-invariant metric on $`SL(n)/SO(n)`$. Now the twistor variety $`X`$ consists of $`\mathrm{log}L^2`$-metrics, that is,
$$\rho (g_0,g)<\mathrm{}.$$
Alternatively, let $`A_x`$ be a self-adjoint (with respect to $`(g_0)_x`$) operator such that $`g_x=g_0(A_x,)`$. Then $`()`$ can be written as
$$_M\mathrm{log}A_x^2𝑑\nu <\mathrm{}.$$
A crucial fact about $`𝒫`$ is a following
Proposition 1.2.$`𝒫`$ is a complete Hilbert Riemannian manifold with nonpositive curvature operator. The action of $`𝒟iff_\nu (M)`$ on $`𝒫`$ is isometric.
Proof.— We will only define a metric, leaving all routine checks to the reader. A tangent space at $`g_0`$ consists of $`L^2`$-sections of $`S^2T^{}M`$, with trace identically zero. If $`A`$ is such a section (so that $`A_x`$ is $`g_0`$-self-adjoint for all $`xM,`$) then we define a square of the length of $`A`$ as
$$_MtrA^2𝑑\nu .$$
This metric is invariant under $`SO(n)`$-valued gauge transformations. Now we define a $`\mathrm{log}L^2`$ $`SL(n)`$-gauge group as a group of measurable sections of $`Aut(TM)`$ such that with respect to $`g_0`$,
$$_M\mathrm{log}(A^{}A)^2𝑑\nu <\mathrm{}.$$
Then $`𝒫`$ is a homogeneous space under the action of this group. We define a metric on $`𝒫`$ as a unique invariant metric, which agrees at $`g_0`$ with the metric just defined.
Now let $`(M^{2n},\omega )`$ be a compact symplectic manifold. Let $`𝔗`$ be the twistor bundle, that is, an $`Sp(2n)/U(n)`$-bundle, associated with the principal $`Sp(2n)`$-bundle, defined by $`\omega `$. A smooth section of $`𝔗`$ is exactly a tamed almost-complex structure. One then defines a space $`Z`$ of $`\mathrm{log}L^2`$-sections of $`𝔗`$ as above (the $`C^{\mathrm{}}`$-version was used in \[Reznikov 4\]).
Prosition 1.3.— The spaces $`X`$ and $`Z`$ are Alexandrov and Busemann nonpositively curved.
Proofs are standard.
### 5.2 A new invariant of smooth volume-preserving dinamical systems
Let $`(M,\nu )`$ be a compact oriented manifold with volume form $`\nu `$. Let $`G`$ be a finitely generated group which acts on $`M`$ by smooth transformations, preserving $`\nu `$. We are going to define a new dinamical invariant which we call $`\mathrm{\Lambda }`$. This is a nonnegative real number. Though it depends on the choice of a system of generators of $`G`$, the crucial fact of whether $`\mathrm{\Lambda }>0`$ or $`\mathrm{\Lambda }=0`$ does not. This relates our $`\mathrm{\Lambda }`$ to Kolmogorov’s entropy \[Kolmogorov 1\]. The invariant $`\mathrm{\Lambda }`$ is highly nontrivial already for $`G=`$, that is, as a new invariant of a volume-preserving diffeomorphism. It is also an invariant under conjugation in $`𝒟iff_\nu (M)`$. A central result of this section is Theorem 2 below stating that if $`G`$ is a Kazhdan group then either $`\mathrm{\Lambda }>0`$ or $`G`$ fixes a $`\mathrm{log}L^2`$-Riemannian metric (again this connects $`\mathrm{\Lambda }`$ to the Kolmogorov’s entropy).
Let $`g_1,\mathrm{},g_n`$ be a system of generators for $`G`$. Let $`X`$ be the twistor variety for $`(M,\nu )`$. Let $`\rho `$ be the distance function for $`X`$, introduced in Section 1. We define $`\mathrm{\Lambda }`$ as the displacement of $`G`$-action:
$$\mathrm{\Lambda }=\underset{zX}{inf}\underset{𝑖}{\mathrm{max}}\rho (g_iz,z).$$
Proposition 2.1.$`\mathrm{\Lambda }`$ is invariant under conjugation in $`𝒟iff_\nu (M)`$.
Proof.$`\rho `$ is $`𝒟iff_\nu `$-invariant.
Proposition 2.2.— Let $`M=(T^n,can)`$ and let $`G=`$ act by iterations of a hyperbolic element of $`SL(n,)`$. Then $`\mathrm{\Lambda }>0`$.
Proof.— The proof is based on an observation about Alexandrov non-positively curved spaces and a trick from \[Reznikov 4\].
Lemma.— Let $`X`$ be an Alexandrov non-positively curved space and let $`\varphi :XX`$ be an isometry which leaves invariant a geodesic $`\gamma `$ of $`X`$. Then the displacement of $`\varphi `$ is realized on the points of $`\gamma `$, that is, for $`y\gamma `$,
$$\rho (y,\varphi y)=\underset{xX}{\mathrm{min}}\rho (x,\varphi x).$$
Proof.— For $`xX`$ let $`y\gamma `$ be a point which realizes the distance from $`x`$ to $`\gamma `$. Then $`\rho (y,\varphi y)\rho (x,\varphi x)`$.
Now let $`X`$ be the twistor space of $`T^n`$ and let $`YX`$ be the space of metrics, invariant under shifts (we view $`T^n`$ as a Lie group). Then $`Y`$ is totally geodesic in $`X`$, because it is a manifold of fixed points of a family of isometries. As a Riemannian manifold, $`YSL(n)/SO(n)`$. Any hyperbolic matrix $`\varphi `$ by definition leaves invariant a geodesic in $`Y`$. The result follows.
A main result in the theory of invariant $`\mathrm{\Lambda }`$ is as follows.
Theorem 2.3.— Let $`G`$ be a Kazhdan group acting on a compact oriented manifold $`(M,\nu )`$ preserving a volume form $`\nu `$. Then either $`\mathrm{\Lambda }>0`$ or $`G`$ fixes a $`\mathrm{log}L^2`$-metric on $`M`$.
Proof.— Consider an isometric action on $`X`$. If the displacement function $`\underset{𝑖}{sup}\rho (g_iz,z)`$ is not bounded away from zero, then either there is a fixed point $`z_0X`$ for $`G`$, or $`G`$ is not Kazhdan, by a result of \[Kovevaar-Schoen 1\]. The result follows.
### 5.3 Non-linear superrigidity alternative
Theorem 3.1.—Let $`G`$ be either a semisimple Lie group of rank$`2`$, or $`Sp(n,1)`$ or $`Iso(a^2)`$. Let $`\mathrm{\Gamma }G`$ be a cocompact lattice. Let $`(M^n,\nu )`$ be a compact oriented manifold, on which $`\mathrm{\Gamma }`$ acts preserving the volume form $`\nu `$. Then either
a) $`\mathrm{\Gamma }`$ preserves a $`\mathrm{log}L^2`$\- metric on $`M`$, or
b) there exists a sequence $`g_0,g_1,\mathrm{}`$ of smooth Riemannian metrics on $`M`$ with volume form $`\nu `$ such that
$$_M\mathrm{log}A_i_{g_0}^2𝑑\nu \mathrm{},$$
where $`g_i=g_0(A_i,)`$ and
$$0<const_1<\underset{𝑗}{sup}_M\mathrm{log}B_{ij}_{g_i}^2d\nu <const_2,(i\mathrm{}),$$
where $`\gamma _j^{}g_i=g_i(B_{ij},)`$, $`\{\gamma _j\}`$ is a fixed finite set of generators for $`\mathrm{\Gamma }`$, or
c) there is a nonconstant totally geodesic $`\mathrm{\Gamma }`$-invariant map
$$\mathrm{\Psi }:G/KX,$$
where $`K`$ is a maximal compact subgroup of $`G`$.
Remarks.-
1) In case b) we say that $`\mathrm{\Gamma }`$ nearly fixes a $`\mathrm{log}L^2`$-metric on $`M`$.
2) the case c) implies, for $`G`$ simple, that $`dimG/KdimSL(n)/SO(n)`$, a so-called Zimmer conjecture.
3) for $`G=SL(m,),m3`$ and $`n=m`$, one deduces in case c) an existence of a measurable frame field $`\widehat{e}(x),\widehat{e}=(e_1,\mathrm{},e_n)`$, such that for almost all $`xM`$,
$$\pi (\gamma )_{}[\widehat{e}(x)]=\gamma \widehat{e}(\pi (\gamma )x)$$
where $`\gamma \mathrm{\Gamma }`$ and $`\pi (\gamma )`$ is an action of $`\gamma `$ on $`M`$.
4) Conversely, a standard action of $`SL(n,)`$ on $`T^n`$ does not satisfy a) (which is well-known) and b). To see this, we notice that $`SL(n,)`$ leaves invariant a totally geodesic space $`Y`$ introduced in the proof of Proposition 2.2. The argument of this proof implies that it is enough to show that the displacement function of the action of $`\mathrm{\Gamma }`$ on $`Y`$ diverges to $`\mathrm{}`$ as one escapes all compact sets of $`Y`$. This follows from the fact that $`Y`$ is a Riemannian symmetric space of non-compact type and $`\mathrm{\Gamma }`$ does not fix a point at infinity of $`Y`$.
5) The statement of Theorem constitutes a definite progress in the nonlinear superrigidity problem. There is still a mystery in the option b) where one would prefer a statement that $`\mathrm{\Gamma }`$ fixes a “point at infinity” of the space of metrics $`X`$, perhaps a measurable distribution of $`k`$-dimension planes, $`kn`$. At the time of writing this chapter (August, 1999) I am unable to make such a reduction.
Proof. follows a long-established tradition \[Siu 1\], \[Corlette 1\], \[Mok-Siu-Yeung 1\], see also a treatment of \[Jost-Yau 1\], in a new infinite-dimensional target context. If neither a) or b) holds then, accounting that $`\mathrm{\Gamma }`$ is Kazhdan, we deduce that the displacement function of $`\mathrm{\Gamma }`$ tends to infinity as one escapes all bounded sets in $`X`$. Let $`F\mathrm{\Gamma }G/K`$ be a flat fibration with fiber $`X`$, corresponding to the action of $`\mathrm{\Gamma }`$ in $`X`$. A theorem of \[Kovevaar-Schoen 1\], or \[Jost 1\] implies that there is a harmonic section of $`F`$. By Propositon 1.2 and main theorem of \[Corlette 1\] and \[Mok-Siu-Yeung 1\], this section must be totally geodesic. The result follows.
In the case of symplectic action of lattice $`\mathrm{\Gamma }`$ on a compact symplectic manifold $`(M,\omega )`$ we have a comletely similar theorem, as follows.
Theorem 3.2.— Let $`G`$ be either a semi-simple Lie group of rank$`2`$, or $`Sp(n,1)`$ or $`Iso(a^2)`$, $`\mathrm{\Gamma }`$ a cocompact lattice in $`G`$ which acts symplectically on a compact symplectic manifold $`(M^{2n},\omega )`$. Then either
a) $`\mathrm{\Gamma }`$ fixes a $`\mathrm{log}L^2`$ tamed almost-complex structure $`J`$, or
b) there exists a sequence of tamed smooth almost-complex structures $`J_iZ`$ with $`\rho (J_0,J_i)\mathrm{}`$ and
$$0<const_2<\underset{𝑗}{sup}\rho (\gamma _jJ_i,J_i)<const_1,\text{or}$$
c) there is a $`\mathrm{\Gamma }`$-invariant totally geodesic map
$$\mathrm{\Psi }:G/KZ.$$
Proof. is exactly as above.
In case c) and $`G`$ simple it follows that $`dimG/KdimSp(2n)/U(n)`$. If $`M=(T^{2n},can),G=Sp(2n,)`$ and case c) one deduces an existence of a measurable symplectic frame $`\widehat{e}(x)=(e_1,\mathrm{},e_{2n}(x))`$, such that for $`\gamma \mathrm{\Gamma },`$
$$\pi (\gamma )_{}[\widehat{e}(x)]=\gamma \widehat{e}(\pi (\gamma )(x)).$$
## Chapter 6 Kähler and quaternionic Kähler groups
In a letter to the author \[Deligne 1\] P.Deligne asked if one can extend the author’s theorem on rationality of secondary characteristic classes of a flat bundle over a projective variety to quasiprojective varieties. In 1994 the author was able to answer this question positively for the special case of noncompact ball quotients using an analytic technique of \[Gromov-Schoen 1\] and the scheme of the original proof for projective varietes. Here we present a full answer to Deligne’s question, Theorem 1.1, using an analytic technique of \[Jost-Zuo 1\], who produced harmonic maps of infinite energy but controlled growth.
We then turn to a well-known open problem of finding restriction on topology of compact quaternionic Kähler manifolds. In case of positive scalar curvature the situation is well-understood, but in case of negative scalar curvature the twistor spaces of \[Solomon 1\] are not Kähler and its technique fails. The only result known was a theorem of \[Corlette 1\] stating that the fundamental group does not have infinite linear representations unless the manifold is locally symmetric. Our result, Theorem 2.2, states that the fundamental group is Kazhdan. This is of course, a severe restriction (Kazhdan groups are rare). As a by-product of our technique, we obtain a new proof of a classical theorem, stating that the lattices in semisimple Lie groups of rank $`2`$, $`Sp(n,1)`$ and $`Iso(a^2)`$ are Kazhdan. We also show using I.1 that the classes in second cohomology space of a Kähler non-Kazhdan group, constructed in \[Reznikov 6\] and shown there nontrivial, are of polynomial growth. This again is very rare for “just a group”, as polynomial growth in cohomology is connected to a polynomial isoperimetric inequality in the Cayley graph, which needs a special reason to hold. This means Kähler groups are rare, too.
### 6.1 Rationality of secondary classes of flat bundler over quasiprojective varietes
A rationality theorem for secondary classes of flat bundles over compact Kähler manifolds (previously known as Bloch conjecture \[Bloch 1\]) has been proved in 1993 in \[Reznikov 3\] and \[Reznikov 5\]. In a letter to the author \[Deligne 1\] P.Deligne asked if one can prove such a statement for local system with logarithmic singularities over a quasiprojective variety. The answer happens to be yes.
Theorem 1.1.— Let $`X`$ be a quasiprojective variety, $`\rho :\pi _1(X)SL(n,)`$ a representation. Let $`b_i(\rho )`$ be the imaginary part and $`ChS_i(\rho )`$ the $`/`$-part of the secondary class $`c_i(\rho )H^{2i1}(X,/)`$ of the flat bundle with monodromy $`\rho `$. Then
A. $`b_i(\rho )=0(i2)`$ (the Vanishing Theorem).
B. $`ChS_i(\rho )H^{2i1}(X,/)`$ (the Rationality Theorem).
Proof.— For any smooth manifold, A implies B, as explained in the above cited papers. So we only prove A. Again it is explained in the above cited papers that we may assume $`\rho `$ to be irreducible. Then by a recent result \[Jost-Zuo 1\] an associated $`SL(n,)/SU(n)`$ flat bundle over $`X`$ possesses a pluriharmonic section $`s`$ which satisfies the Sampson degeneration condition. This means the following. The derivative $`Ds_x`$, $`xX`$ can be viewed as a $``$-linear map to the space $`P`$ of Hermitian matrices. Let $`(Ds_x)_{}^\pm (Y)=(Ds_x(Y)\pm \sqrt{1}Ds_x(\sqrt{1}Y))`$ be a map of $`TX`$ to $`P`$. Then the image of $`(Ds_x)_{}^\pm `$ consists of commuting matrices. Now a first proof of the Main Theorem in \[Reznikov 5\] applies word-to-word and the result follows.
### 6.2 Kazhdan property $`T`$ for Kähler and quaternionic Kähler groups
There are two ways to geometrize group theory. One approach (a time geometry in the terminology of \[Reznikov 7\]) is to consider finitely generated groups which act on a (usually compact) space with some structure (a volume form, a symplectic form, a tree, a circle, a conformal structure, etc). An amazing phenomenon, amply demonstrated in the previous chapters is that these groups tend to be not Kazhdan. Another approach (a space geometry) is to consider groups which are fundamental groups of a compact (or closed to compact) manifold with some structure (like Kähler). It happens that these groups tend to be Kazhdan. Therefore these two families of “geometric” groups are essentically disjoint. A following result is a main theorem of \[Reznikov 6\].
Theorem.— Let $`G`$ be a fundamental group of a compact Kähler manifold. If $`G`$ is not Kazhdan, then $`H^2(G,)0`$. Moreover, if $`H`$ is not Kazhdan and $`\psi :GH`$ is surjective then $`0\psi ^{}:H^2(H,)H^2(G,).`$
I would like to notice an important property, which I overlooked in \[Reznikov 6\].
Proposition 2.1.—Under the conditions of the Theorem, there is a nontrivial class of polynomial growth in $`H^2(G,).`$
Proof.— There is a unitary representation $`\rho :GU()`$ and a class $`lH^1(G,)`$ such that a class $`\gamma `$ in $`H^2(G,)`$ given by $`l,l`$ is nonzero, where $`,`$ is an imaginary part of the scalar product in $``$. This is proved in \[Reznikov 6\]. Now the result follows from Lemma I 1.1.
It is extremely rare for a finitely generated group to have nonzero polynomial cohomology.
We now turn to quaternionic Kähler manifolds. If a scalar curvature is positive, then the topology is very well understood \[Solomon 1\]. On the contrary, if the scalar curvature is negative, the only result known is that the fundamental group satisfy the geometric superrigidity \[Corlette 1\]. This means if $`\pi _1(X)`$ admits a Zariski-dense representation in an algebraic Lie group, then $`\pi _1(X)`$ is a lattice, and $`X`$ a symmetric space of a known type. However, it is a rare occasion for a group to have any finite dimensional linear representation with infinite image. Using a combination of ideas of \[Corlette 1\] and \[Reznikov 6\] which is based on \[Korevaar-Schoen 1\] we now prove a very stong structure theorem.
Theorem 2.2.—Let $`X`$ be a quaternionic Kähler manifold of negative scalar curvature. Then $`\pi _1(X)`$ is Kazhdan.
Proof.—Suppose not. Then by \[Korevaar-Schoen 1\] there exists an affine flat Hilbert bundle $`E`$ over $`X`$ with a nonparallel harmonic section. By a vanishing result of \[Corlette 1\] this section must be totally geodesic. Then $`X`$ must be covered by a flat torus, a contradiction.
Remark.— The same argument provides a new proof of the classical theorem \[Kazhdan 1\], \[Kostant 1\] , that the (cocompact) lattices in semisimple Lie groups of rank$`2`$, $`Sp(n,1)`$ and $`Iso(a^2)`$ are Kazhdan. One uses a vanishing result of \[Mok-Siu-Yeung 1\] ( see also a treatment in \[Jost-Yau 1\]) for lattices in semisimple Lie groups of rank$`2`$, and the above-mentioned result of \[Corlette 1\] for $`Sp(n,1)`$ and $`Iso(a^2)`$. Once established for cocompact lattices, the result follows for all lattices because a Lie group and a lattice in it are or are not Kazhdan at the same time.
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warning/0001/nucl-th0001032.html | ar5iv | text | # Collective modes in asymmetric nuclei
## Chapter 1 Introduction
Giant resonances are currently enjoying a resurgence of attention as a tool for the investigation of many particle effects in finite quantum systems. The study of giant resonances offers the possibility to learn more about nuclear forces. Of special interest are modern experiments with exotic nuclei which broaden our knowledge in a new degree of freedom: the variability in the proton to neutron ratio. The treatment of collective modes in nuclear matter is documented in an enormous literature starting from their discovery up to recent discussions - and citations therein.
We want to investigate the collective excitations in asymmetric nuclear matter by numerical simulation of the Vlasov equation and by comparison with the linear response theory. While the former leads to an insight into finite-size effects like surface, the latter allows to consider collisional correlations in a control-able way. Both methods are contrasted and we demonstrate that the collisional correlations as well as the surface effects are important to describe the experimental damping of giant resonances.
The damping mechanisms of collective motions in excited nuclei are a topic of continuing debate. Mainly two lines of thought are pursued. In one line of thought it is assumed that collisions are the only physical reason for damping which is described via a Fermi liquid approach ,\- . The other line of thought considers new features of the finite nucleus, such as surface oscillations and a level density with finite spacing. The investigations are performed without inertia \- or by including inertia ,\- ; note that inertia is absent in infinite matter.
Both classes of models predict a comparable degree of damping necessary to reproduce the experimental data. Consequently, it is an open question which is the correct physical reason for damping. Of course, the correct description has to assume a finite nucleus consisting of nucleons which are bound via the mean field, through which the nucleons undergo mutual collisions and where the surface is formed by the particles themselves. These features are principally included in Boltzmann-Uehling-Uhlenbeck-(BUU) simulations or in its nonlocal extensions . In full simulations, however, we will not gain a simple insight into the physical origins of the damping mechanism, in particular, how much is contributed by the surface and how much by collisions.
One aim of this article is therefore to compare both pictures in the frame of linear response theory. Within the collision-free Vlasov equation the linear response of finite systems is well known and allows one to calculate the strength function of finite nuclei. The resulting damping, however, does not reproduce the experimental damping of giant resonances since collisions are absent. This motivates us to develop a linear response theory including collisions.
While most of the theoretical treatments of oscillations rely on the linear response method or RPA methods, large amplitude oscillations require methods beyond this level. In particular the question of the appearance of chaos has recently been investigated . The hypothesis was established that the octupole mode is over-damped due to negative curved surface and consequent additional chaotic damping . Here we want to discuss in which conditions one might observe octupole modes at least in Vlasov - simulations of giant resonances which will turn out to be dependent on initial conditions and are consequently an effect beyond linear response.
The outline is as follows: In the second chapter we will give the numerical results of pseudo-particle simulation of the Vlasov equation. This will provide us with some insight into the magnitude of finite size effects on the damping. In the third chapter we present the linear response result of the Vlasov equation including collisional correlations. We will consider a simplified picture of infinite matter response and will consider the Steinwedel Jensen picture. We discuss two possibilities to include surface contributions in the linear response formalism. This allows us to describe the experimental temperature dependence of the damping of giant dipole resonances as well as the structure function of the isovector dipole resonances. The surface consideration allows to describe the isoscalar dipole resonance as a higher order mode. The comparison between linear response results and simulations is performed and finally we discuss nonlinear effects beyond linear response for the example of giant octupole modes.
## Chapter 2 Pseudoparticle simulation
We will describe the giant resonance first by a kinetic equation using a pseudo-particle simulation . The kinetic equation for the quasiclassical distribution function reads for neutrons (for protons analogously)
$`\dot{f}_n(𝐩,𝐫,t)+{\displaystyle \frac{𝐩}{m}}_𝐫f_n(𝐩,𝐫,t)_𝐫(U_n+U_{\mathrm{ext}})_𝐩f_n(𝐩,𝐫,t)=I_{\mathrm{corr}}`$
with the collisional integral $`I_{\mathrm{corr}}`$ discussed later and the self-consistent mean-field potential $`U`$ given by a schematic Skyrme type
$`U_{n/p}(\varrho ,I)`$ $`=`$ $`a(\frac{\varrho }{\varrho _0})+b(\frac{\varrho }{\varrho _0})^s\pm cI(\frac{\varrho }{\varrho _0})`$ (2.2)
with $`a=356\mathrm{M}\mathrm{e}\mathrm{V}`$, $`b=303\mathrm{M}\mathrm{e}\mathrm{V}`$ , $`c=54\mathrm{M}\mathrm{e}\mathrm{V}`$ and $`s=7/6`$. Here the neutron excess is $`I=\frac{\varrho _n\varrho _p}{\varrho _n+\varrho _p}`$, the neutrons feel the +–sign potential, the protons the opposite one.
By multiplying the kinetic equation by $`1,p`$ or $`E=\frac{p^2}{2m}+U`$ respectively one obtains the balance for particle density $`\rho `$, momentum density $`u`$ and energy density $``$. Since the collision integrals vanish for density and momentum balance we get the usual balance equations
$`{\displaystyle \frac{\rho (𝐫,t)}{t}}+{\displaystyle \frac{}{𝐫}}{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}\frac{E}{𝐩}f(𝐩,𝐫,t)}=0`$
$`{\displaystyle \frac{u_i(𝐫,t)}{t}}+{\displaystyle \frac{}{r_j}}{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}(p_i\frac{E}{p_j}f(𝐩,𝐫,t)+(𝐫,t)\delta _{ij})}=0`$ (2.3)
where the mean field energy of the system varies as
$`\delta `$ $`=`$ $`{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}\frac{\delta }{\delta f(𝐩,𝐫,t)}\delta f(𝐩,𝐫,t)}`$ (2.4)
$`=`$ $`{\displaystyle \frac{dp}{(2\pi \mathrm{})^3}(\frac{p^2}{2m}+U)\delta f(𝐩,𝐫,t)}`$
such that from (2.2) follows the total energy density as
$`(𝐫,t)={\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}\frac{p^2}{2m}f(𝐩,𝐫,t)}+{\displaystyle \frac{a}{2}}(\frac{\varrho ^2(𝐫,t)}{\varrho _0})+{\displaystyle \frac{b\rho (𝐫,t)}{s+1}}(\frac{\varrho (𝐫,t)}{\varrho _0})^{s+1}\pm {\displaystyle \frac{c}{2}}I(\frac{\varrho ^2(𝐫,t)}{\varrho _0}).`$ (2.5)
With the help of this quantity the balance of energy density reads from (2)
$`{\displaystyle \frac{(𝐫,t)}{t}}+{\displaystyle \frac{}{𝐫}}{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}E\frac{E}{𝐩}f(𝐩,𝐫,t)}={\displaystyle \frac{}{t}}E_{\mathrm{corr}}(𝐫,t)`$ (2.6)
with the correlation energy arising from the collisional side. The collisional side we will consider later in linear response. Here for the simulation we will neglect the collisions and will restrict ourselves to the Vlasov kinetic equation. In this way we will learn what are the effects of finite size and what are the effects of collisions.
Let us now first look at the input for the simulation. The mass number dependence of binding energy $`/A`$ from (2.5) is shown in figure 2.1. One sees that with increasing asymmetry the binding energy becomes weaker and consequently the compressibility decreases.
The compression modulus is defined as the derivative of the Energy per particle
$`K=9n^2{\displaystyle \frac{(/A)}{n^2}}`$ (2.7)
and took for $`I=0`$ the value $`K=200\mathrm{MeV}`$ according to (2.5). Figure 2.1 shows the asymmetry dependence of K. One sees that for neutron matter no binding energy would occur and a very low compressibility.
We solve now the kinetic equation by representing the distribution function $`f(𝐩,𝐫,t)`$ by a sum of N pseudo-particle distributions
$`f(𝐩,𝐫,t)f_0(𝐩,𝐫,t)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{1}{N}}f_S(𝐩𝐩_𝐢(t),𝐫𝐫_𝐢(t))`$ (2.8)
and use Gaussian pseudo-particles
$`f_S(𝐩𝐩_\mathrm{𝟏},𝐫𝐫_\mathrm{𝟏})=ce^{(𝐩𝐩_\mathrm{𝟏})^2/2\sigma _p^2}e^{(𝐫𝐫_\mathrm{𝟏})^2/2\sigma _r^2}`$ (2.9)
at $`𝐫_\mathrm{𝟏}`$ with momentum $`𝐩_\mathrm{𝟏}`$ . These pseudo-particles follow classical Hamilton equations
$`\dot{𝐩}_i`$ $`=`$ $`U,\dot{𝐫}_i={\displaystyle \frac{𝐩_i}{m}}.`$ (2.10)
The only fluctuations introduced are due to some small unavoidable numerical noise. We are using 300 pseudo-particles per nucleon and a pseudo-particle width of $`\sigma _r=0.466\mathrm{fm}`$. They are adjusted both in such a way that the experimental energy of the $`{}_{}{}^{40}Ca`$ giant monopole mode is reproduced. With these fixed two parameters the experimental behavior of centroid energy with mass number is than reproduced over the full range for giant monopole and giant dipole resonances. We have checked different numbers of test particles. The dependence of observables on the width is discussed in . Numerically the ground states of nuclei are realized by Wood- Saxon shapes of density and Fermi spheres in momentum.
### 2.1 Multipole analysis
Provided we have now solved the kinetic equation we will have the distribution function represented by $`N`$ pseudo-particles. Since we are interested in moments of the distribution function, the density, current and energy, we can expand these moments in the test particle representation as well. The momentum distribution of these moments can be obtained by spatial integration: $`F_a(𝐩)=𝑑𝐫af(𝐩,𝐫)`$ with, for the mass distribution, $`a=1`$, for isospin, $`a=\tau `$, for kinetic energy, $`a=\frac{p^2}{2m}`$, and for kinetic isospin energy, $`a=\frac{\tau p^2}{2m}`$ respectively. Decomposed into spherical coordinates they read
$`F_a(p,\vartheta ,\phi )`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{(2\pi )^3}{N}}a_i{\displaystyle \frac{\delta (pp_i)}{p_i^2}}\delta (\phi \phi _i){\displaystyle \frac{\delta (\vartheta \vartheta _i)}{\mathrm{sin}(\vartheta _i)}}.`$ (2.11)
Radial integration determines now a spherical distribution
$`\overline{F}_a(\vartheta ,\phi )`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{1}{N}}a_i\delta (\phi \phi _i){\displaystyle \frac{\delta (\vartheta \vartheta _i)}{\mathrm{sin}(\vartheta _i)}}`$ (2.12)
which can be decomposed into spherical harmonics
$`\overline{F}_a(\vartheta ,\phi )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}a_{lm}Y_{lm}(\vartheta ,\phi ),`$ (2.13)
$`a_{lm}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{1}{N}}a_iY_{lm}^{}(\vartheta _i,\phi _i).`$ (2.14)
The observable distributions $`\overline{F}_a(\vartheta ,\phi )`$ are normalized to $`\sqrt{4\pi }a_{00}`$, i.e. to mass number $`A`$, total isospin $`T`$, kinetic energy $`E_{kin}`$ and kinetic isospin energy $`E_{kinT}`$, respectively. The polar angular distribution of moments reads now from (2.12) and (2.13)
$`\widehat{F}_a(\vartheta )={\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{1}{N}}a_i{\displaystyle \frac{\delta (\vartheta \vartheta _i)}{\mathrm{sin}(\vartheta _i)}}{\displaystyle _0^{2\pi }}𝑑\phi \overline{F}_a(\vartheta ,\phi )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}a_{l0}\sqrt{\frac{2l+1}{4\pi }}P_l(\mathrm{cos}\vartheta )`$ (2.15)
with Legendre polynomials $`P_l`$. As a measure for the strength of the resonances we obtain now the coefficients $`a_{l0}`$ of the corresponding moment from (2.15) as
$`a_{l0}\sqrt{{\displaystyle \frac{2l+1}{4\pi }}}={\displaystyle \frac{2l+1}{2}}{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \frac{1}{N}}a_i{\displaystyle \frac{P_l(\mathrm{cos}\vartheta _i)}{\mathrm{sin}(\vartheta _i)}}.`$ (2.16)
These amplitudes of multipole moments, $`a_{l0}\sqrt{\frac{2l+1}{4\pi }}`$, are displayed as a function of time in figures 2.2,2.5,2.8 and 2.9. The value $`a_{10}`$ means the dipole moment, vanishing for isoscalar resonances, $`a_{20}`$ is characterizing the quadrupole oscillations and $`a_{30}`$ the octupole ones, etc.
### 2.2 Isoscalar monopoles
The first analyzed mode is the isoscalar giant monopole mode which plays an important role in the determination of nuclear compressibility. The connection of the compressibility and the energy of giant monopole resonances is discussed e.g. in .
The figure 2.2 shows simulation results for a monopole oscillation. Here the excitation has been performed by adding an extra momentum to the test-particles in the direction of the center of mass. We have excited in this way a clear monopole breathing mode which can be seen in the oscillation of the kinetic energy. The corresponding mean field energy performs the opposite oscillations that the total energy is constant. The fact that energy is oscillating between kinetic and correlational ones describes why we have here a compressional or breathing mode. All other modes remain unexcited. The finite value of the isospin and mass for $`l=0`$ reflects the conservation of isospin and particle number.
With the previously chosen fixed width of test-particles we see from figure 2.3 that the experimental mass number dependence of monopole oscillation is well reproduced over the whole range of mass numbers. However, the experimental damping can be seen to be largely underestimated by the Vlasov simulation which yields about $`2`$ MeV. This is a first indication that the mean field, even for a finite system, cannot account for the whole damping and dissipation. Instead we have to take into account collisional correlations which will be performed later.
It is now instructive to examine the isospin dependence of the isoscalar monopole. Since this is the experimental value which determines the nuclear compressibility, the isospin dependence is of direct importance. In figure 2.4 the dependence on asymmetry of the isoscalar monopole energy is plotted for a hard as well as a stiff equation of state. Of course, the hard equation of state leads to a higher monopole energy corresponding to a higher compressibility. Analogously to figure 2.1 the monopole energy decreases with the asymmetry. The hard equation of state leads to a more linear decrease while the soft equation of state remains almost unchanged up to a certain asymmetry and then decreases faster.
### 2.3 Isovector dipole resonances
On the next figure 2.5 the multipole analysis for $`{}_{40}{}^{90}Zr`$ oscillations is performed. The excitation is now chosen as an isovector dipole one created by a shift of proton against neutron spheres. One recognizes that a clear dipole $`a_{10}`$ mode is excited which is seen in the oscillation of the isospin and isospin energy. The kinetic energy now remains constant in contrast to the isoscalar mode because we have no compression mode.
Like the monopole case we now investigate the mass number dependence. With the same parameterization as for the monopoles, the experimental mass dependence of the energy is well reproduced, figure 2.6. The damping is again under-predicted. The Vlasov simulation leads to a nearly constant width of 2 MeV. This is again a hint that collisions cannot be neglected. The isospin dependency of the dipole mode is plotted in figure 2.7 and is characterized by an increasing width and decreasing centroid energy with increasing asymmetry.
### 2.4 Quadrupole resonances
Let us now investigate the quadrupole resonances. For these we distinguish isoscalar and isovector modes. The quadrupole resonances are excited by dividing the spatial distribution into two equal pieces which are accelerated in opposite directions. In the first example the neutrons and protons are in phase giving rise to an isoscalar mode.
We see from the kinetic energy in figure 2.8 that besides the clear isoscalar quadrupole mode there is also a weak isoscalar monopole mode excited. The isospin dependency of the isoscalar quadrupole mode is plotted in figure 2.7. As in the case of the isovector dipole mode one sees a slight decrease of the centroid energy and an increase of damping with increasing asymmetry.
In the second case we excite protons and neutrons out of phase which excites isovector quadrupole oscillations in figure 2.9. The isospin energy and the isospin show nice oscillations in the quadrupole multipole moment. There is a slight excitation of an isovector dipole mode too. We would like to remark that due to the test particle width there occurs a mode coupling which can be seen in the slight excitation of isoscalar modes in the kinetic energy. The isospin dependence shown in figure 2.7 has the same qualitative behaviour as the isoscalar mode, but is stronger in this case.
## Chapter 3 Collective modes in linear response
Now we would like to focus on small amplitude oscillations and we will develop the linear response theory. Large amplitude motions and effects beyond linear response are discussed in chapter 4. The simplest microscopic theory which provides basic experimental features and which allows to include collisional correlations is the Fermi-gas(liquid) model including dissipation. We will compare the results from this linear response with the simulation of finite nuclei, and both will be compared with the experimental data.
The principle of linear response is easily explained. When a system which is described by the kinetic equation (2) is disturbed by an external potential $`U^{\mathrm{ext}}`$ it will react and will create a density change $`\delta \rho `$. The connection of the latter to the external perturbation is called the response function, $`\delta \rho =\chi U^{\mathrm{ext}}`$. Without mean field in (2) we would obtain the polarization $`\mathrm{\Pi }`$ of the system. The mean field introduces a selfconsistency such that the relation between response function and polarization becomes
$`\chi ={\displaystyle \frac{\mathrm{\Pi }}{1\frac{\delta U}{\delta n}\mathrm{\Pi }}}`$ (3.1)
where we have assumed that the mean field $`U`$ is only density dependent. Therefore the potential for the response function is $`V=\frac{\delta U}{\delta n}`$. The collective mode of a system is now characterized by the condition that the denominator in (3.1) vanishes
$`ϵ(𝐪,\omega )=1V(𝐪)\mathrm{\Pi }(𝐪,\omega )=0`$ (3.2)
since then an infinitesimal small external potential can create a finite density oscillation due to the diverging response function. Therefore we will search for complex zeros of the dielectric (or dinuclic) function (DF), $`ϵ(𝐪,\mathrm{\Omega }+i\gamma )=0`$, which provides us with the energy $`\mathrm{\Omega }`$ and the damping $`\gamma `$ of the collective mode. The strength function is then given by
$`S(𝐪,\omega )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{Im}\mathrm{\Pi }}{(1V\mathrm{Re}\mathrm{\Pi })^2+(V\mathrm{Im}\mathrm{\Pi })^2}}.`$ (3.3)
The equation of state like the isothermal compressibility can be expressed by
$`\kappa ={\displaystyle \frac{1}{n^2}}\left({\displaystyle \frac{n}{\mu }}\right)_T={\displaystyle \frac{1}{n^2T}}\underset{q0}{lim}{\displaystyle \frac{d\omega }{\pi }\frac{1}{\mathrm{e}^{\omega /T}1}\mathrm{Im}\chi (q,\omega )}`$
(3.4)
and the fluctuations and diffusion coefficients can also be expressed via the response function.
Two lines of theoretical improvements of the response function can be found in recent publications. The first one starts from TDHF equations and considers the response of nuclear matter described by a non time-reversal Skyrme interaction . The other line tries to improve the response by the inclusion of collisional correlations and by considering multicomponent systems . In both lines of improvements have been combined into one expression. We follow this line and derive the response function from a kinetic equation including mean field (Skyrme) and collisional correlations. Therefore we consider interacting matter which can be described by an energy functional ( mean field) $``$ originally introduced by Skyrme and the residual interaction. The latter we condense in a collisional integral additional to the TDHF equation. Then the response to an external perturbation will contain the effect of Skyrme mean field and additionally the effect of the residual interaction.
### 3.1 Response function for asymmetric matter
The system now consists of a number of different species (neutrons, protons, etc…) interacting with their own kind and with the others. It is important to consider the interaction between different sorts of particles if we want to include friction between different streams of isospin components, etc. In particular the isospin current may not be conserved in this way. Let us start with a set of coupled quantum kinetic equations<sup>1</sup><sup>1</sup>1 The quasiclassical Landau equation (2) follows from gradient expansion as $`{\displaystyle \frac{}{t}}f+_𝐩ϵ_𝐫f_𝐫ϵ_𝐩f={\displaystyle \frac{ff^{\mathrm{l}.\mathrm{e}.}}{\tau }}`$ (3.5) for the reduced density operator $`\widehat{\rho }_a`$ of a certain species denoted by the subscript $`a`$
$`_t\widehat{\rho }_a(t)=i[\widehat{\rho }_a,\widehat{E}_a(t)+\widehat{U}_a^{\mathrm{ext}}]{\displaystyle \underset{b}{}}{\displaystyle \underset{0}{\overset{t}{}}}𝑑t^{}{\displaystyle \frac{\widehat{\rho }_a(t^{})\widehat{\rho }_b^{\mathrm{l}.\mathrm{e}.}(t^{})}{\tau _{ab}(tt^{})}}`$ (3.6)
where $`\widehat{E}=\widehat{P}^2/2m+\widehat{U}`$ denotes the kinetic as well as mean field energy operator and the external field which is assumed to be a nonlinear function of the density. We approximate the collision integral by a non-Markovian relaxation time which is derived in appendix A. This turned out to be necessary to reproduce damping of zero sound . It accounts for the fact that during a two-particle collision a collective mode can couple to the scattering process. Consequently, the dynamical relaxation time represents the physical content of a hidden three-particle process and is equivalent to memory effects.
Furthermore, we assume relaxation towards a local equilibrium
$`f_a^{\mathrm{l}.\mathrm{e}.}(𝐩,𝐑,t)={\displaystyle \underset{q}{}}\mathrm{e}^{i\mathrm{𝐪𝐑}}p+{\displaystyle \frac{q}{2}}|\widehat{\rho }_a^{\mathrm{l}.\mathrm{e}.}|p{\displaystyle \frac{q}{2}}=f_0\left({\displaystyle \frac{\epsilon _a(𝐩𝐐_a(𝐑,t))\mu _a(𝐑,t)}{T_a(𝐑,t)}}\right)`$ (3.7)
with the Fermi distribution $`f_0`$. This local equilibrium will be specified here only by a small deviation of the chemical potential of species $`a`$ ensuring density conservation . The extension of the method including further conservation laws and specifying also the local current and the temperature can be found in .
Linearizing the kinetic equation (3.6) we obtain the matrix equation for the density deviation $`\delta \rho _b`$ due to an external perturbation $`U_{\mathrm{ext}}^b`$
$`{\displaystyle \underset{b}{}}\delta n_b\left\{\delta _{ab}{\displaystyle \frac{i}{\omega \tau _a+i}}\left[\delta _{ab}{\displaystyle \frac{\tau _a}{\tau _b}}C_{ab}\right]\mathrm{\Pi }_a(\omega +{\displaystyle \frac{i}{\tau _a}})\alpha _{ab}\right\}=\mathrm{\Pi }_a(\omega +{\displaystyle \frac{i}{\tau _a}})U_{\mathrm{ext}}^a.`$ (3.8)
The matrix $`C_{ab}`$ is given by
$`C_{ab}={\displaystyle \underset{c}{}}\left\{{\displaystyle \frac{1}{\tau }}\right\}_{ac}{\displaystyle \frac{\mathrm{\Pi }_c\left[(\omega +\frac{i}{\tau _a})\frac{m_a}{m_c}\right]}{\mathrm{\Pi }_c(0)}}\left\{{\displaystyle \frac{1}{\tau }}\right\}_{cb}^1`$ (3.9)
with $`\alpha _{ab}=\frac{}{\rho _b}U_a`$ the linearization of the mean field with respect to the deviation of density from equilibrium value caused by the external perturbation $`U_{\mathrm{ext}}`$. The partial polarization function of species $`a`$ is
$`\mathrm{\Pi }_a(\omega )=2{\displaystyle \frac{dp}{(2\pi )^3}\frac{f_a(p+\frac{q}{2})f_a(p\frac{q}{2})}{ϵ_a(p+\frac{q}{2})ϵ_a(p\frac{q}{2})\omega +i0}}.`$ (3.10)
The factor 2 in front of the integral accounts for the spin degeneracy according to $`\rho _a=2\frac{dp}{(2\pi )^3}f_a(p)`$.
Equation (3.8) represents the complete polarization of the system, because $`\delta \rho =\mathrm{\Pi }U_{\mathrm{ext}}`$. It represents a matrix equation which is solved easily. The collective modes are given by the zeros of the determinant of the matrix on the left hand side of (3.8) because these are the poles of the polarization matrix. Since we took into account relaxation processes between all species we are able to cover current - current friction. For a two component system, e.g. neutrons with density $`\rho _n`$ and protons with density $`\rho _p`$, we write explicitly
$`(1\mathrm{\Pi }_n^\mathrm{M}\alpha _{nn})(1\mathrm{\Pi }_p^\mathrm{M}\alpha _{pp})(D_{np}+\mathrm{\Pi }_n^\mathrm{M}\alpha _{np})(D_{pn}+\mathrm{\Pi }_p^\mathrm{M}\alpha _{pn})=0`$ (3.11)
with the generalization of the Mermin polarization function to a multicomponent system
$`\mathrm{\Pi }_a^\mathrm{M}={\displaystyle \frac{\mathrm{\Pi }_a(\omega +\frac{i}{\tau _a})}{1\frac{i}{\omega \tau _a+i}(1C_{aa})}},`$ (3.12)
and the additional coupling due to asymmetry in the system
$`D_{np}={\displaystyle \frac{\tau _n}{\tau _p}}{\displaystyle \frac{C_{np}}{C_{nn}i\omega \tau _n}}.`$ (3.13)
The $`D_{pn}`$ are given by interchanging species indices. This term does not appear for symmetric matter. Therefore we call this term the asymmetry coupling term further on. The result (3.11) represents the generalization of the dispersion relation (3.2) to a two - component system.
It is illustrative to recover known results for symmetric nuclear matter. This is performed for the case of equal relaxation times $`\tau _p=\tau _n=\tau `$ and equal deviation from the mean field $`\alpha _1=\alpha _{nn}=\alpha _{pp}`$ and $`\alpha _2=\alpha _{np}=\alpha _{pn}`$. Eq. (3.11) takes then the known Mermin form of dispersion relation
$`1(\alpha _1\pm \alpha _2){\displaystyle \frac{\mathrm{\Pi }(\omega +\frac{i}{\tau })}{1\frac{i}{\omega \tau +i}\left[1\frac{\mathrm{\Pi }(\omega +\frac{i}{\tau })}{\mathrm{\Pi }(0)}\right]}}=0`$ (3.14)
with the isovector mode $`\alpha _1\alpha _2`$ and the isoscalar mode $`\alpha _1+\alpha _2`$. Please note that for zero temperature the Mermin expression (3.14) agrees with the response function derived from taking into account multipole expansion in momentum of the disturbed distribution .
We have presented a general dispersion relation for the multicomponent system including known special cases. The dispersion relation (3.11) is similar to the one derived recently in if we neglect the collisional coupling $`D_{np}`$. Also a more general polarization function (3.12) is presented here including collisions within a conserving approximation . In the following we will apply this expression to the damping of giant dipole resonances in symmetric as well as asymmetric nuclear matter.
### 3.2 Application to nuclear matter
Before we can solve the dispersion relation (3.11) we have to specify the wave vector characteristic of the considered mode. This is performed according to the Steinwedel-Jensen model. We assume that the surface remains constant and the density oscillation obeys a wave equation with the boundary condition that the radial velocity vanishes on the spherical surface with radius $`R=1.13A^{1/3}`$. This leads to $`j_l^{}(kR)0`$ with the spherical Bessel function of order $`l=0,1\mathrm{}`$ associated with the monopole, dipole… resonances. From this condition one obtains a connection between the wave vector and the radius of the nuclei or the mass number in the form
$`k_{i,l}={\displaystyle \frac{c_{il}}{1.13A^{1/3}}}`$ (3.15)
for the $`i`$-th mode of multipolarity $`l`$ where $`c_{il}`$ is the $`i`$-th zero of $`j_l^{}(c)0`$.
For the monopole modes as compression modes we have no zero of first order for $`l=0`$. For the dipole modes one has in first order $`k=2.08/R`$ which describes the giant isovector dipole resonance (IVGDR) while the ISGDR is a spurious mode in first order. This would just mean an unphysical oscillation of center of mass motion. However, in second order $`k=5.94/R`$ the isoscalar giant dipole resonance (ISGDR) has been observed recently ( and references therein). This can be considered as a density oscillation inside a sphere as we will discuss in chapter 3.6.
Besides the occurring wave vector we have also to specify the relaxation times which contain the effect of collisions. The dynamic relaxation time has been derived by Sommerfeld expansion in appendix A as
$`{\displaystyle \frac{1}{\tau _{ab}^{\mathrm{gas}}(\omega )}}`$ $`=`$ $`{\displaystyle \frac{1}{\tau _{ab}(0)}}\left[1+{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{\omega }{\pi T}}\right)^2\right]{\displaystyle \frac{1}{\tau _{ab}^{\mathrm{liq}}(\omega )}}={\displaystyle \frac{2}{\tau _{ab}(0)}}\left[1+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\omega }{\pi T}}\right)^2\right]`$ (3.16)
for $`a,b`$ neutrons or protons respectively and dependent on whether we use a Fermi liquid or Fermi gas model. The Markovian relaxation time is given in terms of the cross section $`\sigma _{ab}`$ between species $`a`$ and $`b`$ as $`\tau _{ab}^1=\frac{4m}{3\mathrm{}^3}\sigma _{ab}T^2`$.
We will use as an illustrative example a slightly different mean field potential than (2.2) given by
$`U_a=t_0\left[(1+{\displaystyle \frac{x_0}{2}})(\rho _n+\rho _p)(x_0+{\displaystyle \frac{1}{2}})\rho _a\right]+{\displaystyle \frac{t_3}{4}}\left[(\rho _n+\rho _p)^2\rho _a^2\right]`$ (3.17)
with $`a=n,p`$ the density of neutrons or protons respectively. The Coulomb interaction leads to an additional contribution for the proton mean-field
$`U_p^C(q)={\displaystyle \frac{4\pi e^2}{q^2}}\rho _p(q).`$ (3.18)
The model parameters used here reproduce the Weizsäcker formula
$`{\displaystyle \frac{}{A}}=a_1+{\displaystyle \frac{a_2}{A^{1/3}}}+{\displaystyle \frac{a_3Z^2}{A^{4/3}}}+a_4I^2`$ (3.19)
with the volume energy $`a_1=15.68`$ MeV, Coulomb energy $`a_3=0.717`$ MeV and the symmetry energy $`a_4=28.1`$ MeV with the asymmetry parameter $`I=\frac{n_nn_p}{n_n+n_p}`$.
The energy and damping rates are now determined by the zeros of the (Mermin) polarization function (3.14). First we plot the solution of the dispersion relation (3.11) for symmetric nuclear matter. In figure 3.1 we have plotted the real and imaginary (FWHM) part of complex energy for different approximations with relaxation time (collisions) and with and without Coulomb mean field (3.18). In figure 3.1 (A) we find that the inclusion of Coulomb effects reproduces the experimental shape of the centroid energies at higher mass numbers (dot-dashed line). Taking only collisions into account fails to describe higher mass numbers (solid line). Considering Coulomb together with collisions (dashed line), the centroid energies are reduced towards the data.
The experimental values of the damping rates are also presented versus mass number (figure 3.1 (B)). We recognize that the FWHM is just twice the damping rate $`\mathrm{\Gamma }=2\gamma `$ which has been recently emphasized . It has to be stressed that the experimental data are accessible by this FWHM. Here we have considered only collisions and have a vanishing Landau damping for the infinite matter model .
### 3.3 Comparison between simulation and linear response
We have already stressed that the simulation of finite nuclei neglecting collisional correlations lead to a damping of about $`\mathrm{\Gamma }_{\mathrm{Land}}=2`$MeV for the giant dipole resonance at zero temperature. This was attributed to the surface effects of spherical symmetric nuclei. We consider this as the numerical value of the contribution attributed to the wall at zero temperature. For finite temperature the will be considered an additional contribution due to thermally deformed shapes in section 3.5. Nevertheless it is now interesting to compare the isospin dependences of the simulation without collisions and the linear response including collisions but no surface effects.
In figure 3.2 we see that the damping is slightly increasing with increasing asymmetry for simulations without collisions while the collisional contributions cause a slight decrease of damping. This decrease is much more pronounced in the second order isoscalar dipole mode. The energy in turn shows the same behaviour for simulations without correlations and for linear response with correlations; it slightly decreases the value.
The same qualitative feature can be recognized for the quadrupole modes in figure (3.3). However the importance of collisional damping is in this case much larger than the finite size effects.
### 3.4 Giant resonances in excited nuclei
We consider now the temperature dependence of the GDR. We have found that the centroid energy only slightly decreases with increasing energy. In Fig. 3.4 the theoretical damping rates $`\mathrm{\Gamma }=2\gamma `$ of the IVGDR modes in <sup>120</sup>Sn and <sup>208</sup>Pb are plotted as a function of temperature together with experimental values. The results of the Fermi gas model (A.11) and the Fermi liquid model (A.19) are very close and in good agreement with the data at low temperature. The higher temperature dependence still remains too flat compared to the experiments. The small difference between both models for T=0 vanishes with increasing temperature.
While the figures indicate that the linear response including collisional correlations together with the Steinwedel Jensen model can describe quite well the gross features of giant dipole resonances we should keep in mind that we have neglected here the surface effects. From the simulation results of the second chapter we have seen that there is a damping of about $`\mathrm{\Gamma }_{\mathrm{Land}}=2`$MeV for finite nuclei at zero temperature without collisional correlations. Simply adding now both contributions would overestimate the experimental damping. We would like to point out that the inclusion of only density conservation in the derivation of the response function so far is overestimating the width. In we have shown that the inclusion of momentum conservation reduces this width by approximately $`2`$MeV. Therefore we might anticipate that both the finite size effects together with the collisional contribution lead to the correct value for low temperatures.
In figure 3.5 we see that the temperature increase is too flat compared to the experimental finding if we consider only collisional damping (thin lines). From simulation results of ground state we have already found a contribution from finite size effects. We might now expect that the discrepancy in temperature behaviour is due to finite size effects. This should be explained and understood within a more simplified model.
### 3.5 Nuclear surface contribution - excited nuclei
Besides the $`\mathrm{\Gamma }_{\mathrm{Land}}=2`$MeV of the simulation in the ground state we can have additional damping due to shape deformation if we increase the temperature. In Ref. we have presented a method to include also scattering with the non-spherical nuclear surface. This improves the temperature dependence remarkably (thick lines) in figure 3.5.
The idea is to consider the surface scattering contribution to the damping determined by the additional Lyapunov exponent due to the surface. It has been shown that such Lyapunov exponents appear in the response function additive to the frequency as imaginary shift provided the Lyapunov exponent is small compared to the product of wave vector times Fermi velocity
$`\mathrm{\Pi }_0^{\mathrm{surf}}(q,R,\omega )=\mathrm{\Pi }_0^{\mathrm{inf}}(q,p_f(R),\omega +i\lambda ).`$ (3.20)
The Lyapunov exponent by itself is calculated for different deformations of the surface corresponding to different temperatures
$`R_\lambda (\theta )=R_0\left(1+\alpha _{00}+\alpha _\lambda P_\lambda (\mathrm{cos}(\theta ))\right)`$ (3.21)
with the nuclear radius $`R_0=1.13A^{1/3}`$fm, and where $`\lambda =2`$ corresponds to the quadrupole and $`\lambda =3`$ to the octupole deformation . The corresponding mean deformation <sup>2</sup><sup>2</sup>2For small deviations we found identical Lyapunov exponents for prolate $`\alpha >0`$ and oblate $`\alpha <0`$ deformations and therefore we do not distinguish the sign of $`\alpha `$. is linked to the temperature
$`\alpha ={\displaystyle \frac{𝑑\alpha |\alpha |\mathrm{exp}(E_B(\alpha )/T)}{𝑑\alpha \mathrm{exp}(E_B(\alpha )/T)}}`$ (3.22)
where the surface dependent energy $`E_B(\alpha )`$ is given by the Bethe-Weizsäcker formula <sup>3</sup><sup>3</sup>3Please remember that in principle the Coulomb energy changes with small deformation as well according to the factor $`15(\lambda 1)/(2\lambda +1)^2\alpha _\lambda ^2`$ while the surface term changes as $`1+(\lambda 1)(\lambda +2)/2/(2\lambda +1)\alpha _\lambda ^2.`$ Only the latter correction is considered since the Coulomb energy deformation would lead to corrections of around $`0.3\%`$ and are neglected here.
$`E_B(\alpha )=a_1+{\displaystyle \frac{a_2}{A^{1/3}}}+{\displaystyle \frac{a_3Z^2}{A^{4/3}}}+a_4I^2+a_5A^{2/3}{\displaystyle \frac{S(\alpha )}{S(0)}}`$
with the volume energy $`a_1=15.68`$ MeV, Coulomb energy $`a_3=0.717`$ MeV, the symmetry energy $`a_4=28.1`$ MeV and the surface energy $`a_5=18.56`$ MeV. For the quadrupole, $`S_2`$, and octupole, $`S_3`$, deformations one gets
$`{\displaystyle \frac{S_2(\alpha )}{S(0)}}`$ $`=`$ $`1+{\displaystyle \frac{2}{5}}\alpha ^2+𝒪(\alpha ^3)`$
$`{\displaystyle \frac{S_3(\alpha )}{S(0)}}`$ $`=`$ $`1+{\displaystyle \frac{5}{7}}\alpha ^2+𝒪(\alpha ^3).`$ (3.24)
This represents the lowest order expansion in $`\alpha `$, however the next term gives already corrections in fractions of percent for the highest deformations considered here. By this way, the statistical model (3.22) leads to a connection between temperature and deformation as
$`T=c\pi a_5A^{2/3}<\alpha >^2`$ (3.25)
where the constant $`c`$ is given by the coefficient of $`\alpha ^2`$ in Eq. (3.24) for the corresponding quadrupole or octupole deformation.
Using Eq. (3.25) we can translate the Lyapunov exponent $`\lambda `$ calculated as a function of deformation into a function of the temperature. In figure 3.6 the contribution to the damping of IVGDR for <sup>120</sup>Sn (circles) and <sup>208</sup>Pb (squares) is presented for different shape deformations versus temperature. If we add quadrupole and octupole deformations we come up with a damping curve very similar to . The damping starts at zero and increases rapidly with increasing temperature. We see that the main contribution comes from the quadrupole deformation while the octupole deformation is only sizeable at higher temperature. Let us note that the qualitative difference between Sn and Pb is reproduced by including surface scattering as well as collisional damping.
The collisional contribution $`\mathrm{\Gamma }=2\gamma `$ is plotted as well (Sn: solid line, Pb: dashed line). We recognize that both contributions by themselves, collisional as well as surface scattering, account almost for the same amount required by the experimental values, see figure 3.6. A proper relative weight between both processes is therefore necessary which will be introduced in the following.
Let us note before that we have considered the surface contribution to the damping for a deformed surface by temperature which therefore vanishes at zero temperature. However, for a spherical surface we got already from simulation a zero temperature contribution of $`\mathrm{\Gamma }_{\mathrm{Land}}=2`$MeV. Consequently for the total damping caused by surface both contribute.
So far we have not considered that only particles close to the surface can appreciably contribute to the surface scattering and to this chaotization process, while particles deep inside the nuclei are screened out of this process. Consequently we consider the corresponding collision frequencies as the measure to compare surface collisions with inter-particle collisions. The collision frequency between particles is given by $`1/\tau _0`$ of (3.16). The collision frequency of particles with the deformed surface beyond a sphere, $`\nu _{\mathrm{surf}}`$, is given by the product of the density with the surface increase $`S(\alpha )S(0)=c\alpha ^24\pi R_0^2`$ according to Eq. (3.24) and with the mean velocity in radial direction $`v_r=3/8v_F`$. The result is
$`\nu _{\mathrm{surf}}=1.5Tn_0v_Fr_0^2/a_5`$ (3.26)
where we have used Eq. (3.25) to replace $`\alpha `$. We see that the frequency (3.26) is independent of the size of the nucleus and linearly dependent on the temperature. We use the ratio of these two frequencies to weight properly the two damping mechanisms, the surface collisional, $`\lambda `$, and the inter-particle collisional, $`\gamma `$, contributions. Consequently the full - width - half - maximum (FWHM) reads
$`\mathrm{\Gamma }_{\mathrm{FWHM}}`$ $`=`$ $`2\zeta \gamma +2(1\zeta )(\lambda )+2\gamma _{\mathrm{Land}}`$ (3.27)
$``$ $`\mathrm{\Gamma }_{\mathrm{coll}}+\mathrm{\Gamma }_{\mathrm{surf}}+\mathrm{\Gamma }_{\mathrm{Land}}.`$
With the help of (3.16) and (3.26) the weighting factor $`\zeta `$ is given by
$`\zeta (T)={\displaystyle \frac{\frac{1}{\tau _0(T)}}{\frac{1}{\tau _0(T)}+\nu _{\mathrm{surf}}(T)}}.`$ (3.28)
One sees that for zero and high temperatures $`\zeta =1`$ and due to Eq. (3.27) only the collisional contributions matter. Since $`\nu _{\mathrm{surf}}`$ is linear in the temperature and $`1/\tau `$ depends quadratically on the temperature, the weighting factor $`\zeta `$ has a minimum at temperatures around $`T_c=\frac{\sqrt{3}}{2\pi }\omega `$ for the gas model (3.16) and the surface contributions become important. In the case of the IVGDR this corresponds to a temperature of $`T3.7`$MeV, which is the upper limit of currently achievable experimental temperatures. Therefore we can state that at low and high temperatures the collisional damping is dominant while for temperatures around $`T_c`$ the surface contribution becomes significant.
In figure 3.5 we have compared the effective damping according to Eq. (3.27) with the experimental data. Let us recall that the collisional damping value is reduced here since we include momentum conservation too for the Mermin polarization function. On the other side the surface damping consists now of the deformation contribution by temperature and by the zero temperature contribution obtained by simulation of Vlasov equation. We find a reasonable quantitative agreement combining all these parts. This is illustrated more in detail in Fig. 3.7 where we have plotted the strength function (3.3) for <sup>120</sup>Sn (LHS) and <sup>208</sup>Pb (RHS) within the Fermi gas model (A.11) (dashed lines) and Fermi liquid model (A.19) (solid lines) with the normalized data from Ref. .
The good overall agreement of the shape evolution with the experiment is again accompanied by only minor differences between the Fermi gas and the Fermi liquid model.
### 3.6 Higher order modes - surface modes
The existence of isoscalar giant dipole resonance (ISGDR) in nuclear matter is considered as a spurious mode in most text books since one associates with it a center of mass motion. The more surprising was the experimental justification of a giant resonance carrying the quantum numbers of a isoscalar and dipole mode . Consequently one has to consider higher harmonics as explanation of such a mode . Usually this mode is associated with a squeezing mode analogous to a sound wave .
In this chapter we want to discuss the influence of surface effects on the ISGDR compression mode. We will show that even in the frame of the Fermi liquid model such modes can be understood. Moreover we claim that the surface effects are not negligible for reproducing the strength function. While we have already given a phenomenological approach to include surface effects we want to show now a straightforward possibility to include surface effects in the response function. Consequently we give first a short derivation of response function including surface effects. This will result in a new formula in the temperature-dependent extended Thomas-Fermi approximation .
The starting point is again the semi-classical Vlasov equation (2). Instead of using now a spatial homogeneous equilibrium as done so far we consider now explicitly the spatial dependence. Provided we know the response to the external potential without self-consistent mean field, which is described by the polarization function $`\mathrm{\Pi }`$
$`\delta n(𝐱,t)={\displaystyle 𝑑𝐱^{}\mathrm{\Pi }(𝐱,𝐱^{},t)U^{\mathrm{ext}}(𝐱^{},t)},`$ (3.29)
the response including mean field, $`\chi `$, is given by
$`\chi (𝐱,𝐱^{},t)=\mathrm{\Pi }(𝐱,𝐱^{},t)+{\displaystyle 𝑑𝐱_1𝑑𝐱_2\mathrm{\Pi }(𝐱,𝐱_1,t)\frac{\delta U^{\mathrm{ind}}(𝐱_1,t)}{\delta n(𝐱_2,t)}\chi (𝐱_2,𝐱^{},t)}.`$ (3.30)
Therefore we concentrate first on the calculation of the polarization function $`\mathrm{\Pi }`$ and linearize the Vlasov equation (2) according to $`f(𝐩,𝐑,t)=f_0(𝐩,𝐑)+\delta f(𝐩,𝐑,t)`$ such that the induced density variation $`\delta \rho (𝐑,t)=𝑑𝐩\delta f/(2\pi \mathrm{})^3`$ reads
$`\delta \rho (𝐱,\omega )={\displaystyle \frac{d𝐪}{(2\pi \mathrm{})^3}\mathrm{e}^{i\mathrm{𝐪𝐱}}\frac{d𝐩d𝐱^{}}{(2\pi \mathrm{})^6}\mathrm{e}^{i\mathrm{𝐪𝐱}^{}}\frac{_pf_0(𝐩,𝐱^{})_x^{}U^{\mathrm{ext}}(𝐱^{},t)}{i(\omega \frac{\mathrm{𝐩𝐪}}{m})}}.`$
(3.31)
Here we have employed the Fourier transform of space and time coordinates of (2) to solve for $`\delta f`$ and inverse transform the momentum into the form (3.31). Comparing (3.31) with the definition of the polarization function (3.29) we extract with one partial integration
$`\mathrm{\Pi }(𝐱,𝐱^{},\omega )=_x^{}{\displaystyle \frac{d𝐩d𝐪}{(2\pi \mathrm{})^6}\mathrm{e}^{i𝐪(𝐱𝐱^{})}\frac{_pf_0(𝐩,𝐱^{})}{i(\omega \frac{\mathrm{𝐩𝐪}}{m})}}.`$ (3.32)
With (3.32) and (3.30) we have given the polarization and response functions for a finite system.
In the following we are interested in the gradient expansion since we believe that the first order gradient terms will bear the information about surface effects. Therefore we change to the center of mass and difference coordinates $`𝐑=(𝐱_1+𝐱_2)/2`$, $`𝐫=𝐱_1𝐱_2`$ and retaining only first order gradients we get from (3.32) after Fourier transform of $`𝐫`$ into $`𝐪`$
$`\mathrm{\Pi }(𝐑,𝐪)={\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}\frac{𝐪_pf_0(𝐩,𝐑)}{\omega \frac{\mathrm{𝐩𝐪}}{m}}}+{\displaystyle \frac{i}{2}}_R{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}\left(\frac{_pf_0(𝐩,𝐑)}{\omega \frac{\mathrm{𝐩𝐪}}{m}}\frac{𝐩}{m}\frac{𝐪._pf_0(𝐩,𝐑)}{(\omega \frac{\mathrm{𝐩𝐪}}{m})^2}\right)}`$ (3.33)
$`=`$ $`{\displaystyle \frac{d𝐩}{(2\pi \mathrm{})^3}[𝐪i_R(1+\frac{\omega }{2}_\omega )]\frac{_pf_0(\frac{p^2}{2m},𝐑)}{\omega \frac{\mathrm{𝐩𝐪}}{m}}}`$
where in the last equality we have assumed radial momentum dependence of the distribution function $`f_0`$. We recognize that besides the usual Lindhard polarization function as the first part of (3.33) we obtain a second part which is expressed by a gradient in space. The first part corresponds to the Thomas-Fermi result where we have to use the spatial dependence in the distribution functions and the second part represents the extended Thomas-Fermi approximation. So far we did not assume any special form of the distribution function. Therefore the expression (3.33) is also valid for any high temperature polarization of finite systems.
What remains to be shown is that the response function (3.30) does not contain additional gradients. This is easily confirmed by two equivalent formulations of (3.30), $`\mathrm{\Pi }^1\chi =1+V\chi `$ and $`\chi \mathrm{\Pi }^1=1+\chi V`$, which by adding yield the anticommutator
$`[\mathrm{\Pi }^1,\chi ]_+=2+[V,\chi ]_+.`$ (3.34)
This anticommutator does not contain any gradients up to second order. Therefore we have \[$`V=\delta U^{\mathrm{ind}}/\delta n`$\]
$`\chi (𝐑,𝐪,\omega )={\displaystyle \frac{\mathrm{\Pi }(𝐑,𝐪,\omega )}{1V(𝐑,𝐪,\omega )\mathrm{\Pi }(𝐑,𝐪,\omega )}}+𝒪(_R^2).`$ (3.35)
Equation (3.35) and (3.33) give the response and polarization functions of finite systems in first order gradient approximation.
Now we are ready to derive approximate formulae for spherical nuclei. In this case we can assume $`𝐪||𝐑`$ and we have
$`\mathrm{\Pi }(R,𝐪,\omega )=\mathrm{\Pi }^0(R,𝐪,\omega ){\displaystyle \frac{i}{q}}_R\left[1+{\displaystyle \frac{\omega }{2}}_\omega \right]\mathrm{\Pi }^0(R,𝐪,\omega )`$ (3.36)
where $`\mathrm{\Pi }_0`$ is the usual Lindhard polarization with spatial dependent distributions (chemical potentials, density). We use now further approximations. In the case of giant resonances we are in the regime of small $`q`$ and $`\mathrm{Im}\mathrm{\Pi }^0\omega `$ such that
$`\left[1+{\displaystyle \frac{\omega }{2}}_\omega \right]\mathrm{\Pi }^0=i{\displaystyle \frac{3}{2}}\mathrm{Im}\mathrm{\Pi }^0`$ (3.37)
since $`\mathrm{Re}\mathrm{\Pi }^01c^2q^2/\omega ^2`$. Within the local density approximation we know that the spatial dependence is due to the density $`\rho (R)=\rho _0\mathrm{\Theta }(R_0R)`$. Since we have for zero temperature $`\mathrm{Im}\mathrm{\Pi }^0p_f(\rho )`$ we evaluate
$`_R\mathrm{Im}\mathrm{\Pi }^0=n_0\delta (R_0R)_n\mathrm{Im}\mathrm{\Pi }^0={\displaystyle \frac{1}{3}}\delta (R_0R)\mathrm{Im}\mathrm{\Pi }^0`$
(3.38)
where we assumed the density dependence carried only by the Fermi momentum. Now it is straightforward to average (3.36) over space with the help of (3.38)
$`\mathrm{\Pi }(𝐪,\omega )`$ $`=`$ $`{\displaystyle \frac{3}{R_0^3}}{\displaystyle \underset{0}{\overset{R_0}{}}}𝑑RR^2\mathrm{\Pi }(R,𝐪,\omega )`$ (3.39)
$``$ $`\mathrm{\Pi }^0(𝐪,\omega )+i{\displaystyle \frac{3}{qR_0}}\left[1+{\displaystyle \frac{\omega }{2}}_\omega \right]{\displaystyle \frac{1}{3}}\mathrm{\Pi }^0(𝐪,\omega )`$
$`=`$ $`\mathrm{\Pi }^0+\mathrm{\Pi }^{\mathrm{surf}}.`$
Consequently the surface contribution to the polarization function reads finally with (3.37)
$`\mathrm{\Pi }^{\mathrm{surf}}(𝐪,\omega )={\displaystyle \frac{3}{2qR_0}}\mathrm{Im}\mathrm{\Pi }^0(𝐪,\omega )`$ (3.40)
which is real. With (3.39), (3.40) and (3.35) we obtain now for the structure function
$`S(q,\omega )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{Im}\mathrm{\Pi }^0}{(1V(\mathrm{Re}\mathrm{\Pi }^0+\mathrm{\Pi }^{\mathrm{surf}}))^2+(V\mathrm{Im}\mathrm{\Pi }^0)^2}}.`$ (3.41)
For small $`q`$ expansion we see that the pole of the structure function becomes renormalized similar to what is known from the Mie mode or surface plasmon mode
$`\omega ^2={\displaystyle \frac{\omega _0^2}{1V\mathrm{\Pi }^{\mathrm{surf}}}}.`$ (3.42)
After establishing the structure function including surface contribution we specify the model for actual calculations. We choose as mean field parameterization a Skyrme force (3.17) following Vautherin and Brink which leads to the isoscalar potential
$`V=U_nU_p={\displaystyle \frac{3t_0}{4}}+{\displaystyle \frac{3t_3}{8}}n_0`$ (3.43)
with $`t_0=983.4`$ MeV fm<sup>3</sup>, $`t_3=13106`$ MeV fm<sup>6</sup>, $`x_0=0.48`$ at nuclear saturation density $`n_0=0.16`$ fm<sup>-3</sup> and the incompressibility of $`K=318`$ MeV. Further we employ again the Steinwedel-Jensen model where the basic mode inside a sphere of radius $`R_0`$ is given by a wave vector
$`q_{sp}={\displaystyle \frac{\pi }{2R_0}}.`$ (3.44)
This would correspond to the first order isovector mode . Since this mode is spurious we have to consider the next higher harmonics which is
$`q_{isgdr}={\displaystyle \frac{\pi }{R_0}}.`$ (3.45)
The polarization function with this second order mode contains still contributions from the spurious mode such that we have to subtract this part
$`\mathrm{\Pi }_{\mathrm{ISGDR}}^0(\omega )=\mathrm{\Pi }^0(q_{isgdr},\omega )\mathrm{\Pi }^0(q_{sp},\omega ).`$ (3.46)
In figure 3.8 we have plotted the experimental structure function together with different theoretical estimates according to (3.46) and (3.41). The inclusion of surface corrections (dashed lines) shifts the structure function towards the experimental values. The inclusion of collisions (dot-dashed lines), which should be wrong for isoscalar dipole mode due to cancellation of backscattering leads to really worse results. This supports indirectly that the mode is of isoscalar dipole type and surface dominated.
### 3.7 Simplified model for nuclear matter situation
The resulting wave vectors have very low values compared with the Fermi wave vector in the Steinwedel Jensen model. This allows us to expand the Mermin polarization function (3.12) with respect to small $`qv_c/\omega `$ ratios and $`v_c`$ the sound velocity. We obtain
$`\mathrm{\Pi }_a^\mathrm{M}(\omega )`$ $`=`$ $`{\displaystyle \frac{\rho _a(\mu _a)}{m_a}}{\displaystyle \frac{q^2}{\omega (\omega +i/\tau _a)}},`$ (3.47)
$`\rho _a(\mu _a)`$ $`=`$ $`2\lambda _a^3f_{3/2}(z_a)\zeta _{corr}`$ (3.48)
where the thermal wave length is $`\lambda _a^2=2\pi \mathrm{}^3/(m_aT)`$, $`f_{3/2}`$ the standard Fermi integral and $`z_a=\mathrm{e}^{\mu _a/T}`$ the fugacity. The correction constant $`\zeta _{corr}=1.22`$ is introduced to fit the numerical solution of the dispersion relation (3.2) with the approximative expansion (3.47). With the help of this expansion the dispersion relation (3.11) takes the form
$`0`$ $`=`$ $`\left[\omega \left(\omega +{\displaystyle \frac{i}{\tau _n}}\right)c_{nn}^2q^2\right]\left[\omega \left(\omega +{\displaystyle \frac{i}{\tau _p}}\right)c_{pp}^2q^2\right]`$ (3.49)
$`\left[c_{np}^2+i{\displaystyle \frac{\stackrel{~}{c}_{np}^2}{(\omega +i/\tau _n)\tau _p}}\right]\left[c_{pn}^2+i{\displaystyle \frac{\stackrel{~}{c}_{pn}^2}{(\omega +i/\tau _p)\tau _n}}\right]q^4`$
with the partial sound velocities $`c`$ and $`\stackrel{~}{c}`$
$`c_{ab}^2=\alpha _{ab}{\displaystyle \frac{\rho _a(\mu _a)}{m_a}},\stackrel{~}{c}_{ab}^2={\displaystyle \frac{T}{m_a}}{\displaystyle \frac{\frac{f_{3/2}(z_a)}{f_{1/2}(z_a)}\frac{f_{3/2}(z_b)}{f_{1/2}(z_b)}}{\frac{\tau _{ab}}{\tau _{bb}}\frac{\tau _{aa}}{\tau _{ba}}}}.`$ (3.50)
The dispersion relation (3.11) with the dynamical relaxation times (3.16) is a polynomial of tenth (sixth) order corresponding to the inclusion of memory (in)dependent relaxation times. While most of these solutions will just lead to parasite solutions ($`\mathrm{Re}\omega 0`$) we will get two coupled modes, i.e. the isoscalar and isovector mode. Furthermore a third mode appears at extreme asymmetries and/or strong collisional coupling which we will describe in the next section.
### 3.8 New collective mode
Now we employ the potential (3.17) and assume different neutron and proton densities. In figure 3.9 we plot the isoscalar and isovector modes versus temperature for $`{}_{}{}^{48}Ca`$ with a small asymmetry $`\delta =0.2`$ as well as for $`{}_{}{}^{60}Ca`$ with an asymmetry $`\delta =0.33`$. The kinetic energy is linked to a temperature within the Fermi liquid model via Sommerfeld expansion
$`{\displaystyle \frac{}{A}}_{\mathrm{kin}}`$ $`=`$ $`{\displaystyle \frac{3}{5}}ϵ_f\left[{\displaystyle \frac{(1+\delta )^{5/3}+(1\delta )^{5/3}}{2}}+{\displaystyle \frac{5}{12}}\pi ^2\left({\displaystyle \frac{T}{T_f}}\right)^2{\displaystyle \frac{(1+\delta )^{1/3}+(1\delta )^{1/3}}{2}}\right].`$ (3.51)
This connection between temperature and excitation energy is only valid for a continuous Fermi liquid model. For small nuclei, the concept of temperature is questionable. Some improvement can be obtained by the definition of temperature via the logarithmic derivative of the density of states
$`T^1={\displaystyle \frac{1}{\rho }}{\displaystyle \frac{\rho }{E_{\mathrm{ex}}}}={\displaystyle \frac{5}{4}}E_{\mathrm{ex}}^1+\pi ({\displaystyle \frac{A}{4ϵ_fE_{\mathrm{ex}}}})^{1/2}`$ (3.52)
which provides $`E_{\mathrm{ex}}\frac{1}{4}(E/A)`$ in comparison with (3.51) for small temperatures and small nuclei. We use this temperature to demonstrate possible collective bulk features in an exploratory sense. Of course, the surface energy and shell effects cannot be neglected for realistic calculations of small nuclei.
With increasing temperature the isovector and isoscalar energies decrease and vanish at a certain temperature. At these temperatures the damping becomes twofold because the damping of the spurious mode with negative energy becomes different from the physical mode. We can consider this behavior of damping as a phase transition of isospin demixture. At the same time a very soft mode appears due to collisional coupling which is only present in asymmetric matter . This mode is more pronounced in the next figure 3.9 for $`{}_{}{}^{60}Ca`$ with an asymmetry of $`\delta =0.33`$.
We see that the isovector mode does not disappear but turns over into a flat decrease with increasing temperature. This behaviour is coupled with the pronounced soft mode. In comparison with the more symmetric $`{}_{}{}^{48}Ca`$ we see a different behavior of the damping where the isoscalar mode vanishes. Also the isovector mode appears unique and not two-fold. Now a clear transition of damping behavior for the isovector mode is recognizable which can be considered as a transition from zero to first sound damping.
Besides the standard isovector and isoscalar modes we observe a build up of a very soft mode with a centroid energy around $`1`$ MeV. This mode appears due to the collisional coupling $`\overline{c}_{ab}^2`$ of (3.50). When we turn off the relaxation times, i.e. the collision integral, this mode is vanishing as well as in symmetric nuclear matter, see discussion after (3.13). It shows that this mode appears due to collisional coupling of isovector and isoscalar modes. The corresponding damping of the crossed mode is continuously increasing with temperature.
One may argue whether this third mode can really appear in the system. A simple consideration may convince us about the possible existence of such mode. Let us assume a coupled set of two type of harmonic oscillators (neutrons and protons) interacting between the same sort of particles with strength $`k_n`$ and $`k_p`$, respectively and between different sorts with $`k_{np}`$. Let us choose for simplicity only two neutrons and two protons. Then we obtain the coupled system of harmonic oscillators with frequencies $`\omega _n^2=k_n/m`$, $`\omega _p^2=k_p/m`$ and $`\omega _{np}^2=k_{np}/m`$. The solution yields three basic modes in the system, i.e. $`\omega ^2=2(\omega _n^2+\omega _{np}^2),4\omega _{np}^2,2(\omega _p^2+\omega _{np}^2)`$. If we neglect the different coupling between neutrons and protons $`\omega _{np}`$ we only obtain two modes analogously to isovector and isoscalar ones. We see that the coupling between neutrons and protons can lead to the appearance of a third mode.
Let us now compare the found new mode with the experimental evidence. There are some hints for a soft mode in $`{}_{}{}^{11}Be`$ . The authors have observed a low lying structure at around $`6`$ MeV excitation energy with a damping of around $`1`$ MeV which has not been reproduced yet even within refined coupled channel calculations . A standard explanation would give as the origin a weakly-bound single particle neutron orbital. The observed broad structure at $`6`$ MeV might be explained alternatively as the presented new coupled mode. The centroid energy as well as damping width at least seem to suggest this interpretation.
## Chapter 4 Nonlinear effects beyond linear response
While most of the theoretical treatments of oscillations rely on the linear response method or RPA methods, large amplitude oscillations require methods beyond. In particular the question of the appearance of chaos has recently been investigated . The hypothesis was established that the octupole mode is over-damped due to negative curved surface and consequently additional chaotic damping . Here we want to discuss in which conditions one might observe octupole modes at least in Vlasov - simulation of giant resonances. We will consider different initial conditions of isoscalar giant resonances and will demonstrate that the appearance of octupole modes is dependent on the initial configuration which in turn demonstrates the nonlinear behavior beyond linear response.
As a first initial condition we use the ground state distribution of coordinates while the momentum distribution is deformed anisotropically. We have modified the momenta in a way which corresponds to a giant octupole mode. The local densities and currents remain the same as in ground state. Figure 4.1 shows that at start time $`t=0`$ there is a pure giant octupole, which is damped out and a quadrupole resonance develops instead. The monopole and dipole amplitudes which should remain constant document the stability of simulation. In agreement with the already mentioned hypothesis, the octupole mode is over-damped. The figure shows the nonlinear behavior of mode coupling. Within the linear response the damping rate is expected to be independent of initial conditions. We choose now other initial conditions to show that the result is very much dependent on initial conditions. Therefore we split the nucleus in two parts of mass ratio 3:7 in accordance with the symmetry of the octupole oscillation and accelerate both pieces towards each other. Experimentally this might be realized as a central collision of two nuclei with corresponding masses.
In figure 4.2 a clear quadrupole resonance appears and also a smaller octupole resonance can be seen. Both are damped out. Consequently there is no evidence for an over-damped octupole mode in this case.
In order to understand the different initializations, we split the kinetic energy into a thermal part and a collective part according to $`𝐩^2=(𝐩𝐩)^2+𝐩^2`$. We analyze the time development of the total collective energy in the system
$`E_{coll}(t)`$ $`=`$ $`{\displaystyle \frac{h^2}{2m}}{\displaystyle 𝑑𝐫\varrho (𝐫)𝐩^2(𝐫,t)}`$ (4.1)
with the mean current
$`𝐩(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{1}{\varrho (𝐫)}}{\displaystyle \frac{d𝐩}{(2\pi )^3}𝐩f(𝐩,𝐫,t)}`$ (4.2)
and the density
$`\varrho (𝐫,t)`$ $`=`$ $`{\displaystyle \frac{d𝐩}{(2\pi )^3}f(𝐩,𝐫,t)}.`$ (4.3)
In figure 4.3 the development of collective energy can be seen. There is a background of about 50 MeV due to fixed correlations caused by finite width of pseudo-particles as one can see from the following estimation. Using (3.7) and (3.39) in (4.1) one obtains
$`E_{coll}`$ $``$ $`\frac{1}{\varrho _0}{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \underset{j=1}{\overset{AN}{}}}\frac{1}{N^2}p_ip_j{\displaystyle 𝑑rf_S(rr_i(t),\sigma _r)f_s(rr_j(t),\sigma _r)}`$ (4.4)
$``$ $`\frac{1}{\varrho _0}{\displaystyle \underset{i=1}{\overset{AN}{}}}{\displaystyle \underset{j=1}{\overset{AN}{}}}\frac{1}{N^2}p_ip_jf_S(r_i(t)r_j(t),\sqrt{2}\sigma _r)`$
$``$ $`\frac{1}{\varrho _0N}\frac{1}{(\sqrt{2\pi }\sqrt{2}\sigma _r)^3}{\displaystyle \underset{i=1}{\overset{AN}{}}}\frac{1}{N}p_i.`$
For simulation parameters of $`{}_{}{}^{208}Pb`$, $`\rho _0=0.162\mathrm{fm}^3`$, $`N=75`$, $`\sigma _r=0.53\mathrm{fm}`$, we obtain a basic collective energy of $`54\mathrm{M}\mathrm{e}\mathrm{V}`$. Using more test-particles would diminish this level.
The solid line corresponding to figure 4.1 shows no initial collective energy. The exclusive initial excitation in momentum space without correlation in the spatial domain leads to zero initial collective energy. These correlations are forming during time evolution. Of course, there is no center of mass motion, otherwise we would see just the mean streaming velocity.
This situation is changed if we use the second preparation with simple momentum–space–correlations. The long dashed line in figure 4.3 shows initial collective correlations corresponding to figure 4.2. Since we can deposit enough collective energy in this case, we observe a clear octupole motion.
In order to compare with we calculate the adiabaticity index $`\eta `$ defined in as the ratio of maximum radial surface velocity to the maximum particle speed. A smaller ratio denotes a more adiabatic shape changes in relation to the particle speed. In analogy we define such index as a ratio
$`\eta `$ $`=`$ $`{\displaystyle \frac{\frac{}{t}\sqrt{<r^2>(\vartheta )}}{v_F}}`$ (4.5)
of the root mean square radius speed in forward direction (opening angle $`\vartheta `$) and the Fermi velocity. With opening angle 0.4 rad we obtain a maximum $`\eta =0.12`$ for figure 4.1 and $`\eta =0.30`$ for figure 4.2. This shows that we are essentially still in the adiabatic regime described in .
The nonlinear behavior described so far already documents that we are in a regime of large amplitude oscillations where linear response fails. The corresponding radius elongation in coordinate space varies about 10 %.
## Chapter 5 Summary
We have investigated the giant resonance oscillations by solving the collision-less Vlasov equation as well as by linear response theory.
With the help of numerical solution of the Vlasov equation we could reproduce the mass dependence of energy with one fixed parameter of test-particle width for both giant monopole and giant dipole modes. A multipole analysis was performed which has allowed us to characterize the corresponding excitation. It was shown, that the asymmetry influences the collective behavior. With increasing asymmetry the energy decreases while the damping increases. The damping due to finite systems amount to $`2`$MeV almost constant for all mass numbers which underestimates the experimental values considerably. This motivates to search for additional damping mechanism which was found to be due to the collisional correlations.
The linear response can be used in a simplified liquid drop model to describe the giant resonances in asymmetric matter. We find that the collisional contribution as well as surface contributions are both important to reproduce the experimental damping of giant dipole resonances. While for ground state resonances it is sufficient to add the damping of finite size effects from simulation with the collisional contribution from linear response theory, for the temperature behavior we have needed also the deformation of the surface. The combined model between surface and collisional contribution, weighted properly due to their collision frequencies, is able to reproduce the experimental damping curve with temperature as well as the structure function. A higher order mode like the recently measured ISGDR mode has been described within this simple linear response model.
We have observed that due to correlational coupling there can exist a new mode which appears besides isovector and isoscalar modes in asymmetric nuclear matter due to collisions. We suggest that this mode may be possible to observe as a soft collective excitation in asymmetric systems. The transition from zero sound damping to first sound damping behavior should become possible to observe for isovector modes since they do not vanish at this transition temperature like in symmetric matter. At a certain critical temperature the collective isoscalar mode vanishes and the damping becomes two-fold. This can be considered as a phase transition of demixture.
On the example of octupole resonances in finite nuclei the dependency on the initial configuration is demonstrated. The appearance of an octupole mode was shown by correlating the spatial and momentum initial excitation. It is possible to excite an octupole mode with sufficient collective energy deposited initially. We suggest that isoscalar giant octupole resonances should be possible to observe in nuclear collisions of mass ratio about 3:7 corresponding to octupole symmetry. This effect illustrates a mode beyond linear response.
### 5.1 Acknowledgment
The authors are obliged to M. Vogt who has contributed the calculation of the Lyapunov exponent. M. DiToro is thanked for numerous discussions and the LNS Catania where part of this work was done for hospitality. Valuable hints and comments by J. D. Frankland (GANIL) are also gratefully acknowledged.
## Appendix A Derivation of non-Markovian relaxation time
Here we shortly sketch the derivation of the damping of a collective mode within a Fermi gas and a Fermi liquid model. We will show that we get the latter one from the Fermi gas model with an additional contribution from the quasi-particles.
We start with the Fermi gas where the dispersion relation between momentum and energy is given by $`ϵ=p^2/2m`$ and will show later what has to be changed for a Fermi liquid where $`ϵ`$ is a solution of the quasiparticle dispersion relation. We will see that the contributions from the quasi-particles alone leads to the Landau formula of zero sound damping
$`\gamma \left[1+\left({\displaystyle \frac{\mathrm{\Omega }}{2\pi T}}\right)^2\right].`$ (A.1)
Our considerations use conveniently the Levinson equation for the reduced density matrix $`f`$ which is valid at short time processes compared to inverse Fermi energy $`\mathrm{}/ϵ_f`$ and which collisional side has the form:
$`I_1(t)={\displaystyle \frac{2g}{\mathrm{}^2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau {\displaystyle \frac{d𝐩_2d𝐩_3d𝐩_4}{(2\pi \mathrm{})^6}|T|^2\mathrm{cos}\left(\underset{t}{\overset{t\tau }{}}\mathrm{\Delta }ϵ(\tau )𝑑\tau /\mathrm{}\right)\delta (\mathrm{\Delta }𝐩)\left(\overline{f}_1\overline{f}_2f_3f_4f_1f_2\overline{f}_3\overline{f}_4\right)_{t\tau }}`$ (A.2)
where $`\overline{f}=1f`$, $`\mathrm{\Delta }p=p_1+p_2p_3p_4`$ etc., g is the spin-isospin degeneracy and the transition probability is given by the scattering $`T`$-matrix. In case that the quasiparticle energies $`ϵ(t)`$ become time independent like in the Fermi gas model, the integral in the $`cos`$ function reduces to the familiar expression $`\mathrm{\Delta }ϵ\tau `$. We linearize this collision integral with respect to an external disturbance according to
$`f=n+\delta f`$ (A.3)
where n is the equilibrium distribution. Clearly two contributions have to be distinguished, the one from the quasiparticle energy and the one from occupation factors . First we concentrate on the Fermi gas model where we have only the contribution of the occupation factors and will later add the contribution of the quasiparticle energies for Fermi liquid model. We obtain after Fourier transform of the time
$`\delta I_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{2}}\left[\delta _+(\mathrm{\Delta }ϵ+\mathrm{\Omega })+\delta _{}(\mathrm{\Delta }ϵ\mathrm{\Omega })\right]\left(\delta F_1+\delta F_2\delta F_3\delta F_4\right)(\mathrm{\Omega }).`$ (A.4)
Here we use the abbreviation
$`\mathrm{}`$ $`=`$ $`{\displaystyle \frac{2g}{\mathrm{}^2}}{\displaystyle \frac{d𝐩_2d𝐩_3d𝐩_4}{(2\pi \mathrm{})^6}|T|^2\delta (\mathrm{\Delta }p)\mathrm{}}={\displaystyle \frac{m^3g}{\mathrm{}^2(2\pi \mathrm{})^6}}{\displaystyle 𝑑ϵ_2𝑑ϵ_3𝑑ϵ_4\frac{d\varphi \mathrm{sin}\theta d\theta d\varphi _2}{\mathrm{cos}(\theta /2)}|T|^2\mathrm{}}`$
where the last line appears from standard integration techniques at low temperatures. Further abbreviations are
$`\delta _\pm (x)`$ $`=`$ $`\pi \delta (x)\pm i{\displaystyle \frac{𝒫}{x}}\pi \delta (x)`$
$`\delta F_1`$ $`=`$ $`\delta f_1(\overline{n}_2n_3n_4+n_2\overline{n}_3\overline{n}_4).`$ (A.6)
The approximation used in the first line consists in the neglect of the off-shell contribution from memory effects. This is consistent with the used integration technique (LABEL:int). This terms would lead to divergences which has to be cut off .
Neglecting the backscattering terms $`\delta F_{2/3/4}`$ we obtain from (A.4) a relaxation time approximations with the relaxation time
$`{\displaystyle \frac{1}{\tau (ϵ_1)}}={\displaystyle \frac{3}{4\pi ^2\tau _0}}{\displaystyle \underset{\lambda }{\overset{\mathrm{}}{}}}𝑑x_2𝑑x_3𝑑x_4\left[\delta (\mathrm{\Delta }x+\omega )+\delta (\mathrm{\Delta }x\omega )\right](\overline{n}_2n_3n_4+n_2\overline{n}_3\overline{n}_4)`$ (A.7)
with $`\omega =\mathrm{\Omega }/T`$, $`x=(ϵ\mu )/T`$, $`\lambda =\mu /T`$ and the time
$`{\displaystyle \frac{1}{\tau _0}}={\displaystyle \frac{2gmT^2}{3\mathrm{}^3}}\sigma .`$ (A.8)
Here we have used the definition of cross section $`|T|^2=(4\pi \mathrm{}^2/m)^2d\sigma /d\mathrm{\Omega }`$ and have assumed a constant cross section $`\sigma `$. To calculate (A.7) one needs the standard integrals for large ratios of chemical potentials $`\mu `$ to temperature $`\lambda =\mu /T`$
$`{\displaystyle \underset{\lambda }{\overset{\mathrm{}}{}}}𝑑x_2𝑑x_3𝑑x_4n_2\overline{n}_3\overline{n}_4\delta (\mathrm{\Delta }x\pm \omega )={\displaystyle \frac{1}{2}}\overline{n}_1(x_1\pm \omega )\left[\pi ^2+(x_1\pm \omega )^2\right]`$ (A.9)
to obtain
$`{\displaystyle \frac{1}{\tau (ϵ_1)}}={\displaystyle \frac{3}{8\pi ^2\tau _0}}\left[2\pi ^2+(x_1+\omega )^2+(x_1\omega )^2\right].`$ (A.10)
Further we employ a thermal averaging in order to obtain the mean relaxation time finally<sup>1</sup><sup>1</sup>1One has to use the identities valid up to $`o(\mathrm{exp}[\lambda ])`$
$$\underset{\lambda }{\overset{\mathrm{}}{}}𝑑xn\overline{n}=1;\underset{\lambda }{\overset{\mathrm{}}{}}𝑑xx^2n\overline{n}=\frac{\pi ^2}{3}.$$
$`{\displaystyle \frac{1}{\tau _{\mathrm{gas}}}}={\displaystyle \underset{\lambda }{\overset{\mathrm{}}{}}}𝑑x_1n_1\overline{n}_1{\displaystyle \frac{1}{\tau (ϵ_1)}}={\displaystyle \frac{1}{\tau _0}}\left[1+3\left({\displaystyle \frac{\omega }{2\pi }}\right)^2\right].`$ (A.11)
If we do not use the thermal averaging but take (A.10) at the Fermi energy $`ϵ_1=ϵ_f`$ we will obtain
$`{\displaystyle \frac{1}{\tau (ϵ_f)}}={\displaystyle \frac{1}{\tau _0}}\left[{\displaystyle \frac{3}{4}}+3\left({\displaystyle \frac{\omega }{2\pi }}\right)^2\right].`$ (A.12)
We see that both results disagree with the Landau result of quasiparticle damping (A.1) by factors of 3 at different places . We have point out that the result at fixed Fermi energy will lead to unphysical results for the Fermi liquid case. Therefore we consider the thermal averaged result as the physical one.
We now turn to the Fermi liquid model and replace the free dispersion $`ϵ=p^2/2m`$ by the quasiparticle energy $`ϵ_p`$. Than the variation of the collision integral gives an additional term which comes from the time dependence of the quasiparticle energy on the $`\mathrm{cos}`$-term of (A.2). We have instead of the sum of two complex conjugate exponentials in (A.2) an additional contribution from the linearization of the exponential
$`\delta \mathrm{exp}\left(i\mathrm{\Delta }{\displaystyle \underset{t}{\overset{t\tau }{}}}𝑑\overline{t}ϵ(\overline{t})\right)`$ $`=`$ $`\mathrm{e}^{i\mathrm{\Delta }ϵ\tau }\left(1i\mathrm{\Delta }{\displaystyle \underset{t}{\overset{t\tau }{}}}𝑑\overline{t}[ϵ(\overline{t})ϵ]\right)\mathrm{e}^{i\mathrm{\Delta }ϵ\tau }`$ (A.13)
$`=`$ $`{\displaystyle \frac{iT}{\mathrm{}}}\mathrm{e}^{i\mathrm{\Delta }ϵ\tau /\mathrm{}}\mathrm{\Delta }{\displaystyle \underset{t}{\overset{t\tau }{}}}𝑑\overline{t}{\displaystyle \frac{\delta f(\overline{t})}{n\overline{n}}}.`$
In the last line we have replaced the variation in the quasiparticle energy $`ϵ(t)ϵ`$ by the variation in the distribution function $`\delta f`$ due to the identity
$`\delta f(t)`$ $`=`$ $`f(t)n(ϵ)=f(t)n(ϵ(t))[n(ϵ)n(ϵ(t))]`$ (A.14)
$``$ $`n^{}(ϵ(t)ϵ)={\displaystyle \frac{n\overline{n}}{T}}\left[ϵ(t)ϵ\right]`$
where we assumed within the quasiparticle picture that $`f(t)=n(ϵ(t))`$. This leads now to an additional part in the relaxation time which we write analogously to (A.4)
$`{\displaystyle \frac{1}{\tau _c(ϵ_1)}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{2}}{\displaystyle \frac{\delta _+(\mathrm{\Delta }ϵ\mathrm{\Omega })\delta _+(\mathrm{\Delta }ϵ+\mathrm{\Omega })}{\mathrm{\Omega }}}{\displaystyle \frac{\overline{n}_1\overline{n}_2n_3n_4n_1n_2\overline{n}_3\overline{n}_4}{n_1\overline{n}_1}}.`$ (A.15)
Using again (A.9) we obtain
$`{\displaystyle \frac{1}{\tau _c(ϵ_1)}}={\displaystyle \frac{3}{4\pi ^2\tau _0}}\{{\displaystyle \frac{\overline{n}(x_1+\omega )\left[\pi ^2+(x_1+\omega )^2\right]}{2\omega (\mathrm{e}^\omega 1)}}+[\omega \omega ]\}`$ (A.16)
and get after thermal averaging (A.11)<sup>2</sup><sup>2</sup>2Here one uses \[$`o(\mathrm{exp}[\lambda ])`$\]
$$\underset{\lambda }{\overset{\mathrm{}}{}}𝑑xn=\lambda ;\underset{\lambda }{\overset{\mathrm{}}{}}𝑑xxn=\frac{1}{2}\lambda ^2+\frac{1}{6}\pi ^2;\underset{\lambda }{\overset{\mathrm{}}{}}𝑑xx^2n=\frac{1}{3}\lambda ^3$$
.
$`{\displaystyle \frac{1}{\tau _c}}={\displaystyle \frac{1}{\tau _0}}\left[1+\left({\displaystyle \frac{\omega }{2\pi }}\right)^2\right].`$ (A.17)
Taking instead of thermal averaging the value at Fermi energy $`(ϵ_1=ϵ_f)`$ in (A.16) we find
$`{\displaystyle \frac{1}{\tau _c(ϵ_f)}}={\displaystyle \frac{3(\pi ^2+\omega ^2)}{2\pi ^2\tau _0\omega }}{\displaystyle \frac{\mathrm{e}^\omega 1}{\mathrm{e}^\omega +1}}.`$ (A.18)
Here we like to point out that the Landau result (A.1) appears in (A.17) (see also in Ref. ).
Adding now (A.11) and (A.17) we obtain a final relaxation time for the Fermi liquid model
$`{\displaystyle \frac{1}{\tau _{\mathrm{liq}}}}={\displaystyle \frac{2}{\tau _0}}\left[1+2\left({\displaystyle \frac{\omega }{2\pi }}\right)^2\right].`$ (A.19)
which is the main result in this paper. It contains the typical Landau result of zero sound (A.1) except the factor 2 in front of the frequencies. Comparing (A.19) with the Fermi gas model (A.11) we see that in the limit of vanishing temperature the Fermi liquid value is lower with $`2\mathrm{\Omega }^2`$ compared to the Fermi gas $`3\mathrm{\Omega }^2`$. Further for vanishing frequencies (neglect of memory effects) the Fermi liquid model leads to twice the relaxation rate than the Fermi gas model. The coefficient of temperature increase is than twice larger for the Fermi liquid than for the Fermi gas. If we consider the relaxation times at Fermi energy (no thermal averaging) (A.12) and (A.18) we find the same results as above in the limit of vanishing frequencies. For vanishing temperature only the Fermi gas (A.12) coincides with the result of (A.11) $`3\mathrm{\Omega }^2`$. Expression (A.18) goes to zero for $`T=0`$ and underlines the necessity to thermal average the value. |
warning/0001/hep-th0001174.html | ar5iv | text | # 1 Introduction
## 1 Introduction
It is well known that classical symmetries, expressed in functional form by a set of Ward identities (WI) satisfied by the classical action $`\mathrm{\Gamma }^{\left(0\right)}`$, may be broken at the quantum level . The possible quantum breaking in perturbative QFT of the WI can have two origins. Either it is a truly physical obstruction to the restoration of the WI (and one is then faced with the problem of classifying all possible anomalous terms for the model under investigation) or it is an unwanted effect of the renormalization procedure, needed to handle UV divergences in the full quantum action $`\mathrm{\Gamma }`$. The latter breaking can always be recovered by the choice of a more suited renormalization scheme, or of different renormalization conditions. Since changing the renormalization procedure does not alter the physical content of the theory, this is a spurious breaking of the WI. <sup>3</sup><sup>3</sup>3Though non-physical, these breaking terms nevertheless require a careful treatment. For example, in the Standard Model no regularization scheme is known to preserve all the symmetries of the theory, because of the presence of the $`\gamma _5`$ matrix and of the completely antisymmetric tensor $`ϵ_{\mu \nu \rho \sigma }`$. Thus some procedures are needed to recover the Ward identities broken by the intermediate renormalization: suitably defined finite counter-terms must be added to the regularized quantum action in order to recover the relevant Ward identities.
In gauge theories, the BRST transformations turn out to be a very powerful tool for proving the unitarity of the physical $`S`$ matrix . Thus the restoration of BRST symmetry is an essential step in the perturbative construction of gauge theories. The crucial requirement in the proof of physical unitarity is nilpotency: nilpotency of the BRST transformations is a sufficient condition for a theory to be unitary, provided that the associated WI (known as Slavnov-Taylor identities - STI from now on) hold at the quantum level. The task is then to restore the STI, when this is physically possible (absence of anomalies), and to completely classify the form of the breaking terms in the anomalous case.
The nilpotency of the BRST transformations also allows for an effective cohomological analysis of the ST breaking terms. This is most easily seen in the framework of the field-antifield formalism, an extension of the original BRST formulation .
We consider a general gauge theory with fields $`\varphi _i`$ and ghosts $`c_k`$, introduced by the covariant quantization of the model. The fields $`\varphi _i`$ and the ghosts $`c_k`$ are collectively denoted by $`\mathrm{\Phi }^A`$, $`A=1,\mathrm{},N`$. In the field-antifield formalism, for each field $`\mathrm{\Phi }^A`$ one introduces an antifield $`\mathrm{\Phi }_A^{}`$. The space $``$ of the functionals of $`\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}`$ is endowed with an odd symplectic structure $`(,)`$, the antibracket:
$`(X,Y)={\displaystyle \underset{A=1}{\overset{N}{}}}{\displaystyle d^4x\left(\frac{\delta _rX}{\delta \mathrm{\Phi }^A\left(x\right)}\frac{\delta _lY}{\delta \mathrm{\Phi }_A^{}\left(x\right)}\frac{\delta _rX}{\delta \mathrm{\Phi }_A^{}\left(x\right)}\frac{\delta _lY}{\delta \mathrm{\Phi }^A\left(x\right)}\right)}.`$ (1)
The subscripts $`r`$ and $`l`$ denote right and left differentiation respectively. The classical action $`\mathrm{\Gamma }^{\left(0\right)}[\mathrm{\Phi },\mathrm{\Phi }^{}]`$ is assumed to satisfy the classical master equation
$`(\mathrm{\Gamma }^{\left(0\right)},\mathrm{\Gamma }^{\left(0\right)})=0,`$ (2)
under the condition that $`\mathrm{\Gamma }^{\left(0\right)}[\mathrm{\Phi },\mathrm{\Phi }^{}]|_{\mathrm{\Phi }^{}=0}`$ coincides with the classical gauge-fixed BRST invariant action. The quantization of the theory produces an effective action $`\mathrm{\Gamma }`$
$`\mathrm{\Gamma }=\mathrm{\Gamma }^{\left(0\right)}+{\displaystyle \underset{n1}{}}\mathrm{}^n\mathrm{\Gamma }^{\left(n\right)},`$ (3)
which satisfies the quantum extension of the classical master equation (2) :
$`{\displaystyle \frac{1}{2}}(\mathrm{\Gamma },\mathrm{\Gamma })=\mathrm{}\left(𝒜\mathrm{\Gamma }\right),`$ (4)
where the insertion $`\left(𝒜\mathrm{\Gamma }\right)`$ represents the possible anomalous terms due to the quantum corrections. Using the graded Jacobi identity for the antibracket $`(,)`$
$`((X,Y),Z)+\left(1\right)^{\left(ϵ_X+1\right)\left(ϵ_Y+ϵ_Z\right)}((Y,Z),X)+\left(1\right)^{\left(ϵ_Z+1\right)\left(ϵ_X+ϵ_Y\right)}((Z,X),Y)=0`$ (5)
(where $`ϵ_X=0`$ if $`X`$ obeys Bose statistics and $`ϵ_X=1`$ if $`X`$ obeys Fermi statistics), it is easy to deduce the following identity for $`\mathrm{\Gamma }`$:
$`(\mathrm{\Gamma },(\mathrm{\Gamma },\mathrm{\Gamma }))=0`$ (6)
or, taking into account eq.(4),
$`\mathrm{}(\mathrm{\Gamma },\left(𝒜\mathrm{\Gamma }\right))=0.`$ (7)
At the first order in perturbation theory the previous condition reduces to
$`(\mathrm{\Gamma }^{\left(0\right)},𝒜_1)=0,`$ (8)
where $`𝒜\mathrm{\Gamma }=_{j1}\mathrm{}^{j1}𝒜_j`$. Eq.(8) is the Wess-Zumino consistency condition written in the field-antifield formalism.
Since $`\mathrm{\Gamma }^{\left(0\right)}`$ satisfies the classical master equation (2), the operator $`(\mathrm{\Gamma }^{\left(0\right)},)`$ is nilpotent and eq.(8) gives rise to a cohomological problem. The most general solution of eq.(8) can be cast in the form
$`𝒜_1={\displaystyle \underset{k1}{}}\lambda _k𝒞_k+(\mathrm{\Gamma }^{\left(0\right)},𝒞_0),`$ (9)
where the sum is over the representatives $`𝒞_k`$ of the independent non-trivial cohomology classes of the operator $`(\mathrm{\Gamma }^{\left(0\right)},)`$ and $`(\mathrm{\Gamma }^{\left(0\right)},𝒞_0)`$ is an arbitrary element of the trivial cohomology class with FP-charge $`1`$. The purely algebraic analysis which leads to eq.(9) puts no restrictions on the form of $`𝒞_k`$ and $`𝒞_0`$, apart from saying that they are functionals of $`\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}`$. In particular, $`𝒞_k`$ and $`𝒞_0`$ might very well be non-local functionals of $`\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}`$ and they might contain arbitrarily high powers of $`\mathrm{\Phi }^A,\mathrm{\Phi }_A^{}`$, when expanded on a basis of the space $``$.
However, if the theory is power-counting renormalizable and the quantization is performed by means of a renormalization procedure which satisfies the Quantum Action Principle (QAP) , several strong restrictions are imposed on $`𝒞_k`$ and $`𝒞_0`$: they must be local functionals of $`\mathrm{\Phi }^A`$, $`\mathrm{\Phi }_A^{}`$ and have dimensions that cannot exceed a finite upper limit, predicted by the QAP.
Taking into account these constraints on locality and power counting, one sees from eq.(4) that $`(\mathrm{\Gamma }^{\left(0\right)},𝒞_0)`$ can always be reabsorbed by adding finite counter-terms to $`\mathrm{\Gamma }^{\left(1\right)}`$. Thus $`(\mathrm{\Gamma }^{\left(0\right)},𝒞_0)`$ is a spurious contribution to the anomaly.
If some of the coefficients $`\lambda _k`$ actually turn out to be non-zero (once their calculation has been performed in the intermediate renormalization scheme), the theory is truly anomalous: no matter how one changes the finite part of $`\mathrm{\Gamma }^{\left(1\right)}`$, these terms cannot be reabsorbed. Moreover, it is believed that in every admissible renormalization scheme (compatible with Poincaré invariance and all other exact symmetries of the theory <sup>4</sup><sup>4</sup>4For the sake of definiteness we are supposing that the ST invariance is the only broken invariance of the model.) the calculation of $`\lambda ^{\left(n\right)}`$ will yield a non-zero result.
Some attempts have been made to push forward this kind of analysis of the anomalous terms to higher orders in perturbation theory . The strategy is to find out a suitable higher-order generalization of the Wess-Zumino consistency condition in eq.(8), relying on the consistency condition for the full quantum action $`\mathrm{\Gamma }`$ in eq.(6) and the nilpotency of the operator $`(\mathrm{\Gamma }^{\left(0\right)},)`$.
In this program, the properties of the renormalization procedure and the choice of the renormalization conditions turn out to be as crucial as the algebraic features of the cohomological analysis, dictated by the nilpotency of $`(\mathrm{\Gamma }^{\left(0\right)},)`$.
In this paper we discuss the second order generalization of the Wess-Zumino consistency condition, in the simple framework of the Abelian Higgs-Kibble model . However, the conclusions one can draw from this example are general enough to be of interest for a wide class of gauge theories.
Although the Abelian HK model exhibits spurious anomalous terms only (see for a detailed analysis), we prove that, if the cohomologically trivial contributions to $`𝒜_1`$ are not recovered by suitably chosen finite counter-terms in $`\mathrm{\Gamma }^{\left(1\right)}`$, at the next order the equation for the anomaly is no more the Wess-Zumino consistency condition
$`(\mathrm{\Gamma }^{\left(0\right)},𝒜_2)=0.`$ (10)
Moreover, we show that in this case $`𝒜_2`$ must be non-local, in sharp contrast with the locality of the solutions of eq.(10).
In our discussion we will relax the assumption of nilpotency of the BRST transformations, by adding to the classical action the following mass term
$`{\displaystyle d^4x\left(\frac{M^2}{2}A_\mu ^2+M^2\overline{c}c\frac{M^2}{2\alpha }\left(\varphi _1^2+\varphi _2^2\right)\right)}.`$ (11)
Even though the price of this generalization is the loss of unitarity, in the massive framework it is simpler to appreciate the interplay between algebraic properties and the behavior of the quantum theory under the renormalization procedure. This also allows to discuss the conditions under which strict nilpotency of $`(\mathrm{\Gamma }^{\left(0\right)},)`$ is actually needed to carry out the construction of $`𝒜\mathrm{\Gamma }`$, to higher orders in perturbation theory.
## 2 Consistency conditions in the non-nilpotent case
In the Abelian Higgs-Kibble model the fields $`A_\mu `$ and $`\overline{c}`$ have linear BRST transformations. Thus one can actually avoid to introduce their antifields. Moreover, we work in the on-shell formalism, i.e. we have eliminated the auxiliary Nakanishi-Lautrup field $`B`$ associated with the BRST variation of $`\overline{c}`$ in the off-shell formalism. This in turn allows for a simplification of the Feynman graphs involved in our analysis.
From now on we use the reduced antibracket
$`(X,Y)={\displaystyle d^4x\left[\frac{\delta X}{\delta J_1}\frac{\delta Y}{\delta \varphi _1}+\frac{\delta X}{\delta J_2}\frac{\delta Y}{\delta \varphi _2}\frac{\delta X}{\delta \psi }\frac{\delta Y}{\delta \overline{\eta }}+\frac{\delta X}{\delta \overline{\psi }}\frac{\delta Y}{\delta \eta }\right]}.`$ (12)
All functional derivatives are assumed to act from the left. $`J_1,J_2,\eta ,\overline{\eta }`$ are the antifields of $`\varphi _1,\varphi _2,\overline{\psi },\psi `$ respectively.
The ST identities for the Abelian Higgs-Kibble model then read
$`S\left(\mathrm{\Gamma }\right)=0,`$ (13)
where the ST operator is
$`S\left(\mathrm{\Gamma }\right)={\displaystyle d^4x\left[^\mu c\frac{\delta \mathrm{\Gamma }}{\delta A^\mu }+\left(A+\frac{ev}{\alpha }\varphi _2\right)\frac{\delta \mathrm{\Gamma }}{\delta \overline{c}}\right]}+(\mathrm{\Gamma },\mathrm{\Gamma }).`$ (14)
The complete antibracket $`\frac{1}{2}(\mathrm{\Gamma },\mathrm{\Gamma })`$ in eq.(4) becomes the ST operator in eq.(14).
In the on-shell formalism the ghost equation is:
$`𝒢\mathrm{\Gamma }=\alpha \mathrm{}c+M^2c,`$ (15)
where the ghost operator $`𝒢`$ is defined by:
$`𝒢={\displaystyle \frac{\delta ()}{\delta \overline{c}}}ev{\displaystyle \frac{\delta ()}{\delta J_2}}.`$ (16)
### 2.1 The massless case
In the massless (nilpotent) case, eq. (6) is translated into
$`S_\mathrm{\Gamma }\left(S\left(\mathrm{\Gamma }\right)\right)=0.`$ (17)
$`S_\mathrm{\Gamma }`$ denotes the linearization of the ST operator (14):
$`S_\mathrm{\Gamma }()={\displaystyle d^4x\left[^\mu c\frac{\delta ()}{\delta A^\mu }+\left(A+\frac{ev}{\alpha }\varphi _2\right)\frac{\delta ()}{\delta \overline{c}}\right]}+(\mathrm{\Gamma },)+(,\mathrm{\Gamma }).`$ (18)
The identity (17) is valid for any $`\mathrm{\Gamma }`$ (even for a $`\mathrm{\Gamma }`$ which does not satisfy the STI $`S\left(\mathrm{\Gamma }\right)=0`$).
We denote by $`S_0S_{\mathrm{\Gamma }^{\left(0\right)}}`$ the zero-th order ST linearization. Notice that
$`\{𝒢,S_0\}=0.`$ (19)
We perform a formal expansion for $`S\left(\mathrm{\Gamma }\right)`$ in powers of $`\mathrm{}`$:
$`S\left(\mathrm{\Gamma }\right)={\displaystyle \underset{n0}{}}\mathrm{}^nS\left(\mathrm{\Gamma }\right)^{\left(n\right)}.`$ (20)
At the first order in perturbation theory eq.(17) becomes
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)=0,`$ (21)
which parallels eq.(8) and gives rise to the cohomological analysis of $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$. Thanks to the nilpotency of $`S_0`$ (guaranteed by the invariance of the classical action $`S\left(\mathrm{\Gamma }^{\left(0\right)}\right)=0)`$, it is possible to find the most general form of $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ compatible with condition (21). $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ can be cast in the form
$`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}=Y^{\left(1\right)}+S_0\left(𝒞_0\right),`$ (22)
where $`Y^{\left(1\right)}`$ is characterized as the most general local functional belonging to the kernel of $`S_0`$ and to the orthogonal complement of $`\mathrm{Im}S_0`$. Written on a basis $`\left\{𝒞_k\right\}_{k1}`$ of $`\mathrm{ker}S_0\left(\mathrm{Im}S_0\right)^{}`$, one gets
$`Y^{\left(1\right)}={\displaystyle \underset{k1}{}}\lambda _k𝒞_k,`$ (23)
for some coefficients $`\lambda _k`$. This expansion separates truly anomalous terms ($`Y^{\left(1\right)}`$) from spurious ones ($`S_0\left(𝒞_0\right)`$). The latter can be canceled by a suitable redefinition of the first-order counter-terms entering in the construction of $`\mathrm{\Gamma }^{\left(1\right)}`$.
### 2.2 The massive case
In the massive HK model eq. (17) is modified as follows
$`S_\mathrm{\Gamma }S\left(\mathrm{\Gamma }\right)={\displaystyle d^4x\left(\mathrm{}c+\frac{ev}{\alpha }\frac{\delta \mathrm{\Gamma }}{\delta J_2}\right)\frac{\delta \mathrm{\Gamma }}{\delta \overline{c}}}.`$ (24)
This identity is valid for any functional $`\mathrm{\Gamma }`$, without restrictions as $`S\left(\mathrm{\Gamma }\right)=0`$ or the ghost equation (15).
Taking into account eq.(15) we get
$`S_\mathrm{\Gamma }S\left(\mathrm{\Gamma }\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }}{\delta \overline{c}}}.`$ (25)
The linearized ST operator $`S_0`$ is no more nilpotent; for any functional $`F`$ satisfying the ghost equation (15) we have now
$`S_0^2\left(F\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta F}{\delta \overline{c}}}.`$ (26)
At the classical level the STI are satisfied:
$`S\left(\mathrm{\Gamma }^{\left(0\right)}\right)=0.`$ (27)
At the first order in perturbation theory eq.(25) gives, taking into account eq.(27):
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(1\right)}}{\delta \overline{c}}}.`$ (28)
Noticing that
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)=S_0^2\left(\mathrm{\Gamma }^{\left(1\right)}\right),`$ (29)
we conclude that eq.(26) with $`F=\mathrm{\Gamma }^{\left(1\right)}`$ is embodied in eq.(25), under the assumption of the ST invariance of the classical action in eq.(27).
If we adopt a regularization consistent with locality, power-counting and all other unbroken symmetries of the theory (e.g. Lorentz invariance, C-parity, etc.), by using the Quantum Action Principle (QAP) we conclude that the R.H.S. of eq. (28) is zero. Indeed, from the QAP and the power counting theorem we know that $`d^4xc`$$`\frac{\delta \mathrm{\Gamma }^{\left(1\right)}}{\delta \overline{c}}`$ is a local C-even, Lorentz invariant functional and has dimension less or equal four and FP charge equal two. There are no terms with these properties, so we get the following equation:
$`S_0^2\left(\mathrm{\Gamma }^{\left(1\right)}\right)=0,`$ (30)
i.e. even in the massive (non-nilpotent) case the QAP and the power-counting imply that the linearized ST operator $`S_0`$ is nilpotent on the space of the first-order quantum corrections $`\mathrm{\Gamma }^{\left(1\right)}`$.
Let us come back to the study of eq.(28). Since we know that its R.H.S. is zero, thanks to the QAP and the power counting theorem, we have the following equation for the breaking terms $`𝒜_1=S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$:
$`S_0\left(𝒜_1\right)=0.`$ (31)
Since even in the massive case $`S_0`$ is nilpotent on the action-like functionals, one can apply the same decomposition of eq.(22). In the Abelian HK model there is just one cohomologically non trivial insertion, $`d^4x\overline{c}c_\mu cA^\mu `$. It has the right quantum numbers and the correct exact symmetries of the theory (it is Lorentz-invariant and C-even). However, it can be excluded thanks to the ghost equation (15), the QAP and the power counting. For $`n1`$ the ghost equation (15) can be written in the form
$`𝒢\mathrm{\Gamma }^{\left(n\right)}=0.`$ (32)
Using eq.(32) and (19) we get:
$`𝒢S_0\left(\mathrm{\Gamma }^{\left(1\right)}\right)=S_0\left(𝒢\mathrm{\Gamma }^{\left(1\right)}\right)=0.`$ (33)
By power counting $`S_0\left(\mathrm{\Gamma }^{\left(1\right)}\right)`$ cannot contain the external source $`J_2`$ and from eq. (33) we conclude
$`{\displaystyle \frac{\delta }{\delta \overline{c}}}S_0\left(\mathrm{\Gamma }^{\left(1\right)}\right)=0.`$ (34)
Thus the non-trivial breaking term $`d^4x\overline{c}c_\mu cA^\mu `$ is not present and the HK model turns out to be non-anomalous; suitable counter-terms can be constructed at the first order in perturbation theory (actually at all orders), by which the STI can be restored (see , ).
It is worthwhile noticing that on purely algebraic grounds there are no reasons to exclude the anomalous insertion $`d^4x\overline{c}c_\mu cA^\mu `$. To this extent, the properties of the renormalization procedure, dictated by the QAP and the power counting theorem, are essential.
Suppose now that the STI have been restored up to the $`\left(n1\right)`$-th order in perturbation theory, i.e. we assume that suitable counter-terms have been added iteratively to $`\mathrm{\Gamma }^{\left(j\right)}`$, $`j=1,\mathrm{},k`$, in order to restore the STI till order $`n1`$<sup>5</sup><sup>5</sup>5Of course, this is possible only in the absence of anomalies, as it is in the Abelian HK model; in this case it can be proven that the restoration of the STI to the $`n`$-order doesn’t change the counter-terms needed to recover the STI up to the $`(n1)`$-th order, and can be performed, if the STI are fulfilled till order $`n1`$, by a proper choice of counter-terms at the $`n`$-order only.
$`S\left(\mathrm{I}\mathrm{\Gamma }^{\left(k\right)}\right)=0,k=1,2,\mathrm{},n1`$ (35)
$`\mathrm{I}\mathrm{\Gamma }^{\left(k\right)}`$ denotes the correct symmetric effective action at the $`k`$-th order in perturbation theory.
Then eq. (25) becomes
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(n\right)}\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(n\right)}}{\delta \overline{c}}}`$ (36)
or
$`S_0\left(S_0\left(\mathrm{\Gamma }^{\left(n\right)}\right)+{\displaystyle \underset{j=1}{\overset{n1}{}}}(\mathrm{I}\mathrm{\Gamma }^{\left(nj\right)},\mathrm{I}\mathrm{\Gamma }^{\left(j\right)})\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(n\right)}}{\delta \overline{c}}}.`$ (37)
Taking into account eq.(26) we arrive at the following consistency condition for the lower orders parts of the effective action:
$`S_0\left({\displaystyle \underset{j=1}{\overset{n1}{}}}(\mathrm{I}\mathrm{\Gamma }^{\left(nj\right)},\mathrm{I}\mathrm{\Gamma }^{\left(j\right)})\right)=0.`$ (38)
This consistency condition is a consequence of the form of the linearized ST operator and of the lower order requirements (35). Again, it relies on the use of the QAP to ensure the fulfillment of STI at lower orders in perturbation theory.
## 3 Higher orders
We now consider eq.(25) at the second order in perturbation theory. We do not assume that the STI have been restored at the first order. Then eq.(25) can be written as
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(2\right)}\right)+S_{\mathrm{\Gamma }^{\left(1\right)}}\left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)={\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(2\right)}}{\delta \overline{c}}}.`$ (39)
We show that, if we use a renormalization scheme where the QAP holds, the R.H.S. of eq. (39) is zero. We have to verify that
$`{\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(n\right)}}{\delta \overline{c}}}=0,`$ (40)
for all $`n`$. For $`n=0`$ the classical action $`\mathrm{\Gamma }^{\left(0\right)}`$ (appendix A) satisfies eq.(40), as it can be checked by explicit computation. Suppose now that eq.(40) is verified till order $`n1`$:
$`{\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(k\right)}}{\delta \overline{c}}}=0,k=1,\mathrm{},n1`$ (41)
By using the QAP, at the next order in perturbation theory we get:
$`{\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(n\right)}}{\delta \overline{c}}}={\displaystyle d^4x\mathrm{\Delta }\left(x\right)},`$ (42)
where $`d^4x\mathrm{\Delta }\left(x\right)`$ is an integrated Lorentz invariant local polynomial $`\mathrm{\Delta }\left(x\right)`$ in the fields of the theory. $`\mathrm{\Delta }\left(x\right)`$ has dimension $`4`$, $`FP`$-charge $`+2`$ and it obeys all the exact symmetries of the model. Since there are no terms with these properties ($`d^4xcc=0`$, $`d^4xc\mathrm{}c=0`$, $`d^4xA_\mu c^\mu c`$ is excluded by $`C`$-parity, and so on), we conclude that at the $`n`$-th order
$`{\displaystyle \frac{M^2}{\alpha }}{\displaystyle d^4xc\frac{\delta \mathrm{\Gamma }^{\left(n\right)}}{\delta \overline{c}}}=0.`$ (43)
This in turn implies that eq.(40) holds true for all $`n`$.
This result can be demonstrated by a direct analysis of the Feynman graphs, arising in the perturbative expansion of $`\mathrm{\Gamma }`$.
Moreover, the QAP also implies that $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ cannot depend on external sources. Thus eq.(39) simplifies to
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(2\right)}\right)={\displaystyle \underset{i}{}}{\displaystyle d^4x\frac{\delta \mathrm{\Gamma }^{\left(1\right)}}{\delta J_i\left(x\right)}\frac{\delta \left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)}{\delta \varphi _i\left(x\right)}},`$ (44)
where the sum is extended to all the fields $`\varphi _i`$ whose BRST variation is non linear.
From eq.(44) we see that, if the STI have been restored at the first order (i.e. $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}=0`$), the Wess-Zumino consistency condition holding for $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ is true for $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ too. In particular, $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ is also local, like $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$.
On the contrary, if $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}0`$, the Wess-Zumino consistency condition for $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ is modified by the R.H.S. of eq.(44), which now is non-zero. Moreover, we show in eq.(50) that eq.(44) implies that $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ receives non-local contributions, arising from the insertion of the local functional $`d^4x`$ $`\frac{\delta \left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)}{\delta \varphi _i\left(x\right)}`$ on the non-local quantities $`\frac{\delta \mathrm{\Gamma }^{\left(1\right)}}{\delta J_i\left(x\right)}`$.
In order to simplify the notations, we define $`XS\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$. We also define (in the momentum space) the fourth-order differential operator
$`𝒫(){\displaystyle \frac{\delta ^4()}{\delta \overline{c}\left(q\right)\delta c\left(s\right)c\left(t\right)c\left(r\right)}}`$ (45)
We apply $`𝒫`$ on both sides of eq.(44) and then set the fields (including external sources) to zero. By using $`C`$-parity and the fact that the $`FP`$-charge of $`\mathrm{\Gamma }`$ is zero, we get <sup>6</sup><sup>6</sup>6We use a short hand for functional derivatives: $`Y_{c(p)}`$ denotes $`\frac{\delta Y}{\delta c(p)}`$, $`Y_{c(p)\varphi _1(s)}`$ denotes $`\frac{\delta ^2Y}{\delta c(p)\delta \varphi _1(s)}`$, and so on.
$`𝒫S_0\left(X\right)|_{\phi =0}`$ $`=`$ $`ir_\mu X_{\overline{c}\left(q\right)c\left(s\right)c\left(t\right)A_\mu \left(r\right)}|_{\phi =0}+evX_{\overline{c}\left(q\right)c\left(s\right)c\left(t\right)\varphi _2\left(r\right)}|_{\phi =0}+e^2vX_{c\left(s\right)c\left(t\right)J_1\left(qr\right)}|_{\phi =0}`$ (46)
$`+\mathrm{cycl}.\mathrm{permutations}\mathrm{of}(\mathrm{s},\mathrm{t},\mathrm{r})`$
and
$`𝒫\left[{\displaystyle \underset{i}{}}\left({\displaystyle \frac{d^4p}{\left(2\pi \right)^4}\mathrm{\Gamma }_{J_i\left(p\right)}^{\left(1\right)}S\left(\mathrm{\Gamma }\right)_{\varphi _i\left(p\right)}^{\left(1\right)}}\right)\right]|_{\phi =0}=`$
$`{\displaystyle }{\displaystyle \frac{d^4p}{\left(2\pi \right)^4}}(\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(r\right)c\left(t\right)J_2\left(p\right)}^{\left(1\right)}|_{\phi =0}S\left(\mathrm{\Gamma }\right)_{c\left(s\right)\varphi _2\left(p\right)}^{\left(1\right)}|_{\phi =0}+\mathrm{cycl}.\mathrm{permutations}\mathrm{of}(\mathrm{s},\mathrm{t},\mathrm{r}))`$ (47)
Taking into account the conservation of momenta in 1PI Green functions, we see that the R.H.S. of eq.(47) is not zero for $`s+q+r+t=0`$, for $`t+q+r+s=0`$, or for $`r+s+q+t=0`$.
The amplitude $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)}^{\left(1\right)}|_{\phi =0}`$ gets contributions from the graphs in Figures 1 and 2. All external momenta are understood to be incoming. The relevant Feynman rules are briefly discussed in appendix B.
Working out the explicit form of $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)}^{\left(1\right)}|_{\phi =0}`$ we have (after performing the Wick rotation):
$`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)}^{\left(1\right)}|_{\phi =0}`$ $`=`$ $`i{\displaystyle \frac{e^7v^3}{\alpha ^2}}{\displaystyle }d^4k({\displaystyle \frac{1}{k^2+m_1^2}}{\displaystyle \frac{1}{\left(k+q\right)^2+m_g^2}}{\displaystyle \frac{1}{\left(k+q+r\right)^2+m_1^2}}{\displaystyle \frac{1}{\left(kp\right)^2+m_g^2}}`$ (48)
$`+{\displaystyle \frac{1}{k^2+m_1^2}}{\displaystyle \frac{1}{\left(k+r\right)^2+m_g^2}}{\displaystyle \frac{1}{\left(k+q+r\right)^2+m_1^2}}{\displaystyle \frac{1}{\left(kp\right)^2+m_g^2}}`$
$`(rt))\delta ^{\left(4\right)}(p+q+r+t)`$
All integrals are convergent (no subtraction required). For general momenta $`p,q,r,t`$ the R.H.S. of eq.(48) is non-zero. Moreover, it is non-polynomial in the independent external momenta: by applying Weinberg’s theorem we conclude that for non-exceptional external momenta the amplitude $`\mathrm{\Gamma }_{\overline{c}\left(Qq\right)c\left(Qt\right)c\left(Qr\right)J_2\left(Qp\right)}^{\left(1\right)}|_{\phi =0}`$ behaves as $`Q^4`$ for $`Q\mathrm{}`$ and fixed $`p,q,r,t`$ (in the Euclidean region).
A direct calculation of $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$, obtained by applying the method described in , shows that
$`S\left(\mathrm{\Gamma }\right)_{c\left(s\right)\varphi _2\left(p\right)}^{\left(1\right)}|_{\phi =0}=\left(a+bs^2\right)\delta ^{\left(4\right)}\left(s+p\right)`$ (49)
for some $`c`$-numbers $`a,b`$ (depending on the intermediate renormalization scheme used).
By using eqs.(46), (47) and (49) we obtain the following equation
$`ir_\mu X_{\overline{c}\left(q\right)c\left(s\right)c\left(t\right)A_\mu \left(r\right)}|_{\phi =0}+evX_{\overline{c}\left(q\right)c\left(s\right)c\left(t\right)\varphi _2\left(r\right)}|_{\phi =0}+e^2vX_{c\left(s\right)c\left(t\right)J_1\left(qr\right)}|_{\phi =0}`$
$`+\mathrm{cycl}.\mathrm{permutations}\mathrm{of}(\mathrm{s},\mathrm{t},\mathrm{r})`$
$`=\left(a+bs^2\right)\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(r\right)c\left(t\right)J_2\left(s\right)}^{\left(1\right)}+\mathrm{cycl}.\mathrm{permutations}\mathrm{of}(\mathrm{s},\mathrm{t},\mathrm{r})`$ (50)
Eq.(48) implies that the R.H.S. of eq.(50) is non-polynomial in the variables $`s,q,r,t`$. Thus at least one of the amplitudes $`X_{\overline{c}ccA_\mu },X_{\overline{c}cc\varphi _2},X_{ccJ_1}`$ is non-polynomial in the external momenta $`s,q,r,t`$. This in turn implies that $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ is non-local.
From eq.(50) we learn that some of the terms arising in the expansion of $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ on a basis of the fields $`\mathrm{\Phi }`$ and the external sources $`\mathrm{\Phi }^{}`$ of the theory must contain an arbitrarily high number of derivatives. We will show that the expansion of $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ on a basis of $`\mathrm{\Phi },\mathrm{\Phi }^{}`$ must also contain terms with an arbitrarily high number of fields $`\varphi _1`$.
Consider the insertion of $`k`$ fields $`\varphi _1`$ along the $`\overline{c}c`$-lines or the $`\varphi _1\varphi _1`$-lines in the graphs shown in figures 1 and 2. It is worthwhile performing the construction of these insertions in a recursive way. It is possible to insert a leg $`\varphi _1\left(w_1\right)`$ (carrying momentum $`w_1`$) in the graph on the left of figure 1 by cutting a $`\overline{c}c`$-propagator or by cutting a $`\varphi _1\varphi _1`$-line. The graph on the left of figure 1 thus generates four graphs contributing to the 1-PI amplitude $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)\varphi _1\left(w_1\right)}^{\left(1\right)}|_{\phi =0}`$. They are shown in figure 3.
The same construction can be applied to all other graphs appearing in figures 1 and 2, yielding a family $`^1`$ of graphs contributing to $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)\varphi _1\left(w_1\right)}^{\left(1\right)}|_{\phi =0}`$. Moreover, the graphs in $`^1`$ exhaust all possible graphs contributing to $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)\varphi _1\left(w_1\right)}^{\left(1\right)}|_{\phi =0}`$. Indeed, if $`𝒯`$ is a graph contributing to $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)\varphi _1\left(w_1\right)}^{\left(1\right)}|_{\phi =0}`$, the leg $`\varphi _1\left(w_1\right)`$ must be inserted either on a $`\varphi _1\varphi _1`$-line or on a $`\overline{c}c`$-line. Thus removing the insertion of $`\varphi _1\left(w_1\right)`$ and gluing together the two propagators $`\varphi _1\varphi _1`$ \- $`\varphi _1\varphi _1`$ or $`\overline{c}c`$-$`\overline{c}c`$ of the line with the $`\varphi _1\left(w_1\right)`$ insertion into a single $`\varphi _1\varphi _1`$ or $`\overline{c}c`$ propagator respectively yields a graph contributing to $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)}^{\left(1\right)}|_{\phi =0}`$. The latter must be one of the graphs in figures 1 or 2. Hence $`𝒯`$ belongs to $`^1`$.
From an analytic point of view, the construction of graphs in $`^1`$ amounts to the replacement
$`{\displaystyle \frac{1}{k^2+m_g^2}}{\displaystyle \frac{1}{\left(k+w_1\right)^2+m_g^2}}{\displaystyle \frac{1}{k^2+m_g^2}}`$ (51)
(for a $`\varphi _1\left(w_1\right)`$-insertion on a $`\overline{c}c`$-line) or
$`{\displaystyle \frac{1}{k^2+m_1^2}}{\displaystyle \frac{1}{\left(k+w_1\right)^2+m_1^2}}{\displaystyle \frac{1}{k^2+m_1^2}}`$ (52)
(for a $`\varphi _1\left(w_1\right)`$-insertion on a $`\varphi _1\varphi _1`$-line).
No new cancellations arise from these replacements, as it can be seen from the explicit computation of the associated integrals, leading to a straightforward generalization of eq.(48). So one can conclude that $`\mathrm{\Gamma }_{\overline{c}\left(q\right)c\left(t\right)c\left(r\right)J_2\left(p\right)\varphi _1\left(w_1\right)}^{\left(1\right)}|_{\phi =0}`$ is non-zero.
This construction can be applied to the graphs of $`^1`$ to generate the family $`^2`$, whose elements are graphs contributing to $`\mathrm{\Gamma }_{J_2\left(p\right)\overline{c}\left(q\right)c\left(r\right)c\left(s\right)\varphi _1\left(w_1\right)\varphi _1\left(w_2\right)}^{\left(1\right)}|_{\phi =0}`$. Again, the same remarks as before apply and one finds out that $`\mathrm{\Gamma }_{J_2\left(p\right)\overline{c}\left(q\right)c\left(r\right)c\left(s\right)\varphi _1\left(w_1\right)\varphi _1\left(w_2\right)}^{\left(1\right)}|_{\phi =0}0`$. The recursion can be iterated till the $`k`$ insertions of $`\varphi _1\left(w_j\right)`$, $`j=1,\mathrm{},k`$ have been completed.
We now introduce an extension of the operator $`𝒫`$ defined in eq.(45):
$`𝒫_k(){\displaystyle \frac{\delta ^{4+k}()}{\delta \varphi _1\left(w_k\right)\delta \varphi _1\left(w_{k1}\right)\mathrm{}\delta \varphi _1\left(w_1\right)\delta \overline{c}\left(q\right)\delta c\left(s\right)\delta c\left(t\right)\delta c\left(r\right)}}.`$ (53)
By applying $`𝒫_k`$ to eq.(44) and setting next the fields to zero, one gets the following set of equations:
$`𝒫_k\left[S_0\left(S\left(\mathrm{\Gamma }\right)\right)^{\left(2\right)}\right]_{\phi =0}={\displaystyle \underset{i}{}}{\displaystyle d^4x𝒫_k\left[\frac{\delta \mathrm{\Gamma }^{\left(1\right)}}{\delta J_i\left(x\right)}\frac{\delta \left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)}{\delta \varphi _i\left(x\right)}\right]_{\phi =0}}`$ (54)
For $`k=0`$ we recover eq.(50). The R.H.S. of eq.(54) contains functions of the form
$`S\left(\mathrm{\Gamma }\right)_{c\left(p_1\right)\varphi _2\left(p_2\right)_{a=1}^k\varphi _1\left(q_a\right)}^{\left(1\right)}=Q(p_1,p_2,q_a)\delta ^{\left(4\right)}\left(p_1+p_2+{\displaystyle \underset{a=1}{\overset{k}{}}}q_a\right),`$ (55)
where $`Q(p_1,p_2,q_a)`$ is a polynomial of degree at most $`2`$ in $`p_1,p_2,q_a`$. These functions are zero for $`k>4`$, as it can be seen by a direct calculation applying the method in .
The R.H.S. of eq.(54) has the same feature as the R.H.S. of eq.(50): each configuration of external momenta, compatible with the delta functions contained in the R.H.S. of eq.(54), picks out one and only one of the terms which arise in the expansion of the R.H.S. of eq.(54) in terms of the amplitudes $`\mathrm{\Gamma }_{c\left(p_1\right)\overline{c}\left(p_2\right)c\left(p_3\right)J_2\left(p_4\right)_{a=1}^j\varphi _1\left(q_a\right)}^{\left(1\right)}|_{\phi =0}`$, $`j=1,\mathrm{},k`$. Thus, having shown that
$`\mathrm{\Gamma }_{c\left(p_1\right)\overline{c}\left(p_2\right)c\left(p_3\right)J_2\left(p_4\right)_{a=1}^j\varphi _1\left(q_a\right)}^{\left(1\right)}|_{\phi =0}0,`$ (56)
we conclude that the R.H.S. of eq.(54) is non-zero for every $`k`$. Eq.(54) implies that, in the expansion of $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ on a basis of $`\mathrm{\Phi },\mathrm{\Phi }^{}`$, there are non-zero terms associated with monomials containing an arbitrary number of $`\varphi _1`$ fields.
This has some interesting consequences. We have shown that, if improper finite counter-terms in $`\mathrm{\Gamma }^{\left(1\right)}`$ are chosen, at the second order in perturbation theory the STI
$`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}=𝒜_2`$ (57)
are broken by a non-local functional $`𝒜_2`$. Eq. (57) gives rise upon differentiation with respect to a set of fields and external sources $`\{\mathrm{\Phi }^I,\mathrm{\Phi }_I^{}\}_I`$ (with $`I`$ running in the set of indices $``$) to a number of relations among 1-PI Green functions, once we set $`\mathrm{\Phi }^I=\mathrm{\Phi }_I^{}=0`$ after taking the relevant derivatives.
We can expand $`𝒜_2`$ on a basis of monomials in $`\mathrm{\Phi },\mathrm{\Phi }^{}`$ and their derivatives. Since $`𝒜_2`$ is non-local, an infinite number of monomials appears in this expansion. It may happen that there is a maximum finite number $`O`$ of $`\mathrm{\Phi },\mathrm{\Phi }^{}`$, appearing in every monomial of the expansion. In this case, the expansion is infinite because it contains monomials with arbitrarily high order derivatives. Thus we can differentiate eq.(57) with respect to a number of fields greater than $`O`$, yielding the same result as if $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}=0`$. Only a finite number of relations among 1-PI Green functions, valid in the invariant case $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}=0`$, is altered by this type of non-local breaking terms $`𝒜_2`$. Notice that this is the same behavior one has when the breaking is local.
It may also happen that in the expansion of $`𝒜_2`$ there appear monomials with an arbitrarily high number of $`\mathrm{\Phi },\mathrm{\Phi }^{}`$. Now an infinite set of relations among 1-PI Green functions, derived from eq.(57) upon differentiation, is changed with respect to the invariant case. In this sense, violation of locality by arbitrarily high number of $`\mathrm{\Phi },\mathrm{\Phi }^{}`$ is more severe than violation of locality by arbitrarily high number of derivatives only.
We briefly comment on the results of this section. Had we restored the STI at the first order in perturbation theory, eq.(39) would have read
$`S_0\left(S\left(\mathrm{\Gamma }\right)^{\left(2\right)}\right)=0.`$ (58)
If $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}=0`$, the same Wess-Zumino consistency condition holds true both for $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ (see eq.(21)) and $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$. In particular, $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ is local. Moreover, eq.(50) shows that if $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}0`$, $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ receives non-local contributions. Thus a necessary and sufficient condition for $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ to be local is that $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}=0`$. Notice that one can impose in the Abelian HK model $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}=0`$ because the model does not possess physical anomalies.
This result admits a wider generalization. Suppose that the gauge theory under investigation is truly anomalous. Then at the second order in perturbation theory the consistency condition obeyed by $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ is eq.(44), where now $`S\left(\mathrm{\Gamma }\right)^{\left(1\right)}`$ is non zero for any choice of the first order action-like counterterms. An argument similar to the one leading to eq.(46) thus entails that $`S\left(\mathrm{\Gamma }\right)^{\left(2\right)}`$ must be non local, because of the contributions coming from $`S_{\mathrm{\Gamma }^{\left(1\right)}}\left(S\left(\mathrm{\Gamma }\right)^{\left(1\right)}\right)`$<sup>7</sup><sup>7</sup>7 In it was argued that, if one recovers the spurious contributions to the first order anomaly, $`S(\mathrm{\Gamma })^{(2)}`$ can be made local and chosen in such a way that it satisfies the same Wess-Zumino consistency condition as $`S(\mathrm{\Gamma })^{(1)}`$. To match these requirements, one needs to introduce the first order truly anomalous (local) terms as interaction vertices in the quantum effective action $`\mathrm{\Gamma }^{(1)}`$, by coupling them to external sources of negative dimension. In our opinion, this procedure generates a set of Feynman rules which spoil the validity of the QAP at the next order. In particular, at the next order $`S_0(S(\mathrm{\Gamma })^{(2)})`$ gets non-local contributions, which in our framework are embodied in $`S_{\mathrm{\Gamma }^{(1)}}(\mathrm{\Gamma }^{(1)})`$. Notice that, if the first order physical anomalies had dimension $`4`$ (which is not forbidden by the QAP, saying only that they must have dimension less or equal to $`5`$), no troubles would arise in including them as vertices in $`\mathrm{\Gamma }^{(1)}`$. They could be coupled to external sources with non-negative dimension. Actually, truly anomalous terms have dimension $`5`$ only (at least for a gauge group without Abelian factors) . This in turn implies that they cannot just be thought as new interaction vertices, since these vertices must contain external sources with dimension $`1`$.
## 4 Conclusions
In this paper we have shown that, if the action-like counter-terms entering in $`\mathrm{\Gamma }^{\left(1\right)}`$ are not properly chosen, even a physically non-anomalous theory exhibits a non-local second order anomaly. This anomaly cannot be restored by local second order counterterms. Thus, an improper choice of the finite part of the first order counter-terms renders a first-order physically non-anomalous theory a second order truly anomalous one.
Moreover, we have argued that, if one starts with a truly anomalous theory, locality of the STI breaking terms is satisfied at the first order in perturbation theory only, no matter which renormalization scheme is adopted.
We conclude that locality of the STI breaking terms can be maintained to all orders if and only if there are no truly anomalous terms at the first order in perturbation theory.
Finally, we have shown that strict nilpotency of the BRST transformations (and consequently of the linearized ST operator) is not an essential requirement in order to perform the characterization of the STI breaking terms, independently on the order of the perturbative expansion.
Acknowledgements
We acknowledge Professor R. Ferrari for useful discussions.
## Appendix A Classical action
The classical action for the HK model in the on-shell formalism is
$`\mathrm{\Gamma }^{\left(0\right)}={\displaystyle }d^4x[{\displaystyle \frac{1}{4}}F_{\mu \nu }^2+{\displaystyle \frac{e^2v^2}{2}}A_\mu ^2`$
$`{\displaystyle \frac{\alpha }{2}}A^2+\alpha \overline{c}\mathrm{}c+e^2v^2\overline{c}c+e^2v\overline{c}c\varphi _1`$
$`+{\displaystyle \frac{1}{2}}\left(\left(_\mu \varphi _1\right)^2+\left(_\mu \varphi _2\right)^2\right)\lambda v^2\varphi _1^2{\displaystyle \frac{e^2v^2}{2\alpha }}\varphi _2^2`$
$`+eA_\mu \left(\varphi _2^\mu \varphi _1^\mu \varphi _2\varphi _1\right)+e^2v\varphi _1A^2+{\displaystyle \frac{e^2}{2}}\left(\varphi _1^2+\varphi _2^2\right)A^2`$
$`\lambda v\varphi _1\left(\varphi _1^2+\varphi _2^2\right){\displaystyle \frac{\lambda }{4}}\left(\varphi _1^2+\varphi _2^2\right)^2`$
$`+\overline{\psi }i\overline{)}\psi +Gv\overline{\psi }\psi +{\displaystyle \frac{e}{2}}\overline{\psi }\gamma _\mu \gamma _5\psi A^\mu `$
$`+G\overline{\psi }\psi \varphi _1iG\overline{\psi }\gamma _5\psi \varphi _2`$
$`+J_1\left[ec\varphi _2\right]+J_2ec\left(\varphi _1+v\right)+i{\displaystyle \frac{e}{2}}\overline{\eta }\gamma _5\psi c+i{\displaystyle \frac{e}{2}}c\overline{\psi }\gamma _5\eta `$
$`+{\displaystyle \frac{M^2}{2}}A_\mu ^2+M^2\overline{c}c{\displaystyle \frac{M^2}{2\alpha }}(\varphi _1^2+\varphi _2^2)]`$ (59)
The explicit mass term is in evidence in the last line.
BRST transformations
Off-shell formalism
$`sA_\mu =_\mu c,s\varphi _1=ec\varphi _2,s\varphi _2=ec\left(\varphi _1+v\right)`$
$`s\psi =i{\displaystyle \frac{e}{2}}\gamma _5\psi c,s\overline{\psi }=i{\displaystyle \frac{e}{2}}c\overline{\psi }\gamma _5,s\overline{c}=B,sB={\displaystyle \frac{M^2}{\alpha }}c,sc=0`$ (60)
In the on-shell formalism the $`B`$ field disappears and the BRST transformation of $`\overline{c}`$ becomes
$`s\overline{c}=A+{\displaystyle \frac{ev}{\alpha }}\varphi _2,`$ (61)
## Appendix B Feynman rules
We only recall the Feynman rules needed to evaluate the graphs in Figures 1 and 2.
Propagator for $`\varphi _1\varphi _1`$
$`\mathrm{\Delta }_{\varphi _1\varphi _1}\left(p\right)={\displaystyle \frac{i}{p^2m_1^2+iϵ}},m_1^2=2\lambda v^2+{\displaystyle \frac{M^2}{\alpha }}`$ (62)
The propagator for $`\varphi _1\varphi _1`$ is denoted by a solid line.
Propagator for $`c\overline{c}`$
$`\mathrm{\Delta }_{c\overline{c}}\left(p\right)={\displaystyle \frac{i}{\alpha \left(p^2m_g^2+iϵ\right)}},m_g^2={\displaystyle \frac{e^2v^2+M^2}{\alpha }}`$ (63)
The propagator for $`c\overline{c}`$ is denoted by a dashed line.
Vertices
$`ie\left(2\pi \right)^4\delta ^{\left(4\right)}\left(\text{incoming momenta}\right)`$ (64)
for the vertex $`eJ_2c\varphi _1`$ and
$`ie^2v\left(2\pi \right)^4\delta ^{\left(4\right)}\left(\text{incoming momenta}\right)`$ (65)
for the vertex $`e^2v\overline{c}c\varphi _1`$ (see Figure 4). |
warning/0001/cond-mat0001175.html | ar5iv | text | # Origin of the singular Bethe ansatz solutions for the Heisenberg 𝑋𝑋𝑍 spin chain
## I Introduction
The Bethe ansatz method has been applied to various problems in condensed matter physics. The Heisenberg $`XXZ`$ spin chain, the Hubbard model, the Kondo model, and many other models are solved by the Bethe ansatz or its variants . The Heisenberg $`XXZ`$ spin chain describes various physical systems including interacting fermion systems. Also, it is related to the transfer matrix of the six-vertex model , the asymmetric exclusion process, and a Kardar-Parisi-Zhang-type growth model . For the $`XXZ`$ Hamiltonian, the Bethe ansatz expresses energy eigenvalues and eigenstates as functions of quasi-particle momenta which should satisfy a set of algebraic equations, the Bethe ansatz equations (BAE). Various numerical and analytic methods have been developed to solve the BAE, from which the ground state energy and low-lying excitation spectrum are calculated . Combined with a finite-size-scaling theory, they provide useful tools to study equilibrium and nonequilibrium critical phenomena.
When applying the Bethe ansatz method to the Heisenberg $`XXZ`$ spin chain, one may doubt whether the Bethe ansatz produces the complete set of eigenstates. The completeness of the Bethe ansatz is proved for an inhomogeneous generalization of the $`XXZ`$ chain threaded by the Aharonov-Bohm flux . It implies that the Bethe ansatz for the $`XXZ`$ Hamiltonian is also complete in a general sense. One may obtain the eigenstate of the $`XXZ`$ Hamiltonian from that of the generalized Hamiltonian by taking a limit where the inhomogeneity and the flux vanish.
On the other hand, it has also been investigated for a long time whether the Bethe ansatz produces all eigenstates of the $`XXZ`$ Hamiltonian without the help of the limiting procedure. There have been the arguments for the completeness assuming a string conjecture. For the $`XXX`$ chain, which has isotropic spin-spin interaction, Takahashi showed that the number of the states constructed from the Bethe ansatz using the string conjecture is equal to the total number of the spin states. The string conjecture was also adopted to prove the completeness of the Bethe ansatz for a generalized $`XXZ`$ chain and the Hubbard model . Later, the string conjecture was shown to be invalid by Essler et al. , who analyzed the algebraic Bethe ansatz equations in the two down-spin sector. They found that some solutions expected from the string conjecture are missing and that there exist solutions violating the string conjecture. Nevertheless, they concluded that the Bethe ansatz is complete at least in the two down-spin sector because the numbers of the former and the latter are the same.
However one of the solutions found in Ref. is singular; the corresponding Bethe ansatz wave function is not well-defined. In fact the state $`|\mathrm{\Psi }=_l(1)^l|l,l+1`$, which is an eigenstate of the $`XXZ`$ Hamiltonian, is not produced from the regular (non-singular) solutions of the Bethe ansatz (see below for the definition of the state $`|l,l+1`$). Quasi-particle momenta corresponding to the state are divergent and the Bethe ansatz wave function is ill-defined. An eigenstate that is not produced from the regular solutions of the Bethe ansatz will be referred to as a singular state. Siddharthan reported that there also exist the singular states in the three down-spin sector. It was claimed that a symmetry property might be important since the singular states have a definite symmetry property.
The singular state does not imply the incompleteness of the Bethe ansatz for the $`XXZ`$ Hamiltonian. One may obtain the wave functions for the singular states of the $`XXZ`$ Hamiltonian from those of the generalized Hamiltonian by taking the limit where the inhomogeneity and the flux vanish (see, for example, Ref. ). However it should be examined whether the singular states are relevant to the involved physical quantities of the Heisenberg $`XXZ`$ spin chain since, if so, one must first consider the generalized model and take the suitable limit . Therefore, it is helpful to answer the question why there appear the singular states and how many there are in general down-spin number sectors. In this paper, we investigate the symmetry properties of the $`XXZ`$ Hamiltonian and the Bethe ansatz wave functions to show that the regular solution of the Bethe ansatz with well-defined wave functions does not contain a certain class of eigenstates. Those states are the singular states. We also estimate the number of the singular states approximately at each sector with fixed number of down spins and find that the number of them diverges exponentially with the chain length.
## II Bethe states and symmetry operations
The Heisenberg $`XXZ`$ spin chain of length $`N`$ has the Hamiltonian
$$H=\frac{1}{2}\underset{l=1}{\overset{N}{}}\left\{\sigma _l^+\sigma _{l+1}^{}+\sigma _l^{}\sigma _{l+1}^++\frac{\mathrm{\Delta }}{2}(\sigma _l^z\sigma _{l+1}^z1)\right\},$$
(1)
where $`\sigma _l^{x,y,z}`$ are the Pauli matrices at site $`l`$, $`\sigma _l^\pm =(\sigma _l^x\pm i\sigma _l^y)/2`$. The periodic boundary condition is imposed, $`\sigma _{N+1}^{x,y,z}=\sigma _1^{x,y,z}`$. We assume $`N`$ is even unless stated otherwise. It reduces to the ferromagnetic (antiferromagnetic) $`XXX`$ chain in the case of $`\mathrm{\Delta }=1`$ ($`\mathrm{\Delta }=1`$). Since the Hamiltonian commutes with the magnetization $`M=_l\sigma _l^z`$, one can work in the $`Q`$ down-spin sector in the diagonal basis of $`\sigma ^z`$. A state in the $`Q`$ down-spin sector is spanned by $`{}_{N}{}^{}C_{Q}^{}`$ state vectors $`|x_1,x_2,\mathrm{},x_Q`$ with $`1x_iN`$ ($`x_i<x_j`$ for $`i<j`$) denoting the position of the $`i`$th down spin.
The eigenstate of $`H`$ can be written as $`|\mathrm{\Psi }=_{x_1<\mathrm{}<x_Q}\psi (x_1,\mathrm{},x_Q)|x_1,\mathrm{},x_Q`$. According to the Bethe ansatz, the wave function is given in terms of quasi-particle momenta $`k_m(m=1,\mathrm{},Q)`$ by
$$\psi (x_1,\mathrm{},x_Q)=\underset{P}{}A(P)\mathrm{exp}\left(i\underset{m=1}{\overset{Q}{}}k_{P_m}x_m\right),$$
(2)
where the sum is over all permutations $`P`$ of integers $`\{1,2,\mathrm{},Q\}`$. The quasi-particle momenta are complex numbers with $`\pi <\mathrm{Re}k\pi `$. The amplitudes $`A`$’s for different permutations $`P`$ and $`P^{}`$, which are identical except for two neighboring integers such that $`P_j=P_{j+1}^{}=m,P_{j+1}=P_j^{}=l`$, and $`P_i=P_i^{}`$ for $`ij,j+1`$, are related as
$$\frac{A(P)}{A(P^{})}=\frac{e^{i(k_m+k_l)}2\mathrm{\Delta }e^{ik_m}+1}{e^{i(k_m+k_l)}2\mathrm{\Delta }e^{ik_l}+1}.$$
(3)
The periodic boundary condition implies the Bethe ansatz equations (BAE) for the quasi-particle momenta:
$$e^{iNk_m}=(1)^{Q1}\underset{lm}{}\frac{e^{i(k_m+k_l)}2\mathrm{\Delta }e^{ik_m}+1}{e^{i(k_m+k_l)}2\mathrm{\Delta }e^{ik_l}+1}.$$
(4)
The state constructed from the Bethe ansatz is called the Bethe state. It has the energy $`E=_{m=1}^Q(\mathrm{\Delta }\mathrm{cos}k_m)`$.
In the special case of $`\mathrm{\Delta }=0`$, where the BAE becomes $`e^{ik_mN}=(1)^{Q1}`$, the Bethe ansatz produces all eigenstates of $`H`$. A solution is obtained by selecting $`Q`$ different values among the $`N`$ values of $`(\frac{2\pi }{N})\times `$ integers (half-integers) for odd (even) $`Q`$. There are $`{}_{N}{}^{}C_{Q}^{}`$ different ways for each $`Q`$, which is equal to the possible spin states. Naively one may think that the solutions at $`\mathrm{\Delta }=0`$ evolve continuously as $`\mathrm{\Delta }`$ is turned on. However it will turn out to be false.
The $`XXZ`$ Hamiltonian under the periodic boundary condition commutes with the shift operator $`T`$, which shifts the position of each down spin to the left by one unit. By construction, the Bethe state is the eigenstate of $`T`$ with the eigenvalue $`\mathrm{\Omega }=[_{m=1}^Qe^{ik_m}]`$. Since $`T^N`$ is the identity operator, $`\mathrm{\Omega }^N`$ should be unity. The Hamiltonian also commutes with the lattice inversion operator $`V`$ defined by $`V|\{x_i\}=|\{y_i\}`$ with $`y_i=Nx_{Qi+1}+1`$. Since $`V^2`$ is the identity operator, the eigenvalue of $`V`$, which will be denoted by $`\mathrm{{\rm Y}}`$, takes the value of $`\pm 1`$. The inversion operator does not commute with the shift operator in general. But it is easy to show that they commute with each other in the subspace of $`\mathrm{\Omega }=\pm 1`$; $`[T,V]|\mathrm{\Psi }=0`$ if $`T|\mathrm{\Psi }=\pm |\mathrm{\Psi }`$. Then the eigenstate of $`H`$ can be made as the simultaneous eigenstate of $`T`$ and $`V`$ in the subspace.
The Bethe state with $`\mathrm{\Omega }=\pm 1`$ is not the eigenstate of $`V`$ necessarily. If a Bethe state $`|\mathrm{\Psi }`$ with $`\mathrm{\Omega }=\pm 1`$ with well-defined wave function is a simultaneous eigenstate of $`V`$ ($`V|\mathrm{\Psi }=\mathrm{{\rm Y}}|\mathrm{\Phi }`$), then one can show that $`\mathrm{{\rm Y}}`$ should be equal to $`\mathrm{\Omega }`$ :
$$\mathrm{{\rm Y}}=\mathrm{\Omega }.$$
(5)
It is proved as follows. Suppose that $`|\mathrm{\Psi }`$ is a Bethe state in the $`Q`$ down-spin sector with the quasi-particle momenta $`\{k_1,k_2,\mathrm{},k_Q\}`$ satisfying Eq. (4). The wave function is given as in Eq. (2). When one applies the lattice inversion operator $`V`$ to the state $`|\mathrm{\Psi }`$, the wave function of the state $`V|\mathrm{\Psi }=\stackrel{~}{\psi }(x_1,\mathrm{},x_Q)|x_1,\mathrm{},x_Q`$ is given by
$`\stackrel{~}{\psi }(x_1,\mathrm{},x_Q)`$ $`=`$ $`\psi (Nx_Q+1,\mathrm{},Nx_1+1)`$ (6)
$`=`$ $`\mathrm{\Omega }{\displaystyle \underset{P}{}}\stackrel{~}{A}(P)\mathrm{exp}\left[i{\displaystyle \underset{m}{}}(k_{P_m})x_m\right],`$ (7)
where $`\stackrel{~}{A}(P)=A(P\overline{P})`$ with a parity permutation $`\overline{P}`$ which maps $`\{1,2,\mathrm{},Q\}`$ to $`\{Q,\mathrm{},2,1\}`$, and we used that $`[_{m=1}^Qe^{ik_m}]=\mathrm{\Omega }`$ and $`\mathrm{\Omega }^N=1`$. One can easily verify that $`\{k_m\}`$ satisfies the BAE in Eq. (4) and that the amplitude $`\stackrel{~}{A}(P)`$ satisfies Eq. (3) with $`\{k_m\}`$ replaced by $`\{k_m\}`$. This shows that $`V|\mathrm{\Psi }`$ is a Bethe state with the quasi-particle momenta $`\{k_1,k_2,\mathrm{},k_Q\}`$. Therefore a Bethe state could be an eigenstate of $`V`$ only when the two sets $`\{k_1,k_2,\mathrm{},k_Q\}`$ and $`\{k_1,k_2,\mathrm{},k_Q\}`$ are identical.
We rewrite the wave function in Eq. (7), introducing the permutation $`\stackrel{~}{P}`$ defined by the relation $`k_l=k_{\stackrel{~}{P}_l}`$, as
$`\stackrel{~}{\psi }(\{x_m\})`$ $`=`$ $`\mathrm{\Omega }{\displaystyle \underset{P}{}}\stackrel{~}{A}(P)\mathrm{exp}\left[i{\displaystyle \underset{m}{}}k_{\stackrel{~}{P}P_m}x_m\right]`$ (8)
$`=`$ $`\mathrm{\Omega }{\displaystyle \underset{P}{}}A(\stackrel{~}{P}P\overline{P})\mathrm{exp}\left[i{\displaystyle \underset{m}{}}k_{P_m}x_m\right].`$ (9)
Comparing it with Eq. (2), one can find $`\mathrm{{\rm Y}}`$ from the ratio of $`A(P)`$ and $`\mathrm{\Omega }A(\stackrel{~}{P}P\overline{P})`$ for a certain permutation $`P`$, e.g., the identity permutation $`I`$. If the amplitude $`A`$ is nonzero and finite, $`\mathrm{{\rm Y}}`$ is given by
$$\mathrm{{\rm Y}}=\mathrm{\Omega }A(\stackrel{~}{P}\overline{P})/A(I).$$
(10)
Note that the set $`\{k_m\}`$ is equal to $`\{k_m\}`$. On one hand, unless both $`0`$ and $`\pi `$ are present in the set $`\{k_m\}`$, one can rearrange the momenta in such a way that the permutation $`\stackrel{~}{P}`$ and the parity permutation $`\overline{P}`$ are the same. Then one has $`A(\stackrel{~}{P}\overline{P})=A(\overline{P}^2)=A(I)`$ and hence $`\mathrm{{\rm Y}}=\mathrm{\Omega }`$. On the other hand, if both $`0`$ and $`\pi `$ are present in $`\{k_m\}`$, $`Q`$ should be even and one can rearrange the momenta in the following way:
$$\{k_1,\mathrm{},k_{Q/21},k_{Q/2},k_{Q/2+1},k_{Q/21},\mathrm{},k_1\},$$
with $`k_{Q/2}=0`$ and $`k_{Q/2+1}=\pi `$. With this arrangement, the permutation $`\stackrel{~}{P}\overline{P}`$ maps $`Q/2`$ to $`Q/2+1`$, $`Q/2+1`$ to $`Q/2`$, and $`jQ/2,Q/2+1`$ to $`j`$. Then, using Eq. (3), the ratio between $`A(I)`$ and $`A(\stackrel{~}{P}\overline{P})`$ is given by
$$\frac{A(I)}{A(\stackrel{~}{P}\overline{P})}=\frac{e^{ik_{Q/2+1}+ik_{Q/2}}2\mathrm{\Delta }e^{ik_{Q/2}}+1}{e^{ik_{Q/2+1}+ik_{Q/2}}2\mathrm{\Delta }e^{ik_{Q/2+1}}+1}=1.$$
(11)
Therefore we find that $`V|\mathrm{\Psi }=\mathrm{\Omega }|\mathrm{\Psi }`$ if $`|\mathrm{\Psi }`$ is an eigenstate of $`V`$. This proves Eq. (5).
## III Perturbative construction of singular states
Equation (5) of previous section shows that the Bethe state with well-defined wave function cannot be the simultaneous eigenstate of $`T`$ and $`V`$ with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(\pm 1,1)`$. But there exist many eigenstates of $`H`$ with such symmetry. Suppose an eigenstate of $`H`$ possesses the quantum numbers $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(1,1)`$. The state may be degenerate (having the same energy) with another state which has $`\mathrm{{\rm Y}}=+1`$ in the same $`\mathrm{\Omega }`$ sector. In the degenerate case, the Bethe state may be a linear superposition of the degenerate pair, which does not lead to the contradiction to Eq. (5). On the other hand, if the eigenstate does not have the degenerate partner with $`\mathrm{{\rm Y}}=+1`$ in the same $`\mathrm{\Omega }`$ sector, such a state cannot be a Bethe state, hence is absent in the regular Bethe ansatz solutions. In the remaining part, we show that there exist nondegenerate states with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(\pm 1,1)`$ for nonzero $`\mathrm{\Delta }`$ in the general $`Q`$ down-spin sector.
It is convenient to work with the second quantized form of the Hamiltonian. Under the Jordan-Wigner and Fourier transformations, defined as $`\sigma _l^+=a_l\mathrm{exp}[i\pi _{j=1}^{l1}a_j^{}a_j]`$ and $`a_l=_pa_pe^{ipl}/\sqrt{N}`$, respectively, the $`XXZ`$ Hamiltonian in the $`Q`$ down-spin sector reads as $`H=H_0+H_1`$ where $`H_0=_p(\mathrm{cos}p\mathrm{\Delta })n_p`$ and
$$H_1=\frac{\mathrm{\Delta }}{N}\underset{p_1+p_2=p_3+p_4}{}e^{i(p_1p_4)}a_{p_1}^{}a_{p_2}^{}a_{p_3}a_{p_4}.$$
Here $`a_p`$ and $`a_p^{}`$ are anticommuting fermion operators with the bare momentum $`\pi <p\pi `$ and $`n_p=a_p^{}a_p`$ is the number operator with $`_pn_p=Q`$. The momentum $`p`$ takes the real value of $`(2\pi /N)\times `$ integer (half-integer) for odd (even) $`Q`$.
Consider a symmetric set $`S_{Q2}`$ which consists of $`(Q2)`$ bare momenta satisfying $`_{pS_{Q2}}p=0`$ (mod $`2\pi `$). A set $`S`$ is symmetric if it contains $`p`$ for all $`pS`$. For a given $`S_{Q2}`$, $`R`$ is defined as the set of all momenta in the interval $`\pi /2<p<\pi /2`$ except $`p^{}`$ or $`\pi p^{}`$ (mod $`2\pi `$) for $`p^{}S_{Q2}`$. Note that $`R`$ is also a symmetric set. Then, at $`\mathrm{\Delta }=0`$, the following states
$$|p_\alpha ;S_{Q2}a_{p_\alpha }^{}a_{\pi p_\alpha }^{}\underset{pS_{Q2}}{}a_p^{}|0$$
(12)
with $`p_\alpha R`$ are degenerate eigenstates of $`H`$ in the $`Q`$ down-spin sector having $`\mathrm{\Omega }=1`$. Here $`|0`$ denotes the vacuum, i.e., $`a_p|0=0`$ for all $`p`$. They are not the eigenstates of $`V`$. One can easily show that
$$V|p_\alpha ;S_{Q2}=\mathrm{\Omega }|p_\alpha ;S_{Q2}.$$
(13)
Figure 1 illustrates an example of the set $`S_{Q2}`$ and corresponding $`R`$ with the states in Eq. (12) for $`N=20`$ and $`Q=10`$.
As $`\mathrm{\Delta }`$ is turned on, $`H_1`$ generates the overlap between the degenerate states, which may cause the degeneracy to split. To see the degeneracy splitting, we treat $`H_1`$ as a perturbation and apply the degenerate perturbation theory. In this scheme, the leading order correction to the energy is given by the eigenvalue of the perturbation matrix $`(H_1)_{p_\alpha ,p_\beta }p_\alpha ;S_{Q2}|H_1|p_\beta ;S_{Q2}`$ with $`p_\alpha ,p_\beta R`$. The eigenstates are given by the eigenvectors of the perturbation matrix. A straightforward calculation shows that the matrix element is given by
$$(H_1)_{p_\alpha ,p_\beta }=\frac{4\mathrm{\Delta }}{N}(\mathrm{cos}p_\alpha )(\mathrm{cos}p_\beta ),$$
where we neglect a constant diagonal element. It has the nondegenerate eigenvector
$$|\mathrm{\Psi };S_{Q2}=\underset{p_\alpha R}{}\mathrm{cos}p_\alpha |p_\alpha ;S_{Q2},$$
(14)
with the eigenvalue $`(4\mathrm{\Delta }/N)_{p_\alpha R}\mathrm{cos}^2p_\alpha `$. The other states remain to be degenerate with their common eigenvalue $`0`$. If one applies $`V`$ to the state $`|\mathrm{\Psi };S_{Q2}`$ and uses Eq. (13), one can see that the nondegenerate state is the eigenstate of $`V`$ with $`\mathrm{{\rm Y}}=+1`$ while $`\mathrm{\Omega }=1`$. Therefore, it cannot be produced from the regular solution of the Bethe ansatz by Eq. (5).
The remaining degenerate states may split further in the higher order perturbations. So the number of the nondegenerate states with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(1,+1)`$ is equal to or greater than the number of different realizations of the set $`S_{Q2}`$. For even $`Q`$, $`S_{Q2}`$ consists of $`(Q2)/2`$ values of $`p`$ selected randomly in the interval $`0<p<\pi `$ and their negative partners. For odd $`Q`$, $`S_{Q2}`$ consists of $`p=0`$, $`(Q3)/2`$ $`p`$’s randomly-selected in the interval $`0<p<\pi `$, and their negative partners. So, there are at least $`{}_{N/2}{}^{}C_{(Q2)/2}^{}`$ ($`{}_{N/2}{}^{}C_{(Q3)/2}^{}`$) states which have $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(1,+1)`$ and they are those states which cannot be produced from the Bethe ansatz in the even (odd) $`Q`$ down-spin sector. If we sum over all $`Q`$, the number of such states is of the order of $`2^{N/2}`$.
In the two down-spin sector, the perturbation calculation becomes exact since the states in Eq. (12) with the empty set $`S_{Q2}`$ span the whole two down-spin sector with $`\mathrm{\Omega }=1`$. So the state
$$\underset{\frac{\pi }{2}<p<\frac{\pi }{2}}{}(\mathrm{cos}p)a_p^{}a_{\pi p}^{}|0$$
is the eigenstate of $`H`$ at all $`\mathrm{\Delta }`$, which cannot be produced from the Bethe ansatz. If one performs the inverse Jordan-Wigner and Fourier transformations, the state can be written as $`_l(1)^l|l,l+1`$. This is the example of the singular states mentioned in .
Consider also the states
$$a_{p_\alpha }^{}a_{p_\beta }^{}a_{\pi p_\alpha }^{}a_{\pi p_\beta }^{}\underset{pS_{Q4}}{}a_p^{}|0$$
with $`p_\alpha ,p_\beta R`$, constructed from a symmetric set $`S_{Q4}`$ consisting of $`(Q4)`$ momenta satisfying $`_{pS_{Q4}}p=0`$ (mod $`2\pi `$) and the set $`R`$ of all momenta in the interval $`\frac{\pi }{2}<p<\frac{\pi }{2}`$ except $`p^{}`$ or $`\pi p^{}`$ (mod $`2\pi `$) for $`p^{}S_{Q4}`$. They are degenerate eigenstates of $`H`$ at $`\mathrm{\Delta }=0`$ with $`\mathrm{\Omega }=+1`$. The same perturbation analysis will show that the degeneracy splits for finite $`\mathrm{\Delta }`$ and there appear nondegenerate states with $`\mathrm{{\rm Y}}=1`$. Those states cannot be produced by the Bethe ansatz from Eq. (5) either. There are $`{}_{N/2}{}^{}C_{(Q4)/2}^{}`$ ($`{}_{N/2}{}^{}C_{(Q5)/2}^{}`$) different ways in choosing the set $`S_{Q4}`$ for even (odd) $`Q`$. For each set, there exists at least one state with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(+1,1)`$. So the total number of such states is again of order of $`2^{N/2}`$.
The same analysis can be applied to the states from the symmetric sets $`S_{Q6},S_{Q8},\mathrm{}`$, with $`_{pS}p=0`$ (mod $`2\pi `$) and symmetric sets $`S_{Q1},S_{Q3},\mathrm{}`$ containing $`p=\pi `$ and satisfying $`_{pS}=\pi `$ (mod $`2\pi `$). For each case, there are at least a number of order of $`2^{N/2}`$ singular states.
## IV Numerical Results
The state counting discussed in the previous section is not exact. States from different $`S`$ sets may be degenerate at $`\mathrm{\Delta }=0`$. In that case perturbation calculations should be carried out in an enlarged space and it is not guaranteed that each $`S`$ set generates at least one singular state. And the singular states with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(\pm 1,1)`$ may become degenerate accidentally with a state with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(\pm 1,\pm 1)`$ for a particular value of $`\mathrm{\Delta }`$. Then the Bethe ansatz may produce the linear superpositions of them. We believe that those effects do not affect the leading order of magnitude of the total number of the singular states.
We checked this by diagonalizing the $`XXZ`$ Hamiltonian exactly with $`N=6,8,\mathrm{},16`$. The symmetry property of each eigenstate is checked and the nondegenerate states with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(+1,1)`$ and ($`1,+1`$) are identified as the singular states. Figure 2 shows the energy level, as a function of $`\mathrm{\Delta }`$, of each state with $`\mathrm{\Omega }=\pm 1`$ for $`N=8`$ in the four down-spin sector. The singular states are denoted by dotted lines. Note that each singular state evolves from the degenerate states at $`\mathrm{\Delta }=0`$ as the degeneracy is lifted for nonzero $`\mathrm{\Delta }`$ . At $`\mathrm{\Delta }=\pm \frac{1}{2}`$, many singular states become degenerate with non-singular states and they are not counted as the singular states. In Table 1, we show the total number of the singular states at $`\mathrm{\Delta }=0.2,0.5`$, and $`1.0`$. Due to the accidental degeneracy the number varies with $`\mathrm{\Delta }`$. It is confirmed that the number increases approximately as $`2^{N/2}`$ at the three values of $`\mathrm{\Delta }`$.
## V Discussion and Summary
The singular behaviors of some Bethe ansatz solutions have been noticed for the Heisenberg $`XXZ`$ model (see, e.g., ). So the generalizations to the inhomogeneous $`XXZ`$ model with Aharonov-Bohm flux were considered and their Bethe ansatz solutions were shown to be complete . In those cases the Hamiltonian does not possess the lattice inversion symmetry and the present analysis is not applicable.
The singular states include low-lying states in the half-filling case, $`Q=N/2`$. Figures 1 (b) and (c) illustrate degenerate first excited states at $`\mathrm{\Delta }=0`$ in the $`\mathrm{\Omega }=+1`$ sector. They are not the eigenstates of the inversion operator $`V`$. For nonzero $`\mathrm{\Delta }`$ the degeneracy is removed, and the states evolve into the eigenstates of $`V`$ with $`\mathrm{{\rm Y}}=\pm 1`$. Therefore the Bethe ansatz fails to produce one of them with $`\mathrm{{\rm Y}}=1`$. Those states are important in identifying the operator content of the $`XXZ`$ model in the context of the conformal field theory (see and references therein). They have the scaling dimension $`x_{n,m}`$ with $`n=0`$ and $`m=\pm 1`$. To recover the whole operator content from the Bethe ansatz, one should introduce the Aharonov-Bohm flux and take the limit where the flux vanishes as done in .
So far we have assumed $`N`$ is even. For odd $`N`$, one cannot find a pair of bare momenta $`p`$ and $`p^{}`$ satisfying $`p+p^{}=\pi `$ and the previous analysis cannot be applied. So the above arguments are not valid for odd $`N`$, though Eq. (5) holds good. We have performed the same numerical analysis for odd $`N15`$ and found that the nondegenerate states with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(\pm 1,1)`$ do not exist.
In summary, we have shown that the symmetry properties of the Hamiltonian and the Bethe wave functions prevent the regular solutions of the Bethe ansatz from containing some eigenstates. For the fixed number of down spins, the $`XXZ`$ Hamiltonian can be simultaneously diagonalized with the shift operator $`T`$ and the lattice inversion operator $`V`$ in the subspace where the eigenvalue of $`T`$, $`\mathrm{\Omega }`$, is $`\pm 1`$, and each simultaneous eigenstate has one of the following four pairs of eigenvalues, $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(1,1),(1,1),(1,1),(1,1)`$. But there is no regular Bethe ansatz wave function with $`(\mathrm{\Omega },\mathrm{{\rm Y}})=(1,1),(1,1)`$, which leads to the appearance of the singular Bethe states. We estimated the number of the singular states as of $`O(2^{N/2})`$ for even $`N`$. The number of the singular states is exponentially large though the fraction to the total number of the states, $`2^N`$, is vanishingly small. Moreover, the singular states include low-lying states in the half-filling case where the number of down spins is $`N/2`$.
## ACKNOWLEDGMENTS
This work was supported in part by the BK21 Project of Ministry of Education, Korea. |
warning/0001/cond-mat0001222.html | ar5iv | text | # Localized and Delocalized Charge Transport in Single-Wall Carbon-Nanotube Mats
## Abstract
We measured the complex dielectric constant in mats of single-wall carbon-nanotubes between 2.7 K and 300 K up to 0.5 THz. The data are well understood in a Drude approach with a negligible temperature dependence of the plasma frequency $`\omega _p`$ and scattering time $`\tau `$ with an additional contribution of localized charges. The dielectric properties resemble those of the best ”metallic” polypyrroles and polyanilines. The absence of metallic islands makes the mats a relevant piece in the puzzle of the interpretation of $`\tau `$ and $`\omega _p`$ in these polymers.
PACSnumbers: 71.20.Hk,71.20.Tx,72.80.Le,77.84.Jd
Depending on the wrapping of the graphene sheet, the intrinsic electronic properties of single-wall carbon-nanotubes (SWNTs) are either semiconducting (zigzag and most chiral nanotubes) or metallic (armchair and part of the chiral nanotubes). Single ropes of armchaired SWNTs as well as entangled networks (mats) show a decrease of the dc conductivity, $`\sigma _{\mathrm{dc}}`$, with increasing temperature ($`T`$) for $`T`$ above a critical temperature ($`T^{}`$). For $`T<T^{}`$, $`\sigma _{\mathrm{dc}}`$ decreases with cooling. $`T^{}`$ typically lies between 40 K and 250 K and depends on the morphology and the degree of disorder. The transition from $`d\sigma _{\mathrm{dc}}/dT<0`$ to $`d\sigma _{\mathrm{dc}}/dT>0`$ at $`T^{}`$ is ascribed to structural defects and built-in impurities of the individual nanotubes, or to barriers between the nanotubes or ropes limiting the extension of the charge-carrier states. Additional information about this issue can be obtained from frequency- dependent phase-sensitive permittivity experiments, giving the complex conductivity $`\sigma (\omega )=\sigma ^{}(\omega )+i\sigma ^{\prime \prime }(\omega )`$ or dielectric constant, $`ϵ(\omega )=ϵ^{}(\omega )iϵ^{\prime \prime }(\omega )`$. Different frequency dependencies of $`\sigma ^{}`$ and $`ϵ^{}`$ are expected for the two limiting cases of localized and delocalized charge transport and also the sign of $`ϵ^{}`$ changes with the two models. Here we apply this method to mats of SWNTs. The data do not need Kramers-Kronig analysis, which is a great advantage for the analysis of the response of the delocalized charge carriers. Like heavily doped polymers SWNTs are shown to be an example of a system with exceptional long scattering times $`\tau `$ and low plasma frequencies $`\omega _p`$, but in contrast to the polymers metallic islands can be excluded as an explanation. For that reason SWNTs form an important new element in the not-yet understood physics behind $`\tau `$ and $`\omega _p`$.
Weak inter-nanotube contacts or strong intra-nanotube defects might lead to charge localization. The motion of the localized charge carriers will be diffusion controlled and the frequency-dependent conductivity $`\sigma (\omega )=(ne^2/k_BT)D(\omega )`$ is given by linear response theory as
$$D(\omega )=\frac{1}{2d}\omega ^2_0^{\mathrm{}}(r(t)r(0))^2e^{i\omega t}𝑑t$$
(1)
with $`d`$ the dimensionality of the transport system, $`r`$ the charge-carrier position and $``$ the configurational average. Eq. (1) reduces to Fick’s law for a frequency independent $`D`$. With increasing frequency $`\sigma ^{}`$ will increase while the positive $`ϵ^{}`$ decreases. In the delocalized case charge transport is expected to follow the scheme of a Drude electron-gas with a conductivity:
$$\sigma (\omega )=ϵ_0\omega _p^{\mathrm{\hspace{0.25em}2}}\tau \frac{1}{1+i\omega \tau }$$
(2)
with $`\omega _p^2=nq^2/(ϵ_0m^{})`$. For most metals $`\omega _p10^{15}`$ s<sup>-1</sup>. If charge transport is governed by (anomalous) diffusion (Eq.(1)), $`d\sigma ^{}/d\omega 0`$ and $`ϵ^{}1`$. In contrast, the Drude model of delocalized charge transport (Eq.(2)) predicts $`ϵ^{}<0`$ for $`\omega \omega _p`$ and $`d\sigma ^{}/d\omega 0`$ and $`dϵ^{}/d\omega 0`$ for $`\omega \tau ^1`$.
SWNT mats were prepared by vacuum-filtering a suspension of SWNT’s in water with approx. 0.5 % Triton X-100, a non-ionic surfactant, through filter paper with a pore size of 1 $`\mu `$m. The SWNT’s were produced using laser-ablation. Purification of the mats was performed by washing the filter paper with the attached SWNT mat with deionized water to remove the Triton X-100 and with methanol to remove residual NaOH. In this way, mats with a diameter of 34 mm and a thickness of typically 10 $`\mu `$m were obtained. Some of these SWNT mats were investigated with the filter paper attached to it. Other mats could be peeled off the filter paper. Up to now SWNT mats are always mixtures of chiral, zig-zag and armchaired nanotubes. The fraction of metallic nanotubes is estimated to be of the order of 0.1-0.5. For these strand-like materials such numbers are sufficiently high to be well above the percolation threshold for dc-conduction.
The dc conductivity ($`\sigma _{\mathrm{dc}}`$) of the mats was measured with the four-probe technique. The $`T`$-dependent dc measurements were performed in an Oxford flow (down to 4 K) and a <sup>3</sup>He cryostat (down to 0.4 K). Complex $`\sigma (\omega )`$ or $`ϵ(\omega )`$-data in the GHz regime were obtained by running-wave transmission-measurements with the electric-field vector parallel to the plane of the mat. The transmission and the phase shift introduced by the sample were directly measured with an ABmm millimeter-wave network analyzer. The complex dielectric constant at the given $`\omega `$ could be obtained by fitting the data to first principles formulae without the need of a Kramers-Kronig analysis. The attenuation of the samples was approximately 45 dB. For the transmission experiments at 10 GHz and at 15 GHz, resp. X-band and P-band rectangular waveguides were used. The sample was mounted onto a choke flange of the waveguides. Between 40 GHz and 500 GHz a free-space electromagnetic wave, focused with polyethylene lenses was transmitting the SWNT mat. For the $`T`$-dependent measurements at 285 GHz the SWNT mat was mounted in an optical He-flow cryostat. The mm-wave reached the sample by passing two quartz and capton windows with a diameter of 40 mm. The sample space of the cryostat was filled with liquid He for measurements below 4.2 K. The resulting change of $`ϵ`$ of the medium surrounding the sample was considered in the analysis.
In Fig. 1 the $`T`$ dependences of $`\sigma _{285\mathrm{G}\mathrm{H}\mathrm{z}}^{}`$ and $`\sigma _{\mathrm{dc}}`$, normalized to the 300 K value of $`6.9\times 10^4`$ S/m for the cleaned and $`2\times 10^4`$ S/m for the uncleaned SWNT mats are shown.
Clearly visible is the change from a positive $`d\sigma ^{}/dT`$ at low $`T`$ to a negative $`d\sigma ^{}/dT`$ at high $`T`$. The transition temperatures between the two regimes is close to $`T^{}=120`$ K in agreement with recently reported values for similar prepared SWNT mats (for the uncleaned mat $`T^{}200`$ K). The measured frequency dependences of $`ϵ^{}`$ and $`\sigma ^{}`$ show the following main features, see Fig 2: the conductivity almost keeps its dc value up to about 10 GHz and decreases at higher frequencies. Saturation is observed close to 1 THz. The dielectric constant increases from $`ϵ^{}10^4`$ in the 10 GHz-regime to $`ϵ^{}100`$ close to 1 THz, apparently approaching the regime of $`ϵ^{}>0`$ at still higher frequencies.
The dielectric constant at 285 GHz remains negative down to 2.7 K, see the temperature-dependent data in Fig. 3. For $`T>T^{}`$, $`\sigma _{285\mathrm{G}\mathrm{H}\mathrm{z}}^{}(T)`$ is proportional to $`\sigma _{\mathrm{dc}}(T)`$. For $`T<T^{}`$, $`\sigma _{285\mathrm{G}\mathrm{H}\mathrm{z}}^{}(T)`$ remains constant, while $`\sigma _{\mathrm{dc}}(T)`$ decreases with decreasing $`T`$, see Fig. 1.
To study the influence of inter-rope barriers, desiccated SWNT mats were prepared by evaporating the water after casting the suspension onto a quartz substrate. After desiccation, the remaining surfactant covers the SWNT-ropes. Fig. 4 shows the dielectric constant to be positive for such an unpurified SWNT mat. After rinsing the mat with deionized water and methanol, the charge-transport properties turned from a localized-carrier dominated regime with $`ϵ^{}>0`$, $`dϵ^{}/d\omega <0`$ and $`d\sigma ^{}/d\omega >0`$ for the unpurified sample to a delocalized-carrier dominated (metallic) regime with an increased $`\sigma _{\mathrm{dc}}`$, $`ϵ^{}<0`$, $`dϵ^{}/d\omega <0`$ and $`d\sigma ^{}/d\omega <0`$, see Fig. 4.
The peculiar $`T`$ dependence of $`\sigma _{\mathrm{dc}}`$ has been explained as a transition from metallic charge-transport at high temperatures to a non-metallic regime at low temperatures caused by charge-transport barriers of the order of some 10 meV. Such a model is expected to give a thermally activated behavior of the conductivity at very low temperatures. In our data (see Fig. 1) $`\sigma _{\mathrm{dc}}(T)`$ flattens below 10 K showing the model to be too crude. The metallic nature at room temperature is supported by the negative dielectric constant obtained from the high-frequency dielectric measurements, and $`d\sigma ^{}/d\omega 0`$ and $`dϵ^{}/d\omega 0`$ for $`\omega \tau ^1`$, see Fig. 2. A qualitative discrepancy between the observed frequency dependencies and the Drude behavior lies in the saturation of $`\sigma ^{}`$ for $`\omega >\tau ^1`$. The high-frequency conductivity remains at approx. 40 % of $`\sigma _{\mathrm{dc}}`$. This points to a background-conductivity, $`\sigma _b`$, due to localized charge carriers present in the system. A possible increase with frequency of this background (here neglected) will enhance the Drude response at the lowest frequencies at most by a factor of 2. The positive contribution to $`ϵ^{}`$ (also neglected) is usually well below $`10^3`$ in the GHz-regime and decreases with $`\omega `$.
A fit of the Drude model including background conductivity is shown in Fig. 2. At low frequencies the value of $`\sigma ^{}`$ is not accurate and the experimental error of $`ϵ^{}`$ in the 10 GHz-range is considerable. Although discrepancies in the shape of the frequency dependencies remain, the principal characteristics as described above can be reproduced. The observed more stretched frequency dependence is expected for a distribution of $`\tau `$ and $`\omega _p`$. Given the experimental inaccuracies, the order of magnitude of the fit parameters is correct. The scattering time is estimated as $`\tau =25\times 10^{12}`$ s, the plasma frequency as $`\omega _p=2.55.5\times 10^{13}`$ s<sup>-1</sup> and the background conductivity as $`\sigma _b=23\times 10^4`$ S/m.
Using the Fermi velocity of graphite, $`v_F=8\times 10^5`$ m/s, the scattering time $`\tau `$ gives a mean free path of $`\mathrm{\Lambda }3`$ $`\mu `$m. A similar value has been estimated from ESR measurements at 100 K on SWNT mats. From DC measurements on an isolated SWNT at a few mK $`\mathrm{\Lambda }=3`$ $`\mu `$m is suggested as a lower limit. Although theoretical arguments predict a very low scattering probability with acoustic phonons inside a tube, which might allow similar values also at higher temperatures, the $`\mathrm{\Lambda }`$ found here refers to 3D transport, for which such a value seems (far) too high. We will return to this problem below.
The obtained plasma frequency is about one percent of $`\omega _p`$ for normal metals. Assuming $`m^{}=m_e`$, the value of $`\omega _p`$ implies a charge-carrier density $`n=39\times 10^{23}`$ m<sup>-3</sup>. Correcting for the lower density of the mats ($`0.65`$ g/cm<sup>3</sup>) compared to that of a SWNT ($`2`$ g/cm<sup>3</sup>) and assuming a fraction of $`50\%`$ metallic tubes would give $`n4\times 10^{24}`$ m<sup>-3</sup>, comparable to graphite. However, for SWNT $`n`$ is predicted to be about $`10^2\times `$ higher. Fischer et al. found after mass-density corrections a more than $`10\times `$ higher $`\sigma _{\mathrm{dc}}`$ for single ropes than for mats. It is plausible that in a mat only a fraction of the charge carriers present participate in the delocalized (metallic) charge transport, while due to localization the remaining charge carriers have a smaller contribution to $`\sigma _{\mathrm{dc}}`$. In the model proposed here localized charge carriers are incorporated and contribute to $`\sigma _b`$.
At $`T=300`$ K $`\sigma _{285\mathrm{G}\mathrm{H}\mathrm{z}}^{}`$ is mainly due to $`\sigma _b`$, i.e. localized charge carriers (see fit in Fig.2). At such a high frequency $`\sigma ^{}(\omega )`$ might well be determined by photon- in stead of phonon-assisted hopping, which explains the constant value of $`\sigma _{285\mathrm{G}\mathrm{H}\mathrm{z}}^{}(T)`$ for $`T10^2`$ K . However, $`ϵ^{}`$ is still dominated by the delocalized charges over the whole temperature range ($`ϵ^{}<0`$ at 285 GHz, Fig. 3). The decrease by 40 % of $`|ϵ^{}|`$ between 70 K and 2.7 K can be accounted for by the decrease of carriers from the semiconducting tubes and a growing contribution of localized states. It shows that the metallic part of $`ϵ^{}(\omega )`$ for $`\omega >1/\tau `$ has no strong $`T`$ dependence, which implies an almost $`T`$ independence of $`\omega _p`$. Also $`\sigma _{\mathrm{dc}}`$ at 0.4 K is only a factor 0.7 lower than at 300 K, which indicates that not only $`\omega _p`$ but also $`\tau `$ have not changed appreciably with $`T`$. An opening of a gap due to twistons in the order of 20 meV as suggested by Kane and Mele seems therefore unlikely. Based on the measured metallic low-temperature behavior down to 4.2 K of the thermoelectric power in SWNT ropes, Hone et al. also excluded the opening of a gap at low $`T`$.
The room temperature data presented in Fig. 4 confirm the importance of inter-rope contacts at low temperatures. The purification procedure removed the surfactant and other impurities from the surface of the SWNT ropes allowing better contacts between ropes. Intra-tube or intra-rope transport are likely not changed by the purification, meaning that charge-localization effects due to defects in the graphene-sheet pattern of the tubes or bending of them should be unaffected. The effect of the purification on $`\sigma _{\mathrm{dc}}(T)`$, see Fig. 1, supports the picture. The transition temperature $`T^{}`$ for the uncleaned sample is higher than for the cleaned one. Also, $`d\sigma _{\mathrm{dc}}/dT`$ below $`T^{}`$ is larger for the latter. Both indicate that inter-rope barriers limit $`\sigma _{\mathrm{dc}}`$ at low temperatures. These findings are consistent with the higher $`\sigma _{\mathrm{dc}}`$, lower $`T^{}`$ and weaker $`d\sigma _{\mathrm{dc}}/dT`$ below $`T^{}`$ of single rope data, where inter-rope barriers are eliminated.
Highly conducting polymers like doped polyaniline (PAN) and polypyrrole (PPy) and SWNT-mats show analogous dielectric behavior. In the metallic polymers $`\sigma _{\mathrm{dc}}`$ typically has a maximum value of order $`10^4`$ S/m (around 200 K) and decreases to lower temperatures. The values of $`ϵ^{}(\omega )`$ are strongly frequency and $`T`$ dependent. Let us take one of the best conducting materials, PAN doped with $`d`$,1-camphorsulfonic acid (PAN-CSA), as an example. Around 1 meV ($`\omega `$ 1.5 THz) at 200 K $`ϵ^{}(\omega )`$ is a few times $`10^3`$, and becomes less negative at lower $`T`$. For the same samples at room temperature $`ϵ^{}(\omega )`$ starts negative, becomes positive around 30 meV, returns negative around 0.1 eV and finally comes close to zero in the optical regime. In the SWNT-mats the maximum value of $`\sigma _{\mathrm{dc}}`$ is almost $`10^5`$ S/m, and $`ϵ`$(285 GHz) reaches a value of $`10^3`$ and decreases in absolute value with decreasing temperature (by a factor of two at 4.2 K). At room temperature below 0.5 meV $`ϵ^{}(\omega )`$ is negative, likely becomes positive at higher energies and returns negative again around 10 meV. The comparison shows that best conducting polymers are essentially behaving like well-rinsed mats of single wall nanotubes (the same similarity exists between not-rinsed mats and the slightly less conducting polymers). Like in the mats, in these polymers values for $`\mathrm{\Lambda }`$ will be of the order of 100 nm for $`v_F=5\times 10^5`$ m/s. In the metallic polymers homogeneous and inhomogeneous disorder models are frequently used to explain these extreme values. For the nanotubes crystalline regions are excluded (TEM pictures show the nanotube mats to be completely entangled and disordered), which at least for the SWNT-mats requires an alternative for the inhomogeneous disorder model.
In summary, we have shown that the delocalized properties of nanotube mats can be completely determined by sub-THz measurements and are well described in a Drude picture with a negligible temperature dependence of the plasma frequency and scattering time. The values of $`\omega _p`$ and $`\tau `$ resemble those found in well conducting doped polymers. When modeled with the usually chosen Fermi velocities, unusually large values of the mean free path result for both systems. This finding and the absence of crystalline regions in the mats underline the need for a better description of these disordered strand-like systems.
Acknowledgements. We like to acknowledge Hubert Martens for enlightening discussions and critical reading of the manuscript and Roel Smit for support in the experiments with the <sup>3</sup>He cryostat. This investigation is part of the research program of FOM-PPM with financial support from NWO. M. Ahlskog was supported by the Academy of Finland. |
warning/0001/astro-ph0001203.html | ar5iv | text | # Disc instability models for X-ray transients: evidence for evaporation and low limit-from𝛼-viscosity ?
## 1 Introduction
Low Mass Black Hole Binaries (LMBHBs) and Low Mass Neutron Stars Binaries (LMNSBs) are close binary systems in which a Roche-lobe filling, low mass (main sequence or sub-giant) secondary star transfers mass to a compact primary. All known LMBHBs and several LMNSBs are transients (see e.g. Tanaka & Lewin 1995, van Paradijs & McClintock 1995, White, Nagase & Parmar 1995); these transient systems are known as soft X-ray transients (SXTs), or X-ray novae. Their outbursts last from several weeks to several months during which X–ray luminosities can approach the Eddington limit. The outbursts are separated by long periods of quiescence lasting from one to tens of years (see e.g. Chen, Shrader & Livio 1997; hereafter CSL).
SXTs are in some respects similar to dwarf novae (DN), which exhibit outbursts of 4 – 6 magnitudes in optical. These outbursts typically last days and are separated by periods of quiescence of a few weeks. DN are a subclass of Cataclysmic Variables (CVs; see Warner 1995 for a review), close binary systems similar to LMBHBs and LMNSBs, except for the compact primary which is a white dwarf.
Is is widely believed that dwarf nova outbursts are due to a thermal-viscous instability in the accretion disc of these systems. For effective temperatures between about 5000 K and 8000 K, accretion discs in which viscosity is described by the “$`\alpha `$prescription” (Shakura & Sunyaev 1973) are thermally and viscously unstable because of large changes in the opacity when hydrogen recombines. The disc instability model (DIM; see Cannizzo 1993b for a review; Hameury et al. 1998 and Menou, Hameury & Stehle 1999a for the most recent version of the model) accounts for a number of general properties of dwarf novae. Similarities between the properties of DN and SXTs led to the suggestion that their outbursts have the same origin (van Paradijs & Verbunt 1984; Cannizzo, Ghosh & Wheeler 1982; Huang & Wheeler 1989; Mineshige & Wheeler 1989).
Several important features of DN remain unexplained by the standard version of the DIM, which assumes that the disc extends down to the white dwarf’s surface, that the mass transfer rate from the secondary star is constant, and that effects of disc and secondary irradiation can be neglected (see e.g. Lasota & Hameury 1998; Lasota 1999; Warner 1998; Smak 1999, Hameury. Lasota & Warner 1999b for a discussion of these effects). This version cannot account for the quiescent accretion rates onto the white dwarf (they exceed the model predictions by more than two orders of magnitude), and for the very long recurrence times of WZ Sge–type systems (Lasota 1996a,b). The observed delay between the rise in optical and in EUV (e.g. in SS Cyg; Mauche 1996) was also long considered as a problem, but Smak (1998) recently argued that it can be naturally reproduced by the DIM if appropriate outer boundary conditions for the disc are used.
Similar difficulties, but more pronounced, appear for SXTs, since the standard DIM fails to reproduce a number of essential properties of these systems. In particular, an early attempt to explain the properties of BH SXT outbursts with the DIM (Mineshige & Wheeler 1989) was unable to reproduce the observed long recurrence times (comparable to those of WZ Sge systems) and the high accretion rates required in outburst. In addition, the weak X-ray flux detected from quiescent BH SXTs, which was mistakenly thought by some to be an argument in favour of the DIM versus the competing mass transfer instability model, is many orders of magnitude larger than any flux predicted by the DIM \[1996a\].
In the DN case, the problems can been solved, at least in part, if a “hole” is present in the inner regions of the disc. The “hole” can either be caused by evaporation of the disc , or due to the extended magnetospheric cavity of the WD . A model in which a “hole” is included can account for both the EUV delay and the quiescent X-ray luminosity (if the hole is filled with a tenuous, X–ray emitting gas; Meyer & Meyer-Hofmeister 1994; Hameury, Lasota & Dubus 1999a).
The case of WZ Sge requires additional modifications of the standard DIM. The very long recurrence time of this system, and the large amount of mass accreted during the outburst, can be reproduced either by using a value of the viscosity parameter $`\alpha _{\mathrm{cold}}`$ several orders of magnitude lower than in other DN (Smak 1993; Osaki 1995), or by making the “hole” sufficiently large for the remaining disc to be globally stable. In the latter case, the outbursts would have to be triggered by enhanced mass transfer (Lasota, Hameury & Huré 1995; Hameury, Lasota & Huré 1997b). A marginally stable, truncated disc, as proposed by Warner, Livio & Tout , is equivalent to a globally stable disc from this point of view \[1997b\]. If the truncated disc is globally stable, and $`\alpha _{\mathrm{cold}}0.02`$, however, there is not enough mass in the disc to account for the outburst’s total energy, so that enhanced mass transfer is also required during the outburst. Hameury et al. \[1997b\] find that if this enhanced mass transfer is due to irradiation of the secondary, the model reproduces the shape of the observed lightcurve.
The presence of “holes” in the accretion discs of quiescent SXTs seems unavoidable if their outbursts are due to the thermal-viscous instability (Lasota 1996a,b). In the BH case, such holes can only result from the evaporation of the accretion disc since black holes are not magnetized. Narayan, McClintock & Yi (1996) proposed a model in which, as a result of evaporation, the inner parts of the accretion flow in quiescent BH SXTs form an advection-dominated accretion flow (ADAF; see also Lasota, Narayan & Yi 1996; Narayan, Barret & McClintock 1997). The model explains the observed spectra and solves the problem of incompatibility between the accretion rates predicted by the DIM and the observed flux.
Esin, McClintock & Narayan (1997; see also Narayan 1996; Esin et al. 1998) extended the idea of two-component accretion flows a step further. They showed that the various spectral states of black hole X-ray binaries can be modeled as a sequence of physical states with varying accretion rates and sizes of the two components of the accretion flow. According to this model, in outburst, the disc extends down to the black hole, while the inner disc is gradually replaced by an ADAF as a system goes to lower luminosity levels. A gradual evaporation of the inner disc during the decline is also qualitatively consistent with the accretion geometry inferred by Życki, Done & Smith (1999) from the X-ray reprocessing properties of several BH X-ray transients.
Hameury et al. (1997a) used the two component (inner ADAF + outer disc) model to describe the quiescent state and the rise to outburst of the BH SXT GRO J1655-40. The model predictions agree with existing constraints on the spectrum of the system in quiescence, the observed 6 day delay between the rise to outburst of the optical and X-ray fluxes, and the rise-times of these two fluxes. Recently, Esin, Lasota & Hynes (1999) showed that the whole, unusual outburst can be described in the framework of the DIM thus validating Hameury et al. (1997a) results.
In this paper, we consider the role of disc evaporation on the predictions of the DIM for SXT outbursts. We shall refer to a DIM that includes evaporation as a Truncated–Disc Instability Model (hereafter TDIM). Cannizzo (1998) already constructed TDIMs of the BH SXT A0620-00 with a weak evaporation and the prescription $`\alpha =50(H/R)^{3/2}`$ for the viscosity parameter. Cannizzo’s simulations show that models without evaporation lead to unreasonably short recurrence times, and that evaporation acts to increase the predicted recurrence times to values of several tens of years, as required. This is expected, since evaporation removes the inner, most unstable parts of the accretion disc. However, the outburst rise time predicted by Cannizzo’s model is an order of magnitude longer than the observed one. This is due to the use of the prescription $`\alpha =50(H/R)^{3/2}`$ according to Lasota & Hameury (1998).
There are two additional reasons to go beyond Cannizzo’s calculations. First, his calculations do not allow the disc outer radius to vary with time. Smak (1984, 1998; see also Hameury et al. 1998) showed that including this effect is crucial for reliable DIM predictions of DN outbursts, and it is likely to be important for SXT models as well. Second, in the models of Cannizzo (1998), the effect of including evaporation and using a very small value of $`\alpha `$ in the quiescent disc are not clearly separated. In particular, in his work, it is unclear if evaporation alone can be held responsible for the long recurrence times, if values of $`\alpha `$ similar to those inferred from applications of the DIM to DN (typically $`\alpha _{\mathrm{hot}}10^1`$ and $`\alpha _{\mathrm{cold}}0.02`$) are used for the discs of SXTs. In our view, this is an important question because the evidence for disc evaporation, even indirect, is stronger than for any specific $`\alpha `$prescription.
In § 2, we recall some important observational properties of BH and NS SXTs, later used to constrain our models. In § 3, we first describe the numerical truncated disc instability model used. We then show that, in the limit where illumination can be neglected, TDIMs with values of $`\alpha `$ considered as standard for DN discs are unable to reproduce the long recurrence times of SXTs, no matter how strong the evaporation is. In § 4, we present other plausible models of SXTs. We find that a TDIM with a small value of $`\alpha _{\mathrm{cold}}`$ ($`5\times 10^3`$) and strong evaporation does reproduce the long recurrence times and the high maximum values of the accretion rate in the discs of SXTs. We also present models in which the disc of SXTs are globally stable and the outbursts are triggered by a slow variation of the mass transfer rate of the secondary. In § 5, we discuss important consequences, limitations and possible extensions of this work. In § 6, we summarize our main results.
## 2 Main observational characteristics of SXTs
BH and NS X-ray Transients show a rather complex behaviour (see, e.g., Tanaka & Shibazaki 1996 and CSL for reviews). In this section, we identify several important properties of these systems which may be considered as characteristic of their class and therefore must be explained by a successful model of SXT outbursts. We will also often compare SXTs with DN, since the same underlying thermal-viscous instability is supposed to be responsible for the outbursts in the two classes of systems.
As a guide, we use the observed properties of SS Cyg, Aql X-1 and A0620-00 as indicative of the general properties of U Gem-type DN, NS SXTs and BH SXTs respectively, keeping in mind that significant differences exist even between two outbursts in a same system (e.g. at least four types of outbursts in SS Cyg; Warner 1995).
### 2.1 Timing properties
The timescales of evolution of an unsteady disc are simple and robust predictions of the DIM, which provide powerful tests of the model when compared to the observations. This is because the time evolution of the disc is due to the existence of a well defined limit cycle in the DIM, and because the observed variability is ‘ready-to-interpret’ (free of noise, systematic errors, etc..).
Most of the properties of BH and NS SXTs mentioned in this section are taken from the compilation of CSL. We identify several important timescales for X-ray transients: the recurrence time $`t_{\mathrm{rec}}`$ between two successive outbursts, the e-folding rise timescale $`\tau _{\mathrm{rise}}`$ and the total rise time $`t_{\mathrm{rise}}`$ of an outburst (which can be wavelength-dependent), the e-folding decay timescale $`\tau _{\mathrm{decay}}`$ (which is wavelength dependent), and the total duration of the outburst $`t_{\mathrm{dur}}`$. In our view, a model which fails to reproduce all these characteristic times is not fully satisfactory. None of the previous models of SXTs passed this test (e.g. Mineshige & Wheeler 1989; Cannizzo 1998).
We focus our discussion on a sample of SXTs with clearly identified primaries (see Garcia et al. 1998): the BH systems have firm dynamical lower limits on the mass of the primary, and the NS SXTs showed type I X-ray bursts. We also focus on systems which have exhibited standard Fast Rise Exponential Decay (FRED) lightcurves. This excludes the long period systems GRO J1655-40 and V 404 Cyg and reduces the sample to four BH SXTs (A0620-00, GS 2000+25, GRS 1124-683=Nova Mus 91 and GRO J0422+32 – “best examples of FREDs” according to CSL) and two NS SXTs (Cen X-4 and Aql X-1; here we chose only their outbursts of the FRED type : it is not clear how the standard DIM could account for the other types of lightcurves observed in SXTs).
#### 2.1.1 Recurrence times
Aql X-1 experiences an outburst every year approximately, while the recurrence time of Cen X-4 is $`10`$ years. A0620-00 experienced two outbursts separated by $`58`$ years (note that no X-ray data are available for the presumed first outburst). Most BH SXTs did not experience a second outburst since their discovery by X-ray satellites, so that there is good evidence for the recurrence times of BH SXTs generally exceeding $`2030`$ years. As mentioned above, SXT recurrence times are similar to those of WZ Sge-type DN. For comparison, the recurrence time of SS Cyg is $`t_{\mathrm{rec}}50`$ days (e.g. Cannizzo & Mattei 1992) and the recurrence time of other U Gem type DN is generally of the order of one month (e.g. Szkody & Mattei 1984).
#### 2.1.2 Rise times
The characteristic X-ray rise timescales ($`\tau _{\mathrm{rise}}`$) of FRED-type SXTs vary from 0.3 to 2 days (for the period of fastest flux increase – CSL). The total rise time $`t_{\mathrm{rise}}`$ from quiescence to outburst peak is difficult to estimate because it is instrument dependent, and because the coverage of the rising phase is, for obvious reasons, rather poor. For SXTs, $`t_{\mathrm{rise}}`$ is $``$ 5 to 10 days in X–rays, or perhaps longer. For comparison, the total optical rise time in SS Cyg is $`t_{\mathrm{rise}}`$ 3 days (for a ‘type B’ outburst - see below). The optical rise time of several U Gem-type DN in the sample of Szkody & Mattei (1984) is $`t_{\mathrm{rise}}2`$ days.
In several DN, there is evidence for at least two types of outbursts. ‘Type A’ outbursts correspond to FREDs and are usually interpreted as ‘outside-in’ outbursts, i.e. outbursts triggered in the outer regions of the disc (Smak 1984). This interpretation seems, in a few cases, consistent with the observed delay between the rise to outburst of optical and EUV (or UV) fluxes (Warner 1995). ‘Type B’ outbursts have a more symmetric shape with a slower rise phase, and are usually interpreted as ‘inside-out’ outbursts (Smak 1984; see also Menou et al. 1999a).
In the case of SXTs, there is no observational evidence for outside-in outbursts. The delay between the rise of the optical and X-ray fluxes (observed only during a non-FRED outburst of the BH SXT GRO J1655-40) is not an indication of an outside-in outburst if the disc is truncated. Indeed, Hameury et al. (1997a) obtain correct rise times and the X-ray/optical delay for an inside-out outburst. In this case the ‘inside’ corresponds to the inner edge of the truncated disc, far from the central object.
#### 2.1.3 Decay times
The X-ray decay timescales $`\tau _{\mathrm{dec}}`$ of SXTs with FRED light-curves are typically 25 - 40 days (CSL). The optical decay timescales of BH SXTs are $``$ 50 to 200 days, while the optical decay timescales of NS SXTs are $``$ 10 to 30 days (but the available information is poor). We note that long decay timescales are also observed in WZ Sge–type DN (Kuulkers, Howell & van Paradijs 1996; Kuulkers 1998).
#### 2.1.4 Duration times
The reported duration times for the FRED–type outbursts of BH SXTs range from 170 to 260 days in X-rays (and even longer in optical), while $`t_{\mathrm{dur}}70`$ days in X-rays for FRED–type outbursts of NS SXTs. The typical duration time of U Gem–type DN in the sample of Szkody & Mattei (1984) is $`t_{\mathrm{dur}}10`$ days (optical).
One has to remember, however, that the reported durations for SXTs depend on the detection threshold of the instrument and that the coverage of some outbursts is incomplete; sometimes the reported duration is only a lower limit on the true duration. In some cases, the decay timescale $`\tau _{\mathrm{dec}}`$ may be more appropriate than the duration time $`t_{\mathrm{dec}}`$ for comparison of the model predictions with the observations.
### 2.2 Other important observational constraints
An estimate of the maximum luminosity reached in outburst by a SXT depends on the distance to the system and, when compared to the Eddington luminosity, on the mass of the primary. According to Garcia et al. (1998), the NS SXTs considered here approximately reached the Eddington luminosity at outburst peak, as did some of the BH SXTs. According to CSL, the maximum luminosities reached are closer to 10% of the Eddington luminosity, but their values could be underestimated. For example, CSL obtain for Nova Muscae 1991 a peak luminosity (0.4 – 10 keV) of 0.08 in Eddington units, whereas according to Esin et al. (1997), who used the entire available spectral information, the peak luminosity in this system was almost exactly Eddington. In our view, this means that models of SXTs have to be able to reach the Eddington accretion rate. The Eddington accretion rate is defined here as
$$\dot{M}_{\mathrm{Edd}}=1.39\times 10^{18}\left(\frac{M_1}{M_{}}\right)\mathrm{g}\mathrm{s}^1,$$
(1)
the rate at which an accretion flow around a central object of mass $`M_1`$ reaches the Eddington luminosity for a standard $`10\%`$ radiative efficiency. The Eddington luminosity, however, seems to be the maximum luminosity for SXTs. Combined with outburst durations this implies that the mass accreted during outbursts by the central body is $``$ a few $`10^{24}`$ g (roughly the same as accreted onto the white dwarf during WZ Sge eruptions). This constraint must be satisfied by SXT models as well.
Observations of broadened $`H_\alpha `$ emission lines provide an upper limit on the value of the disc inner radius $`R_{\mathrm{in}}`$ in quiescence for several BH SXTs (e.g. Orosz et al. 1994; Narayan et al. 1996; Menou, Narayan & Lasota 1999b). In most systems, the constraints indicate a value of $`R_{\mathrm{in}}`$ less than $`10^4`$ to a few $`10^4R_S`$, where $`R_s=2.95\times 10^5(M_1/M_{})`$ cm is the Schwarzschild radius of the black hole primary of mass $`M_1`$. The limits provide a useful constraint on the strength of evaporation to be included in the TDIMs of SXTs (see below).
For some systems, other constraints on $`R_{\mathrm{in}}`$ exist at various stages of the outburst from the modeling of X-ray reflection properties (Życki et al. 1998), as well as spectral fits which involve $`R_{\mathrm{in}}`$ as a parameter (e.g. Sobczak et al. 1999). However, these two methods are model dependent and somewhat controversial, because the relevant spectral features can be modeled differently (M. Nowak, private communication). We choose to ignore them here.
## 3 “Standard” truncated disc instability models of SXTs
In this section, we investigate a “standard” version of the TDIM for SXTs. As mentioned above, by “standard”, we mean a model with values of the viscosity parameter $`\alpha `$ that are usually considered as standard for DN discs: $`\alpha _{\mathrm{hot}}0.1`$ and $`\alpha _{\mathrm{cold}}0.02`$. Non-standard models, in particular models using significantly smaller values of $`\alpha _{\mathrm{cold}}`$, are considered in § 4.
### 3.1 Numerical specifications
#### 3.1.1 General
We use the numerical code described in Hameury et al. (1998) to simulate the evolution of discs around NSs and BHs. The equations of conservation of mass, angular momentum and energy for a thin Keplerian disc are solved on an adaptive grid which resolves narrow structures in the disc (transition fronts in particular). A fully implicit numerical scheme is used for time evolution, to avoid any Courant-type limiting condition on the timestep of integration.
A grid of disc vertical structures, which give the local cooling rate of the disc as a function of its surface density, central temperature and the vertical gravity is precalculated before running the disc evolution. The effect of changing the mass of the central object (from a NS SXT to a BH SXT) enters the time-dependent calculations both through the radial conservation equations solved, and through the grid of vertical structures with different gravities.
The outer radius of the disc is allowed to vary with time under the action of the torque applied by the secondary star on the outermost regions of the disc. Including this proper boundary condition (as opposed to a rigid boundary condition with a fixed value of the outer radius) has been found to drastically affect the predictions of the DIM for DN discs (see Hameury et al. 1998 for details). The effect of including this proper boundary condition for discs of SXTs is discussed in § 3.2.1.
#### 3.1.2 Formation of a hole in the disc
The formation of a hole in the discs of SXTs is probably determined by the physics of the evaporation of gas from the disc, possibly into an ADAF. This physics is poorly understood. Several mechanisms have been proposed in the literature to explain the transition (see, e.g., Meyer & Meyer-Hofmeister 1994; Liu, Meyer, Meyer-Hofmeister 1997; Narayan & Yi 1995; Honma 1996; Shaviv, Wickramasinghe & Wehrse 1999), but none of these models has reached a sufficient level of detail to allow robust quantitative predictions.
Disc evaporation has already been included in the calculations of Hameury et al. (1997a) and Hameury et al. (1999a). In these studies, a ‘reasonable’ formula was used (close to the one proposed by Meyer & Meyer-Hofmeister 1994) and evaporation was included in the numerical models as a sink of matter in the mass conservation equation (see also the models with evaporation of Cannizzo 1998).
Here, we chose a slightly different approach. We assume an ad hoc law for the rate of evaporation as a function of radius, as in Hameury et al. (1997a; 1999a), but in addition we assume that evaporation acts only in a very narrow region at the inner edge of the disc. This is a good approximation if the rate of evaporation of the disc material is a steep function of radius.
To be more specific, the inner radius of the disc, $`R_{\mathrm{in}}`$, varies according to:
$$\dot{M}_{\mathrm{disc}}(R_{\mathrm{in}})=\dot{M}_{\mathrm{evap}}(R_{\mathrm{in}}),$$
(2)
where $`\dot{M}_{\mathrm{disc}}(R_{\mathrm{in}})`$ is the time-dependent accretion rate at the inner edge $`R_{\mathrm{in}}`$ of the disc, and $`\dot{M}_{\mathrm{evap}}(R)`$ is the evaporation law with an arbitrary functional form. This evaporation law is chosen so that, in outburst, $`\dot{M}_{\mathrm{evap}}(R)`$ is much less than the accretion rate in the disc at any radius, and the disc is fully extended. The radius of truncation of the disc in quiescence depends on the functional form of $`\dot{M}_{\mathrm{evap}}(R)`$ and the profile of $`\dot{M}_{\mathrm{disc}}`$ in the disc. We find that this numerical recipe is much more stable than that used in Hameury et al. (1997a), while it gives very similar results for the evolution of $`R_{\mathrm{in}}`$ and the outburst cycles in general.
Following Hameury et al. (1997a), we could have assumed that the disc is evaporated into an ADAF which accretes everywhere at the maximum rate allowed, $`\dot{M}_{\mathrm{crit}}^{\mathrm{ADAF}}(R)`$, and take for the evaporation law $`\dot{M}_{\mathrm{evap}}(R)=\dot{M}_{\mathrm{crit}}^{\mathrm{ADAF}}(R)`$. Recently, however, the estimates of $`\dot{M}_{\mathrm{crit}}^{\mathrm{ADAF}}(R)`$ have been revised to higher values (Esin et al. 1997; Menou et al. 1999b). Using these recent estimates for the evaporation law (i.e. a stronger evaporation than, e.g., in Hameury et al. 1997a) would imply that truncated discs are globally stable (and stationary) for all mass transfer rates below $`10^{17}`$ g s<sup>-1</sup> (this is demonstrated by the Fig. 5 of Menou et al. 1999b), whereas typical SXT mass–transfer rates are much lower: $`10^{15}10^{16}`$ g s<sup>-1</sup> (van Paradijs 1996; CSL; Menou et al. 1999b). In addition, for mass transfer rates $`10^{17}`$ g s<sup>-1</sup>, the stable discs would be truncated at radii $`3\times 10^4`$ $`R_s`$, which is inconsistent with the upper limits derived from $`H_\alpha `$ emission lines (§ 2.2).
To circumvent these difficulties, and ensure that discs in our models experience outbursts even at lower mass transfer rates, $`\dot{M}_{\mathrm{evap}}(R)`$ must be less than $`\dot{M}_{\mathrm{crit}}^{\mathrm{ADAF}}(R)`$. We use the following evaporation law:
$$\dot{M}_{\mathrm{evap}}(R)=\frac{0.08\dot{M}_{\mathrm{Edd}}}{\left(\frac{R}{R_s}\right)^{1/4}+\left(\frac{R}{800R_s}\right)^2}.$$
(3)
where $``$ is a constant ensuring that a maximally truncated disc is still unstable. The power laws and scalings in Eq. (3) are such that the profile of $`\dot{M}_{\mathrm{evap}}(R)`$ is initially slowly decreasing with radius, and then decreases much more steeply beyond a radius $`800R_s\times ^{1/2}`$. In all the models with evaporation described below, we take $`=20`$ (except in section 4.1, in which $`=30`$). This corresponds to a situation where the evaporation rate during quiescence is significantly smaller than the mass transfer rate from the secondary, thus ensuring that the disc is not marginally unstable.
This expression for the evaporation law is ad hoc but possesses several advantages. First, it is a continuously decreasing function of radius, which seems quite natural, and is appropriate for the numerical implementation of evaporation described above. Second, the maximum value of $`\dot{M}_{\mathrm{evap}}(R)`$, as defined in Eq. (3) when $`RR_S`$, is consistent with the high value required in the ADAF models of Esin et al. (1997) to fit the high luminosity states of the BH SXT Nova Muscae 91. Third, it leads to disc truncation radii in quiescence which are smaller than the upper limits inferred from the broadening of $`H_\alpha `$ emission lines. The resulting discs are unsteady and experience limit cycles for reasonable mass transfer rates ($`10^{15}`$ g s<sup>-1</sup>), as required for the DIM to be the explanation of SXT outbursts.
Naturally, one could worry that the predictions of the time-dependent models depend on the evaporation law used. This is true about the strength of evaporation, but we find that our conclusions do not crucially depend on the precise shape of the evaporation law. We discuss the effects of changing the evaporation law on the model predictions in § 5.
Note that in the models presented here, the presence of a tenuous region in the inner accretion flow, presumably an ADAF, is only taken into account through disc evaporation (i.e. as a sink of mass): we neglect any other physical effect that this inner region of the accretion flow could have on the disc evolution. Disc irradiation by the central X-ray source during high luminosity states is also neglected, although it is a potentially important effect (see § 5 for a discussion). For simplicity, we use the same evaporation law for BH and NS SXTs, except that $`\dot{M}_{\mathrm{Edd}}`$ and $`R_S`$ are scaled with the actual primary mass $`M_1`$.
### 3.2 Results
We present in this section two SXT models, one with parameters relevant for the BH SXT A0620-00, and the other for the NS SXT Aql X-1. In both models, $`\alpha _{\mathrm{hot}}=0.1`$, $`\alpha _{\mathrm{cold}}=0.02`$, the disc inner radius is limited during the evolution to the minimum value $`R_{\mathrm{min}}=5\times 10^8`$ cm (see § 5 for a discussion of this numerical limitation) and Eq. (3), scaled with $`M_1`$, is used for the evaporation law.
The model of A0620-00 has the following parameters: primary mass $`M_1=6M_{}`$, time-averaged value of the disc outer radius $`<R_{\mathrm{out}}>=1.24\times 10^{11}`$ cm and mass transfer rate $`\dot{M}_T=3\times 10^{15}`$ g s<sup>-1</sup>. The model of Aql X-1 has the following parameters: primary mass $`M_1=1.4M_{}`$, time-averaged value of the disc outer radius $`<R_{\mathrm{out}}>=1.8\times 10^{11}`$ cm and mass transfer rate $`\dot{M}_T=2\times 10^{16}`$ g s<sup>-1</sup>. The main properties of these models are summarized in Figures 1 to 5, and are discussed in more detail below.
For comparison, we also show in Fig. 3 and 4 results from a model with parameters relevant for A0620-00, but without evaporation. Results are similar, except for a significantly shorter recurrence time, a larger number of reflares, and the existence of two type of outbursts. Successive small and large outbursts are also predicted in models of dwarf novae (see e.g. Cannizzo 1993a; Menou et al. 1999a). During small outbursts, the outward propagating heating front does not reach the outer disc radius. In the DIM, the inner regions of quiescent accretion discs are more subject to instabilities than the outer ones, so that the truncation of the inner disc has a stabilizing effect. This explains the absence of ‘small’ outbursts and the reduced number of reflares in models that include evaporation (see below in § 3.2.6).
The panel representing the total mass in Fig. 1 shows that the disc evolution is periodic and therefore relaxed (i.e. does not depend on the arbitrary properties of the initial disc used at the beginning of the simulation). This is true of every model discussed in this paper. Strictly periodic outburst cycles are produced in SXT models, like for DN, while observed cycles show, in general, only some regularity.
In the models presented here (as well as in all the other sections), the outbursts are triggered in the inner regions of the truncated disc (inside-out type). Such outbursts never start exactly at the disc inner edge, so that there is always a heating front traveling some distance inwards. Its effect is seen as a ‘spike’ in the light-curve in Fig. 4, which lasts for less than one minute and is therefore not observable. In a truncated disc, these short-lived fronts are absent, because the innermost regions where they would propagate no longer exist. Heating fronts only marginally reach the outer edge of the disc in the model with parameters relevant for A 0620-00. In the model with parameters relevant for Aql X-1, the disc is so extended that the heating fronts never go beyond one third of the disc outer radius. This would of course not be the case for shorter period systems such as GRO J0422+32.
In the following discussion of the rise times and decay times predicted by the models, we will assume for simplicity that the X-ray emission from the system is proportional to the accretion rate $`\dot{M}_{\mathrm{in}}`$ at the inner edge of the disc. The real situation is likely to be more complicated, in particular if some of the X-rays observed are emitted by an ADAF whose luminosity is not simply proportional to $`\dot{M}`$ (in that case, the timescales could be even shorter because the radiative efficiency increases with $`\dot{M}`$). Similarly, the V magnitude shown are not supposed to directly represent the V magnitude observable from a system but only the intrinsic optical luminosity variations of the disc. In observed systems, the V band magnitude emission has also contributions from the secondary, from reprocessing of X-ray photons by the disc and possibly from an ADAF. We do not attempt to describe all these processes here, and we concentrate on the basic properties of the models.
#### 3.2.1 Recurrence times
As for DN discs, it is important to include a correct boundary condition at the outer edge of the disc of SXTs. We find that the recurrence times predicted by models in which $`R_{\mathrm{out}}`$ is allowed to vary with time are typically a factor of a few (say 2-3) times shorter than those predicted if $`R_{out}`$ is kept fixed. Indeed, when a heating front reaches $`R_{\mathrm{out}}`$, mass is spread over a more extended region if $`R_{\mathrm{out}}`$ is allowed to vary. Consequently the critical surface density $`\mathrm{\Sigma }_{\mathrm{min}}`$ below which a cooling front appears in the disc is reached by the disc earlier, less mass is accreted during the outburst, and overall it takes less time for the disc to reach the state of critical surface density $`\mathrm{\Sigma }_{\mathrm{max}}`$ at which the next outburst is triggered (Hameury et al. 1998). Therefore, all existing models of BH SXTs with fixed values of $`R_{\mathrm{out}}`$ have systematically overestimated the recurrence times $`t_{\mathrm{rec}}`$ by a factor of a few when the heating front was able to reach the disc’s outer edge. In the present paper, we only consider models with varying $`R_{\mathrm{out}}`$.
The recurrence time in the model of A0620-00 is $`t_{\mathrm{rec}}5`$ years, i.e. shorter than the known recurrence time by more than one order of magnitude (§ 2.1.1). The recurrence time in the model of Aql X-1 is $`t_{\mathrm{rec}}1`$ year, in very good agreement with the know recurrence time.
#### 3.2.2 quiescence
During quiescence, all of the accreted mass is processed through the ADAF and is therefore evaporated, so that the mass accretion rate shown in figs. 1 and 5 is the actual evaporation rate except when the disc radius reaches its minimum value during outbursts. As can be seen, the evaporation rate remains significantly smaller than the mass transfer rate from the secondary, varying between 0.15 and 0.60 $`\dot{M}_\mathrm{T}`$ in the case of the A0620-00 model, and between 0.10 and 0.15 $`\dot{M}_\mathrm{T}`$ in the Aql X-1 case. These are close to the maximum evaporation rate for which the disc is unstable but they do not imply fine tuning.
#### 3.2.3 Rise times
In the model of A0620-00, once an outburst is triggered, it takes the disc $`9`$ days to reach the compact object (the time required to go from $`5\times 10^8`$ cm to $`10^7`$ cm, not included in our simulations, is small because of the short timescales in the inner disc, so that we neglect it). It also takes $`14`$ days for the accretion rate at the disc inner edge to reach its maximum value. Overall, the total rise time in X-rays is therefore $`t_{\mathrm{rise}}5`$ days, which is in reasonable agreement with the available information on FRED-type light-curves. The total rise time in optical is much shorter because the outbursts are triggered in a truncated disc, i.e. effectively in the outer regions of the accretion flow. Our results are similar to those of Hameury et al. (1997a) who addressed in more details the issue of time delays between the rise in optical and X-ray.
The rise timescale for the period of fastest flux increase is $`\tau _{\mathrm{rise}}4`$ days in X-rays. This is somewhat longer than the observed rise timescales of 0.3-2 days (§ 2.1.2). For completeness, we also give the rise timescale in optical, $`\tau _{\mathrm{rise}}1`$ hour, for which there is very little available observational data (GRO J1655-40 being the only system in which the rise in optical was observed; see Orosz et al. 1997). This timescale is approximately the thermal time at the disc inner edge, since this is where most of the optical light comes from at the beginning of an outburst. It is not easily comparable to the observations, because of light dilution by both the secondary and the ADAF.
In the model of Aql X-1 shown in Fig. 5, it takes the disc $`3`$ days to reach the compact object once an outburst is triggered. The fact that this time is shorter than in the model of A0620-00 can be partly attributed to a less efficient evaporation in this model (because $`\dot{M}_{\mathrm{Edd}}`$ and $`R_S`$, in Eq. , scale with the mass $`M_1`$ of the primary). It also takes $`12`$ days for the accretion rate at the disc inner edge to reach its maximum value. Overall, the total rise time is therefore $`t_{\mathrm{rise}}9`$ days, which again is in reasonable agreement with the available data on FRED-type light-curves.
The total rise time in optical is not so short: $`t_{\mathrm{rise}}9`$ days. This is because the disc is less strongly truncated than in the model of A0620-00, which results in much more “inside-out/symmetric” outburst shapes (e.g. Smak 1984). In addition, because the disc is large, the inner parts do not contribute as much as in the case of A0620-00. The rise timescale for the period of fastest flux increase is $`\tau _{\mathrm{rise}}1`$ days in X-rays, which is consistent with available observational data. It is $`\tau _{\mathrm{rise}}1.7`$ hours in optical.
#### 3.2.4 Decay and duration times
In the model of Aql X-1, the total duration time of outbursts is $`t_{\mathrm{dur}}25`$ days. This is short compared to the duration $`70`$ days reported in the literature for FRED-type outbursts. Again, we note that the outbursts produced in this model are not really of the FRED-type, but have a rather symmetric “triangle” shape.
The decay timescale early in the decline is $`\tau _{\mathrm{dec}}3.5`$ days in X-rays and $`\tau _{\mathrm{dec}}4`$ days in optical. The timescales become smaller later on because the light-curves steepen. These values are clearly shorter than the observed values of 25-40 days in X-rays and 10-30 days in optical.
In the model of A0620-00, the total duration time of an outburst is $`t_{\mathrm{dur}}55`$ days. This is quite short compared to the typical duration times for BH SXTs with FRED-type light-curves, $`170260`$ days in X-rays and even more in optical. Note that $`t_{\mathrm{dur}}`$ would be shorter if reflares did not occur (see below for a discussion of reflares).
Estimating the decay timescales is complicated by the presence of reflares in this model. For completeness, we estimate two types of $`\tau _{\mathrm{dec}}`$. The overall $`\tau _{\mathrm{dec}}`$ corresponds to the decay timescale for the global light-curve, ignoring the “fluctuations” induced by the reflares. It is $`\tau _{\mathrm{dec}}20`$ days in X-rays and $`\tau _{\mathrm{dec}}22`$ days in optical. The actual decay timescales, in between two successive reflares, are shorter: $`\tau _{\mathrm{dec}}9`$ days in X-rays and $`\tau _{\mathrm{dec}}11`$ days in optical. We find it difficult to compare these timescales with observed values, since observed light-curves do not resemble the light-curves shown in Fig. 2 (see § 5 for a discussion)
One should keep in mind that we did not take into account outer disc X–ray irradiation which will increase significantly both the duration and the decay timescale (Dubus 1999a,b).
#### 3.2.5 Outburst amplitude
In the model of A0620-00, the maximum accretion rate reached at the inner edge of the disc is substantially sub–Eddington ($`0.015\dot{M}_{\mathrm{Edd}}`$; Fig. 2). This is far too small. In the model of Aql X-1, on the other hand, the Eddington limit is reached without difficulty (Fig. 5).
The total amount of mass accreted during an outburst is also quite small in the model of A0620-00. This is not surprising because the recurrence time is more than 10 times shorter than observed, and not much mass is accumulated in quiescence. On the other hand, the total accreted mass is comparable to what is deduced from the observations in the case of Aql X-1. Again, this is not a surprise since the recurrence time is about right.
#### 3.2.6 Reflares
A striking feature of the models described here is the possible occurrence of reflares during the decline of an outburst (Fig. ). These reflares correspond to multiple reflections of cooling and heating fronts during their propagation in the disc. For the specified set of parameters, we find that they are present in models with a BH primary, but absent when the primary is a NS. The number of reflares is reduced when the disc is truncated because increasing the inner disc radius gradually removes the regions of the disc where reflections occur.
Reflares are in fact a natural outcome of the DIM. They have been observed in many simulations over the years, in particular in models of DN discs in which a same value $`\alpha _{\mathrm{hot}}=\alpha _{\mathrm{cold}}`$ was used (e.g. Smak 1984). In that case, a cooling front develops, propagates over some fraction of the disc radial extent, is reflected into a heating front which goes back all the way up to the disc outer edge. This is repeated indefinitely and leads to small amplitude variations of the disc luminosity. This failure of the model has been overcome by using a higher value of $`\alpha `$ for the hot ionized disc than for the cold neutral disc, which results in large amplitude outbursts as those observed.
The fundamental reason why a reflare occurs is because the surface density just behind a cooling front reaches the critical value $`\mathrm{\Sigma }_{\mathrm{max}}`$ sometime during its propagation. This happens typically when the front is located at $`R_{\mathrm{front}}\mathrm{}<10^{10}`$ cm in our simulations. An outward propagating heating front then appears which destroys the inner cooling front as mass starts being accreted efficiently in the cooling region. Later on, this heating front is reflected when the surface density just behind the front reaches the critical value $`\mathrm{\Sigma }_{\mathrm{min}}`$. A cooling front develops which shuts off the outward transport of angular momentum at the origin of the heating front propagation. A sequence of heating and cooling front reflections corresponds to the multiple reflares observed in the simulations.
What was probably not previously realized is that the occurrence of reflares depends not only on the values of $`\alpha `$ but also on the mass of the central object, as shown by our simulations. We note that reflares were also found in the simulations of the decline phase of BH SXT outbursts of Cannizzo, Chen & Livio (1995). Cannizzo et al. interpreted the reflares as numerical artifacts. This is not the case in our models: reflares are a generic property of the DIM and the TDIM for some range of masses and viscosity parameters.
One can rather easily see how the presence or absence of reflares depends on the mass of the accreting object. As noticed by Menou et al. (1999a) the surface density at the cooling front is
$$\mathrm{\Sigma }\left(R_{\mathrm{front}}\right)=\mathrm{\Sigma }_{\mathrm{min}}\left(R_{\mathrm{front}}\right),$$
(4)
where $`\mathrm{\Sigma }_{\mathrm{min}}(R)`$ is the ‘minimum’ critical surface density in the disc. Vishniac (1997) emphasized that the cold, outer regions of the disc behind the cooling front are essentially frozen during the cooling front propagation. Therefore, $`\mathrm{\Sigma }(R)`$ behind the front is constant with time and equal to $`K(M_1,\alpha )\mathrm{\Sigma }_{\mathrm{min}}`$, where $`K`$ depends on the viscosity parameter $`\alpha `$ and the mass of the accreting object $`M_1`$. According to the disc models of Menou et al. (1999a), for a central mass $`M_1=1.2`$ M, $`K4`$, while for $`M_1=7`$ M, $`K67`$. The ratio of the ‘maximum’ to ‘minimum’ surface densities is (Hameury et al. 1998)
$$\frac{\mathrm{\Sigma }_{\mathrm{max}}\left(\alpha _{\mathrm{cold}}=0.02\right)}{\mathrm{\Sigma }_{\mathrm{min}}\left(\alpha _{\mathrm{hot}}=0.1\right)}6.4,$$
(5)
so that for $`M_1=1.2`$ M the post-cooling-front surface density $`K\mathrm{\Sigma }_{\mathrm{min}}<\mathrm{\Sigma }_{\mathrm{max}}`$, and no reflares are expected. For $`M_1=7`$ M, the post-cooling-front surface density $`K\mathrm{\Sigma }_{\mathrm{min}}\mathrm{}>\mathrm{\Sigma }_{\mathrm{max}}`$ and it is not surprising that reflares are present.
To explain the dependence of $`K`$ on the mass of the accreting object, one can use the cooling front model of Vishniac & Wheeler (1996). Although some of the predictions of this model were found to be somewhat inaccurate by Menou et al. (1999a), we expect the scalings derived here to be valid. In this model the front speed can be written as:
$$v_{\mathrm{front}}\left(\frac{H}{R}\right)^{0.7}.$$
(6)
From mass conservation the post-cooling front density can be written as (see e.g. Vishniac 1997)
$$\mathrm{\Sigma }_{\mathrm{pf}}\mathrm{\Sigma }_{\mathrm{min}}\frac{v_r(\mathrm{front})}{v_{\mathrm{front}}},$$
(7)
where $`v_r(\mathrm{front})`$ is the speed of matter at the cooling front. Using Vishniac & Wheeler (1996) model, one obtains
$$\frac{\mathrm{\Sigma }_{\mathrm{pf}}}{\mathrm{\Sigma }_{\mathrm{min}}}M_1^{0.35}.$$
(8)
This scaling, $`KM_1^{0.35}`$, is in good agreement with the numerical results of Menou et al. (1999a).
#### 3.2.7 Consequences
Our models confirm the results of Cannizzo (1998) showing that a stronger evaporation increases $`t_{\mathrm{rec}}`$ because it removes the inner, most unstable parts of the disc. The recurrence times are also increased if the mass transfer rate $`\dot{M}_T`$ in the model is reduced, since more time is then needed to accumulate mass in the disc up to the critical surface density $`\mathrm{\Sigma }_{\mathrm{max}}`$ at which the next outburst is triggered. An exploration of the parameter space of the models shows that it is not possible, however, to reproduce recurrence times of ten to tens of years in TDIMs with ‘standard’ values of $`\alpha `$. This is because to obtain these long recurrence times, the increase in the strength of evaporation or the reduction of $`\dot{M}_T`$ required lead to globally stable discs.
To demonstrate this, we computed a series of models of NS and BH SXTs with fixed values of $`R_{\mathrm{in}}`$. The various values of $`R_{\mathrm{in}}`$ are equivalent to various strengths of disc evaporation. For example, a model with parameters relevant for A0620-00, with a fixed value of $`R_{\mathrm{in}}=6\times 10^9`$ cm and $`\dot{M}_T=10^{16}`$ g s<sup>-1</sup> has a recurrence time $`t_{\mathrm{rec}}3`$ years. This is slightly smaller than the recurrence time in the model with evaporation shown in Fig. 1, in which $`R_{\mathrm{in}}6\times 10^9`$ cm before the onset of an outburst, but $`\dot{M}_T`$ is slightly smaller. This type of comparison shows that the main effect of evaporation on the outburst cycles is to affect the radial extent of the disc in quiescence, which in turn determines $`t_{\mathrm{rec}}`$.
Figure 6 shows the recurrence times predicted for three mass transfer rates ($`\dot{M}=10^{16}`$, $`3\times 10^{16}`$ and $`10^{17}`$ g s<sup>-1</sup>) and several values of the fixed inner radius $`R_{\mathrm{in}}`$. At a given $`\dot{M}_T`$, the same trend is observed: the recurrence time initially increases rather slowly with $`R_{\mathrm{in}}`$, but then the disc rapidly becomes globally stable at some critical value of $`R_{\mathrm{in}}`$. For reasonable values of $`\dot{M}_T`$, this critical value of $`R_{\mathrm{in}}`$ corresponds to recurrence times well below ten to tens of years. It is possible to adjust the inner disc radius so that the recurrence time will be arbitrary long (see e.g. the two upper points in the lower panel of Fig. 6) but the resulting disc is marginally stable. Keeping it in this state for the entire quiescence period requires a mass-transfer rate $`\dot{M}_T`$ strictly constant, which is clearly unrealistic.
This result on short recurrence times is confirmed by an analytical estimate of the recurrence times of SXTs.
An outburst starts when the surface density reaches the critical value $`\mathrm{\Sigma }_{\mathrm{max}}`$. In the case of outside-in (type A) outbursts, this happens when enough mass has accumulated at the disc outer edge. If the accumulation time at the outer edge is longer than the viscous time (as it is the case for low accretion rates), matter transferred from the secondary diffuses inwards and accumulates at shorter radii, giving rise to an inside-out (type B) outburst. Since, for a given set of parameters, the recurrence time of a type A outburst is shorter than the corresponding time for a type B outburst ($`t_\mathrm{A}<t_\mathrm{B}`$), we consider here only type B outbursts (in any case the only type obtained in our models).
Following Smak (1993), the recurrence time of a type B outburst is
$$t_\mathrm{B}\left(\frac{\mathrm{ln}\mathrm{\Sigma }}{t}\right)^1,$$
(9)
which is the characteristic growth-time of a surface density contrast. From the standard diffusion equation describing the evolution of $`\mathrm{\Sigma }`$, one obtains the relation between $`t_\mathrm{B}`$ and $`\nu \mathrm{\Sigma }`$, where $`\nu `$ is the kinematic viscosity coefficient. The relation between $`\nu \mathrm{\Sigma }`$ and the effective temperature $`T_{\mathrm{eff}}`$ is
$$\sigma _{\mathrm{SB}}T_{\mathrm{eff}}^4=\frac{9}{8}\frac{GM_1\xi }{R^3}\nu \mathrm{\Sigma },$$
(10)
where $`\sigma _{\mathrm{SB}}`$ is the Stefan-Boltzmann constant, $`G`$ is the gravitational constant, $`R`$ is the distance from the central object of mass $`M_1`$, and $`\xi 35`$ is a constant factor that accounts for the non-stationarity of the quiescent disc (Idan et al . 1999). Using the values of $`\mathrm{\Sigma }_{\mathrm{max}}`$ calculated by Hameury et al. (1998) and assuming a constant effective temperature, which is a good approximation in quiescence (Smak 1984; Hameury et al. 1998), one gets
$`t_\mathrm{B}3\left({\displaystyle \frac{\xi }{3}}\right)\left({\displaystyle \frac{M_1}{M_{}}}\right)^{0.62}\left({\displaystyle \frac{R}{10^{10}\mathrm{cm}}}\right)^{0.14}\times `$
$`\left({\displaystyle \frac{\alpha _{\mathrm{cold}}}{0.02}}\right)^{0.83}\left({\displaystyle \frac{T_{\mathrm{eff}}}{3000\mathrm{K}}}\right)^4`$ $`\mathrm{yr}.`$ (11)
This characteristic time $`t_\mathrm{B}`$ is expected to be an upper limit to the recurrence time (e.g. $`T_{\mathrm{eff}}`$ may be larger than 3000 K). Again, this rough estimate is much less than tens of years. It also suggests that one way to obtain long enough recurrence times is to reduce the value of $`\alpha _{\mathrm{cold}}`$. Interestingly, we expect the maximum accretion rate reached by the disc during outburst to increase as well if $`\alpha _{\mathrm{cold}}`$ is reduced. This is because more mass will be present in the disc at the time an outburst is triggered ($`\mathrm{\Sigma }_{\mathrm{max}}\alpha _{\mathrm{cold}}^{0.83}`$; e.g. Hameury et al. 1998). Models of SXTs with smaller values of $`\alpha _{\mathrm{cold}}`$ are considered in the following section.
## 4 “Non-standard” truncated disc models of SXTs
In the following, we present two non-standard truncated disc models of SXTs. The first model is a non-standard TDIM because the value of the viscosity parameter $`\alpha _{\mathrm{cold}}`$ used is significantly smaller than the standard value $`0.02`$. The second model explores the possibility that quiescent, truncated discs in SXTs are globally stable. This model is not a Disk Instability Model, since the origin of an outburst has to be a variation of the mass transfer rate in the system.
### 4.1 TDIMs with a smaller $`\alpha `$ in quiescence
The model described in this section is similar to the model with parameters relevant for the BH SXT A0620-00 presented in § 3, except for $`\alpha _{\mathrm{cold}}`$ which is set to $`5\times 10^3`$ here. We have also taken a slightly larger $`=30`$. The predictions for the outburst cycles are shown in Fig. 7, and a details on the evolution of important quantities during an outburst are shown in Fig. 8.
As expected, the recurrence time predicted in this model is much longer ($`t_{\mathrm{rec}}53`$ years) and the maximum accretion rate reached in outburst is much larger ($`0.1\dot{M}_{\mathrm{Edd}}`$) than in the model with $`\alpha _{\mathrm{cold}}=0.02`$. In our view, this is a success of the model that such long $`t_{\mathrm{rec}}`$ and high $`\dot{M}`$ in outburst can be reached with smaller values of $`\alpha _{\mathrm{cold}}`$.
There are other important differences caused by the reduction of $`\alpha _{\mathrm{cold}}`$. The disc mass is much larger in this model, as is the total mass accreted during an outburst. The variations of $`R_{\mathrm{out}}`$ are also larger. Finally, the disc optical emission in quiescence is quite reduced compared to what it is in the model with $`\alpha _{\mathrm{cold}}=0.02`$.
The most obvious difference, however, is probably the disappearance of reflares in the decline phase of the outbursts when $`\alpha _{\mathrm{cold}}=5\times 10^3`$. This is in fact consistent with the explanation for the occurrence of reflares given in § 3.2.6. The ratio $`\mathrm{\Sigma }_{\mathrm{max}}/\mathrm{\Sigma }_{\mathrm{min}}`$ is indeed increased from 7 to $`20`$ when $`\alpha _{\mathrm{cold}}`$ is reduced from 0.02 to $`5\times 10^3`$ (for the same $`\alpha _{\mathrm{hot}}`$), which makes it more difficult for the post-cooling-front surface density to reach $`\mathrm{\Sigma }_{\mathrm{max}}`$ and trigger a reflare.
The total duration of an outburst predicted in the model with $`\alpha _{\mathrm{cold}}=5\times 10^3`$ is $`t_{\mathrm{dur}}50`$ days (similar in optical and X-rays). This is the same as in the case $`\alpha _{\mathrm{cold}}=0.02`$ and is still fairly short compared to the observed duration times (§ 2.1.4).
Once an outburst is triggered, it takes the disc $`2`$ days to reach the compact object. The rise phase of the outburst (both in X-rays and in optical) can be clearly separated in an early phase of very rapid increase and a late phase of much slower increase (see Fig. 8). The total rise time for the early phase in X-rays (up to $`\dot{M}_{\mathrm{in}}2\times 10^{17}`$ g s<sup>-1</sup>; see Fig. 8) is $`t_{\mathrm{rise}}3`$ days, while the slower phase takes an additional $`15`$ days to reach the outburst peak. The total rise time for the early phase in optical (up to $`V=13`$; see Fig. 8) is $`t_{\mathrm{rise}}1`$ day, while the slower phase takes an additional $`10`$ days to reach the outburst peak.
The rise timescales (we only consider the early phases of rapid increase here) are $`\tau _{\mathrm{rise}}0.4`$ days in X-rays (short but consistent with observed values), and $`\tau _{\mathrm{rise}}0.1`$ days in optical. The decay timescales (early in the decline phase, before it steepens) are $`\tau _{\mathrm{dec}}10`$ days in X-rays and $`\tau _{\mathrm{dec}}13`$ days in optical. These values are too small, and not consistent with the lowest observed values.
### 4.2 Models with globally stable truncated discs
For completeness, we also constructed a model with parameters relevant for the BH SXT A0620-00 in which the disc is truncated (same evaporation law as before), the mass transfer rate $`\dot{M}_T`$ is low enough for the disc to be globally stable during a long ‘quiescent’ period, and an outburst is triggered by a slow variation of $`\dot{M}_T`$. (Note that this is different from a burst of mass, or an enhanced mass transfer during the outburst.) More specifically, $`\dot{M}_T`$ was initially set to the value $`10^{15}`$ g s<sup>-1</sup> and was linearly increased over a period of several years until it reaches a value sufficient to trigger an outburst in the disc. The values of $`\alpha `$ chosen for this model are $`\alpha _{\mathrm{hot}}=0.1`$ and $`\alpha _{\mathrm{cold}}=0.02`$, i.e. similar to the ‘standard’ values of the models described in § 3.2.
We find that in most respects, the predictions of this model for the outburst properties are similar to those of the TDIM of A0620-00 described in § 3.2 (which is why we do not show a figure for an outburst in this model; Fig. 2 is basically applicable): the maximum accretion rate reached is slightly smaller, and two reflares instead of three are predicted (as well as a duration a bit shorter).
Of course, there is no prediction for the recurrence time in this model, since one has to refer to some activity cycle in the secondary to obtain outburst cycles for the disc.
## 5 Discussion
We have constructed models of SXT outbursts that include evaporation of the disc inner regions during quiescence. The only numerical limitation in our calculations is the minimum radius $`R_{\mathrm{min}}=5\times 10^8`$ cm allowed for the disc inner radius $`R_{\mathrm{in}}`$ during evolution. Our working assumption is that, in outburst, the disc extends down to the compact object with a nearly constant accretion rate in its inner regions. It seems a reasonable assumption given the results of existing calculations on such extended discs (e.g. Cannizzo et al. 1995), and it allows us not to worry about modeling the innermost regions of the disc during the outbursts.
One of the advantage of this simplification is that we do not have to handle the problematic radiation-pressure dominated regions of the disc, located close to the compact object in outburst. Our results concerning the recurrence times should not be affected by the numerical limitation $`R_{\mathrm{in}}R_{\mathrm{min}}`$ since it does not affect the propagation of the cooling wave and is irrelevant in quiescence (during which $`R_{\mathrm{in}}R_{\mathrm{min}}`$).
There are, however, important physical effects that are not included in our simulations. Disc irradiation by the central X-ray source is not taken into account. Irradiation is likely to increase the duration times of the outbursts, as well as increase the recurrence times and the maximum accretion rates reached in outburst, as suggested by King & Ritter (1998) and confirmed by numerical calculations of Dubus (1999a,b). Irradiation is also likely to increase the decay timescales, both in X-rays and optical. If the optical emission from the system is dominated by reprocessing of X-rays by the disc during outburst, the predicted optical light-curves could be markedly different from those found here in models without irradiation. Nevertheless, it is possible that irradiation is unimportant in some SXTs (e.g. Shahbaz, Charles & King 1998). Our results should be applicable to these systems.
Another effect that has been neglected in our simulations is the fact that the disc could reach the 3:1 resonance during its expansion in outburst, at least in short orbital period systems, where heating fronts can reach the outer edge of the disc. Tidal effects could then significantly affect the outburst properties (Osaki 1996). Similarly, enhanced mass transfer due to the irradiation of the secondary might play an important role in SXTs, as it has been argued for WZ Sge-type DN (Smak 1993; Hameury et al. 1997b; Augusteijn, Kuulkers & Shaham 1993).
One of the advantages of our calculations is that they are free of all these complications, and isolate the effects of including the evaporation of the disc and of reducing the value of $`\alpha _{\mathrm{cold}}`$. We acknowledge that none of the models presented in this paper matches the requirements that we defined in § 2 for a fully successful model of SXTs. This suggests that at least one ingredient of the model is missing (the most obvious candidate being irradiation). In future applications of the DIM to SXTs, additional physical effects will have to be included if ones wishes to reach a better agreement between the model predictions and the observations.
One of the predictions of the DIM that was clearly identified in our models is the possible occurrence of reflares during the decline from outburst. Some of the BH SXT light-curves show ‘reflares’ during decline from outburst but these observed features are totally different from the reflares appearing in our models. According to CSL, reflares (which they call ‘secondary maxima’) observed in X–ray transients are of three morphological types: ‘glitches’, ‘bumps’ and ‘mini-outbursts’. Glitches are upward inflections superposed on a smooth exponential decay. Neither the shape, nor the amplitude of reflares obtained in our models correspond to glitches. Bumps and mini-outbursts are more heterogeneous classes of features but it is difficult to find in the CSL compilation a light–curve morphology similar to the one seen in Fig. (2). The reflares are only present in models with $`\alpha 10^3`$, i.e. in models with recurrence times too short to correspond to the observations. It is, therefore, a success of the model that the set of parameters ensuring both long recurrence times and absence of unobserved reflares are consistent. Some of the observed reflares may be due to disc irradiation (G. Dubus, private comm. 1999), which means that this effect has to be included before any detailed model of SXT with reflares can be constructed.
The fact that, in our simulations, reflares occur only when the mass of the accreting compact object is several $`M_{}`$ (the boundary being $`3M_{}`$ for standard values of $`\alpha `$) is interesting because it is apparently consistent with the fact that reflares have been observed only in SXTs containing BHs, and not in SXTs containing NSs (e.g. CSL). However, since our reflares are different from the observed ones and absent in the favoured model, this might be a coincidence. On the other hand, since a complete, self-consistent model of X–ray transients is still to be constructed it could be useful to keep in mind this property differentiating neutron star from black holes.
By using the strongest evaporation compatible with existing upper limits on the value of $`R_{\mathrm{in}}`$ in quiescent BH SXTs, we tested if evaporation alone can be responsible for the long recurrence times of SXTs, or if very small values of $`\alpha _{\mathrm{cold}}`$ in their quiescent discs are also required. Our calculations in § 3.2.7 show that a strong evaporation is not sufficient, and that either a viscosity parameter $`\alpha _{\mathrm{cold}}`$ at the level of $`10^3`$ is required to reach agreement with recurrence times of tens of years, or that disc outbursts must be drastically affected by X-ray illumination of the disc.
The very low values of $`\alpha _{\mathrm{cold}}`$ required to model BH SXTs would suggest that more than one viscosity mechanism drives accretion in the disc. The viscosity in the hot state arises probably from a well developed MHD turbulence with a universal viscosity parameter $`\alpha _{\mathrm{hot}}\mathrm{}>0.1`$. The much lower value of $`\alpha `$ in quiescent discs is consistent with the disappearance of MHD turbulence in quiescent discs proposed by Gammie & Menou (1998) (see also Meyer & Meyer-Hofmeister 1999). We note that if the relevant mechanisms in quiescence were non-local in nature, the use of a local $`\alpha `$prescription would be questionable and a major revision of the DIM might be required.
## 6 Conclusion
In this paper, we investigated the effect of disc evaporation on the predictions of the DIM for SXTs. As for DN discs, we show that the effect of allowing the disc outer radius $`R_{\mathrm{out}}`$ to vary with time is crucial to obtain robust predictions for the outburst cycles.
We use the strongest evaporation still compatible with available upper limits on the value of the disc inner radius in quiescent BH SXTs to test if evaporation alone can be responsible for the long recurrence times of SXTs.
We find that such a strong evaporation does increase the recurrence times predicted (in agreement with earlier calculations by Cannizzo 1998), but is not sufficient to reproduce the longest known recurrence times of SXTs when standard values of the viscosity parameter $`\alpha `$ are used in the disc ($`\alpha _{\mathrm{hot}}0.1`$, $`\alpha _{\mathrm{cold}}0.02`$).
We show that models with strong evaporation and a significantly smaller value of $`\alpha _{\mathrm{cold}}`$ ($``$ a few $`10^3`$) do reproduce the long recurrence times of tens of years of BH SXTs, as well as the high X-ray luminosities reached by these systems in outburst. The value of $`\alpha _{\mathrm{cold}}`$ needed to reproduce the long recurrence times would be even smaller if evaporation was weaker than what we assumed. Another possibility is of course that X-ray illumination drastically alters the outburst cycle. We defer the examination of this effect to a future work.
We argue that the requirement for a smaller value of $`\alpha _{\mathrm{cold}}`$ in at least some SXTs is consistent with the disappearance of MHD turbulence in quiescent discs (Gammie & Menou 1998) and suggests that another mechanism, perhaps non-local, could be responsible for accretion in the disc during this low luminosity phase.
Finally, since the version of the DIM used in this paper fails to reproduce several important properties of SXTs (such as duration and shape of the light-curve), future work must include in the model the observed missing ingredient: X–ray irradiation of the disc.
## Acknowledgments
JPL thanks Guillaume Dubus for illuminating discussions. This work was supported in part by the National Science Foundation under Grant No. PHY94-07194 and by NASA grant 5-2837, and in part by ASPS/CNRS. KM was supported by a SAO predoctoral fellowship and a French Higher Education Ministry grant. |
warning/0001/cs0001021.html | ar5iv | text | # REFINEMENT OF A STRUCTURED LANGUAGE MODEL11footnote 1This work was funded by the NSF IRI-19618874 grant STIMULATE
## 1 INTRODUCTION
The task of a speech recognizer is to automatically transcribe speech into text. The most successful approach to speech recognition so far is a statistical one : given the observed string of acoustic features $`A`$, find the most likely word string $`\widehat{W}`$ among those that could have generated $`A`$:
$`\widehat{W}=argmax_WP(W|A)=argmax_WP(A|W)P(W)`$ (1)
This paper is concerned with the estimation of the language model probability $`P(W)`$. We will first describe current modeling approaches to the problem, followed by a detailed explanation of our model. A few preliminary experiments that show the potential of our approach for language modeling will then be presented.
## 2 BASIC LANGUAGE MODELING
The language modeling problem is to estimate the source probability $`P(W)`$ where $`W=w_1,w_2,\mathrm{},w_n`$ is a sequence of words. This probability is estimated from a text training corpus. Usually the model is parameterized: $`P_\theta (W),\theta \mathrm{\Theta }`$ where $`\mathrm{\Theta }`$ is referred to as the parameter space. Due to the sequential nature of an efficient search algorithm, the model operates left-to-right, allowing the computation
$`P(w_1,w_2,\mathrm{},w_n)=P(w_1){\displaystyle \underset{i=2}{\overset{n}{}}}P(w_i/w_1\mathrm{}w_{i1})`$ (2)
We thus seek to develop parametric conditional models:
$`P_\theta (w_i/w_1\mathrm{}w_{i1}),\theta \mathrm{\Theta },w_i𝒱`$ (3)
where $`𝒱`$ is the vocabulary chosen by the modeler. Currently most successful is the *n-gram language model*:
$`P_\theta (w_i/w_1\mathrm{}w_{i1})=P_\theta (w_i/w_{in+1}\mathrm{}w_{i1})`$ (4)
### 2.1 LANGUAGE MODEL QUALITY
All attempts to derive an algorithm that would estimate the model parameters so as to minimize the word error rate have failed. As an alternative, a statistical model is evaluated by how well it predicts a string of symbols $`W_t`$ — commonly named *test data* — generated by the source to be modeled.
#### 2.1.1 Perplexity
Assume we compare two models $`M_1`$ and $`M_2`$; they assign probability $`P_{M_1}(W_t)`$ and $`P_{M_2}(W_t)`$, respectively, to the sample test string $`W_t`$. “Naturally”, we consider $`M_1`$ to be a better model than $`M_2`$ if $`P_{M_1}(W_t)>P_{M_2}(W_t)`$. The test data is not seen during the model estimation process.
A commonly used quality measure for a given model $`M`$ is related to the entropy of the underlying source and was introduced under the name of perplexity (PPL) :
$`PPL(M)=exp(1/|W_t|{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{ln}[P_M(W_t)])`$ (5)
### 2.2 SMOOTHING
Assume that our model $`M`$ is faced with the prediction $`w_i|w_1\mathrm{}w_{i1}`$ and that $`w_i`$ has not been seen in the training corpus in context $`w_1\mathrm{}w_{i1}`$ which itself has possibly not been encountered in the training corpus. If $`P_M(w_i|w_1\mathrm{}w_{i1})=0`$ then $`P_M(w_1\mathrm{}w_N)=0`$ thus forcing a recognition error; good models are smooth, in the sense that
$`ϵ(M)>0`$ s.t. $`P_M(w_i|w_1\mathrm{}w_{i1})>ϵ,w_i𝒱`$, $`(w_1\mathrm{}w_{i1})𝒱^{i1}`$.
One standard approach that ensures smoothing is the deleted interpolation method . It interpolates linearly among contexts of different order $`h_n`$:
$`P_\theta (w_i|w_{in+1}\mathrm{}w_{i1})={\displaystyle \underset{k=0}{\overset{k=n}{}}}\lambda _kf(w_i/h_k)`$ (6)
where: $`h_k=w_{ik+1}\mathrm{}w_{i1}`$ is the context of order $`k`$ when predicting $`w_i`$; $`f(w_i/h_k)`$ is the relative frequency estimate for the conditional probability $`P(w_i/h_k)`$; $`\lambda _k,k=0\mathrm{}n`$ are the interpolation coefficients satisfying $`\lambda _k>0,k=0\mathrm{}n`$ and $`_{k=0}^{k=n}\lambda _k=1`$.
The model parameters $`\theta `$ then are: the counts $`C(h_n,w_i)`$ — lower order counts are inferred recursively by: $`C(h_k,w_i)=_{w_{ik}𝒱}C(w_{ik},h_k,w_i)`$ — and the interpolation coefficients $`\lambda _k,k=0\mathrm{}n`$.
A simple way to estimate the model parameters involves a two stage process:
1. gather counts from *development data* — about 90% of training data;
2. estimate interpolation coefficients to minimize the perplexity of *check data* — the remaining 10% of the training data.
Different smoothing techniques are also used e.g., maximum entropy or back-off .
## 3 DESCRIPTION OF THE STRUCTURED LANGUAGE MODEL
The model we present is closely related to the one investigated in , however different in a few important aspects:
* our model operates in a left-to-right manner, thus allowing its use directly in the hypothesis search for $`\widehat{W}`$ in (1);
* our model is a factored version of the one in , thus enabling the calculation of the joint probability of words and parse structure; this was not possible in the previous case due to the huge computational complexity of the model.
### 3.1 THE BASIC IDEA AND TERMINOLOGY
Consider predicting the word `after` in the sentence:
`the contract ended with a loss of 7 cents after trading as low as 89 cents`.
A 3-gram approach would predict `after` from `(7, cents)` whereas it is intuitively clear that the strongest word-pair predictor would be `contract ended` which is outside the reach of even 7-grams. Our assumption is that what enables humans to make a good prediction of `after` is the syntactic structure of its sentence prefix. The linguistically correct *partial parse* of this prefix is shown in Figure 1.
A binary branching parse for a string of words is a binary tree whose leaves are the words. The headword annotation makes the tree an oriented graph: at each node we have two children; the current node receives a *headword* from either child; one arrow suffices to describe which of the children — left or right — is percolated to become the headword of the parent. It was found that better parse trees are generated when using tags: part-of-speech(POS) tags for the leaves and non-terminal(NT) tags for the intermediate nodes in the parse tree. Any subtree identifies a *constituent*. The word `ended` is called the *headword* of the *constituent* `(ended (with (...)))` and `ended` is an *exposed headword* when predicting `after` — topmost headword in the largest constituent that contains it. The syntactic structure in the past filters out irrelevant words and points to the important ones, thus enabling the use of long distance information when predicting the next word.
Our model will attempt to build the syntactic structure incrementally while traversing the sentence left-to-right; it will assign a probability $`P(W,T)`$ to every sentence $`W`$ with every possible POStag assignment, binary branching parse, non-terminal tag and headword annotation for every constituent of the parse tree $`T`$.
Let $`W`$ be a sentence of length $`n`$ words to which we have prepended `<s>` and appended `</s>` so that $`w_0=`$`<s>` and $`w_{n+1}=`$`</s>`. Let $`W_k`$ be the word k-prefix $`w_0\mathrm{}w_k`$ of the sentence and $`W_kT_k`$ the *word-parse k-prefix*. A word-parse k-prefix contains — for a given parse — only those binary subtrees whose span is completely included in the word k-prefix, excluding $`w_0=`$`<s>`. Single words along with their POStag can be regarded as root-only subtrees. Figure 2 shows a word-parse k-prefix; `h_0 .. h_{-m}` are the *exposed heads*, each head being a pair (headword, non-terminal tag), or (word, POStag) in the case of a root-only tree.
A *complete parse* — Figure 3 — is a binary parse of the
`(<s>, SB)` $`(w_1,t_1)\mathrm{}(w_n,t_n)`$ `(</s>, SE)` sequence with the following two restrictions:
1. $`(w_1,t_1)\mathrm{}(w_n,t_n)`$ `(</s>, SE)` is a constituent, headed by `(</s>, TOP’)`;
2. `(</s>, TOP)` is the only allowed head. Note that $`((w_1,t_1)\mathrm{}(w_n,t_n))`$ *needn’t* be a constituent, but for the parses where it is, there is no restriction on which of its words is the headword or what is the non-terminal tag that accompanies the headword.
Our model can generate all and only the complete parses for a string
`(<s>, SB)` $`(w_1,t_1)\mathrm{}(w_n,t_n)`$ `(</s>, SE)`.
The model will operate by means of three modules:
* WORD-PREDICTOR predicts the next word $`w_{k+1}`$ given the word-parse k-prefix $`W_kT_k`$ and then passes control to the TAGGER;
* TAGGER predicts the POStag $`t_{k+1}`$ of the next word given the word-parse k-prefix and the newly predicted word $`w_{k+1}`$ and then passes control to the PARSER;
* PARSER grows the already existing binary branching structure by repeatedly generating the transitions: `(adjoin-left, NTtag)` or `(adjoin-right, NTtag)` until it passes control to the PREDICTOR by taking a `null` transition. `NTtag` is the non-terminal tag assigned to each newly built constituent and `{left,right}` specifies from where the new headword is inherited. The parser operates always on the two rightmost exposed heads, starting with the newly tagged word $`w_{k+1}`$.
The operations performed by the PARSER are illustrated in Figures 6-6 and they ensure that all possible binary branching parses with all possible headword and non-terminal tag assignments for the $`w_1\mathrm{}w_k`$ word sequence can be generated.
It is easy to see that any given word sequence with a possible parse and headword annotation is generated by a unique sequence of model actions.
### 3.2 PROBABILISTIC MODEL
The probability $`P(W,T)`$ of a word sequence $`W`$ and a complete parse $`T`$ can be broken into:
$`P(W,T)=`$ (7)
$`{\displaystyle \underset{k=1}{\overset{n+1}{}}}[P(w_k/W_{k1}T_{k1})P(t_k/W_{k1}T_{k1},w_k){\displaystyle \underset{i=1}{\overset{N_k}{}}}P(p_i^k/W_{k1}T_{k1},w_k,t_k,p_1^k\mathrm{}p_{i1}^k)]`$
where:
* $`W_{k1}T_{k1}`$ is the word-parse $`(k1)`$-prefix
* $`w_k`$ is the word predicted by WORD-PREDICTOR
* $`t_k`$ is the tag assigned to $`w_k`$ by the TAGGER
* $`N_k1`$ is the number of operations the PARSER executes at position $`k`$ of the input string before passing control to the WORD-PREDICTOR (the $`N_k`$-th operation at position k is the `null` transition); $`N_k`$ is a function of $`T`$
* $`p_i^k`$ denotes the i-th PARSER operation carried out at position k in the word string;
$`p_i^k\{`$ `(adjoin-left, NTtag)`, `(adjoin-right, NTtag)`$`\},1i<N_k`$ ,
$`p_i^k=`$`null`$`,i=N_k`$
Each $`(W_{k1}T_{k1},w_k,t_k,p_1^k\mathrm{}p_{i1}^k)`$ is a valid word-parse k-prefix $`W_kT_k`$ at position $`k`$ in the sentence, $`i=\overline{1,N_k}`$.
To ensure a proper probabilistic model certain PARSER and WORD-PREDICTOR probabilities must be given specific values:
* $`P(`$`null`$`/W_kT_k)=1`$, if `h_{-1}.word = <s>` and `h_{0}` $``$ `(</s>, TOP’)` — that is, before predicting `</s>` — ensures that `(<s>, SB)` is adjoined in the last step of the parsing process;
* $`P(`$`(adjoin-right, TOP)`$`/W_kT_k)=1`$,
if `h_0 = (</s>, TOP’)` and `h_{-1}.word = <s>` and
$`P(`$`(adjoin-right, TOP’)`$`/W_kT_k)=1`$,
if `h_0 = (</s>, TOP’)` and `h_{-1}.word` $``$ `<s>` ensure that the parse generated by our model is consistent with the definition of a complete parse;
* $`ϵ>0,W_{k1}T_{k1},P(w_k`$=`</s>`$`/W_{k1}T_{k1})ϵ`$ ensures that the model halts with probability one.
In order to be able to estimate the model components we need to make appropriate equivalence classifications of the conditioning part for each component, respectively.
The equivalence classification should identify the strong predictors in the context and allow reliable estimates from a treebank. Our choice is inspired by :
$`P(w_k/W_{k1}T_{k1})=`$ $`P(w_k/[W_{k1}T_{k1}])`$ $`=P(w_k/h_0,h_1)`$ (8)
$`P(t_k/w_k,W_{k1}T_{k1})=`$ $`P(t_k/w_k,[W_{k1}T_{k1}])`$ $`=P(t_k/w_k,h_0.tag,h_1.tag)`$ (9)
$`P(p_i^k/W_kT_k)=`$ $`P(p_i^k/[W_kT_k])`$ $`=P(p_i^k/h_0,h_1)`$ (10)
It is worth noting that if the binary branching structure developed by the parser were always right-branching and we mapped the POStag and non-terminal tag vocabularies to a single type then our model would be equivalent to a trigram language model.
### 3.3 SMOOTHING
All model components — WORD-PREDICTOR, TAGGER, PARSER — are conditional probabilistic models of the type $`P(y/x_1,x_2,\mathrm{},x_n)`$ where $`y,x_1,x_2,\mathrm{},x_n`$ belong to a mixed bag of words, POStags, non-terminal tags and parser operations ($`y`$ only). For simplicity, the smoothing method we chose was deleted interpolation among relative frequency estimates of different orders $`f_n()`$ using a recursive mixing scheme:
$`P(y/x_1,\mathrm{},x_n)=`$ (11)
$`\lambda (x_1,\mathrm{},x_n)P(y/x_1,\mathrm{},x_{n1})+(1\lambda (x_1,\mathrm{},x_n))f_n(y/x_1,\mathrm{},x_n),`$
$`f_1(y)=uniform(vocabulary(y))`$ (12)
The $`\lambda `$ coefficients are tied based on the range into which the count $`C(x_1,\mathrm{},x_n)`$ falls. The approach is a standard one .
### 3.4 PRUNING STRATEGY
Since the number of parses for a given word prefix $`W_k`$ grows exponentially with $`k`$, $`|\{T_k\}|O(2^k)`$, the state space of our model is huge even for relatively short sentences. We thus have to prune most parses without discarding the most likely ones for a given sentence $`W`$. Our pruning strategy is a synchronous multi-stack search algorithm.
Each stack contains hypotheses — partial parses — that have been constructed by *the same number of predictor and the same number of parser operations*. The hypotheses in each stack are ranked according to the $`\mathrm{ln}(P(W_k,T_k))`$ score, highest on top. The width of the search is controlled by two parameters:
* the maximum stack depth — the maximum number of hypotheses the stack can contain at any given time;
* log-probability threshold — the difference between the log-probability score of the top-most hypothesis and the bottom-most hypothesis at any given state of the stack cannot be larger than a given threshold.
### 3.5 WORD LEVEL PERPLEXITY
Attempting to calculate the conditional perplexity by assigning to a whole sentence the probability:
$`P(W/T^{})={\displaystyle \underset{k=0}{\overset{n}{}}}P(w_{k+1}/W_kT_k^{}),`$ (13)
where $`T^{}=argmax_TP(W,T)`$ — the search for $`T^{}`$ being carried according to our pruning strategy — is not valid because it is not causal: when predicting $`w_{k+1}`$ we would be using $`T^{}`$ which was determined by looking at the entire sentence. To be able to compare the perplexity of our model with that resulting from the standard trigram approach, we need to factor in the entropy of guessing the prefix of the final best parse $`T_k^{}`$ *before predicting* $`w_{k+1}`$, based solely on the word prefix $`W_k`$.
To maintain a left-to-right operation of the language model, the probability assignment for the word at position $`k+1`$ in the input sentence was made using:
$`P(w_{k+1}/W_k)={\displaystyle \underset{T_kS_k}{}}P(w_{k+1}/W_kT_k)\rho (W_k,T_k),`$ (14)
$`\rho (W_k,T_k)=P(W_kT_k)/{\displaystyle \underset{T_kS_k}{}}P(W_kT_k)`$ (15)
where $`S_k`$ is the set of all parses present in our stacks at the current stage $`k`$.
Note that if we set $`\rho (W_k,T_k)=\delta (T_k,T_k^{}|W_k)`$ — 0-entropy guess for the prefix of the parse $`T_k`$ to equal that of the final best parse $`T_k^{}`$— the two probability assignments (13) and (14) would be the same, yielding a lower bound on the perplexity achievable by our model when using a given pruning strategy.
A second important observation is that the next-word predictor probability $`P(w_{k+1}/W_kT_k)`$ in (14) need not be the same as the WORD-PREDICTOR probability (8) used to extract the structure $`T_k`$, thus leaving open the possibility to estimate it separately.
### 3.6 PARAMETER REESTIMATION
#### 3.6.1 First Model Reestimation
Our parameter re-estimation is inspired by the usual EM approach. Let $`(W,T^{(k)}),k=1,2,\mathrm{},N`$ denote the set of parses of $`W`$ that survived our pruning strategy. Each parse was produced by a unique sequence of model actions: predictor, tagger, and parser moves. The collection of these moves will be called a derivation. Each of the $`N`$ members of the set is produced by exactly the same number of moves of each type. Each move is uniquely specified by identifiers $`(y^{(m)},\underset{¯}{x}^{(m)}),`$ where $`m\{`$WORD-PREDICTOR, TAGGER, PARSER$`\}`$ denotes the particular model, $`y^{(m)}`$ is the specification of the particular move taken (e.g., for $`m=`$PARSER, the quantity $`y^{(m)}`$ specifies a choice from $`\{left,right,null\}`$ and the exact tag attached), and $`\underset{¯}{x}^{(m)}`$ specifies the move’s context (e.g., for $`m=`$PARSER, the two heads).
For each possible value $`(y^{(m)},\underset{¯}{x}^{(m)})`$ we will establish a counter which at the beginning of any particular iteration will be empty. For each move $`(y^{(m)},\underset{¯}{x}^{(m)})`$ present in the derivation of $`(W,T^{(j)})`$ we add to the counter specified by $`(y^{(m)},\underset{¯}{x}^{(m)})`$ the amount
$$\rho (W,T^{(k)})=\frac{P(W,T^{(k)})}{_{j=1}^NP(W,T^{(j)})}$$
where $`P(W,T^{(j)})`$ are evaluated on the basis of the model’s parameter values established at the end of the preceding iteration. We do that for all $`(W,T^{(j)}),j=1,2,\mathrm{},N`$ and for all sentences $`W`$ in the training data. Let $`C^{(m)}(y^{(m)},\underset{¯}{x}^{(m)})`$ be the counter contents at the end of this process. The corresponding relative frequency estimate will be
$$f(y^{(m)}|\underset{¯}{x}^{(m)})=\frac{C^{(m)}(y^{(m)},\underset{¯}{x}^{(m)})}{_{z^{(m)}}C^{(m)}(z^{(m)},\underset{¯}{x}^{(m)})}$$
The lower order frequencies needed for the deleted interpolation of probabilities in the next iteration are derived in the obvious way from the same counters.
It is worth noting that because of pruning (which is a function of the statistical parameters in use), the sets of surviving parses $`(W,T^{(k)}),k=1,2,\mathrm{},N`$ for the same sentence $`W`$ may be completely different for different iterations.
#### 3.6.2 First Pass Initial Parameters
Each model component — WORD-PREDICTOR, TAGGER, PARSER — is initialised from a set of hand-parsed sentences, after each parse tree $`(W,T)`$ is decomposed into its $`derivation(W,T)`$. Separately for each $`m`$ model component, we:
* gather joint counts $`C^{(m)}(y^{(m)},\underset{¯}{x}^{(m)})`$ from the derivations that make up the “development data” using $`\rho (W,T)=1`$;
* estimate the deleted interpolation coefficients on joint counts gathered from “check data” using the EM algorithm . These are the initial parameters used with the reestimation procedure described in the previous section.
#### 3.6.3 Language Model Refinement
In order to improve performance, we develop a model to be used in (14), different from the WORD-PREDICTOR model (8). We will call this new component the L2R-WORD-PREDICTOR.
The key step is to recognize in (14) a hidden Markov model (HMM) with fixed transition probabilities — although dependent on the position in the input sentence $`k`$ — specified by the $`\rho (W_k,T_k)`$ values.
The Expectation-step of the EM algorithm for gathering joint counts $`C^{(m)}(y^{(m)},\underset{¯}{x}^{(m)})`$, $`m=`$ L2R-WORD-PREDICTOR-MODEL, is the standard one whereas the Maximization-step uses the same count smoothing technique as that descibed in section 3.6.1.
The second reestimation pass is seeded with the $`m=`$ WORD-PREDICTOR model joint counts $`C^{(m)}(y^{(m)},\underset{¯}{x}^{(m)})`$ resulting from the first parameter reestimation pass (see section 3.6.1).
## 4 EXPERIMENTS
We have carried out the reestimation technique described in section 3.6 on 1 Mwds of “development” data. For convenience we chose to work on the UPenn Treebank corpus — a subset of the WSJ (Wall Stree Journal) corpus. The vocabulary sizes were:
word vocabulary: 10k, open — all words outside the vocabulary are mapped to the `<unk>` token; POS tag vocabulary: 40, closed; non-terminal tag vocabulary: 52, closed; parser operation vocabulary: 107, closed. The development set size was 929,564wds (sections 00-20), check set size 73,760wds (sections 21-22), test set size 82,430wds (sections 23-24).
Table 1 shows the results of the reestimation techniques presented in section 3.6; `E?` and `L2R?` denote iterations of the reestimation procedure described in sections 3.6.1 and 3.6.3, respectively. A deleted interpolation trigram model had perplexity 167.14 on the same training-test data.
Simple linear interpolation between our model and the trigram model:
$`Q(w_{k+1}/W_k)=\lambda P(w_{k+1}/w_{k1},w_k)+(1\lambda )P(w_{k+1}/W_k)`$
yielded a further improvement in PPL, as shown in Table 2. The interpolation weight was estimated on check data to be $`\lambda =0.36`$. An overall relative reduction of 11% over the trigram model has been achieved.
As outlined in section 3.5, the perplexity value calculated using (13) is a lower bound for the achievable perplexity of our model; for the above search parameters and E3 model statistics this bound was 99.60, corresponding to a relative reduction of 41% over the trigram model.
## 5 CONCLUSIONS AND FUTURE DIRECTIONS
A new source model that organizes the prefix hierarchically in order to predict the next symbol is developed. As a case study we applied the source model to natural language, thus developing a new language model with applicability in speech recognition.
We believe that the above experiments show the potential of our approach for improved language modeling for speech recognition. Our future plans include:
* experiment with other parameterizations for the word predictor and parser models;
* evaluate model performance as part of an automatic speech recognizer (measure word error rate improvement).
## 6 Acknowledgments
This research has been funded by the NSF IRI-19618874 grant (STIMULATE).
The authors would like to thank to Sanjeev Khudanpur for his insightful suggestions. Also to Harry Printz, Eric Ristad, Andreas Stolcke, Dekai Wu and all the other members of the dependency modeling group at the summer96 DoD Workshop for useful comments on the model, programming support and an extremely creative environment. Also thanks to Eric Brill, Sanjeev Khudanpur, David Yarowsky, Radu Florian, Lidia Mangu and Jun Wu for useful input during the meetings of the people working on our STIMULATE grant. |
warning/0001/hep-lat0001022.html | ar5iv | text | # Strings and Branes in Nonabelian Gauge Theory CERN-TH/99-198
## 1 Introduction
It is an old speculation, motivated by work of Wilson , ‘t Hooft , Polyakov and many others, that non-abelian gauge theories can alternatively be formulated as string theories. Recently this subject has been revived, in the wake of the discovery of Dirichlet-branes and the investigation of their properties. In particular, based on work by Polyakov and Klebanov , Maldacena has argued that at least $`N=4`$ super-Yang-Mills theory has such a string representation; the argument has been extended in , and in many other papers.
The following is a summary of some of these developments, intended especially for non-string theorists. This is mainly a review, although I will use various new arguments and simplified derivations. For a more complete review and an extensive list of references, see .
We consider $`SU(N)`$ gauge theory with gauge field $`A_\mu ^aT_a`$, where $`T_a`$ are the generators of $`SU(N)`$ in the adjoint representation. The Lagrangean is
$$=\frac{1}{g^2}d^4x\text{tr}\frac{1}{4}F_{\mu \nu }F^{\mu \nu },$$
where $`F_{\mu \nu }=_{[\mu }A_{\nu ]}+i[A_\mu ,A_\nu ]`$ is the Yang-Mills field strength and $`g`$ is the coupling constant. There might also be a certain number $`n_F`$ of massless flavors in the fundamental representation of the gauge group.
An important property of the coupling constant $`g`$ is that it flows under scale transformations:
$$\frac{d}{d\mathrm{log}\mu }g^2(\mu )=\beta [g^2(\mu )],$$
where $`\mu `$ is the scale at which, e.g., scattering experiments are performed. The one-loop-beta function
$$\beta (g^2)(\frac{11}{3}N\frac{2}{3}n_F)(g^2)^2$$
has a coefficient that is negative as long as the number of flavors is not too big. Then the theory is asymptotically free, flowing to weak coupling in the UV and to strong coupling in the IR. More precisely, in the UV
$$g^2(\mu )\frac{1}{\mathrm{log}\frac{\mu }{\mathrm{\Lambda }}},$$
where $`\mathrm{\Lambda }`$ is the scale at which $`g`$ becomes of order 1 ($``$ 350 MeV in QCD).
This behavior of the coupling constant means that perturbation theory in $`g`$ is o.k. at short distances (say, much smaller than the size of a nucleus), but is useless in the IR where the theory is strongly coupled. But of course it is often the IR properties of gauge theories that we are particularly interested in. E.g., we would like to see that QCD accounts for quark confinement. In the simpler case of the pure gauge theory, we would like to see that the Wilson loop obeys an area law.
When at least two flavors are present, we would like to see whether and how the global chiral symmetry is spontaneously broken to its diagonal subgroup. And of course we would like to compute baryon masses and compare them with experiment, or (already in the pure gauge theory) compute glueball masses. Also, we would like to study QCD at finite temperature, i.e. under conditions present in the early universe. Is there a phase transition at some critical temperature, above which confinement is lost and chiral symmetry is restored? In order to adress any of these questions, we need methods to study nonabelian gauge theories at strong coupling.
In this talk I will focus on the case $`n_f=0`$ without flavors. Then the coefficient of the quadratic beta function coefficient is simply proportional to $`N`$, and it is useful to redefine the coupling constant to $`\lambda =g^2N`$, such that $`N`$ disappears from the flow equation:
$$\dot{\lambda }\lambda ^2\text{with}\lambda =g^2N.$$
## 2 Lattice Gauge Theory
One way of studying strongly coupled euclidean gauge theory is to regularize the theory by putting it on lattice with lattice spacing $`a`$, using, e.g., the Wilson lattice action, and to then take the continuum limit. If the theory is considered at large bare coupling constant $`\lambda _{bare}`$, then the partition function and all correlation functions can be expanded in a power series in $`\lambda _{bare}^1`$. It is well-know that in this expansion e.g. the partition function can be written as a sum over closed surfaces on the lattice (Fig. 2), with some surface tension $`\sigma `$, weighted by the ‘t Hooft factor $`N^{22g}`$ where $`g`$ is the genus of the surface:
$$Z\underset{cl.surf.\mathrm{\Sigma }}{}\mathrm{exp}\{\sigma Area(\mathrm{\Sigma })\}N^{22g}.$$
In particular, in the “planar limit”
$$g0,N\mathrm{},\lambda =g^2N\text{fixed},$$
only surfaces of spherical topology (genus 0) survive. The question naturally arises, whether there is a continuum limit of this picture: if we go to the continuum limit of the gauge theory by taking the lattice spacing $`a`$ to zero while adjusting the bare coupling constant $`\lambda `$ such that the renormalized coupling constant remains fixed, is then the partition function of the gauge theory given by a sum over closed continuous surfaces? I.e., is SU(N) gauge theory described by a string theory with string tension $`\sigma `$ and string coupling constant $`\kappa =\frac{1}{N^2}`$? In particular, is SU(N) gauge theory in the large-N limit described by some classical string theory?
The same questions can be asked for Wilson loops, i.e. for the trace of the path-ordered exponential of the gauge field, integrated along a closed contour $`C`$:
$$W(C)=<\text{tr}Pe^{i_CA}>.$$
In the lattice theory at strong bare coupling, $`W(C)`$ is given by a sum over surfaces on the lattice bounded by $`C`$ (Fig. 3). Is in the continuum theory $`W(C)`$ given by a sum over continuous surfaces with some finite surface tension $`\sigma `$, and can this be used to derive an area law for the Wilson loop?
There is, of course, a problem. The strong coupling expansion is an expansion in the inverse bare coupling constant. But the coupling constant runs. It runs in such a way that in order to keep the renormalized coupling fixed, we have to take the bare coupling constant (the coupling constant at scale $`\mu \frac{1}{a}`$) to zero:
$$\lambda _{bare}\frac{1}{|\mathrm{log}a|}.$$
Thus we cannot trust the strong coupling expansion any more, and we cannot prove from the lattice regularization that the partition function is given by a sum over surfaces.
One possible way out is to consider instead the $`𝒩=4`$ supersymmetric version of the theory. One of the nice properties of $`𝒩=4`$ $`SU(N)`$ super-Yang-Mills theory is that the coupling constant does not run: $`\dot{\lambda }=0`$. So we can have a continuum theory that is simultaneously strongly coupled in the IR and in the UV. We will return to this case below and identify the string theory by which it indeed seems to be described, following the work mentioned above. As for the non-supersymmetric theory, we will say in the end which problem in two-dimensional conformal field theory needs to be solved in oder to decide whether it is a simple string theory or not.
But for the moment let us stick with the non-supersymmetric theory, let us simply assume that it has a string representation, and let us study what consequences this would have. That is, we parametrize the surfaces that are bounded by the loop $`C`$ by world-sheet coordinates $`\xi _1,\xi _2`$. The embedding coordinates of the surface are $`X^\mu (\xi )`$. So we assume that the Wilson loop is given by a path integral
$`W(C)`$ $``$ $`{\displaystyle _{X^\mu |_C=\widehat{X}^\mu (\xi _1)}}[dX^\mu ]\times `$ (2)
$`\mathrm{exp}\{\sigma {\displaystyle d^2\xi \sqrt{det_\alpha X^\mu _\beta X_\mu }}\}`$
Here we have used the Nambu-Goto world-sheet action, which expresses the area of the world-sheet in embedding space as the integral over the square root of the determinant of the induced metric on the world-sheet. In string theory, the tension $`\sigma `$ is conventionally called $`\frac{1}{\alpha ^{}}`$. The boundary values $`\widehat{X}^\mu (\xi _1)`$ are the embedding coordinates of the loop $`C`$.
We would now like to review that if such a string theory of SU(N) gauge theory exists, it must have at least two perhaps unexpected properties: first, the string theory must live in five – rather than four – dimensions; and second, it cannot be a bosonic string theory but is probably a superstring theory.
## 3 D-branes in gauge theory
Why does the QCD string have to live in more than four dimensions? Here we follow an argument due to Polyakov . The nonpolynomial Nambu-Goto action is difficult to handle in the quantum theory. As is usual in string theory, we rewrite it with the help of an auxiliary two-dimensional world-sheet metric $`h_{\alpha \beta }`$:
$`e^{\sigma {\scriptscriptstyle d^2\xi \sqrt{det_\alpha X^\mu _\beta X_\mu }}}`$ (3)
$``$ $`{\displaystyle [dh_{\alpha \beta }(\xi )]e^{\sigma {\scriptscriptstyle d^2\xi \sqrt{h}h^{\alpha \beta }_\alpha X^\mu _\beta X_\mu }}}.`$ (4)
Classically, $`h`$ is a non-propagating field. It is easy to check that the path integral over $`h`$ has a saddle point where the auxiliary metric is equal to the world-sheet metric,
$`h_{\alpha \beta }=_\alpha X^\mu _\beta X_\mu .`$ (5)
The saddle point value of the integrand indeed reproduces the Nambu-Goto action. The embedding coordinates $`X^\mu `$ are now free fields coupled to the metric $`h`$. We will return to the boundary conditions for $`h`$ in a moment.
A two-dimensional metric has three components, and there are two diffeomorphisms. Thus there is one gauge invariant degree of freedom which can be chosen to be the conformal factor. I.e., every two-dimensional metric, at least on a surface with the topology of a disc, can be written as
$$h_{\alpha \beta }=e^{\varphi (\xi )}\widehat{h}_{\alpha \beta },$$
where $`\widehat{h}`$ is an arbitrarily chosen background metric on the world-sheet that nothing physical can depend on.
Now, classically, even the conformal factor $`\varphi `$ drops out of the action in (4). Quantum mechanically, there is the conformal anomaly:
$$\frac{\delta S_{eff}}{\delta \varphi (x)}\sqrt{h}<T_\mu ^\mu >c\sqrt{h}R^{(2)}=2\mathrm{}\varphi ,$$
where $`c`$ is the “central charge”. Integrating this equation shows that the effective action contains a piece proportional to
$$d^2\xi \sqrt{\widehat{h}}\widehat{h}^{\alpha \beta }_\alpha \varphi _\beta \varphi .$$
Comparing with (4), we see that the conformal factor enters exactly as if it were another embedding coordinate: the string effectively lives in the five-dimensional space $`(X^\mu ,\varphi )`$.
Where in this five-dimensional space is the four-dimensional space in which the Wilson loop lives? To this end we consider the issue of boundary conditions for the metric $`h_{\alpha \beta }`$. In oder to be consistent with the saddle point equation (5) for $`h`$, we choose
$$h_{11}|_C=_1X^\mu _1X_\mu |_C.$$
(More precisely, there is a zero-mode corresponding to constant shifts of $`\varphi `$, which is now also fixed by this boundary condition.) Now, the background metric $`\widehat{h}`$ could be arbitrary, but it is convenient to choose it such that at the boundary:
$$\widehat{h}_{11}|_C=_1X^\mu _1X_\mu |_C.$$
This implies that, at the boundary,
$$\varphi |_C=0.$$
This is a Dirichlet boundary condition. It means that, although the bulk of the world-sheet can move in five dimensions, its boundary is restricted to lie on a four-dimensional hyperplane (see Fig. 4). In string theory, such $`(p+1)`$-dimensional hyperplanes on which otherwise closed strings can end are called “Dirichlet $`p`$-branes”. So at best we can hope to recover SU(N) gauge theory as the world-brane theory that lives on a Dirichlet-3-brane in a higher-dimensional string theory. This was the first remark.
## 4 Fermionic Strings
The second remark is that the string theory cannot be a bosonic string theory . To see this, let us first consider the simpler Ising model. It consists of spins $`\sigma _i`$ whose values can be either +1 or –1, sitting on the sites of a lattice. The energy of a configuration is the sum over all links $`<ij>`$ of the product of the neighboring spins $`\sigma _i`$ and $`\sigma _j`$. The partition function is the sum over all spin configurations, weighted with the Boltzmann factor:
$$Z_{Ising}=\underset{\{\sigma _i\}}{}\mathrm{exp}\{\beta \underset{<ij>}{}\sigma _i\sigma _j\},$$
where $`\beta `$ plays the role of $`1/g^2`$ in the gauge theory.
Just like the partition function of the gauge theory can be written, in the strong coupling expansion, as a sum over closed surfaces, the partition function of the Ising model can be written as a sum over closed paths:
$$Z_{Ising}=\underset{cl.paths}{}e^{M(\beta )L},$$
where $`L`$ is the total length of the paths and $`M=\mathrm{log}(2\mathrm{tanh}\beta )`$ can be interpreted as the mass of a particle moving on the lattice. So at first sight it looks as if in the continuum limit the Ising model was equivalent to a theory of free scalar particles, whose world-lines are these paths.
But we know that this is not true: instead, the Ising model in the continuum limit is equivalent to a theory of free fermions. An intuitive way to see this is the following. Consider the configuration of links drawn in Fig. 5. There are three ways to represent this lattice configuration in terms of closed paths, as shown on the right-hand side. So if we describe the Ising model by free scalar particles, we make a mistake: we count one configuration three times. We can correct the mistake by weighing each configuration with a relative factor
$$(1)^n,$$
where $`n`$ is the number of 360<sup>o</sup> rotations that the tanget vectors to the paths make as they go around the paths. This corrects the mistake because it subtracts the middle path, but now we see that the particle behaves as a fermion: its wave function flips its sign when the particle is rotated by 360 degrees.
There is no similarly clear argument for the case of surfaces on a lattice, but it is clear that there are plaquette configurations that we overcount if we interpret the sum over surfaces on the lattice as a sum over bosonic string world-sheets. And it is reasonable to expect that we can correct this mistake by introducing fermions on the string world-sheet.
One can now try to guess or derive from first principles just what kind of fermionic string it is that we need. This has in fact been tried for a long time without clear result. But there is a simple string theory that contains fermions on the world-sheet and admits Dirichlet-3-branes: this is the type IIB string theory. Type IIB string theory lives in 10 - not in 5 - dimensions. But it can be turned into a 5-dimensional string theory by Kaluza-Klein “compactification” on a compact 5-dimensional manifold $`K^5`$, such as the 5-sphere $`S^5`$. So in the following let us simply “try out” these 5-dimensional superstring theories: that is, instead of starting with a gauge theory and trying to construct a dual string theory out of it, we start with these string theories and try to see what kind of gauge theories on the 3-brane they describe.
To admit it right away: following Maldacena , we will get not the standard $`SU(N)`$ Yang-Mills theory but - depending on what the compactification manifold $`K^5`$ is - various of its conformally invariant cousins, including the $`N=4`$ supersymmetric theory. So the target will be missed, but not by too much, and we will mention in the end how one might be able to get to realistic gauge theories.
## 5 The open bosonic string
The purpose of the next three sections is to explain that the picture in figure 4 is oversimplified in the following respect. It might seem as if the 3-brane was just a fictitious object that sits in flat, empty 5-dimensional embedding space without affecting it. This is not true. First, it turns out that the 3-brane is electrically charged, in a sense that I will explain; in the case of $`SU(N)`$ gauge theory, the charge is proportional to $`N`$. And second, the 5–dimensional space around the D–brane is curved. It turns out to be anti-de Sitter space $`AdS_5`$.
In order to explain these points, we need to spend a few minutes reviewing some basic facts of string theory. To keep things simple, it is useful to first consider the open bosonic string without D-branes (Fig. 6). That is, the boundary of the string world-sheet is allowed to fluctuate all over the embedding space, which in the case of the bosonic string is 26-dimensional.
We consider the two-dimensional field theory that lives on the string world-sheet. There are 26 world-sheet fields $`X^m(\xi )`$, all of which obey Neumann boundary conditions. The world-sheet action is
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle d^2\xi _\alpha X^m^\alpha X^ng_{mn}(X)}`$ (6)
$`+`$ $`{\displaystyle \frac{1}{\sqrt{\alpha ^{}}}}{\displaystyle _\alpha X^mA_m(X)d\xi ^\alpha }.`$ (7)
$`g_{mn}(X)`$ is the embedding space metric, which – from the point of view of the world-sheet theory – represents the possibility of adding generally complicated interactions between the fields $`X^m`$ (there is also a tachyon, a dilaton and an antisymmetric tensor field which we suppress for now). The boundary interactions of the world-sheet theory are similarly represented in embedding space by the 26-dimensional gauge field $`A_m(X)`$.
A consistency condition in classical string theory is that the world-sheet theory must be conformally invariant. One way to understand this is as follows: In section 3, we wrote the world-sheet metric as $`h_{\alpha \beta }=\widehat{h}_{\alpha \beta }e^\varphi `$. $`\varphi `$ became one of the embedding coordinates. $`\widehat{h}`$ was chosen arbitrarily, and we have already said that nothing physical can depend on this choice. In particular, the world-sheet theory must be completely invariant under rescaling the background metric $`\widehat{h}`$.
Scale invariance means that all the beta functions in the theory must vanish. If $`A_m=0`$, the beta function(al) for the metric $`g_{mn}`$ in (6) is well–known to be to lowest order in $`\alpha ^{}`$:
$$\frac{d}{d\mathrm{log}\mu }g_{mn}R_{mn},$$
where $`R_{mn}`$ is the Ricci tensor of $`g_{mn}`$. So we get the Einstein equations (in vacuum). More generally, the $`\beta =0`$ conditions of $`2d`$ field theory are, by definition, the string equations of motion. The string effective action is defined as the action whose variation yields these equations of motion. In the low–energy limit, it is given by :
$$S=S_{Gravity}+S_{Gauge}+o(\alpha ^{})$$
with
$`S_{Gravity}`$ $``$ $`{\displaystyle d^{26}x\sqrt{g}\frac{1}{\kappa ^2}R^{(26)}[g]}`$ (8)
$`S_{Gauge}`$ $``$ $`{\displaystyle d^{26}x\sqrt{g}\frac{1}{\kappa }F_{mn}F^{mn}},`$ (9)
where $`\kappa `$ is the string coupling constant, $`R^{(26)}`$ is the Ricci scalar and $`F_{mn}`$ is the field strength of the gauge field $`A_m`$. We see that the gauge coupling constant is related to the string coupling constant by
$`g_{YM}^2=\kappa .`$ (10)
The above action is, to lowest oder in $`\alpha ^{}`$, the action of Einstein gravity coupled to Maxwell theory. The $`\alpha ^{}`$ corrections denote the difference between Einstein gravity + Maxwell theory and exact conformal invariance of the world-sheet theory.
In the above, we have considered world-sheets with the topology of a disc. We can also allow topology fluctuations in the bulk or at the boundary of the world–sheet (Fig. 7). Those correspond to string loop corrections – i.e. loop corrections to the gravity action and the gauge action, respectively. Formally, the classical action must then be replaced by the full effective action. For the bosonic string, this is of course not well-defined but for the superstring it is.
For the moment, though, let us stay with classical string theory and now restrict the boundary of the world-sheet to lie on a Dirichlet 3-brane. This is done by switching the boundary conditions for all but four of the world–sheet fields $`X^m`$ from “Neumann” to “Dirichlet”. Then the Maxwell theory is also restricted to live on the 3-brane. $`S_{Gauge}`$ becomes a 4-dimensional action, delta-function restricted to the brane:
$`S`$ $`=`$ $`S_{Gravity}^{(26)}+S_{Gauge}^{(4)}\times \delta ^{22}(x)`$ (11)
The fields that $`S_{Gauge}^{(4)}`$ contains are the 26 components of the gauge field $`A_m`$ as before, but now they split up into the 4 components parallel to the brane which make up a 4-dimensional Maxwell field $`A_\mu `$, plus the 22 components transverse to the brane. The latter just become scalar fields that live on the brane.
So we learn two interesting things from bosonic string theory: first of all, the embedding space metric cannot be anything. It must obey - to lowest order in $`\alpha ^{}`$ \- the Einstein equations derived from the action (11). Away from the brane, this implies
$$R_{mn}=0.$$
Second, there is a dynamical gauge field $`A_\mu `$ that lives on the 3-brane. This helps answering a question: in the first part of the talk we have started with a 4-dimensional gauge theory and looked for its “dual” string theory. Given such a string theory, how can we recover the original gauge field? We would like to identify it with the 3-brane gauge field we have just found. However, so far this is only an abelian gauge field. To get a nonabelian gauge field we must consider a generalization: Instead of a single D-brane we formally consider a set of $`N`$ D-branes sitting on top of each other (figure 8). The Wilson loop now carries a “color” index running from 1 to $`N`$ (labelling the D-brane on which it ends), and one can argue that the Maxwell field is replaced by a SU(N) gauge field .
## 6 Superstrings
How is the discussion modified if we consider type IIB superstring theory instead of bosonic string theory? One modification is quite expected: the gauge theory that lives on the brane also becomes supersymmetric. The other crucial difference between D–branes in bosonic string theory and D–branes in superstring theory is that the latter are charged. Let us begin with the first modification.
In the case of the superstring, there are only 10 (and not 26) embedding coordinates $`X^m`$. Viewed as world–sheet fields, four of them (say $`X^\mu `$ with $`\mu \{0,1,2,3\}`$) obey Neumann boundary conditions, and six of them (denoted by $`X^t,t\{4,\mathrm{},9\}`$) obey Dirichlet boundary conditions. We will denote by $`X^4`$ the radial coordinate of this transverse space, and by $`X^5,\mathrm{},X^9`$ its angles.
If the world-sheet boundary could lie anywhere in the embedding space (i.e. if all coordinates had Neumann boundary conditions), the gauge theory would be 10-dimensional super-Yang Mills theory which contains the gauge field $`A_m`$ and a gaugino $`\psi `$ with 16 spinor components. But since the Yang–Mills theory is restricted to the 3-brane, $`A_m`$ splits up into a 4-dimensional gauge field $`A_\mu `$ plus six scalars $`A_i`$, while the 10-dimensional gaugino splits up into four 4-dimensional gauginos $`\psi ^a`$ (since 4-dimensional gauginos have only four spinor components). This results in precisely the field content of $`𝒩=4`$ supersymmetric Yang-Mills theory. So this is why one expects a dual gauge theory with up to 4 supersymmetries. So much for the first modification.
Next, let us discuss in what sense the 3–branes are “charged”. In the case of the bosonic string, the bulk theory was 26–dimensional bosonic gravity. In the case of type IIB superstring theory, the bulk theory is what is called “type IIB supergravity” in the 10-dimensional embedding space. Apart from the metric $`g_{mn}`$, type IIB supergravity contains a variety of fields. But the only one that we need to mention here is the so-called self-dual Ramond-Ramond 4-form gauge field $`C^{(4)}`$. This is an antisymmetric tensor field with 5-form field strength $`F^{(5)}`$. The gauge invariance consists of adding to $`C^{(4)}`$ the total derivative of a 3–form $`\mathrm{\Lambda }^{(3)}`$. “Self-duality” means that $`F^{(5)}`$ is equal to its dual $`F^{(5)}`$:
$`F^{(5)}`$ $`=`$ $`dC^{(4)}`$ (12)
$`\delta C^{(4)}`$ $`=`$ $`d\mathrm{\Lambda }^{(3)}`$ (13)
$`F^{(5)}`$ $`=`$ $`F^{(5)}`$ (14)
The 10–dimensional “gauge field” $`C^{(4)}`$ must of course not be confused with the 1-form $`A_\mu `$ that lives only on the 3–brane.
It is this 4-form gauge field $`C^{(4)}`$, with respect to which – as realized by Polchinski – the 3–brane is charged: 3–branes are sources of “electric field strength” $`F_{01234}`$ (again, “4” denotes the radial direction of the transverse space). Each 3-brane carries one unit of Ramond-Ramond charge:
$$_{S^5}F^{(5)}N$$
where the 5–sphere $`S^5`$ surrounds the 3–brane in 10 dimensions. Moreover, because of self-duality, there is also a component $`F_{56789}`$: each brane also carries one unit of magnetic charge:
$$_{S^5}F^{(5)}N.$$
Here it is useful to think of a D-brane in 10–dimensional type IIB superstring theory as an an analog of a charged particle in ordinary 4–dimensional Maxwell theory. In this analogy, $`F^{(5)}`$ in 10 dimensions is replaced by the Maxwell field strength $`F^{(2)}`$ in 4 dimensions, and the $`S^5`$ is replaced by an $`S^2`$ surrounding the particle. The D–brane is then analogous to a dyon, carrying simultaneously electric and magnetic charge.
Let us conclude our brief review of the relevant facts of string theory with the 10–dimensional bulk action that replaces $`S_{Gravity}^{(26)}`$ in (11). It contains the terms
$`S^{(10)}`$ $``$ $`{\displaystyle d^{10}x\sqrt{g}\frac{1}{\kappa ^2}R^{(10)}[g]}`$ (15)
$`+`$ $`{\displaystyle d^{10}x\sqrt{g}\frac{1}{25!}F_{mnpqr}F^{mnpqr}}.`$ (16)
As a result, the Einstein equations away from the brane change to
$`R_{mn}\kappa ^2F_{mabcd}F_n^{abcd}.`$ (17)
In other words, the electric flux $`F^{(5)}`$ emanating from the branes carries energy-momentum and therefore curves the 5-dimensional space-time. Let us next discuss what the curved space looks like.
## 7 Anti-de Sitter space
From now on we will consider the 10-dimen-sional string theory, “compactified” down to five dimensions. These five dimensions will now be denoted by $`x^m`$ with $`m\{0,1,2,3,\varphi \}`$. The remaining coordinates $`x^i`$ with $`i\{5,6,7,8,9\}`$ parametrize the compact manifold $`K^5`$, whose size and shape we take to be constant, i.e. independent of the $`x^m`$. All fields $`g_{mn},C^{(4)}`$ are assumed to be independent of the compact coordinates $`x^i`$. The 3-branes are again parametrized by $`x^\mu ,\mu \{0,1,2,3\}`$ and are assumed to be “smeared out” over $`K^5`$.
The 3–branes are then analogous to charged capacitor plates: the electric charge of the 4–dimensional branes simply creates a constant density of electric flux $`F_{0123\varphi }`$ in the 5-dimensional space-time. This density is proportional to $`N`$, because there are $`N`$ branes. It is also inversely proportional to the Volume Vol<sub>K</sub> of the 5-dimen-sional compactification space $`K^5`$, because – upon replacing the $`S^5`$ by $`K^5`$ in the previous section – the integral of $`F^{(5)}`$ over $`K^5`$ gives one unit of charge:
$$\frac{F_{0123\varphi }}{\sqrt{g^{(5)}}}\frac{N}{\text{Vol}_K}.$$
Here, $`g^{(5)}`$ is the determinant of the 5-dimensional metric $`g_{mn}`$. From the Einstein equation (17) we see that the curved space is anti-de Sitter space, $`R_{mn}g_{mn}`$: it has constant negative curvature scalar
$$R\kappa ^2F_{0123\varphi }F^{0123\varphi }\frac{\kappa ^2N^2}{\text{Vol}_K^2}.$$
To determine the volume of $`K^5`$, we write the metric on $`K^5`$ as
$$ds_K^2=L^2d\stackrel{}{\theta }^2,$$
where $`\stackrel{}{\theta }`$ are angular coordinates on $`K^5`$ (e.g. $`d\stackrel{}{\theta }^2=d\mathrm{\Omega }_5^2`$ for $`K^5=S^5`$). Because the 5–form $`F^{(5)}`$ is self–dual, the Einstein equation (17) tells us that $`K^5`$ is an Einstein manifold, $`R_{ij}+g_{ij}`$, whose Ricci scalar
$$R_K\frac{\kappa ^2N^2}{\text{Vol}_K^2}$$
has the opposite sign but the same magnitude as that of the anti-de Sitter space-time. Since the Ricci scalar is $`L^2`$ and Vol$`{}_{K}{}^{}L^5`$ we conclude that the curvature radius is
$`{\displaystyle \frac{L}{\sqrt{\alpha ^{}}}}(\kappa N)^{\frac{1}{4}}.`$ (18)
Here we have reinstated $`\alpha ^{}`$ to make $`L`$ dimensionless.
To summarize, the $`N`$ branes carry “Ramond-Ramond charge”, and the resulting flux curves the five-dimensional space, thereby turning it into anti-de Sitter space. One way of writing the metric on anti-de Sitter space is
$`ds^2=d\varphi ^2+\mathrm{exp}\{{\displaystyle \frac{2}{L}}(\varphi \varphi _0)\}dx_{||}^2.`$ (19)
Here, $`x_{||}`$ denotes the coordinates $`x^\mu `$ parallel to the D-branes, and we are assuming that the branes are to the left. Note however that there is a free parameter $`\varphi _0`$; it will be given an interpretation later.
We have derived this result to lowest order in $`\alpha ^{}`$ and for a surface with disc topology. Let us discuss corrections to this result . We have already mentioned three types of corrections. First, there are corrections in $`\alpha ^{}`$. Those come out to be proportional to
$$\alpha ^{}R^{(10)}\frac{\alpha ^{}}{L^2}\frac{1}{\sqrt{\kappa N}}\frac{1}{\sqrt{\lambda }},$$
where $`\lambda =g^2N=\kappa N`$ is the coupling constant of the super-Yang-Mills theory, using (10,18). So the anti-de Sitter space solution that we have found in the supergravity approximation to superstring theory can be trusted (away from the brane) as long as the super-Yang-Mills coupling constant $`\lambda `$ is strong. As $`\lambda `$ is decreased, one has to replace the supergravity equations of motion by the conditions of exact conformal invariance of the world-sheet theory.
The second type of correction comes from topology fluctuations on the world-sheet boundary, as in figure 7. Those correspond to loop corrections to the super-Yang-Mills theory, which are power series in $`\lambda `$ and therefore large. As a result, the classical action $`S_{Gauge}^{(4)}`$ is replaced by the full quantum effective action $`\mathrm{\Gamma }_4`$. So our supergravity solution contains information about this effective action in the limit of large $`\lambda `$ – precisely the limit we were interested in! On the other hand, the supergravity solution is not appropriate in the perturbative regime of small $`\lambda `$.
The third type of corrections is due to topology fluctuations in the bulk of the world-sheet. Those correspond to loop corrections to supergravity. They come out to be proportional to
$$\kappa ^2(\alpha ^{}R^{(10)})^4\frac{\kappa ^2\alpha ^4}{L^8}\frac{1}{N^2}.$$
This is precisely what one expects from a string theory that describes a gauge theory: non-planar diagrams are suppressed by powers of $`\frac{1}{N^2}`$! In particular, the limit $`N\mathrm{}`$ can be taken, in which the classical supergravity solution can be trusted and string loop corrections can be neglected.
## 8 Scales and Wilson loops
We would like to conlude this review with the project that is suggested by figure 3: to compute a Wilson loop by computing the minimal area of the string world-sheet which it bounds. What we have learned in the meantime is that this world-sheet lives in a 5-dimensional curved space with anti-de Sitter metric (19), while its boundary (the Wilson loop) lives on a flat four-dimensional hypersurface located at $`\varphi =0`$.
Before the Wilson loop can be computed, the free parameter $`\varphi _0`$ in the metric (19) must be interpreted, because obviously the result will depend on it. The interpretation is that $`\varphi _0`$ defines the scale in the gauge theory in the following sense: a scale transformation
$`x_{||}`$ $``$ $`e^\tau x_{||}`$ (20)
of the 4-dimensional physical space parametrized by $`x_{||}`$ is equivalent in (19) to a shift of $`\varphi _0`$:
$`\varphi _0`$ $``$ $`\varphi _0+L\tau .`$ (21)
It is useful to redefine
$$z=Le^{\frac{\varphi \varphi _0}{L}},$$
so that the metric becomes independent of the free parameter $`\varphi _0`$:
$`ds^2={\displaystyle \frac{L^2}{z^2}}(dz^2+dx_{||}^2).`$ (22)
We see that anti-de Sitter space has a boundary at $`z=0`$. $`\varphi _0`$ now enters in terms of the location of the branes. The branes sit at $`\varphi =0`$, i.e. at
$$z=z_0Le^{\frac{\varphi _0}{L}}.$$
From (20,21) we then see that translating the D-branes in $`z`$ corresponds to probing different scales in the four-dimensional gauge theory. And in particular, going to the ultraviolet limit $`\varphi _0\mathrm{}`$ of the gauge theory corresponds to placing the branes at the boundary $`z=0`$ of anti-de Sitter space (compare with ). In the following computation, $`z_0`$ will be regarded as an ultraviolet cutoff. The gauge theory on the brane is then the bare gauge theory, while properties of the renormalized gauge theory will be read off from the interior of anti-de Sitter space. As $`z_0`$ is taken to zero, the bare parameters will be adjusted to keep these renormalized parameters fixed.
So let us illustrate these remarks at the example of the Wilson loop in the large-$`N`$, large-$`\lambda `$ limit, outlining a calculation by Maldacena and by Rey and Yee .
In Minkowski space, it is convenient to let the “loop” consist of two lines of length $`T`$ parallel to the time axis, separated by a spatial distance $`l`$. We define $`l`$ to be measured in the metric $`dx_{||}^2`$ (rather than $`\frac{L^2}{z^2}dx_{||}^2`$). A constant time slice through the world-sheet is a line $`z(x)`$ in the $`xz`$plane, as drawn in figure 9. We denote by $`z^{}`$ the maximal value of $`z`$,
$$z^{}=z_{max},$$
and pick the origin of the $`x`$–axis such that $`x=0`$ for $`z=z^{}`$. The Nambu-Goto action (the area of the world-sheet) becomes $`TE`$, where
$`E={\displaystyle 𝑑x\frac{L^2}{z^2}\sqrt{1+(z^{}(x))^2}}`$ (23)
is the energy of a pair of oppositely charged particles, separated by the distance
$$l=2x(z_0).$$
The equations of motion of this action can be seen to yield the following minimal area world-sheet:
$`|x(z)|=z^{}{\displaystyle _{z/z^{}}^1}{\displaystyle \frac{d\zeta }{\sqrt{\frac{1}{\zeta ^4}1}}}.`$ (24)
Let us now define the continuum limit as follows: as discussed above, we think of the gauge theory on the brane as the bare gauge theory with ultraviolet cutoff $`z_0`$. We take this cutoff $`z_0`$ to zero, thus placing the branes at the boundary of anti-de Sitter space. Keeping the renormalized physics fixed means keeping $`z^{}`$ fixed. This requires adjusting $`l`$, which has a finite limit:
$$l(z_0)=2x(z_0)\frac{(2\pi )^{\frac{3}{2}}}{\mathrm{\Gamma }(\frac{1}{4})^2}z^{}\text{as}z_00$$
(On top of this, $`l`$ has a dimension, but this has already been taken care of by measuring $`l`$ using the metric $`dx_{||}^2`$.)
Next, we can compute the area of the minimal area world-sheet by differentiating (24) and plugging the result into (23). In the limit $`z_00`$, the area diverges because
$$x(z)z^3+\text{const.}\text{near}z=0E\frac{1}{z_0}.$$
The divergence is equal to the length of the two dashed lines the figure. Its interpretation is the following: The energy $`E`$ consists of two contributions: one from the electrostatic attraction of the charged particles, and the other one from their mass. In computing the Wilson loop, the latter must be subtracted; it is equal to the energy of massive uncharged particles. The dashed lines just represent the world–sheets of such uncharged particles. Subtracting this divergence and performing the integral, one finds :
$`E`$ $`=`$ $`{\displaystyle \frac{4\sqrt{2}\pi ^2}{\mathrm{\Gamma }(\frac{1}{4})^4}}{\displaystyle \frac{\sqrt{\lambda }}{l}}.`$ (25)
What do we learn from this result? First, we see that the energy obeys a Coulomb law, rather than being linear in $`l`$ as would have been expected in a confining theory. Indeed, the $`𝒩=4`$ theory is conformally invariant and not confining.
So the Wilson loop does not obeys an area law. In the above calculation this arises essentially because the world-sheet drops towards the center of AdS so quickly that the bulk of the world-sheet gives a vanishing contribution to its area. A confining theory would have to be described by a different 5–dimensional geometry (not $`AdS_5`$), in which the world–sheet cannot drop to the center (see ).
What is unexpected in (25) is the factor $`\sqrt{\lambda }`$: perturbatively one would have expected a factor of $`\lambda =g^2N`$. So this is an example of a nontrivial claim of the string representation of non-perturbative N=4 supersymmetric gauge theory .
## 9 Outlook
To conclude, at least for the conformally invariant N=4 supersymmetric version of large–N Yang–Mills theory we now seem to have a strong-coupling description in terms of a particular 5-dimensional string theory, expanded around an $`AdS_5`$ background. Previously, explicit string representations for Yang–Mills theories have been known only in the two–dimensional case .
Much work has followed this “discovery”, of which I will mention just a few examples. First, one can also compute the correlation functions of local operators, following Gubser, Klebanov, Polyakov and Witten . This has been used to predict the scaling dimensions of various operators in the large-$`N`$, large-$`\lambda `$ limit .
Second, replacing the compactification manifold $`S^5`$ by a more general Einstein manifold $`K^5`$, the dual string theories have been identified as certain exotic conformally invariant, but not always supersymmetric four-dimensional gauge theories . To identify them, it helps to compare global symmetries. E.g., the $`SO(6)`$ symmetry of the $`S^5`$ calls for being identified with the SU(4) $``$ SO(6) R-symmetry of the $`𝒩=4`$ super-Yang-Mills theory that rotates the four gauginos into each other . All these theories must be conformal, as the $`SO(2,3)`$ symmetry of $`AdS_5`$ reflects the $`SO(2,3)`$ conformal group in four dimensions. So there appears to be a rich set of RG fixed points and flows between them in four-dimensional gauge theory. The author’s own work in this context can be found in .
Third, it has been pointed out that the 5-dimensional AdS-space is in fact populated with all kinds of elementary and solitonic objects. E.g., instantons in the Yang-Mills theory have been identified with instantons of the 5-dimensional string theory, with their position in the fifth dimension just representing the scale modulus of the instanton solution . Furthermore, it is an old conjecture that baryons appear as solitons in the QCD string theory. In our case there is no chiral matter, and there are no baryons - but in the case of $`SO(N)`$ gauge theory there are analogs of baryons, the “Pfaffian particles”. Those have indeed been identified with string theory solitons – namely five-branes, wrapped over the $`K^5`$ . These and other objects are symbolically drawn in figure 11.
So far we have been discussing the $`𝒩=4`$ supersymmetric or other conformally invariant versions of SU(N) gauge theory. But what about the standard, asymptotically free non-supersymmetric Yang-Mills theory we started with? Can we draw a similar picture as in figure 10?
There is a suggestion due to Witten how to generalize the string representation to non-supersymmetric gauge theory in principle . Using this prescription, some qualitative features such as confinement have been argued to emerge as expected . There is even a quantitative computation of glueball masses (see e.g. ). However, the problem is that this defines a theory with strong bare coupling $`\lambda `$. As mentioned in the beginning, what we need in the non-supersymmetric case is weak bare coupling and strong renormalized coupling. Just as in lattice gauge theory at strong bare coupling, there is no reason to expect the results for glueball masses, etc., to be universal . Universal results are only expected in the continuum limit
$$a0,\lambda _{bare}\frac{1}{|\mathrm{log}a|}.$$
Can this limit be studied in the string theory representation? Since $`\lambda 0`$, the supergravity solution can no longer be trusted. Instead one has to sum up all the $`\alpha ^{}`$-corrections. In other words, one needs an exact two-dimensional conformal field theory.
Many exact conformal field theories in two dimensions are known, but none of them includes the ‘Ramond-Ramond’ background $`F^{(5)}`$ that we need here. So here is the problem that needs to be solved before one can decide whether bosonic SU(N) gauge theory is “dual” to a simple string theory: extend the current algebra methods of conformal field theory to the case of a self-dual Ramond-Ramond background!
Whether this can be done or whether it is as intractable as the perturbative approach to QCD remains to be seen.
NOTE:
One aspect in which this review deviates from the standard treatment of the subject should be mentioned. Usually, the 3-branes are introduced in 10 dimensions. Then one goes to the near–horizon limit of the resulting geometry, which is $`AdS_5\times S^5`$. Here we basically first “compactify” on the $`S^5`$ and then indroduce D–branes in 5 dimensions (section 7). Then there is no need to go to a near-horizon limit: the $`5d`$ geometry is automatically $`AdS_5`$ everywhere. The embarassment of having to explain six new coordinates (instead of only one) also disappears. Moreover, the branes are now “inside the geometry” before the continuum limit of the gauge theory is taken. The continuum limit automatically moves them to the “correct” end of $`AdS_5`$ – the boundary (section 8). |
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