Stable-Diffusion-XL-Base-1.0

How to run

Model Description

Stable Diffusion XL Base 1.0 (SDXL 1.0) is a foundation text-to-image model released by Stability AI.
It is the flagship successor to Stable Diffusion 2.1, designed for photorealism, artistic flexibility, and high-resolution generation.

SDXL 1.0 is a latent diffusion model trained on a broad dataset of images and captions. Compared to prior versions, it improves prompt alignment, visual coherence, and output quality, especially in complex scenes and detailed compositions.

Features

High fidelity image generation: sharper details and improved realism.
Flexible style range: from photorealistic renders to artistic illustration.
Better prompt alignment: improved understanding of nuanced or multi-concept prompts.
High resolution support: natively trained for 1024×1024 images.
Compositional strength: more accurate handling of multiple subjects and fine object placement.

Use Cases

Creative content generation (illustrations, art, concept design)
Product mockups and marketing visuals
Character and environment ideation
Storyboarding and visual storytelling
Research in generative imaging

Inputs and Outputs

Input:

Text prompts (descriptions, concepts, artistic directions)
Optional negative prompts to avoid undesired elements

Output:

Generated image(s) matching the prompt
Default resolution: 1024×1024 pixels

How to use

1) Install Nexa-SDK

Download and follow the steps under "Deploy Section" Nexa's model page: Download Windows SDK

2) Get an access token

Create a token in the Model Hub, then log in:

nexa config set license '<access_token>'

3) Run the model

Running:

nexa infer NexaAI/sdxl-base

License

Licensed under: CreativeML Open RAIL++-M License

References

Model card: https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0