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1111.1671
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# Construction of wedge-local nets of observables through Longo-Witten
endomorphisms. II
Marcel Bischoff 111Supported in part by the ERC Advanced Grant 227458 OACFT
“Operator Algebras and Conformal Field Theory”. and Yoh Tanimoto∗
e-mail: bischoff@mat.uniroma2.it, tanimoto@mat.uniroma2.it
Dipartimento di Matematica, Università di Roma “Tor Vergata”
Via della Ricerca Scientifica, 1 - I–00133 Roma, Italy.
###### Abstract
In the first part, we have constructed several families of interacting wedge-
local nets of von Neumann algebras. In particular, there has been discovered a
family of models based on the endomorphisms of the ${\rm U(1)}$-current
algebra ${{\mathcal{A}}^{(0)}}$ of Longo-Witten.
In this second part, we further investigate endomorphisms and interacting
models. The key ingredient is the free massless fermionic net, which contains
the ${\rm U(1)}$-current net as the fixed point subnet with respect to the
${\rm U(1)}$ gauge action. Through the restriction to the subnet, we construct
a new family of Longo-Witten endomorphisms on ${{\mathcal{A}}^{(0)}}$ and
accordingly interacting wedge-local nets in two-dimensional spacetime. The
${\rm U(1)}$-current net admits the structure of particle numbers and the
S-matrices of the models constructed here do mix the spaces with different
particle numbers of the bosonic Fock space.
Dedicated to Roberto Longo on the occasion of his 60th birthday
###### Contents
1. 1 Introduction
2. 2 Preliminaries
1. 2.1 Fermi nets on $S^{1}$
2. 2.2 Subnets and the character argument
3. 2.3 Scattering theory of waves in ${\mathbb{R}}^{2}$ (revisited)
4. 2.4 Restriction of wedge-local nets
3. 3 Examples of fermi nets
1. 3.1 ${\rm U(1)}$-current net ${{\mathcal{A}}^{(0)}}$
2. 3.2 The free complex fermion net $\mathrm{Fer}_{\mathbb{C}}$
3. 3.3 ${\rm U(1)}$-current net as a subnet of $\mathrm{Fer}_{\mathbb{C}}$
4. 4 A new family of Longo-Witten endomorphisms on ${\rm U(1)}$-current net
5. 5 Interacting wedge-local net with particle production
1. 5.1 Construction of scattering operators
2. 5.2 Action of the S-matrix on the 1+1 particle space
6. 6 Conclusion and outlook
## 1 Introduction
As already explained in Part I [Tan11a], construction of interacting models of
Quantum Field Theory in (physical) four spacetime dimensions has been a long-
standing open problem, and recently the algebraic approach had several
progress [Lec08, GL07, GL08, BS08, BLS11, Lec11] and two dimensional cases
work particularly well: these works constructed models of QFT with weaker
localization property, and in some case such models turned out to be strictly
local and fully interacting [Lec08]. One should recall, however, that the
models in [Lec08] allow a complete interpretation in terms of particles
(asymptotic completeness) and the particle number is preserved under the
scattering operator. On the other hand, it is known that in four dimensions an
interacting model inevitably involves particle production [Aks65]. In the
present paper, we construct a further new family of interacting wedge-local
two-dimensional massless models and find that their S-matrices mix the spaces
with different particle numbers.
In fact, the requirement to involve particle production non-perturbatively is
already not simple. On the one hand, an asymptotically complete model must
behave like the free theory and hence must be compatible with the Fock space
structure at asymptotic time. On the other hand, a particle production process
properly means a violation of the Fock structure at physical time. To overcome
this difficulty, one would have to “deform” the free theory in a somewhat
involved way (cf. [Lec11]) or should rely on a nice trick. Here we take the
second way. Standard examples and techniques from Conformal Field Theory
provide such a trick.
Conformal Field Theory has been well studied particularly on the circle, which
can be seen as a chiral part of 1+1 dimensional theory. There are many
important examples of such models, or nets in operator-algebraic terms, and
both field-theoretic and operator-algebraic techniques allow one to analyze
their interrelationships. Our trick can be briefly summarized as follows: we
consider the free complex fermionic field $\psi$ on the circle; the field
$\psi$ admits a gauge group action by ${\rm U(1)}$, and the fixed point with
respect to this action is known to be isomorphic to the algebra of the ${\rm
U(1)}$-current $J$. Both fields are free fields acting naturally on the Fock
space (fermionic and bosonic, respectively) but the correspondence between the
spaces is quite involved. The passage to 1+1 dimensional models is simply the
tensor product of two such chiral parts. Now, we can easily “deform” the two-
dimensional Dirac field (built up from the chiral parts
$\psi\otimes{\mathbbm{1}}$ and ${\mathbbm{1}}\otimes\psi$) in such a way that
it commutes with the product action of the gauge group ${\rm U(1)}\times{\rm
U(1)}$. Hence the deformation restricts to the algebra of the conserved
current $J^{\mu}=(J^{0},J^{1})=(J\otimes{\mathbbm{1}}+{\mathbbm{1}}\otimes
J,{\mathbbm{1}}\otimes J-J\otimes{\mathbbm{1}})$, and this deformation is
sufficiently complicated so that the resulting S-matrix does not preserve the
bosonic Fock structure, thanks to the involved fermion-boson correspondence.
In Part I, we have constructed a family of two-dimensional massless models
based on the free current $J^{\mu}$ or more precisely its net
${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ of von Neumann algebras of
observables. The main ingredient was endomorphisms of the algebra
${{\mathcal{A}}^{(0)}}({\mathbb{R}}_{+})$ of observables localized in the
positive half-line ${\mathbb{R}}_{+}$ commuting with the translations. A
family of such endomorphisms has been studied first by Longo and Witten [LW11]
in order to construct Quantum Field Theory with boundary. We used those
endomorphisms to construct two-dimensional models without boundary. In the
present article, we study the fermi net $\mathrm{Fer}_{\mathbb{C}}$ generated
by the free complex fermionic field $\psi$ and its Longo-Witten endomorphisms.
We construct endomorphisms of $\mathrm{Fer}_{\mathbb{C}}$ which commute with
the gauge action of ${\rm U(1)}$, hence restrict to the fixed point subnet of
${{\mathcal{A}}^{(0)}}$. It turns out that the restricted endomorphisms cannot
be implemented by second quantization operators, hence are different from the
ones considered in [LW11]. We again knit them up to construct S-matrices and
wedge-local nets.
Then the fixed point with respect to the action of ${\rm U(1)}\times{\rm
U(1)}$ is considered. We find that its asymptotic behaviour is the same as the
free (bosonic) current $J^{\mu}$ and the S-matrix does not preserve the space
of $1+1$ particles ($1$ left and $1$ right moving particle) in the sense of
the Fock space structure. We stress that the Fock space particle number has no
intrinsic meaning as particles, because we are in a massless case where just a
scattering between two waves is considered. One has to pass to the massive
case to talk about particle production. We will discuss in more detail the
implication of this phenomenon at the end of Section 5.
This paper is organized as follows. In Section 2 we recall the standard
notions of algebraic QFT and the scattering theory of two-dimensional massless
models [Buc75, DT11]. Some simple observations are given about subtheories and
inner symmetries. Main examples of nets, the free complex fermionic net
$\mathrm{Fer}_{\mathbb{C}}$ and the ${\rm U(1)}$-current net
${{\mathcal{A}}^{(0)}}$, are introduced in Section 3. Although it is well-
known [KR87, Kac98, Reh98] that the fixed point subnet
$\mathrm{Fer}_{\mathbb{C}}$ with respect to ${\rm U(1)}$ is
${{\mathcal{A}}^{(0)}}$ at the field-theoretic level, we prove it in the
framework of algebraic approach. Section 4 is devoted to the construction of
new Longo-Witten endomorphisms on ${{\mathcal{A}}^{(0)}}$. They are used in
Section 5 to construct new interacting wedge-local nets. Outlook and open
problems are summarized in Section 6.
## 2 Preliminaries
### 2.1 Fermi nets on $S^{1}$
Here we give a summary of one-dimensional nets, since they will be our
building blocks of the construction of two-dimensional interacting models. In
the first part, we considered local nets of von Neumann algebras on $S^{1}$.
Since we need to exploit the free fermionic field in this second part, a
generalized concept of nets is recalled.
We follow the definition in [CKL08] and denote by ${\rm
M\ddot{o}b}^{(2)}(\cong{\rm SL}(2,{\mathbb{R}})\cong{\rm SU}(1,1))$ the double
cover of the Möbius group ${\rm PSL}(2,{\mathbb{R}})$. We denote by
$\mathcal{I}$ the set of proper intervals $I\subset S^{1}$, where proper means
that $I$ is open and connected and neither dense nor the empty set. A (Möbius
covariant) fermi net is an assignment of von Neumann algebras
${\mathcal{F}}_{0}(I)$ on ${\mathcal{H}}_{{\mathcal{F}}_{0}}$ to intervals
$I\in\mathcal{I}$ on $S^{1}$ satisfying the following conditions:
1. (1)
Isotony. If $I_{1}\subset I_{2}$, then
${\mathcal{F}}_{0}(I_{1})\subset{\mathcal{F}}_{0}(I_{2})$.
2. (2)
Möbius covariance. There exists a strongly continuous unitary representation
$U_{0}$ of the group ${\rm M\ddot{o}b}^{(2)}$ such that for any interval
$I\in\mathcal{I}$ it holds that
$U_{0}(g){\mathcal{F}}_{0}(I)U_{0}(g)^{*}={\mathcal{F}}_{0}(gI),\mbox{ for
}g\in{\rm M\ddot{o}b}^{(2)},$
where the action of ${\rm M\ddot{o}b}^{(2)}\cong{\rm SU}(1,1)$ on $S^{1}$ is
defined through linear fractional transformation.
3. (3)
Positivity of energy. The generator of the one-parameter subgroup of the lift
of rotations in $\mathrm{M\ddot{o}b}$ in the representation $U_{0}$ is
positive.
4. (4)
Existence of the vacuum. There is a unique (up to a phase) unit vector
$\Omega_{0}$ in ${\mathcal{H}}_{{\mathcal{F}}_{0}}$ which is invariant under
the action of $U_{0}$, and cyclic for $\bigvee_{I\Subset
S^{1}}{\mathcal{F}}_{0}(I)$.
5. (5)
${\mathbb{Z}}_{2}$-grading. There is a unitary operator $\Gamma_{0}$ with
$\Gamma_{0}^{2}={\mathbbm{1}}$ such that $\Gamma_{0}\Omega_{0}=\Omega_{0}$ and
${\hbox{\rm Ad\,}}\Gamma_{0}({\mathcal{F}}_{0}(I))={\mathcal{F}}_{0}(I)$.
6. (6)
Graded locality. If $I_{1}\cap I_{2}=\emptyset$, then
$[{\mathcal{F}}_{0}(I_{1}),{\hbox{\rm
Ad\,}}Z_{0}({\mathcal{F}}_{0}(I_{2}))]=0$, where
$Z_{0}:=\frac{{\mathbbm{1}}-\mathrm{i}\Gamma_{0}}{1-\mathrm{i}}$.
If the grading operator is trivial: $Z_{0}={\mathbbm{1}}$, then the net
${\mathcal{F}}_{0}$ is said to be local.
Among the consequences of these conditions are (see [CKL08]):
1. (7)
Reeh-Schlieder property. The vector $\Omega_{0}$ is cyclic and separating for
each ${\mathcal{F}}_{0}(I)$.
2. (8)
Additivity. If $I=\bigcup_{i}I_{i}$, then
${\mathcal{F}}_{0}(I)=\bigvee_{i}{\mathcal{F}}_{0}(I_{i})$.
3. (9)
Twisted Haag duality on $S^{1}$. For an interval $I\in\mathcal{I}$, it holds
that ${\mathcal{F}}_{0}(I)^{\prime}={\hbox{\rm
Ad\,}}Z_{0}({\mathcal{F}}_{0}(I^{\prime}))$, where $I^{\prime}$ is the
interior of the complement of $I$ in $S^{1}$.
4. (10)
Bisognano-Wichmann property. The modular group
$\Delta_{{\mathcal{F}}_{0}({\mathbb{R}}_{+})}^{\mathrm{i}t}$ of
${\mathcal{F}}_{0}({\mathbb{R}}_{+})$ with respect to $\Omega_{0}$ is equal to
$U_{0}(\delta(-2\pi t))$, where $S^{1}$ is identified as the one-point
compactification of ${\mathbb{R}}$ as below and $\delta$ is the one-parameter
group of dilations.
5. (11)
Irreducibility. It holds that
$\bigvee_{I\in\mathcal{I}}{\mathcal{F}}_{0}(I)=B({\mathcal{H}}_{{\mathcal{F}}_{0}})$.
Each algebra ${\mathcal{F}}_{0}(I)$ is referred to as a local algebra (even
for a fermi net). Note that if the grading operator $\Gamma_{0}$ is trivial,
then the definition of fermi net coincides with the one of local Möbius-
covariant net. We identify the circle $S^{1}$ and the compactified real line
${\mathbb{R}}\cup\\{\infty\\}$ through the Cayley transform
$t=-\mathrm{i}\frac{z-1}{z+1}\Longleftrightarrow
z=-\frac{t-\mathrm{i}}{t+\mathrm{i}},\phantom{...}t\in{\mathbb{R}},\phantom{..}z\in
S^{1}\subset{\mathbb{C}}$
and refer to the algebra ${\mathcal{F}}_{0}(I)$ for an interval
$I\subset{\mathbb{R}}$. The representation $U_{0}$ of ${\rm
M\ddot{o}b}^{(2)}\cong\mathrm{SL}(2,{\mathbb{R}})$ restricts indeed to a
projective unitary representation of ${\rm PSL}(2,{\mathbb{R}})$ [CKL08]. Let
$\rho$ be the ($4\pi$ periodic) lift of the rotations in $\mathrm{PSU}(1,1)$
(acting by $\rho(\theta)z=\mathrm{e}^{\mathrm{i}\theta}z$) to
$\mathrm{M\ddot{o}b}^{(2)}$ and let us denote
$R_{0}(\theta)=U_{0}(\rho(\theta))=\mathrm{e}^{\mathrm{i}\theta L_{0}}$. Under
the identification between $S^{1}$ and ${\mathbb{R}}\cup\\{\infty\\}$, one can
talk about the translations and dilations of ${\mathbb{R}}$, which are
included in $\mathrm{M\ddot{o}b}$. In particular, the representation of
translations (which we denote by $\tau$) plays a crucial role. Let us denote
$T_{0}(t)=U_{0}(\tau(t))$.
A Longo-Witten endomorphism of a fermi net ${\mathcal{F}}_{0}$ is an
endomorphism of the algebra ${\mathcal{F}}_{0}({\mathbb{R}}_{+})$ implemented
by a unitary $V_{0}$ which commutes with the translation $T_{0}(t)$. A family
of Longo-Witten endomorphisms has been found for the ${\rm U(1)}$-current net
and the real free fermion net [LW11]. The examples will be explained later in
detail.
Note that a Longo-Witten endomorphism is uniquely implemented up to scalar.
Indeed, since it commutes with translation, ${\hbox{\rm Ad\,}}V_{0}$ is an
endomorphism of ${\mathcal{F}}_{0}({\mathbb{R}}_{+}+t)$ for any
$t\in{\mathbb{R}}$. If there is another unitary $W_{0}$ which satisfies
${\hbox{\rm Ad\,}}W_{0}(x)={\hbox{\rm Ad\,}}V_{0}(x)$ for any
$x\in{\mathcal{F}}_{0}({\mathbb{R}}_{+}+t)$, $t\in{\mathbb{R}}$, then by the
irreducibility $W_{0}^{*}V_{0}$ must be scalar.
### 2.2 Subnets and the character argument
Let ${\mathcal{F}}_{0}$ be a fermi (or local) net on
${\mathcal{H}}_{{\mathcal{F}}_{0}}$. Another assignment ${\mathcal{A}}_{0}$ of
von Neumann algebras $\\{{\mathcal{A}}_{0}(I)\\}_{I\in\mathcal{I}}$ on
${\mathcal{H}}_{{\mathcal{F}}_{0}}$ is called a subnet of ${\mathcal{F}}_{0}$
if it satisfies isotony, Möbius covariance with respect to the same $U_{0}$
for ${\mathcal{F}}_{0}$ and it holds that
${\mathcal{A}}_{0}(I)\subset{\mathcal{F}}_{0}(I)$ for every interval
$I\in\mathcal{I}$. We simply write
${\mathcal{A}}_{0}\subset{\mathcal{F}}_{0}$. In this case, let us denote
${\mathcal{H}}_{{\mathcal{A}}_{0}}=\overline{\bigvee_{I\in\mathcal{I}}{\mathcal{A}}_{0}(I)\Omega_{0}}$.
Then it is immediate to see that ${\mathcal{A}}_{0}(I)$ and $U_{0}$ restrict
to ${\mathcal{H}}_{{\mathcal{A}}_{0}}$, and by this restriction
${\mathcal{A}}_{0}|_{{\mathcal{H}}_{{\mathcal{A}}_{0}}}$ becomes a fermi net
with the representation of covariance
$U_{0}|_{{\mathcal{H}}_{{\mathcal{A}}_{0}}}$. This restriction is also said to
be a subnet of ${\mathcal{F}}_{0}$ if no confusion arises.
For a fermi net ${\mathcal{F}}_{0}$ on $S^{1}$, a gauge automorphism
$\alpha_{0}$ is a family of automorphisms $\\{\alpha_{0,I}\\}$ of local
algebras which satisfies the consistency condition
$\alpha_{0,I_{2}}|_{{\mathcal{A}}_{0}(I_{1})}=\alpha_{0,I_{1}}\,\,\mbox{ for
}I_{1}\subset I_{2}\,.$
If a gauge automorphism $\alpha_{0}$ preserves the vacuum state
$\langle\Omega_{0},\cdot\,\Omega_{0}\rangle$, it is said to be an inner
symmetry. An inner symmetry $\alpha_{0}$ can be unitarily implemented by the
formula $V_{\alpha_{0}}x\Omega_{0}=\alpha_{0}(x)\Omega_{0}$, where $x$ is an
element of some local algebra ${\mathcal{F}}_{0}(I)$. We say that a compact
group $G$ acts on the net ${\mathcal{F}}_{0}$ when there is automorphisms
$\\{\alpha_{0,g}\\}_{g\in G}$ which satisfy the composition law when
restricted to local algebras. The fixed point subnet with respect to this
action of $G$ is the subnet defined by
${\mathcal{F}}_{0}^{G}(I):={\mathcal{F}}_{0}(I)^{G}$.
Let ${\mathcal{F}}_{0}$ be a fermi net and ${\mathcal{A}}_{0}$ be a subnet.
Recall that, for a Möbius covariant fermi net, the Bisognano-Wichmann property
is automatic. As a consequence, for each interval there is a conditional
expectation $E_{0,I}:{\mathcal{F}}_{0}(I)\to{\mathcal{A}}_{0}(I)$ which
preserves the vacuum state $\langle\Omega_{0},\cdot\,\Omega_{0}\rangle$ and
implemented by the projection $P_{{\mathcal{A}}_{0}}$ onto
${\mathcal{H}}_{{\mathcal{A}}_{0}}$ (see [Tak03, Theorem IX.4.2]). This
projection $P_{{\mathcal{A}}_{0}}$ contains much information of
${\mathcal{A}}_{0}$.
Consider the case where ${\mathcal{A}}_{0}={\mathcal{F}}_{0}^{G}$ is the fixed
point subnet with respect to an action $\alpha_{0}$ of a compact group $G$ by
inner symmetry. Then we have a unitary representation $V_{\alpha_{0}}$ of $G$
on ${\mathcal{H}}_{{\mathcal{F}}_{0}}$. If we write the set of invariant
vectors with respect to $V_{\alpha_{0}}$ by
${\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}$, it holds that
${\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}={\mathcal{H}}_{{\mathcal{A}}_{0}}$.
Indeed, the inclusion
${\mathcal{H}}_{{\mathcal{A}}_{0}}\subset{\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}$
is obvious. On the other hand, for $x\in{\mathcal{F}}_{0}(I)$, we have
$\left(\int_{G}\alpha_{0}(x)\,\mathrm{d}g\right)\Omega_{0}=\int_{G}\left(V_{\alpha_{0}}(g)x\Omega_{0}\right)\,\mathrm{d}g,$
which implies that any vector in ${\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}$ can
be approximated from ${\mathcal{H}}_{{\mathcal{A}}_{0}}$ by the Reeh-Schlieder
property.
For the later use, we put here a simple observation.
###### Proposition 2.1.
In the situation above, if a Longo-Witten endomorphism is implemented by
$W_{0}$ and $W_{0}$ commutes with $V_{\alpha_{0}}$, then ${\hbox{\rm
Ad\,}}W_{0}$ restricts to a Longo-Witten endomorphism of the fixed point
subnet ${\mathcal{A}}_{0}$.
###### Proof.
The unitary $W_{0}$ commutes with the projection $P_{{\mathcal{A}}_{0}}$,
hence also with the conditional expectation $E_{0}$ onto ${\mathcal{A}}_{0}$.
∎
Let ${\mathcal{F}}_{0}$ be fermi (or local) net on
${\mathcal{H}}_{{\mathcal{F}}_{0}}$. The Hilbert space
${\mathcal{H}}_{{\mathcal{F}}_{0}}$ is graded by the action of the rotation
subgroup $R_{0}(\theta)=\mathrm{e}^{\mathrm{i}\theta L_{0}}$:
${\mathcal{H}}_{{\mathcal{F}}_{0}}={\mathbb{C}}\Omega_{0}\oplus\bigoplus_{r\in\frac{1}{2}{\mathbb{N}}}{\mathcal{H}}_{r}=\bigoplus_{r\in\frac{1}{2}{\mathbb{N}}_{0}}{\mathcal{H}}_{r}$
with
${\mathcal{H}}_{r}=\\{\xi\in{\mathcal{H}}_{{\mathcal{F}}_{0}}:R_{0}(\theta)\xi=\mathrm{e}^{\mathrm{i}r\theta}\xi\\}$
and the sum only going over ${\mathbb{N}}_{0}$ for a local net. The conformal
character of the net ${\mathcal{F}}_{0}$ is given as a formal power series of
$t=\mathrm{e}^{-\beta}$:
$\operatorname{tr}_{{\mathcal{H}}_{{\mathcal{F}}_{0}}}(\mathrm{e}^{-\beta
L_{0}})=\sum_{r\in\frac{1}{2}{\mathbb{N}}_{0}}^{\infty}{\hbox{dim}\,}{\mathcal{H}}_{r}\cdot
t^{r}\,.$
Let us assume that there is an action of $G={\rm U(1)}$ by inner symmetry. We
denote by $V_{0}(\theta)$ the implementing unitary. Then $V_{0}$ and $U_{0}$
commute and ${\mathcal{H}}_{F_{0}}$ is graded also by the gauge action
$V_{0}(\theta)=\mathrm{e}^{\mathrm{i}\theta Q_{0}}$:
${\mathcal{H}}_{{\mathcal{F}}_{0}}={\mathbb{C}}\Omega_{0}\oplus\bigoplus_{r\in\frac{1}{2}{\mathbb{N}},q\in{\mathbb{Z}}}{\mathcal{H}}_{r,q}=\bigoplus_{q\in{\mathbb{Z}}}{\mathcal{H}}_{\,\cdot\,,q},\qquad\text{with}\qquad{\mathcal{H}}_{\,\cdot\,,q}:=\bigoplus_{r\in\frac{1}{2}{\mathbb{N}}_{0}}{\mathcal{H}}_{r,q}$
and the character is given as a formal power series in $t=\mathrm{e}^{-\beta}$
and $z=\mathrm{e}^{-E}$:
$\operatorname{tr}_{{\mathcal{H}}_{{\mathcal{F}}_{0}}}(\mathrm{e}^{-\beta
L_{0}-EQ_{0}})=\sum_{r\in\frac{1}{2}{\mathbb{N}}_{0},q\in{\mathbb{Z}}}{\hbox{dim}\,}{\mathcal{H}}_{r,q}\cdot
t^{r}z^{q}\,.$
Recall that it holds that
${\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}={\mathcal{H}}_{{\mathcal{A}}_{0}}$. The
operator $Q_{0}$ acts by $0$ on ${\mathcal{H}}_{{\mathcal{F}}_{0}}^{G}$, hence
we can obtain the conformal character of ${\mathcal{A}}_{0}$ just by taking
the coefficient of $z^{0}$ in
$\operatorname{tr}_{{\mathcal{H}}_{{\mathcal{F}}_{0}}}(\mathrm{e}^{-\beta
L_{0}-EQ_{0}})$.
Later in this paper we need to compare the size of two subnets. Let
${\mathcal{A}}_{0}\subset{\mathcal{B}}_{0}\subset{\mathcal{F}}_{0}$ be an
inclusion of three fermi nets. If the conformal characters of
${\mathcal{A}}_{0}$ and ${\mathcal{B}}_{0}$ coincide, then this means that the
subspaces ${\mathcal{H}}_{{\mathcal{A}}_{0}}$ and
${\mathcal{H}}_{{\mathcal{B}}_{0}}$ coincide, since we have already an
inclusion
${\mathcal{H}}_{{\mathcal{A}}_{0}}\subset{\mathcal{H}}_{{\mathcal{B}}_{0}}$
and the coefficients of the conformal character are the dimensions of
eigenspaces of $L_{0}$. This in turn implies that two subnets
${\mathcal{A}}_{0}$ and ${\mathcal{B}}_{0}$ are the same since the conditional
expectations which are implemented by
$P_{{\mathcal{A}}_{0}},P_{{\mathcal{B}}_{0}}$ are the same. We will see such
an argument in an example.
### 2.3 Scattering theory of waves in ${\mathbb{R}}^{2}$ (revisited)
Here we just collect some basic notions regarding scattering theory of two-
dimensional massless models. As recalled in Part I [Tan11a], this theory has
been established by Buchholz [Buc75] and extended to the wedge-local case
[DT11]. A Borchers triple on a Hilbert space ${\mathcal{H}}$ is a triple
$({\mathcal{M}},T,\Omega)$ of a von Neumann algebra ${\mathcal{M}}\subset
B({\mathcal{H}})$, a unitary representation $T$ of ${\mathbb{R}}^{2}$ on
${\mathcal{H}}$ and a vector $\Omega\in{\mathcal{H}}$ such that
* •
${\hbox{\rm Ad\,}}T(t_{0},t_{1})({\mathcal{M}})\subset{\mathcal{M}}$ for
$(t_{0},t_{1})\in
W_{\mathrm{R}}=\\{(x_{0},x_{1})\in{\mathbb{R}}^{2}:x_{1}>|x_{0}|\\}$, the
standard right wedge.
* •
The joint spectrum ${\rm sp}\,T$ is contained in the closed forward lightcone
$\overline{V}_{+}=\\{(p_{0},p_{1})\in{\mathbb{R}}_{2}:p_{0}\geq|p_{1}|\\}$.
* •
$\Omega$ is a unique (up to scalar) invariant vector under $T$, and cyclic and
separating for ${\mathcal{M}}$.
We recall that one interprets ${\mathcal{M}}$ as the algebra assigned to the
wedge $W_{\mathrm{R}}$. Let ${\mathcal{W}}$ be the set of wedges, i.e. the set
of all $W=gW_{\mathrm{R}}$ where $g$ is a Poincaré transformation, then we
define the wedge-local net ${\mathcal{W}}\ni W\mapsto{\mathcal{M}}(W)$
associated with the Borchers triple $({\mathcal{M}},T,\Omega)$ by
${\mathcal{M}}(W_{R}+a)=T(a){\mathcal{M}}T(a)^{\ast}$ and
${\mathcal{M}}(W_{R}^{\prime}+a)=T(a){\mathcal{M}}^{\prime}T(a)^{\ast}$. With
the help of the modular objects one can define a representation of the
Poincaré group extending the one of translations $T$ [Bor92]. For details we
refer to the first part.
Take a Borchers triple $({\mathcal{M}},T,\Omega)$ and $x\in B({\mathcal{H}})$.
We write $x(a)={\hbox{\rm Ad\,}}T(a)(x)$ for $a\in{\mathbb{R}}^{2}$ and
consider observables sent to lightlike directions with parameter
${\mathcal{T}}$:
$x_{\pm}(h_{\mathcal{T}}):=\int h_{\mathcal{T}}(t)x(t,\pm t)\,\mathrm{d}t,$
where
$h_{\mathcal{T}}(t)=|{\mathcal{T}}|^{-\varepsilon}h(|{\mathcal{T}}|^{-\varepsilon}(t-{\mathcal{T}}))$,
$0<\varepsilon<1$ is a constant, ${\mathcal{T}}\in{\mathbb{R}}$ and $h$ is a
nonnegative symmetric smooth function on ${\mathbb{R}}$ such that $\int
h(t)\,\mathrm{d}t=1$. Then for $x\in{\mathcal{M}}$, the limits
$\Phi^{\mathrm{out}}_{+}(x):=\underset{{\mathcal{T}}\to+\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,x_{+}(h_{\mathcal{T}})$
and
$\Phi^{\mathrm{in}}_{-}(x):=\underset{{\mathcal{T}}\to-\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,x_{-}(h_{\mathcal{T}})$
exist. Furthermore we set
$\Phi_{+}^{\mathrm{in}}(y^{\prime}):=J_{\mathcal{M}}\Phi_{+}^{\mathrm{out}}(J_{\mathcal{M}}y^{\prime}J_{\mathcal{M}})J_{\mathcal{M}},\,\,\Phi_{-}^{\mathrm{out}}(y^{\prime}):=J_{\mathcal{M}}\Phi_{-}^{\mathrm{in}}(J_{\mathcal{M}}y^{\prime}J_{\mathcal{M}})J_{\mathcal{M}}$
for $y^{\prime}\in{\mathcal{M}}^{\prime}$, where $J_{\mathcal{M}}$ is the
modular conjugation of ${\mathcal{M}}$ with respect to $\Omega$. The
properties of these asymptotic fields are summarized in [DT11, Tan11a]. For
example, it holds that
$\Phi_{+}^{\mathrm{in}}(y^{\prime})=\underset{{\mathcal{T}}\to-\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,y^{\prime}_{+}(h_{\mathcal{T}})$
and
$\Phi_{-}^{\mathrm{out}}(y^{\prime})=\underset{{\mathcal{T}}\to+\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,y^{\prime}_{-}(h_{\mathcal{T}})$.
Let ${\mathcal{H}}_{+}$ (respectively by ${\mathcal{H}}_{-}$) be the space of
the single excitations with positive momentum, (respectively with negative
momentum), i.e. ${\mathcal{H}}_{+}=\\{\xi\in{\mathcal{H}}:T(t,t)\xi=\xi\mbox{
for }t\in{\mathbb{R}}\\}$ (respectively
${\mathcal{H}}_{-}=\\{\xi\in{\mathcal{H}}:T(t,-t)\xi=\xi\mbox{ for
}t\in{\mathbb{R}}\\}$). For $\xi_{+}\in{\mathcal{H}}_{+}$,
$\xi_{-}\in{\mathcal{H}}_{-}$, there are sequences of local operators
$\\{x_{n}\\},\\{y_{n}\\}\subset{\mathcal{M}}$ and
$\\{x^{\prime}_{n}\\},\\{y^{\prime}_{n}\\}\subset{\mathcal{M}}^{\prime}$ such
that
$\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,P_{+}x_{n}\Omega=\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,P_{+}x^{\prime}_{n}\Omega=\xi_{+}$
and
$\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,P_{-}y_{n}\Omega=\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,P_{+}y^{\prime}_{n}\Omega=\xi_{-}$.
We define collision states as in [DT11]:
$\xi_{+}{\overset{{\mathrm{in}}}{\times}}\xi_{-}=\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,\Phi^{\mathrm{in}}_{+}(x^{\prime}_{n})\Phi^{\mathrm{in}}_{-}(y_{n})\Omega,\quad\xi_{+}{\overset{{\mathrm{out}}}{\times}}\xi_{-}=\underset{n\to\infty}{{{\mathrm{s}\textrm{-}\lim}}}\,\Phi^{\mathrm{out}}_{+}(x_{n})\Phi^{\mathrm{out}}_{-}(y^{\prime}_{n})\Omega\,.$
We denote by ${\mathcal{H}}^{\mathrm{in}}$ (respectively
${\mathcal{H}}^{\mathrm{out}}$) the subspace generated by
$\xi_{+}{\overset{{\mathrm{in}}}{\times}}\xi_{-}$ (respectively
$\xi_{+}{\overset{{\mathrm{out}}}{\times}}\xi_{-}$). The isometry
$S:{\mathcal{H}}^{\mathrm{out}}\ni\xi_{+}{\overset{{\mathrm{out}}}{\times}}\xi_{-}\longmapsto\xi_{+}{\overset{{\mathrm{in}}}{\times}}\xi_{-}\in{\mathcal{H}}^{\mathrm{in}}$
is called the scattering operator or the S-matrix of the Borchers triple
$({\mathcal{M}},T,\Omega)$. We say that the Borchers triple
$({\mathcal{M}},T,\Omega)$ is interacting if $S$ is not equal to the identity
operator on ${\mathcal{H}}^{\mathrm{out}}$ and asymptotically complete (with
respect to waves) if it holds that
${\mathcal{H}}^{\mathrm{in}}={\mathcal{H}}^{\mathrm{out}}={\mathcal{H}}$.
We have studied the general structure of asymptotically complete local and
wedge-local nets (using Borchers triple) in [Tan11a, Section 3]. The point was
that for a given (strictly local) $({\mathcal{M}},T,\Omega)$ we can construct
the chiral net, and the original object ${\mathcal{M}}$ can be recovered from
the chiral net and a single operator $S$. Here we rephrase this observation
from the point of view of constructing examples based on chiral components.
See also the general structure of asymptotically complete strictly local nets
[Tan11a, Section 3]
$x$$t$$W_{\mathrm{R}}$$\scriptstyle\mathrm{Ad\,}S({\mathbbm{1}}\otimes{\mathcal{F}}_{-}({\mathbb{R}}_{+}))$$\scriptstyle{\mathcal{F}}_{+}({\mathbb{R}}_{-})\otimes{\mathbbm{1}}$$W_{\mathrm{R}}+(t_{1},x_{1})$$\scriptstyle\mathrm{Ad\,}S({\mathbbm{1}}\otimes{\mathcal{F}}_{-}({\mathbb{R}}_{+}+\frac{t_{1}+x_{1}}{\sqrt{2}}))$$\scriptstyle{\mathcal{F}}_{+}({\mathbb{R}}_{-}+\frac{t_{1}-x_{1}}{\sqrt{2}})\otimes{\mathbbm{1}}$$x$$t$
Figure 1: On the definition of the wedge-local net
###### Proposition 2.2.
Let ${\mathcal{F}}_{\pm}$ be two fermi nets on $S^{1}$ defined on
${\mathcal{H}}_{\pm}$ and assume that there is a unitary operator $S$ on
${\mathcal{H}}_{+}\otimes{\mathcal{H}}_{-}$ commuting with $T_{+}\otimes
T_{-}$, leaving ${\mathcal{H}}_{+}\otimes\Omega_{-}$ and
$\Omega_{+}\otimes{\mathcal{H}}_{-}$ pointwise invariant, such that
$x\otimes{\mathbbm{1}}$ commutes with ${\hbox{\rm
Ad\,}}S(x^{\prime}\otimes{\mathbbm{1}})$ where
$x\in{\mathcal{F}}_{+}({\mathbb{R}}_{-})$ and $x^{\prime}\in{\hbox{\rm
Ad\,}}Z_{+}({\mathcal{F}}_{+}({\mathbb{R}}_{+}))$, and ${\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes y)$ commutes with ${\mathbbm{1}}\otimes
y^{\prime}$ where $y\in{\mathcal{F}}_{-}({\mathbb{R}}_{+})$ and
$y^{\prime}\in{\hbox{\rm Ad\,}}Z_{-}({\mathcal{F}}_{-}({\mathbb{R}}_{-}))$.
Then the triple
* •
${\mathcal{M}}_{S}:=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes
y):x\in{\mathcal{F}}_{+}({\mathbb{R}}_{-}),y\in{\mathcal{F}}_{-}({\mathbb{R}}_{+})\\}^{\prime\prime}$,
* •
$T(t,x):=T_{+}(\frac{t-x}{\sqrt{2}})\otimes T_{-}(\frac{t+x}{\sqrt{2}})$,
* •
$\Omega:=\Omega_{+}\otimes\Omega_{-}$
is an asymptotically complete Borchers triple with the S-matrix $S$.
###### Proof.
As in Part I [Tan11a], the conditions on $T$ and $\Omega$ are automatic
because they are just tensor products of objects for fermi nets. Similarly,
the condition that ${\hbox{\rm
Ad\,}}T(a){\mathcal{M}}_{S}\subset{\mathcal{M}}_{S}$ for $a\in W_{\mathrm{R}}$
is easily seen from the assumption that $T$ commutes with $S$ and the
covariance of fermi nets.
What remains is the cyclicity and separating property of $\Omega$ for
${\mathcal{M}}_{S}$. Cyclicity is immediate because we have
$\displaystyle{\mathcal{M}}_{S}\Omega$ $\displaystyle\supset$
$\displaystyle\\{x\otimes{\mathbbm{1}}\cdot S({\mathbbm{1}}\otimes
y)S^{*}\cdot\Omega:x\in{\mathcal{F}}_{+}({\mathbb{R}}_{-}),y\in{\mathcal{F}}_{-}({\mathbb{R}}_{+})\\}$
$\displaystyle=$ $\displaystyle\\{x\otimes
y\cdot\Omega:x\in{\mathcal{F}}_{+}({\mathbb{R}}_{-}),y\in{\mathcal{F}}_{-}({\mathbb{R}}_{+})\\}$
by the assumed property of $S$, and the latter set is total in
${\mathcal{H}}_{+}\otimes{\mathcal{H}}_{-}$ by the Reeh-Schlieder property for
fermi nets. As for the separating property, we define:
${\mathcal{M}}_{S}^{1}:=\\{{\hbox{\rm
Ad\,}}S(x^{\prime}\otimes{\mathbbm{1}}),{\mathbbm{1}}\otimes
y^{\prime}:x^{\prime}\in{\hbox{\rm
Ad\,}}Z_{+}({\mathcal{F}}_{+}({\mathbb{R}}_{+})),y^{\prime}\in{\hbox{\rm
Ad\,}}Z_{-}({\mathcal{F}}_{-}({\mathbb{R}}_{-}))\\}^{\prime\prime}.$
By an analogous proof, one sees that $\Omega$ is cyclic for
${\mathcal{M}}_{S}^{1}$. Furthermore, ${\mathcal{M}}_{S}$ and
${\mathcal{M}}_{S}^{1}$ commute by assumption. Hence $\Omega$ is separating
for ${\mathcal{M}}_{S}$. In other words, $({\mathcal{M}}_{S},T,\Omega)$ is a
Borchers-triple.
It is immediate that
$\Phi^{\mathrm{out}}_{+}(x\otimes{\mathbbm{1}})=x\otimes{\mathbbm{1}}$ and
$\Phi^{\mathrm{in}}_{-}({\hbox{\rm Ad\,}}S({\mathbbm{1}}\otimes y))={\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes y)$ (the latter follows since $S$ commutes with
$T$). Similarly, we have $\Phi^{\mathrm{in}}_{+}({\hbox{\rm
Ad\,}}S(x^{\prime}\otimes{\mathbbm{1}}))={\hbox{\rm
Ad\,}}S(x^{\prime}\otimes{\mathbbm{1}})$ and
$\Phi^{\mathrm{out}}_{-}({\mathbbm{1}}\otimes y^{\prime})={\mathbbm{1}}\otimes
y^{\prime}$. From this, one concludes that the Borchers triple
$({\mathcal{M}}_{S},T,\Omega)$ is asymptotically complete and its S-matrix is
$S$. ∎
We remark that we see $({\mathcal{M}}_{S},T,\Omega)$ as a fermi (i.e. twisted
local) net defined by ${\mathcal{M}}(W_{\mathrm{R}}^{\prime})={\hbox{\rm
Ad\,}}Z_{+}\otimes Z_{-}({\mathcal{M}})$ and that the scattering theory of
waves [Buc75] is considered to be an analogue of the Haag-Ruelle scattering
theory and it is not intended to be applied to fermionic nets. But we will not
pay much attention to this restriction, since our result is a construction of
wedge-local nets with a free massless bosonic net as the asymptotic net, and
fermionic nets appear as auxiliary objects.
### 2.4 Restriction of wedge-local nets
We consider a Borchers triple $({\mathcal{M}},T,\Omega)$. It is in some cases
interesting to consider a subalgebra ${\mathcal{N}}$ of ${\mathcal{M}}$. Let
us denote ${\mathcal{H}}_{\mathcal{N}}:=\overline{{\mathcal{N}}\Omega}$.
###### Proposition 2.3.
If the subspace ${\mathcal{H}}_{\mathcal{N}}$ is invariant under $T$ and
${\hbox{\rm Ad\,}}T(a)({\mathcal{N}})\subset{\mathcal{N}}$ for $a\in
W_{\mathrm{R}}$. Then
$({\mathcal{N}}|_{{\mathcal{H}}_{\mathcal{N}}},T|_{{\mathcal{H}}_{\mathcal{N}}},\Omega)$
is a Borchers triple on ${\mathcal{H}}_{\mathcal{N}}$.
###### Proof.
The components ${\mathcal{N}},T$ and $\Omega$ naturally restricts to
${\mathcal{H}}_{\mathcal{N}}$. The conditions on $T$ and $\Omega$ are trivial,
even restricted to ${\mathcal{H}}_{\mathcal{N}}$. The cyclicity of $\Omega$ is
immediate from the definition of ${\mathcal{H}}_{\mathcal{N}}$. Since $\Omega$
is already separating for ${\mathcal{M}}$, so is also for ${\mathcal{N}}$.
Endomorphic action of $T$ on ${\mathcal{N}}$ is in the hypothesis. ∎
We call a triple $({\mathcal{N}},T,\Omega)$ a (Borchers) subtriple of
$({\mathcal{M}},T,\Omega)$ if ${\mathcal{N}}$ is a subalgebra of
${\mathcal{M}}$, ${\mathcal{H}}_{\mathcal{N}}$ is invariant under $T(a)$,
${\hbox{\rm Ad\,}}T(a)({\mathcal{N}})\subset{\mathcal{N}}$ for $a\in
W_{\mathrm{R}}$, and ${\mathcal{N}}$ is invariant under ${\hbox{\rm
Ad\,}}\Delta_{\mathcal{M}}^{\mathrm{i}t}$, where $\Delta_{\mathcal{M}}$ is the
modular operator of ${\mathcal{M}}$ with respect to $\Omega$.
Recall that a Borchers triple $({\mathcal{M}},T,\Omega)$ gives rise to a
strictly local net if $\Omega$ is cyclic for ${\mathcal{M}}\cap{\hbox{\rm
Ad\,}}T(a)({\mathcal{M}})^{\prime}$ for any $a\in W_{\mathrm{R}}$. We call
such a triple therefore strictly local. The following proposition shows that
the concept of Borchers subtriple corresponds to the one of a local subnet.
###### Proposition 2.4.
If a Borchers triple $({\mathcal{M}},T,\Omega)$ is strictly local, any
subtriple $({\mathcal{N}},T,\Omega)$ is again strictly local when restricted
on $\overline{{\mathcal{N}}\Omega}$.
###### Proof.
Since ${\mathcal{N}}$ is invariant under the modular automorphism ${\hbox{\rm
Ad\,}}\Delta_{\mathcal{M}}^{\mathrm{i}t}$, there is a conditional expectation
$E$ from ${\mathcal{M}}$ onto ${\mathcal{N}}$ which preserves the state
$\langle\Omega,\,\cdot\,\Omega\rangle$ and is implemented by the projection
$P_{\mathcal{N}}$ (see [Tak03, Theorem IX.4.2] for the original reference and
[Tan11b, Appendix A] for an application to nets).
We have to show that $\Omega$ is cyclic for the relative commutant
${\mathcal{N}}\cap{\hbox{\rm Ad\,}}T(a)({\mathcal{N}})^{\prime}$ on the
subspace ${\mathcal{H}}_{\mathcal{N}}$. Let us denote
${\mathcal{M}}_{0,a}:={\mathcal{M}}\cap{\hbox{\rm
Ad\,}}T(a)({\mathcal{M}})^{\prime}$. We claim that $E({\mathcal{M}}_{0,a})$ is
contained in ${\mathcal{N}}\cap{\hbox{\rm Ad\,}}T(a)({\mathcal{N}})^{\prime}$.
Indeed, by the definition of $E$, the image of $E$ is contained in
${\mathcal{N}}$. Furthermore, if $x\in{\mathcal{M}}_{0,a}$, $y\in{\hbox{\rm
Ad\,}}T(a)({\mathcal{N}})\subset{\hbox{\rm Ad\,}}T(a)({\mathcal{M}})$, then
$E(x)y=E(xy)=E(yx)=yE(x),$
hence they commute and the image $E({\mathcal{M}}_{0,a})$ lies in the relative
commutant. Now we have
$\overline{\left({\mathcal{N}}\cap{\hbox{\rm
Ad\,}}T(a)({\mathcal{N}})^{\prime}\right)\Omega}\supset\overline{E({\mathcal{M}}_{0,a})\Omega}\supset\overline{P_{\mathcal{N}}{\mathcal{M}}_{0,a}\Omega}={\mathcal{H}}_{\mathcal{N}}$
by the assumed strict locality of $({\mathcal{M}},T,\Omega)$. ∎
Let $({\mathcal{B}},U,\Omega)$ be an asymptotically complete local Poincaré
covariant net on ${\mathbb{R}}^{2}$ fulfilling the Bisognano-Wichmann property
(see [Tan11a] for related definitions). We recall that one can define the
(out-) asymptotic algebras ${\mathcal{B}}_{+}\otimes{\mathcal{B}}_{-}$ and the
scattering operator $S$ which is a unitary operator, and that it is possible
to recover the original net by the formula
${\mathcal{B}}(W_{\mathrm{R}})=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes
y):x\in{\mathcal{B}}_{+}({\mathbb{R}}_{-}),y\in{\mathcal{B}}_{-}({\mathbb{R}}_{+})\\}^{\prime\prime}.$
Note that $({\mathcal{B}}(W_{\mathrm{R}}),U|_{{\mathbb{R}}^{2}},\Omega)$ is an
asymptotically complete, strictly local Borchers triple. Here we exhibit a
simple way to construct subtriples. Let ${\mathcal{A}}_{+},{\mathcal{A}}_{-}$
be (Möbius covariant) subnets of ${\mathcal{B}}_{+},{\mathcal{B}}_{-}$,
respectively. If we define
${\mathcal{N}}=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes
y):x\in{\mathcal{A}}_{+}({\mathbb{R}}_{-}),y\in{\mathcal{A}}_{-}({\mathbb{R}}_{+})\\}^{\prime\prime},$
then $({\mathcal{N}},U|_{{\mathbb{R}}^{2}},\Omega)$ is a Borchers subtriple of
$({\mathcal{B}}(W_{\mathrm{R}}),U|_{{\mathbb{R}}^{2}},\Omega)$. Indeed,
conditions regarding ${\mathcal{N}},U|_{{\mathbb{R}}^{2}},\Omega$ are
immediate. As for the invariance of ${\mathcal{N}}$ under ${\hbox{\rm
Ad\,}}\Delta_{\mathcal{M}}^{\mathrm{i}t}$, it suffices to note that $S$ and
$\Delta_{\mathcal{M}}^{\mathrm{i}t}$ commute ([Tan11a, Lemma 2.4], cf.
[Buc75]) and that ${\mathcal{A}}_{+}({\mathbb{R}}_{-})$ and
${\mathcal{A}}_{-}({\mathbb{R}}_{+})$ are preserved by ${\hbox{\rm
Ad\,}}\Delta_{\mathcal{M}}^{\mathrm{i}t}$ because of Bisognano-Wichmann
property.
The trouble is, however, that such Borchers triples constructed as above are
not necessarily asymptotically complete in general. Indeed, the out-asymptotic
states span the subspace
$\overline{{\mathcal{A}}_{+}({\mathbb{R}}_{-})\Omega}\otimes\overline{{\mathcal{A}}_{-}({\mathbb{R}}_{+})\Omega}$.
It is easy to see that this coincides with the full space
$\overline{{\mathcal{N}}\Omega}$ if and only if it is invariant under $S$.
Since a clear-cut scattering theory is so far available only for
asymptotically complete nets, it is worthwhile to give a general condition to
assure that subnets are asymptotically complete. For simplicity, we consider
the following situation: let ${\mathcal{A}}_{0}$ be a fermi net on
${\mathcal{H}}_{0}$ with an action of a compact group $G$ by inner symmetry
implemented by $V_{g}$. Suppose that there is a unitary operator $S$ on
${\mathcal{H}}_{0}\otimes{\mathcal{H}}_{0}$ such that
$({\mathcal{M}}_{S},T,\Omega)$ is a Borchers triple where
* •
${\mathcal{M}}_{S}:=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes
y):x\in{\mathcal{A}}_{0}({\mathbb{R}}_{-}),y\in{\mathcal{A}}_{0}({\mathbb{R}}_{+})\\}^{\prime\prime}$,
* •
$T(t,x):=T_{0}(\frac{t-x}{\sqrt{2}})\otimes T_{0}(\frac{t+x}{\sqrt{2}})$,
* •
$\Omega:=\Omega_{0}\otimes\Omega_{0}$,
as in Proposition 2.2.
###### Proposition 2.5.
If $S$ commutes with $V_{g}\otimes V_{g^{\prime}}$, $g,g^{\prime}\in G$, then
the triple
* •
${\mathcal{N}}_{S}:=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S({\mathbbm{1}}\otimes
y):x\in{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-}),y\in{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{+})\\}^{\prime\prime}$
(restricted to $\overline{{\mathcal{N}}_{S}\Omega}$),
* •
$T(t,x):=T_{0}(\frac{t-x}{\sqrt{2}})\otimes T_{0}(\frac{t+x}{\sqrt{2}})$,
(restricted to $\overline{{\mathcal{N}}_{S}\Omega}$)
* •
$\Omega:=\Omega_{0}\otimes\Omega_{0}$
is an asymptotically complete Borchers triple with asymptotic algebra
${\mathcal{A}}_{0}^{G}\otimes{\mathcal{A}}_{0}^{G}$ and scattering operator
$S|_{\overline{{\mathcal{N}}_{s}\Omega}}$.
###### Proof.
As remarked above, $({\mathcal{N}}_{S},T,\Omega)$ is a Borchers triple on
${\mathcal{H}}_{{\mathcal{N}}_{S}}$, hence the only thing to be proven is
asymptotic completeness. We show that the subspace
$\overline{{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-})\otimes{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{+})\Omega}$
is invariant under $S$.
We claim that $\overline{{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-})\Omega_{0}}$
coincides with the subspace ${\mathcal{H}}_{0}^{G}$ of invariant vectors under
$\\{V_{g}\\}_{g\in G}$. Indeed, for any $x\in{\mathcal{A}}_{0}$, the averaging
$\int_{g}V_{g}x\Omega_{0}\,\mathrm{d}g=\left(\int_{g}\alpha_{g}(x)\,\mathrm{d}g\right)\Omega_{0}$
gives a projection onto ${\mathcal{H}}_{0}^{G}$. By the Reeh-Schlieder
property, any vector in ${\mathcal{H}}_{0}^{G}$ can be approximated by vectors
in ${\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-})\Omega_{0}$. The converse inclusion
is obvious.
Now it is easy to see that
$\overline{{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-})\otimes{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{+})\Omega}={\mathcal{H}}_{0}^{G}\otimes{\mathcal{H}}_{0}^{G}$.
This is the space of invariant vectors under the action $\\{V_{g}\otimes
V_{g^{\prime}}:g,g^{\prime}\in G\\}$. Since $S$ commutes with $V_{g}\otimes
V_{g^{\prime}}$ by assumption, this subspace is preserved under $S$. Then, as
remarked before, $\overline{{\mathcal{N}}_{S}\Omega}$ coincides with
$\overline{{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{-})\otimes{\mathcal{A}}_{0}^{G}({\mathbb{R}}_{+})\Omega}$
and we obtain the asymptotic completeness.
The statement on S-matrix is immediate from the definition and by Proposition
2.2. ∎
## 3 Examples of fermi nets
### 3.1 ${\rm U(1)}$-current net ${{\mathcal{A}}^{(0)}}$
Let $U_{1}$ be the irreducible unitary positive-energy representation of
$\operatorname{M\ddot{o}b}$ with lowest weight $1$ on a Hilbert space denoted
by ${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$, which can be
identified with the one-particle space of the ${\rm U(1)}$-current. This has
the following concrete realization: consider $C^{\infty}(S^{1},{\mathbb{R}})$,
where we write the periodic function $f\in C^{\infty}(S^{1},{\mathbb{R}})$ as
a Fourier series
$f(\theta)=\sum_{k\in{\mathbb{Z}}}\hat{f}_{k}\mathrm{e}^{\mathrm{i}k\theta},\quad\hat{f}_{k}=\int_{0}^{2\pi}\mathrm{e}^{-\mathrm{i}k\theta}f(\theta)\frac{\,\mathrm{d}\theta}{2\pi}=\overline{\hat{f}_{-k}}\,.$
We introduce a semi-norm
$\|f\|^{2}=\sum_{k=1}^{\infty}k\cdot|\hat{f}_{k}|^{2}$
and a complex structure, i.e. an isometry $\mathcal{J}$ w.r.t. $\|\,\cdot\,\|$
satisfying $\mathcal{J}^{2}=-1$, by
$\mathcal{J}:\hat{f}_{k}\mapsto-\mathrm{i}\operatorname{sign}(k)\hat{f}_{k}$
and finally we get the Hilbert space
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}=\overline{C^{\infty}(S^{1},{\mathbb{R}})/{\mathbb{R}}}^{\|\,\cdot\,\|}$
by completion with respect to the norm $\|\,\cdot\,\|$, where ${\mathbb{R}}$
is identified with the constant functions. By abuse of notation we denote also
the image of $f\in C^{\infty}(S^{1},{\mathbb{R}})$ in
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ by $f$. The scalar
product (linear in the second component) and the sesquilinear form
$\omega(\,\cdot\,,\,\cdot\,)\equiv\Im\langle\,\cdot\,,\,\cdot\,\rangle$ are
given by
$\displaystyle\langle f,g\rangle$
$\displaystyle=\sum_{k=1}^{\infty}k\hat{f}_{k}\hat{g}_{-k}\,,$
$\displaystyle\omega(f,g)$
$\displaystyle=\frac{-\mathrm{i}}{2}\sum_{k\in{\mathbb{Z}}}k\hat{f}_{k}\hat{g}_{-k}=\frac{1}{2}\int_{0}^{2\pi}f(\theta)g^{\prime}(\theta)\frac{\,\mathrm{d}\theta}{2\pi}=\frac{1}{4\pi}\int
f\,\mathrm{d}g\,,$
respectively. The unitary action $U_{1}$ of $\operatorname{M\ddot{o}b}$ on
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ is induced by the action
on $C^{\infty}(S^{1},{\mathbb{R}})$
$(U_{1}(g)f)=(g_{\ast}f)(\theta):=f(g^{-1}(\theta))$.
For $I\in\mathcal{I}$ we denote by $H(I)$ the closure of the subspace of real
functions with support in $I$. This space is standard (i.e.
$H(I)+\mathrm{i}H(I)$ is dense in
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ and
$H(I)\cap\mathrm{i}H(I)=\\{0\\}$) and the family
$\\{H(I)\\}_{I\in\mathcal{I}}$ is a local Möbius covariant net of standard
subspaces [Lon08, LW11].
We explain briefly the bosonic second quantization procedure in general. Let
${\mathcal{H}}^{1}$ be a separable Hilbert space, the one-particle space, and
$\omega(\,\cdot\,,\,\cdot\,)=\Im\langle\,\cdot\,,\,\cdot\,\rangle$ the
sesquilinear form. There are unitaries $W(f)$ for $f\in{\mathcal{H}}^{1}$
fulfilling
$W(f)W(g)=\mathrm{e}^{-\mathrm{i}\omega(f,g)}W(f+g)=\mathrm{e}^{-2\mathrm{i}\omega(f,g)}W(g)W(f)$
and acting naturally on the bosonic Fock space
$\mathrm{e}^{{\mathcal{H}}^{1}}$ over ${\mathcal{H}}^{1}$. This space is given
by
$\mathrm{e}^{{\mathcal{H}}^{1}}=\oplus_{n=0}^{\infty}P_{n}({\mathcal{H}}^{1})^{\otimes
n}$, where $P_{n}$ is the projection
$P_{n}(\xi_{1}\otimes\cdots\otimes\xi_{n})=1/n!\sum_{\sigma}\xi_{\sigma(1)}\otimes\cdots\otimes\xi_{\sigma(n)}$
and the sum goes over all permutations. The set of coherent vectors
$\mathrm{e}^{h}:=\oplus_{n=0}^{\infty}h^{\otimes n}/\sqrt{n!}$ with
$h\in{\mathcal{H}}^{1}$ is total in $\mathrm{e}^{{\mathcal{H}}^{1}}$ and it
holds $\langle\mathrm{e}^{f},\mathrm{e}^{h}\rangle=\mathrm{e}^{\langle
f,h\rangle}$. The vacuum is given by $\Omega=\mathrm{e}^{0}$ and the action of
$W(f)$ is given by
$W(f)\mathrm{e}^{0}=\mathrm{e}^{-\frac{1}{2}\|f\|^{2}}\mathrm{e}^{f}$, in
other words the vacuum representation is characterized by
$\phi(W(f))=\mathrm{e}^{-\frac{1}{2}\|f\|^{2}}$, where
$\phi(\,\cdot\,)=\langle\Omega,\,\cdot\,\Omega\rangle$. For a real subspace
$H\subset{\mathcal{H}}^{1}$, we define the von Neumann algebra
$\displaystyle R(H)=\\{W(f):f\in H\\}^{\prime\prime}\subset
B(\mathrm{e}^{{\mathcal{H}}^{1}})\,.$
Let $U$ be a unitary on the one-particle space ${\mathcal{H}}^{1}$ then
$\mathrm{e}^{U}:=\oplus_{n=0}^{\infty}U^{\otimes n}$ acts on coherent states
by $\mathrm{e}^{U}\mathrm{e}^{h}=\mathrm{e}^{Uh}$ and is therefore a unitary
on $\mathrm{e}^{{\mathcal{H}}^{1}}$, the second quantization unitary.
We obtain the ${\rm U(1)}$-current net ${{\mathcal{A}}^{(0)}}$ on
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}:=\mathrm{e}^{{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}}$
with $\Omega_{0}=\mathrm{e}^{0}$ by defining
${{\mathcal{A}}^{(0)}}(I):=R(H(I))$ which is covariant with respect to
$U(g):=\mathrm{e}^{U_{1}(g)}$. For $f\in C^{\infty}(S^{1},{\mathbb{R}})$ we
consider a self-adjoint operator $J(f)$ given by the generator of the unitary
one-parameter group $W(t\cdot f)=\mathrm{e}^{\mathrm{i}t\cdot J(f)}$ with
$t\in{\mathbb{R}}$. This defines the usual current (field operator) smeared
with the real test function $f$, which fulfills
$J(f)\Omega_{0}=f\in{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ and
$\displaystyle[J(f),J(g)]$
$\displaystyle=2\mathrm{i}\omega(f,g)=\sum_{k}k\hat{f}_{k}\hat{g}_{-k}=\frac{\mathrm{i}}{2\pi}\int
f\,\mathrm{d}g\,.$
It can be extended to complex test functions via
$J(f+\mathrm{i}g)=J(f)+\mathrm{i}J(g)$, and one obtains the usual operator
valued ($z$-picture) distribution $J(z)$ with the relations
$\displaystyle J(f)$
$\displaystyle=\sum_{n\in{\mathbb{Z}}}\hat{f}_{n}J_{n}=\oint_{S^{1}}f(z)J(z)\frac{\,\mathrm{d}z}{2\pi\mathrm{i}},$
$\displaystyle J(z)=\sum_{n}J_{n}z^{-n-1}$ $\displaystyle[J_{m},J_{n}]$
$\displaystyle=m\delta_{m+n,0}\,,$
where the modes $J_{n}=J(e_{n})$ with
$e_{n}(\theta)=\mathrm{e}^{\mathrm{i}n\theta}$ satisfy $J_{n}\Omega_{0}=0$ for
$n\geq 0$.
The space ${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}$ is spanned by
vectors of the form $\xi=J_{-n_{1}}\cdots J_{-n_{k}}\Omega_{0}$ with
$0<n_{1}\leq\cdots\leq n_{k}$ with “energy” $N=\sum_{m}n_{m}$, i.e.
$R(\theta)\xi=\mathrm{e}^{\mathrm{i}N\theta}\xi$. Therefore it is graded with
respect to the rotations
$\displaystyle{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}$
$\displaystyle={\mathbb{C}}\Omega_{0}\oplus\bigoplus_{n\in{\mathbb{N}}}{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}},n}$
$\displaystyle{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}},n}$
$\displaystyle=\bigoplus_{k=1}^{n}\bigoplus_{\begin{subarray}{c}0<n_{1}\leq\cdots\leq
n_{k}\\\ n_{1}+\cdots+n_{k}=n\end{subarray}}{\mathbb{C}}J_{-n_{1}}\cdots
J_{-n_{k}}\Omega_{0}$
and ${\hbox{dim}\,}{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}},n}$ is the
number of partitions of $n$ elements, whose generating function $p(t)$ is the
inverse of Euler’s function $\phi(t)=\prod_{k=1}^{\infty}(1-t^{k})$ and
therefore the conformal character of the ${\rm U(1)}$-current net is given by
($t=\mathrm{e}^{-\beta}$):
$\operatorname{tr}_{{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}}(\mathrm{e}^{-\beta
L_{0}})=\sum_{n=0}^{\infty}{\hbox{dim}\,}{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}},n}\cdot
t^{n}=\prod_{n\in{\mathbb{N}}}(1-t^{n})^{-1}$
(a conformal character is defined as a formal power series, but it is often
convergent for $|t|<1$ and here we used the formula
$(1-z)^{-1}=1+z+z^{2}\cdots$). It will be convenient to use the real
parametrization $x\in{\mathbb{R}}\cong S^{1}\setminus\\{-1\\}$ of the cut
circle and use the conventions
$\displaystyle
f(s)=\int_{\mathbb{R}}\mathrm{e}^{-\mathrm{i}sp}\hat{f}(p)\,\mathrm{d}p.$
By writing $f(s)=f_{0}(\theta(s))$ for $f_{0}\in
C^{\infty}(S^{1},{\mathbb{R}})$ where $\theta(s)=2\arctan(s)$, the space
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ above can be identified
with the space $L^{2}({\mathbb{R}}_{+},p\,\mathrm{d}p)$ in which the space
${\mathcal{S}}({\mathbb{R}},{\mathbb{R}})$ embeds by restriction of the
Fourier transformation to ${\mathbb{R}}_{+}$. In other words
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ can be seen as the
closure of the space ${\mathcal{S}}({\mathbb{R}},{\mathbb{R}})$ with complex
structure $\mathcal{J}\hat{f}(p)=\mathrm{i}\operatorname{sign}(p)\hat{f}(p)$
and the scalar product and sesquilinear form given by:
$\displaystyle\langle f,g\rangle$
$\displaystyle=\int_{{\mathbb{R}}_{+}}\hat{f}(-p)\hat{g}(p)p\,\mathrm{d}p\,,$
$\displaystyle\omega(f,g)$
$\displaystyle=\frac{-\mathrm{i}}{2}\int_{{\mathbb{R}}}\hat{f}(-p)\hat{g}(p)p\,\mathrm{d}p=\frac{1}{4\pi}\int_{\mathbb{R}}f(x)g^{\prime}(x)\,\mathrm{d}x\,.$
Using the above identification we denote for
$f\in{\mathcal{S}}({\mathbb{R}},{\mathbb{R}})$ by $J(f)$ the smeared current
with $J(f)\Omega_{0}=f\in{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$.
In this parametrization commutation relations read:
$\displaystyle[J(f),J(g)]$
$\displaystyle=\frac{\mathrm{i}}{2\pi}\int_{\mathbb{R}}f(x)g^{\prime}(x)\,\mathrm{d}x=\int_{{\mathbb{R}}}\hat{f}(-p)\hat{g}(p)p\,\mathrm{d}p\,.$
### 3.2 The free complex fermion net $\mathrm{Fer}_{\mathbb{C}}$
We construct the net of the free complex fermion on the circle, which can be
seen as the chiral part of the net of the free massless Dirac (or complex)
fermion on two dimensional Minkowski space. The notations of this section are
basically in accordance with [Was98], but we use a different convention of
positive-energy, which leads to the conjugated complex structure. For giving a
simple description of the one-particle space, we consider first the Hilbert
space $L^{2}(S^{1})$ and the Hardy space $H^{2}(S^{1})$, namely
$H^{2}(S^{1}):=\left\\{f:\mbox{ analytic on the unit disk }D,\sup_{0\leq
r<1}\int_{0}^{2\pi}|f(r\mathrm{e}^{\mathrm{i}\theta})|^{2}\,\mathrm{d}\theta<\infty\right\\}\,.$
Any function in $H^{2}(S^{1})$ has a $L^{2}$-boundary value and can be
considered as an element of $L^{2}(S^{1})$. In this sense, $H^{2}(S^{1})$ is a
subspace of $L^{2}(S^{1})$. Furthermore, it holds that
$H^{2}(S^{1})=\\{f\in L^{2}(S^{1}):\hat{f}_{n}=0\mbox{ for }n<0\\},$
where $\hat{f}_{n}$ is the $n$-th Fourier component of $f$. We denote the
orthogonal projection onto $H^{2}(S^{1})$ by $P$.
The group
$\mathrm{SU}(1,1):=\left\\{\left(\begin{matrix}\alpha&\beta\\\
\overline{\beta}&\overline{\alpha}\end{matrix}\right)\in
M_{2}({\mathbb{C}}):|\alpha|^{2}-|\beta|^{2}=1\right\\}$
acts on the circle $S^{1}$ by $g\cdot z=\frac{\alpha
z+\beta}{\overline{\beta}z+\overline{\alpha}}$ and there is a unitary action
of $\mathrm{SU}(1,1)$ on $L^{2}(S^{1})$ by
$(U(g)f)(z):=(V_{g}f)(z)=\frac{1}{-\overline{\beta}z+\alpha}f(g^{-1}\cdot
z)\,.$
One sees that the projection $P$ commutes with $V_{g}$ , since $V_{g}f$ is
still an analytic function for $|\alpha|>|\beta|$.
Then one defines a new Hilbert space
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}=\overline{PL^{2}(S^{1})}\oplus({\mathbbm{1}}-P)L^{2}(S^{1})$:
namely, ${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$ is identical with
$L^{2}(S^{1})$ as a real linear space and the multiplication by $\mathrm{i}$
is given by $-\mathrm{i}(2P-{\mathbbm{1}})$, or in other words, by
$-\mathrm{i}$ on $PL^{2}(S^{1})$ and $\mathrm{i}$ on
$({\mathbbm{1}}-P)L^{2}(S^{1})$. Because $P$ and $U(g)$ commute, the action of
$\mathrm{SU}(1,1)$ remains unitary on
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$.
Then for $I\in\mathcal{I}$ one takes real Hilbert subspaces $K(I):=L^{2}(I)$
of ${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$. This subspaces turn out
to be standard [Was98, Theorem (p. 497)]. If $I_{1}$ and $I_{2}$ are disjoint
intervals, $K(I_{1})$ are real orthogonal to $K(I_{2})$, in other words
$K(I_{1})\subset K(I_{2})^{\perp}$, where
$K^{\perp}=\\{\xi\in{\mathcal{H}}:\Re\langle\xi,K\rangle=0\\}$. It turns out
that $I\mapsto K(I)$ is a twisted-local Möbius covariant net of standard
subspaces.
We briefly explain the fermionic second quantization in general. Let
${\mathcal{H}}^{1}$ be a complex Hilbert space and
${\mathcal{H}}=\Lambda({\mathcal{H}}^{1})$ the antisymmetric (fermionic) Fock
space obtained by completing the exterior algebra with the inner product. For
$A\in B({\mathcal{H}}^{1})$ with $\|A\|\leq 1$ we define $\Lambda(A)$ to be
$A^{\otimes k}$ on
${\mathcal{H}}^{k}:=\Lambda^{k}({\mathcal{H}}^{1})\subset({\mathcal{H}}^{1})^{\otimes
k}$. The space is ${\mathbb{Z}}_{2}$ graded by
$\Gamma:=\Lambda(-{\mathbbm{1}})$. We define
$Z=\frac{{\mathbbm{1}}-\mathrm{i}\Gamma}{1-\mathrm{i}}$ and note that
$Z^{2}=\Gamma$. For $f\in{\mathcal{H}}^{1}$ let $a(f)$ be the bounded operator
obtained by continuing the exterior multiplication $f\wedge\cdot$. The
operators fulfill the complex Clifford relations
$a(f)^{\ast}a(g)+a(g)a(f)^{\ast}=\langle f,g\rangle$ and
$\\{a(f),a(g)\\}=\\{a(f)^{\ast},a(g)^{\ast}\\}=0$ for all
$f,g\in{\mathcal{H}}^{1}$. For a standard subspace $K\subset{\mathcal{H}}^{1}$
we define the von Neumann algebra
$C(K)=\\{c(f):f\in K\\}^{\prime\prime}\subset B(\Lambda{\mathcal{H}}^{1})$
where $c(f)=a(f)+a(f)^{\ast}$, which fulfills the real Clifford relations
$c(f)c(g)+c(g)c(f)=2\Re\langle f,g\rangle$. By $\Omega=1\in\Lambda^{0}$ we
denote the vacuum which is cyclic and separating for $C(K)$ for every standard
subspace $K\subset{\mathcal{H}}^{1}$. Further it holds Haag-Araki duality,
i.e. $C(K^{\perp})$ equals $C(K)^{\perp}:=ZC(K)^{\prime}Z^{\ast}$, the twisted
commutant of $C(K)$. For a unitary $U$ on ${\mathcal{H}}^{1}$ it holds
$\Lambda(U)c(f)\Lambda(U^{\ast})=c(Uf)$, which implies that $C$ is covariant
with respect to the unitaries $U({\mathcal{H}}^{1})$, i.e.
$\Lambda(U)C(K)\Lambda(U)^{\ast}=C(UK)$.
We note that in the case like the complex fermion the one-particle space is
obtained from a Hilbert space ${\mathcal{H}}^{1}$ (the space of test
functions) and a projection $P$ by
${\mathcal{H}}^{1}_{P}=P{\mathcal{H}}^{1}\oplus\overline{P^{\perp}{\mathcal{H}}^{1}}$
and one gets a new representation of the complex Clifford algebra on
$\Lambda({\mathcal{H}}^{1}_{P})$ by
$a_{P}(f)=a(Pf)+a(\overline{P^{\perp}f})^{\ast}$ where $a(f)$ is the creation
operator. For a standard subspace $K\subset{\mathcal{H}}^{1}_{P}$ which is
invariant under the multiplication of $\mathrm{i}_{{\mathcal{H}}^{1}}$ in
${\mathcal{H}}^{1}$, the von Neumann algebra $C(K)$ on
$\Lambda({\mathcal{H}}_{P}^{1})$ coincides with the von Neumann algebra
$\\{a_{P}(f),a_{P}(f)^{\ast}:f\in K\\}^{\prime\prime}$. Indeed, the one
inclusion follows from $c(f)=a_{P}(f)+a_{P}(f)^{\ast}$ and the other follows
from Araki-Haag duality and $\\{a_{P}(f),c(g)\\}=\langle
g,f\rangle_{{\mathcal{H}}^{1}}=\Re\langle
g,f\rangle_{{\mathcal{H}}^{1}_{P}}-\mathrm{i}\Re\langle
g,\mathrm{i}_{{\mathcal{H}}^{1}}f\rangle_{{\mathcal{H}}^{1}_{P}}=0$ for $f\in
K$ and $g\in K^{\perp}$. We further note that the space
$\Lambda({\mathcal{H}}^{1}_{P})$ is as a real Hilbert space the same as
$\Lambda({\mathcal{H}}^{1})$ and can be identified canonically with
$\Lambda(P{\mathcal{H}}^{1})\otimes\Lambda(\overline{P^{\perp}{\mathcal{H}}^{1}})$.
We turn to the concrete case where ${\mathcal{H}}^{1}=L^{2}(S^{1})$ and define
the net
$\mathrm{Fer}_{\mathbb{C}}(I):=C(K(I))=\\{a_{P}(f),a_{P}(f)^{\ast}:f\in
L^{2}(I)\\}^{\prime\prime}$ (where here
$a_{P}(f):=a(\overline{Pf})+a(P^{\perp}f)^{\ast}$) on
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}=\Lambda({\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1})\cong\Lambda(\overline{PL^{2}(S^{1})})\otimes\Lambda(P^{\perp}L^{2}(S^{1}))$
which is isotonic by definition and fulfills twisted duality, namely by Haag-
Araki duality
$\mathrm{Fer}_{\mathbb{C}}(I^{\prime})=C(K(I)^{\perp})=C(K(I))^{\perp}=\mathrm{Fer}_{\mathbb{C}}(I)^{\perp}$.
In addition, the net $\mathrm{Fer}_{\mathbb{C}}$ is Möbius covariant. Indeed,
we can take the representation $\Lambda U(\cdot)$ by promoting the one-
particle representation $U$ to the second quantization operator. It is easy to
see that the covariance of this net $\mathrm{Fer}_{\mathbb{C}}$ follows from
the covariance of the net of standard spaces $K$. The representation $\Lambda
U$ has positive energy since so does the representation $U$, and leaves
invariant the vacuum vector $\Omega_{0}$ of the Fock space. Summing up, the
net $\mathrm{Fer}_{\mathbb{C}}$ is a fermi net (cf. [Was98]). This net is
referred to as the free complex fermi net on $S^{1}$. The scalar
multiplication by a constant phase $\mathrm{e}^{-\mathrm{i}\vartheta}$ in the
original structure of the one-particle space is still a unitary operator in
the new structure. Its promotion by the second quantization $V(\theta)$
implements an action of ${\rm U(1)}$ on $\mathrm{Fer}_{\mathbb{C}}$ by inner
symmetry. This will be referred to as the ${\rm U(1)}$-gauge action.
For $r\in\frac{1}{2}+{\mathbb{Z}}$ let $\psi_{r}=a_{P}(e_{-r-\frac{1}{2}})$
and $\bar{\psi}_{r}=a_{P}(e_{r-\frac{1}{2}})^{\ast}$ where $e_{r}\in
L^{2}(S^{1})$ with $e_{r}(\theta)=\mathrm{e}^{\mathrm{i}\theta r}$. The
$\psi_{r},\bar{\psi}_{r}$ are the modes of the free complex fermion, namely
$\displaystyle\\{\psi_{n},\psi_{m}\\}$
$\displaystyle=\\{\bar{\psi}_{m},\bar{\psi}_{n}\\}=0$
$\displaystyle\\{\bar{\psi}_{n},\psi_{m}\\}$ $\displaystyle=\delta_{m+n,0}$
$\displaystyle\psi_{n}^{\ast}$ $\displaystyle=\bar{\psi}_{-n}$
and it holds that $\psi_{r}\Omega_{0}=\bar{\psi}_{r}\Omega_{0}=0$ for
$r\in\frac{1}{2}+{\mathbb{N}}_{0}$. Each of $\psi_{r}$ or $\bar{\psi}_{r}$ has
norm $1$ following from the commutation relation. We can introduce the usual
fields ($f,g\in L^{2}(S^{1})$) and operator valued distributions in the
$z$-picture:
$\displaystyle\Psi(f)$
$\displaystyle=\sum_{r\in\frac{1}{2}+{\mathbb{Z}}}\hat{f}_{r}\Psi_{r}=\oint_{S^{1}}f(z)z^{-\frac{1}{2}}\Psi(z)\frac{\,\mathrm{d}z}{2\pi\mathrm{i}}\,,$
$\displaystyle\Psi(z)$
$\displaystyle=\sum_{r\in\frac{1}{2}+{\mathbb{Z}}}\Psi_{r}z^{-r-\frac{1}{2}}\,,$
$\displaystyle\bar{\psi}(f)$
$\displaystyle=\psi(\overline{f})^{\ast}=a_{P}(e_{-\frac{1}{2}}f)^{\ast}\,,$
$\displaystyle\\{\bar{\psi}(f),\psi(g)\\}$ $\displaystyle=\oint
f(z)g(z)\frac{\,\mathrm{d}z}{2\pi\mathrm{i}z}\,,$
where $\Psi$ is either $\psi$ or $\bar{\psi}$. The fields $\psi,\bar{\psi}$
are covariant, e.g. $U(g)\psi(f)U(g)^{\ast}=\psi(f_{g})$ with
$f_{g}(z)=\frac{1}{|\alpha-\overline{\beta}z|}f(g^{-1}z)$ for
$g=\left(\begin{matrix}\alpha&\beta\\\
\overline{\beta}&\overline{\alpha}\end{matrix}\right)\in\mathrm{SU(1,1)}$.
We note that vectors of the form
$\xi=\psi_{-r_{1}}\cdots\psi_{-r_{k}}\bar{\psi}_{-s_{1}}\cdots\bar{\psi}_{-s_{\ell}}\Omega_{0}$
with $0<r_{1}<\cdots<r_{k}$ and $0<s_{1}<\cdots<s_{\ell}$ form a basis of
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}=\Lambda({\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1})$
and that such a $\xi$ is an eigenvector for the rotations,
$R(\theta)\xi\equiv\mathrm{e}^{\mathrm{i}\theta
L_{0}}\xi=\mathrm{e}^{\mathrm{i}N\theta}\xi$ with
$N=\sum_{j=1}^{k}r_{j}+\sum_{j=1}^{\ell}s_{j}$ and of the gauge action
$V(\theta)\xi\equiv\mathrm{e}^{\mathrm{i}\theta
Q}\xi=\mathrm{e}^{\mathrm{i}(k-\ell)\theta}\xi$. In each vector of this basis
the $r$-th energy level can either be empty, be occupied by $\psi_{-r}$ or
$\bar{\psi}_{-r}$ or occupied by both. The contribution of this level to the
character
$\operatorname{tr}_{{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}}(\mathrm{e}^{-\beta
L_{0}-EQ})$ is then $1$, $zt^{r}$, $z^{-1}t^{r}$ or $t^{2r}$, respectively,
where $t=\mathrm{e}^{-\beta}$ and $z=\mathrm{e}^{-E}$. By summing over all
possibilities one gets that the character of $\mathrm{Fer}_{\mathbb{C}}$ is
given by (cf. [Kac98, Reh98]):
$\displaystyle\operatorname{tr}_{{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}}(\mathrm{e}^{-\beta
L_{0}-EQ})=\operatorname{tr}_{{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}}(t^{L_{0}}z^{Q})$
$\displaystyle=\prod_{r\in{\mathbb{N}}_{0}+\frac{1}{2}}(1+zt^{r}+z^{-1}t^{r}+t^{2r})$
$\displaystyle=\prod_{r\in{\mathbb{N}}_{0}+\frac{1}{2}}(1+zt^{r})(1+z^{-1}t^{r})$
$\displaystyle=p(t)\sum_{q\in{\mathbb{Z}}}z^{q}t^{\frac{q^{2}}{2}}\,,$
where the last equality follows directly from the Jacobi triple product
formula (see [Apo76, Theorem 14.6])
$\displaystyle\prod_{r\in{\mathbb{N}}}(1+zw^{2r-1})(1+z^{-1}w^{2r-1})(1-w^{2r})=\sum_{q\in{\mathbb{Z}}}z^{q}w^{q}$
by setting $2r-1=2n$ and $t=w^{2}$. In particular, for the local net
$\mathrm{Fer}_{\mathbb{C}}^{{\rm U(1)}}$ the character is given by
$\operatorname{tr}_{{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{{\rm
U(1)}}}(\mathrm{e}^{-\beta L_{0}})=p(t)$, since it is the fixed point with
respect to the ${\rm U(1)}$-gauge action and the conformal character is the
coefficient of $z^{0}$.
### 3.3 ${\rm U(1)}$-current net as a subnet of $\mathrm{Fer}_{\mathbb{C}}$
In this section we use the well-known fact that the Wick product
$:\\!\bar{\psi}\psi\\!:$ of the complex fermion $\psi$ equals the ${\rm
U(1)}$-current and give an analogue of the boson-fermion correspondence (see
e.g. [Kac98, 5.2]) in the operator algebraic setting. Let us denote by
${\mathcal{D}}_{0}$ the subspace of $\Lambda({\mathcal{H}}_{P}^{1})$ of
vectors with finite energy:
${\mathcal{D}}_{0}:=\mathrm{span}\left\\{\psi_{-r_{1}}\cdots\psi_{-r_{k}}\bar{\psi}_{-s_{1}}\cdots\bar{\psi}_{-s_{l}}\Omega_{0}:k,l\in{\mathbb{N}}_{0},r_{i},s_{j}\in{\mathbb{N}}+\frac{1}{2}\right\\}\,.$
Then we define the unbounded operators on the domain ${\mathcal{D}}_{0}$:
$\displaystyle J_{n}=\sum_{r+s=n}{:\\!\bar{\psi}_{r}\psi_{s}\\!:}$
$\displaystyle=$
$\displaystyle\sum_{r<0}\bar{\psi}_{r}\psi_{n-r}-\sum_{r>0}\psi_{n-r}\bar{\psi}_{r}$
$\displaystyle=$
$\displaystyle\sum_{r}\left(\bar{\psi}_{r}\psi_{n-r}-\langle\Omega_{0},\bar{\psi}_{r}\psi_{n-r}\Omega_{0}\rangle\right)$
with $r,s\in\frac{1}{2}+{\mathbb{Z}}$. Note that any vector in
${\mathcal{D}}_{0}$ is annihilated by $\psi_{r}$ for sufficiently large $r$,
thus the action of $J_{n}$ on such a vector can be defined and remains in
${\mathcal{D}}_{0}$. In particular, we have $J_{n}\Omega_{0}=0$ for
$n\in{\mathbb{N}}_{0}$.
###### Lemma 3.1.
On ${\mathcal{D}}_{0}$ it holds that
1. 1.
$[J_{n},\psi_{k}]=-\psi_{n+k}$ and $[J_{n},\bar{\psi}_{k}]=\bar{\psi}_{n+k}$
2. 2.
$[J_{m},J_{n}]=m\delta_{m+n,0}$
###### Proof.
Using $[ab,c]=a\\{b,c\\}-\\{a,c\\}b$, one obtains
$[\bar{\psi}_{r}\psi_{n},\psi_{k}]=-\delta_{r+k,0}\psi_{n}$ and
$[\psi_{n}\bar{\psi}_{r},\psi_{k}]=\delta_{r+k,0}\psi_{n}$ from which directly
follows
$[J_{n},\psi_{k}]=\sum_{r<0}[\bar{\psi}_{r}\psi_{n-r},\psi_{k}]-\sum_{r>0}[\psi_{n-m}\bar{\psi}_{r},\psi_{k}]=-\psi_{n+k}$.
Analogously one shows $[J_{n},\bar{\psi}_{k}]=\bar{\psi}_{n+k}$.
From the Jacobi identity, it follows immediately that $[J_{n},J_{m}]$ commutes
with all $\psi_{k}$ and $\bar{\psi}_{k}$ and hence $[J_{n},J_{m}]$ is a
multiple of the identity, therefore
$[J_{n},J_{m}]=\langle\Omega_{0},[J_{n},J_{m}]\Omega_{0}\rangle{\mathbbm{1}}$.
It is
$\displaystyle[J_{n},J_{p}]$
$\displaystyle=\sum_{r<0}[J_{n},\bar{\psi}_{r}\psi_{p-r}]-\sum_{r>0}[J_{n},\psi_{p-r}\bar{\psi}_{r}]$
$\displaystyle=-\sum_{r<0}\left(\bar{\psi}_{r}\psi_{p-r+n}-\bar{\psi}_{r+n}\psi_{p-r}\right)-\sum_{r>0}\left(\psi_{p-r}\bar{\psi}_{r+n}-\psi_{p-r+n}\bar{\psi}_{r}\right)$
and in the case $p\neq-n$ we get
$\langle\Omega_{0},[J_{n},J_{p}]\Omega_{0}\rangle=0$, and otherwise
$\displaystyle\langle\Omega_{0},[J_{n},J_{-n}]\Omega_{0}\rangle$
$\displaystyle=\begin{cases}\sum_{r<0}\langle\Omega_{0},\bar{\psi}_{r+n}\psi_{-r-n}\Omega_{0}\rangle=\sum_{r=\frac{1}{2}}^{n-\frac{1}{2}}\langle\Omega_{0},\\{\bar{\psi}_{r},\psi_{-r}\\}\Omega_{0}\rangle&n>0\\\
-\sum_{r>0}\langle\Omega_{0},\psi_{-r-n}\bar{\psi}_{r+n}\Omega_{0}\rangle=-\sum_{r=\frac{1}{2}}^{-n-\frac{1}{2}}\langle\Omega_{0},\\{\psi_{r},\bar{\psi}_{-r}\\}\Omega_{0}\rangle&n<0\end{cases}$
$\displaystyle=n\,,$
which completes the proof. ∎
Let $L_{0}$ be the generator of the rotation:
$R(\theta)=\mathrm{e}^{\mathrm{i}\theta L_{0}}$. From its action (see the end
of Section 3.2), one verifies that ${\mathcal{D}}_{0}$ is a core for $L_{0}$.
###### Lemma 3.2 (Linear energy bounds).
It holds that $[L_{0},J_{n}]=-nJ_{n}$ on ${\mathcal{D}}_{0}$. For a
trigonometric polynomial $f=\sum_{n}\hat{f}_{n}e_{n}$ where the sum is finite
and $\xi\in{\mathcal{D}}_{0}$, we have
$\displaystyle\|J(f)\xi\|$ $\displaystyle\leq c_{f}\|(L_{0}+1)\xi\|$
$\displaystyle\|[L_{0},J(f)]\xi\|$ $\displaystyle\leq
c_{\partial_{\theta}f}\|(L_{0}+1)\xi\|,$
where $c_{f}$ depends only on $f$.
###### Proof.
For the commutation relation, it is enough to choose an energy eigenvector
$\xi\in{\mathcal{D}}_{0}$, i.e. $L_{0}\xi=N\xi$. It is
$J_{n}L_{0}\xi=NJ_{n}\xi$ and
$L_{0}J_{n}\xi=L_{0}\left(\sum_{r<0}\bar{\psi}_{r}\psi_{n-r}\xi-\sum_{r>0}\psi_{n-r}\bar{\psi}_{r}\xi\right)=(N-n)J_{n}\xi,$
and the first statement follows.
We have seen that $\psi_{r}$ and $\bar{\psi}_{r}$ have norm $1$ in Section
3.2. First we claim that $\|J_{n}\xi\|\leq\|(2(L_{0}+1)+|n|)\xi\|$. Let $\xi$
be again an eigenvector of $L_{0}$, i.e. $L_{0}\xi=N\xi$. From the defining
sum of $J_{n}$, one sees that only $2N+|n|+2$ terms contribute to $J_{n}\xi$.
Hence we have $\|J_{n}\xi\|\leq(2N+|n|+2)\|\xi\|=\|2(L_{0}+1)+|n|\xi\|$. If
the inequality holds for eigenvectors, then for $\\{\xi_{r}\\}$ with different
eigenvalues, we have $\xi_{r}\perp\xi_{s}$ and $J_{n}\xi_{r}\perp
J_{n}\xi_{s}$, and hence
$\displaystyle\left\|J_{n}\sum_{r}\xi_{r}\right\|^{2}$ $\displaystyle=$
$\displaystyle\sum_{r}\|J_{n}\xi_{r}\|^{2}$ $\displaystyle\leq$
$\displaystyle\sum_{r}\|(2(L_{0}+1)+|n|)\xi_{r}\|^{2}$ $\displaystyle=$
$\displaystyle\left\|(2(L_{0}+1)+|n|)\sum_{r}\xi_{r}\right\|^{2}$
and the general case follows.
For a smeared field, we have
$\|J(f)\xi\|=\left\|\sum_{n}\hat{f}_{n}J_{n}\xi\right\|\leq
2\tilde{c}_{f}\|(L_{0}+1)\xi\|+\tilde{c}_{\partial_{\theta}f}\|\xi\|\leq(2\tilde{c}_{f}+\tilde{c}_{\partial_{\theta}f})\|(L_{0}+1)\xi\|,$
where $\tilde{c}_{f}=\sum_{n}|\hat{f}_{n}|$. By defining
$c_{f}=2\tilde{c}_{f}+\tilde{c}_{\partial_{\theta}f}$, we obtain the first
inequality of the statement. The rest follows by noting that
$[L_{0},J(f)]=J(\mathrm{i}\partial_{\theta}f)$. ∎
For a smooth function $f=\sum_{n\in{\mathbb{Z}}}\hat{f}_{n}e_{n}\in
C^{\infty}(S^{1})$, its Fourier coefficients $\hat{f}_{n}$ are strongly
decreasing and, in particular, it is summable:
$\sum_{n}|\hat{f}_{n}|=\tilde{c}_{f}<\infty$. Hence we can naturally extend
the definition of the smeared current to smooth functions using the above
estimate by
$J(f)=\sum_{n\in{\mathbb{Z}}}f_{n}J_{n}=\sum_{r,s\in\frac{1}{2}+{\mathbb{Z}}}f_{r+s}:\\!\psi_{r}\bar{\psi}_{s}\\!:,$
and the same inequality in Lemma 3.2 holds. The operator is closable since we
have $J(f)\subset J(\overline{f})^{*}$ and we still denote the closure by
$J(f)$. We note that from the above definition it follows that $J(f)$ is
obtained by a limit $\sum_{n}:\\!\psi(h_{n})\bar{\psi}(k_{n})\\!:$ with
suitable functions such that $\sum_{n}h_{n}(\theta)k_{n}(\vartheta)\to 2\pi
f(\theta)\delta(\theta-\vartheta)$. This implies covariance of the “field”,
i.e. $U(g)J(f)U(g)^{\ast}=J(f\circ g^{-1})$.
Recall that $\|\psi_{r}\|=1$, hence the smeared field is still bounded:
$\|\psi(g)\|\leq\tilde{c}_{g}$. We claim that, for $f,g\in C^{\infty}(S^{1})$
and $\xi\in{\mathcal{D}}_{0}$, $\psi(g)\xi$ is in the domain of $J(f)$.
Indeed, for a trigonometric polynomial $g$, we have the estimate
$\displaystyle\|J(f)\psi(g)\xi\|$ $\displaystyle\leq$ $\displaystyle
c_{f}\|(L_{0}+1)\psi(g)\xi\|$ $\displaystyle\leq$ $\displaystyle
c_{f}(\tilde{c}_{g}\|\xi\|+\|[L_{0},\psi(g)]\xi+\psi(g)L_{0}\xi\|)$
$\displaystyle\leq$ $\displaystyle
c_{f}(\tilde{c}_{g}(\|\xi\|+\|L_{0}\xi\|)+\tilde{c}_{\partial_{\theta}g}\|\xi\|)\,.$
Then if we have a sequence of trigonometric polynomial $g_{n}$ converging to a
smooth function $g\in C^{\infty}(S^{1})$, the sequence
$\\{J(f)\psi(g_{n})\xi\\}$ is also converging.
###### Lemma 3.3.
For $\xi,\eta\in{\mathcal{D}}_{0}$, it holds that
$\displaystyle[J(f),\psi(g)]\xi$ $\displaystyle=-\psi(f\cdot g)\xi$
$\displaystyle[J(f),\bar{\psi}(g)]\xi$ $\displaystyle=\bar{\psi}(f\cdot g)\xi$
$\displaystyle\langle J(\bar{f})\xi,J(g)\eta\rangle$ $\displaystyle=\langle
J(\bar{g})\xi,J(f)\eta\rangle+2\mathrm{i}\omega(f,g)\langle\xi,\eta\rangle.$
###### Proof.
For trigonometric polynomials $f,g$, the statements can be proved easily from
Lemma 3.1. The general case is shown by approximating first $f$ by
polynomials, then $g$, according to the convergence considered above (as for
the third statement, obviously the order of limits does not matter). ∎
We need the following well-known result [DF77, Theorem 3.1]:
###### Theorem 3.4 (The commutator theorem).
Let $H$ be a positive self-adjoint operator and $A,B$ symmetric operators
defined on a core ${\mathcal{D}}_{0}$ for $(H+{\mathbbm{1}})^{2}$. Assume that
there is a constant $C$ such that
$\displaystyle\|A\xi\|\leq C\|(H+{\mathbbm{1}})\xi\|,\,\,\,\|B\xi\|\leq
C\|(H+{\mathbbm{1}})\xi\|,$ $\displaystyle\|[H,A]\xi\|\leq
C\|(H+{\mathbbm{1}})\xi\|,\,\,\,\|[H,B]\xi\|\leq C\|(H+{\mathbbm{1}})\xi\|,$
$\displaystyle\langle A\xi,B\eta\rangle=\langle B\xi,A\eta\rangle\mbox{ for
any }\xi,\eta\in{\mathcal{D}}_{0}.$
Then $A$ and $B$ are essentially self-adjoint on any core of $H$ and any
bounded functional calculus of $A$ and $B$ commute.
###### Remark 3.5.
In the original literature [DF77], this Theorem is proved under the assumption
of certain operator inequalities. In fact, what is really used in the proof of
commutativity of bounded functions is the norm estimates
$\|A(H+{\mathbbm{1}})^{-1}\|<C,\|[H,A](H+{\mathbbm{1}})^{-1}\|<C$ etc. and
they follow from the assumptions here. The essential self-adjointness of $A$
and $B$ can be proved by [RS75, Theorem X.37]. An analogous application of
this theorem with norm estimates can be found in [BSM90].
By the commutator theorem, we get that $J(f)$ is self-adjoint for $f\in
C^{\infty}(S^{1},{\mathbb{R}})$ and that all bounded functions of $J(f)$
commute with all bounded functions of $J(g)$ for $f,g\in
C^{\infty}(S^{1},{\mathbb{R}})$ with disjoint support.
Let $I$ be a proper interval and let us define the von Neumann algebra
${\mathcal{B}}(I)=\\{\mathrm{e}^{\mathrm{i}J(f)}:{\rm supp}f\subset
I\\}^{\prime\prime}.$
The local net ${\mathcal{B}}(I)$ restricted to
$\overline{{\mathcal{B}}(I)\Omega_{0}}$ can be identified with the ${\rm
U(1)}$-current net ${{\mathcal{A}}^{(0)}}$ on
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}$, in particular we can
identify
$\overline{{\mathcal{B}}(I)\Omega_{0}}\cong{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}$.
###### Proposition 3.6.
Let $I$ be a proper interval, then
${\mathcal{B}}(I)\subset\mathrm{Fer}_{\mathbb{C}}^{{\rm U(1)}}(I)$.
###### Proof.
We see that ${\mathcal{B}}(I)$ commutes with
$\mathrm{Fer}_{\mathbb{C}}(I^{\prime})=\\{c(g):g\in
L^{2}(I^{\prime})\\}^{\prime\prime}$ because, for $f,g$ with disjoint
supports, $c(g)$ commutes with $J(f)$ on a core by Lemma 3.3 and therefore any
spectral projection of $c(g)$ commutes with $J(f)$, and hence with any bounded
functions of $J(f)$.
Further because $J(f)$ commutes by construction with the gauge action $V(t)$
and is in particular even because $V(\pi)=\Gamma$, it follows that
${\mathcal{B}}(I)$ lies in the twisted commutant
$\mathrm{Fer}_{\mathbb{C}}(I^{\prime})^{\perp}$. By twisted Haag duality it is
${\mathcal{B}}(I)\subset\mathrm{Fer}_{\mathbb{C}}(I^{\prime})^{\perp}=\mathrm{Fer}_{\mathbb{C}}(I)$
and therefore ${\mathcal{B}}(I)={\mathcal{B}}(I)^{{\rm
U(1)}}\subset\mathrm{Fer}_{\mathbb{C}}^{{\rm U(1)}}(I)$. ∎
Since the covariance has been seen, we have the following.
###### Corollary 3.7.
${\mathcal{B}}$ is a subnet of $\mathrm{Fer}_{\mathbb{C}}^{{\rm U(1)}}$.
Now the following is straightforward.
###### Proposition 3.8.
The ${\rm U(1)}$-fixed point subnet of the complex free fermion net
$\mathrm{Fer}_{\mathbb{C}}$ is the ${\rm U(1)}$-current net, i.e.
$\mathrm{Fer}_{\mathbb{C}}^{{\rm
U(1)}}={\mathcal{B}}\cong{{\mathcal{A}}^{(0)}}$.
###### Proof.
Let us see ${\mathcal{B}}$ as a subnet of the fermi net
$\mathrm{Fer}_{\mathbb{C}}^{{\rm U(1)}}$ on
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{{\rm
U(1)}}\equiv{\mathcal{H}}_{\,\cdot\,,0}$. Further
$\overline{{\mathcal{B}}(I)\Omega}$ does not depend on $I$ by the same proof
of the Reeh-Schlieder property and is clearly a subspace of
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{{\rm
U(1)}}\equiv{\mathcal{H}}_{\,\cdot\,,0}$.
In fact they coincide, since we have confirmed that
$\operatorname{tr}_{{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}}(\mathrm{e}^{-\beta
L_{0}})=\operatorname{tr}_{{\mathcal{H}}_{\,\cdot\,,0}}(\mathrm{e}^{-\beta
L_{0}})=p(t)$, where $\mathrm{e}^{-\beta}=t$, namely, their conformal
characters coincide (see also Section 2.2). ∎
We finish this section by giving the parametrization in $x$-picture, where the
action of the translation is more natural. With
$f(x)=\frac{1}{\sqrt{2\pi}}\sqrt{\left|\frac{\partial\theta(x)}{\partial
x}\right|}\mathrm{e}^{\mathrm{i}\theta(x)/2}f_{0}(\theta(x))$
we identify $L^{2}({\mathbb{R}})=L^{2}({\mathbb{R}},\,\mathrm{d}x)$ with
$L^{2}(S^{1})=L^{2}([0,2\pi],\,\mathrm{d}\theta/(2\pi))$ and therefore the
space ${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$ is given by
$PL^{2}({\mathbb{R}})\oplus\overline{P^{\perp}L^{2}({\mathbb{R}})}$ with
$P:\hat{f}(p)\mapsto\Theta(p)\hat{f}(p)$ and it can be identified in “momentum
space” with $L^{2}({\mathbb{R}}_{+},2\pi\,\mathrm{d}p)\oplus
L^{2}({\mathbb{R}}_{+},2\pi\,\mathrm{d}q)$ by
$\displaystyle f(x)$
$\displaystyle\longmapsto\widehat{Pf}(p)\oplus\overline{\widehat{P^{\perp}f}(-q)}$
$\displaystyle p,q>0\,.$
The field operators are defined for $f\in L^{2}({\mathbb{R}})$ by
$\psi(f)=a_{P}(f)$ and $\bar{\psi}(f)=a_{P}(\overline{f})^{\ast}$. For
$\Psi\in{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}$ we write its components
$\Psi_{m,n}\in{\mathcal{H}}_{m,n}:=L^{2}({\mathbb{R}}_{+}^{m+n},(2\pi)^{m+n}\,\mathrm{d}p_{1}\cdots\,\mathrm{d}p_{m}\,\mathrm{d}q_{1}\cdots
q_{n})_{-},$ (1)
where $-$ means the antisymmetrization within $p_{1},\ldots,p_{m}$ and
$q_{1},\ldots,q_{n}$. By this notation
$(\psi(f)\Omega_{0})_{1,0}(p)=\hat{f}(p)$ and
$(\bar{\psi}(f)\Omega_{0})_{0,1}(q)=\hat{f}(q)$. Further the bi-field
${:\\!\bar{\psi}(f)\psi(g)\\!:}=\bar{\psi}(f)\psi(g)-\langle\Omega_{0},\bar{\psi}(f)\psi(g)\Omega_{0}\rangle{\mathbbm{1}}$
creates from the vacuum $\Omega_{0}$ a fermionic 1+1 particle state
$\Psi_{f,g}:={:\\!\bar{\psi}(f)\psi(g)\\!:}\Omega_{0}$ with
$(\Psi_{f,g})_{1,1}(p,q)=-\hat{f}(q)\hat{g}(p)$ and it follows for $h\in
C^{\infty}({\mathbb{R}},{\mathbb{R}})$ that for the ${\rm U(1)}$-current $J$,
it holds $(J(h)\Omega_{0})_{1,1}(p,q)=-\frac{1}{2\pi}\hat{h}(p+q)$ which is
obtained by taking a limit $\sum_{n}\Psi_{f_{n},g_{n}}$ with test functions
$\sum_{n}f_{n}(x)g_{n}(y)\to h(x)\delta(x-y)$. We make the important
observation that the $J(f)\Omega_{0}$ generate the one-particle space which we
can identify with ${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ and
this is obviously a proper subspace of the fermionic 1+1-particle space
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1,1}$.
## 4 A new family of Longo-Witten endomorphisms on ${\rm U(1)}$-current net
We use the description of
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}=\overline{PL^{2}(S^{1})}+P^{\perp}L^{2}(S^{1})$
which equals $L^{2}(S^{1})$ as a real Hilbert space and is described in the
beginning of Section 3.2. First we decompose
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$ into irreducible
representations of $\mathrm{SU}(1,1)$ in a compatible way with $K(I)$. Let us
define
$\displaystyle{\mathcal{H}}_{\Re}$ $\displaystyle:=$
$\displaystyle\\{f\in{\mathcal{H}}^{1}_{\mathrm{Fer}_{\mathbb{C}}}:z^{\frac{1}{2}}f(z)\mbox{
is real}\\},$ $\displaystyle{\mathcal{H}}_{\Im}$ $\displaystyle:=$
$\displaystyle\\{f\in{\mathcal{H}}^{1}_{\mathrm{Fer}_{\mathbb{C}}}:z^{\frac{1}{2}}f(z)\mbox{
is pure imaginary}\\}.$
By their definition, it is clear that ${\mathcal{H}}_{\Re}$ and
${\mathcal{H}}_{\Im}$ are real Hilbert subspaces of $L^{2}(S^{1})$. In fact,
they are complex subspaces with respect to the new complex structure. To see
this, we take another description of ${\mathcal{H}}_{\Re}$: in terms of
Fourier components, it holds that $f\in{\mathcal{H}}_{\Re}$ if and only if
$f_{n}=\overline{f_{-n-1}}$. Recall that, on $L^{2}(S^{1})$, the new scalar
multiplication by $\mathrm{i}$ is given by $\mathrm{i}\cdot
f_{n}=-\mathrm{i}f_{n}$, $\mathrm{i}\cdot f_{-n-1}=\mathrm{i}f_{-n-1}$ for
$n\geq 0$. Hence this condition is preserved under the multiplication by
$\mathrm{i}$ and ${\mathcal{H}}_{\Re}$ is a complex subspace. An analogous
argument holds for ${\mathcal{H}}_{\Im}$. Next we see that
${\mathcal{H}}_{\Re}$ and ${\mathcal{H}}_{\Im}$ are orthogonal. Note that
because of the change of the complex structure, for $f(z)=\sum_{n}f_{n}z^{n}$
and $h(z)=\sum_{n}h_{n}z^{n}$ the inner product is written as follows:
$\langle f,h\rangle=\sum_{n\geq
0}f_{n}\overline{h_{n}}+\sum_{n<0}\overline{f_{n}}h_{n}.$
Now $f\in{\mathcal{H}}_{\Re}$ implies $f_{n}=\overline{f_{-n-1}}$ and
$h\in{\mathcal{H}}_{\Im}$ implies $h_{n}=-\overline{h_{-n-1}}$ for non-
negative $n$, hence it is easy to see that
$\langle f,h\rangle=\sum_{n\geq
0}f_{n}\overline{h_{n}}+\sum_{n<0}\overline{f_{n}}h_{n}=-\sum_{n\geq
0}\overline{f_{-n-1}}{h_{-n-1}}+\sum_{n<0}\overline{f_{n}}{h_{n}}=0.$
In other words, these two complex subspaces are mutually orthogonal.
Furthermore, ${\mathcal{H}}_{\Re}$ and ${\mathcal{H}}_{\Im}$ are invariant
under the action of $\mathrm{SU}(1,1)$. We recall that the action is given by
$(V_{g}f)(z)=\frac{1}{-\overline{\beta}z+\alpha}f\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)$.
Then if $z^{\frac{1}{2}}f(z)$ is real then it holds that
$\displaystyle z^{\frac{1}{2}}(V_{g}f)(z)$ $\displaystyle=$
$\displaystyle\frac{1}{{\bar{z}}^{\frac{1}{2}}(-\overline{\beta}z+\alpha)}\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)^{-\frac{1}{2}}\cdot\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)^{\frac{1}{2}}f\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)$
$\displaystyle=$
$\displaystyle\frac{1}{(-\overline{\beta}+\alpha\bar{z})^{\frac{1}{2}}(\overline{\alpha}z-\beta)^{\frac{1}{2}}}\cdot\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)^{\frac{1}{2}}f\left(\frac{\overline{\alpha}z-\beta}{-\overline{\beta}z+\alpha}\right)$
and both factors are real. Similarly one shows that ${\mathcal{H}}_{\Im}$ is
preserved under $V_{g}$. It is obvious that these two representations are
intertwined by the multiplication by $\mathrm{i}$ in the old complex
structure. This is still a unitary map, thus they are unitarily equivalent.
One can see that each representation is indeed irreducible, and when
restricted to $\mathrm{PSU}(1,1)={\rm PSL}(2,{\mathbb{R}})$, it is the
projective positive energy representation with lowest weight $\frac{1}{2}$.
It is easy to see that $e^{\Re}_{n}:=\\{e_{n}+e_{-n+1},n\geq 0\\}$ and
$e^{\Im}_{n}:=\\{\mathrm{i}(e_{n}+e_{-n+1}),n\geq 0\\}$ form bases of
${\mathcal{H}}_{\Re}$ and ${\mathcal{H}}_{\Im}$, respectively, where
$e_{n}(z)=z^{n}$ and the multiplication by $\mathrm{i}$ is given in the old
structure. Now we describe the gauge action in terms of this basis. By the
definition, for a given complex number $\alpha$ with modulus $1$, the action
is given by the multiplication in the old structure. Hence if
$\alpha=\cos\theta+\mathrm{i}\sin\theta$, we have
$U_{\alpha}e^{\Re}_{n}=\cos\theta e^{\Re}_{n}+\sin\theta e^{\Im}_{n}$ and
$U_{\alpha}e^{\Im}_{n}=-\sin\theta e^{\Re}_{n}+\cos\theta e^{\Im}_{n}$. This
means that $U_{\alpha}$ acts as the real rotation by $\theta$ in this basis.
#### Construction of endomorphisms
We construct Longo-Witten endomorphisms on the free fermion net
$\mathrm{Fer}_{\mathbb{C}}$ commuting with the gauge action. The key is the
following theorem. We remind that a standard pair $(\tilde{H},\tilde{T})$ is a
standard subspace $\tilde{H}\subset\tilde{\mathcal{H}}$ of a Hilbert space
$\tilde{\mathcal{H}}$ and a positive energy representation $\tilde{T}$ of
${\mathbb{R}}$ on $\tilde{\mathcal{H}}$, such that
$\tilde{T}(a)\tilde{H}\subset\tilde{H}$ for $a\geq 0$. If $\tilde{T}$ is
maximally abelian, the standard pair is said to be irreducible and there is a
(up to unitary equivalence) unique irreducible standard pair.
###### Theorem 4.1 ([LW11, Theorem 2.6]).
Let $(\tilde{H},\tilde{T})$ be a standard pair with multiplicity $n$, i.e. it
decomposes into $n$-fold direct sum of irreducible standard pairs, each
unitarily equivalent to the unique standard pair $(H,T)$ and
$T(t)=\mathrm{e}^{\mathrm{i}tP}$. Then a unitary $\tilde{V}$ commuting with
the translation $\tilde{T}$ preserves $\tilde{H}$ if and only if $\tilde{V}$
is a $n\times n$ matrix $(V_{hk})$ (with respect to the decomposition of
$\tilde{\mathcal{H}}$ into $n$ direct sum as above) such that
$V_{hk}=\varphi_{hk}(P)$, where $\varphi_{hk}:{\mathbb{R}}\to{\mathbb{C}}$ are
complex Borel functions such that $(\varphi_{hk})$ is a unitary matrix for
almost every $p>0$, each $\varphi_{hk}$ is the boundary value of a function in
${\mathbb{H}}(\SS_{\infty})$ and is symmetric, i.e.
$\varphi_{hk}(-p)=\overline{\varphi_{hk}(p)}$.
Consider the one-particle space
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$ for
$\mathrm{Fer}_{\mathbb{C}}$. The pair of the standard space
$K({\mathbb{R}}_{+})=L^{2}({\mathbb{R}}_{+})$ defined in Section 3.2 (under
the identification of $S^{1}$ and ${\mathbb{R}}\cup\\{\infty\\})$ and the
natural translation has multiplicity 2. If we take a matrix-valued function
$(\varphi_{hk})$ as above and take the second quantization operator
$\Lambda(V)$ of the (matrix-valued) operator $(V_{kh})=(\varphi_{hk}(P))$,
then it implements a Longo-Witten endomorphism of $\mathrm{Fer}_{\mathbb{C}}$
(see [LW11]).
As the gauge group acts by real rotation
$\left(\begin{matrix}\cos\theta&-\sin\theta\\\
\sin\theta&\cos\theta\end{matrix}\right)$, any matrix-valued function of $p$
which commute with them must have the form
$\left(\begin{matrix}a(p)&\mathrm{i}b(p)\\\
-\mathrm{i}b(p)&a(p)\end{matrix}\right)$. If each component is symmetric, then
$a$ is symmetric and $b$ is antisymmetric. Such a matrix-valued function can
be diagonalized by the matrix $\left(\begin{matrix}1&\mathrm{i}\\\
\mathrm{i}&1\end{matrix}\right)$ and becomes
$\left(\begin{matrix}a(p)+b(p)&0\\\ 0&a(p)-b(p)\end{matrix}\right)$. We claim
that such $a$ and $b$ exist. Indeed, let $\varphi$ be a inner function (not
necessarily symmetric), namely the boundary value with modulus $1$ of a
bounded analytic function on the upper half-plane ${\mathbb{H}}$, and define
$a(p)=\frac{1}{2}(\varphi(p)+\overline{\varphi(-p)})$,
$b(p)=\frac{1}{2}(\varphi(p)-\overline{\varphi(-p)})$. Then it is obvious that
$a$ is symmetric and $b$ is antisymmetric. In addition, $a(p)+b(p)=\varphi(p)$
and $a(p)-b(p)=\overline{\varphi(-p)}$, hence the diagonalized matrix is
unitary for almost every $p$. By the theorem of Longo-Witten, the operator
$\varphi(P_{1}):=\left(\begin{matrix}a(P)&\mathrm{i}b(P)\\\
-\mathrm{i}b(P)&a(P)\end{matrix}\right)$ preserves the real Hilbert space
$\widetilde{H}:=K({\mathbb{R}}_{+})$, where $P_{1}$ is the generator of the
translation in ${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}$ which has
multiplicity $2$ and $P$ is the generator of $T$ of the irreducible standard
pair $(H,T)$.
It is easy to see that the above diagonalization is given exactly by the
decomposition
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}=\overline{PL^{2}(S^{1})}\oplus({\mathbbm{1}}-P)L^{2}(S^{1})$.
We remind that ${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}^{1}$ can be
identified with $L^{2}({\mathbb{R}})$ as a real space. In
$L^{2}({\mathbb{R}})$ the function $\varphi(P_{1})f$ is the function with
Fourier transform $\varphi(p)\hat{f}(p)$ and we remark that it also follows
directly from the Paley-Wiener theorem that $\varphi(P_{1})$ leaves
$L^{2}({\mathbb{R}}_{+})\subset L^{2}({\mathbb{R}})$ invariant for $\varphi$
inner. Further using that the space
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}$ decomposes in
${\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}=\bigoplus_{m,n\in{\mathbb{N}}_{0}}{\mathcal{H}}_{m,n}$
like in (1) with the gauge action given by
$V(\theta)\Psi_{m,n}=\mathrm{e}^{\mathrm{i}(m-n)\theta}\Psi_{m,n}$, the action
of the Longo-Witten unitary $V_{\varphi}=\Lambda(\varphi(P_{1}))$ is given by
$\displaystyle(V_{\varphi}\Psi)_{m,n}(p_{1},\cdots,p_{m},q_{1},\cdots,q_{n})$
$\displaystyle\quad=\varphi(p_{1})\cdots\varphi(p_{m})\overline{\varphi(-q_{1})}\cdots\overline{\varphi(-q_{n})}\Psi_{m,n}(p_{1},\cdots,p_{m},q_{1},\cdots,q_{n})\,.$
###### Lemma 4.2.
Let $\iota:\Psi\in
L^{2}({\mathbb{R}}_{+},p\,\mathrm{d}p)\equiv{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\hookrightarrow{\mathcal{H}}_{1,1}\subset{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}$
be the embedding given by $\iota(\Psi)_{1,1}(p,q)=-\frac{1}{2\pi}\Psi(p+q)$. A
second quantization Longo-Witten unitary $V_{\varphi}$ commuting with the
gauge action $V(\,\cdot\,)$ satisfies
$V_{\varphi}\iota{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\subset\iota{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$
if and only if $V_{\varphi}=V(\theta)T(t)$ with $t\geq 0$.
###### Proof.
The translations commute with the gauge action and it follows immediately that
they leave ${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ invariant. We
note that $\varphi(p)\overline{\varphi(-q)}\Psi(p+q)$ belongs to
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$ only if it can be written
as a function of $g(p+q)$. This means that
$\varphi(p)\overline{\varphi(-q)}=\widetilde{\varphi}(p+q)$ for $p,q\geq 0$,
where $\widetilde{\varphi}$ is another function. Then, putting $q=0$ and $p=0$
respectively, we see that
$\varphi(p)\overline{\varphi(0)}=\widetilde{\varphi}(p)$ for $p\geq 0$ and
$\varphi(0)\overline{\varphi(-q)}=\widetilde{\varphi}(q)$ for $q\geq 0$, in
particular $\widetilde{\varphi}(0)=1$. Multiplying the each side of these
equations, one sees that
$\widetilde{\varphi}(p+q)=\widetilde{\varphi}(p)\widetilde{\varphi}(q)$
because $|\varphi(0)|=1$. Then it follows that
$\widetilde{\varphi}(p)=\mathrm{e}^{\mathrm{i}\kappa p}$ for some $\kappa\geq
0$, and $\varphi(p)=\mathrm{e}^{\mathrm{i}(\kappa p+\theta)}$ for some
$\theta\in{\mathbb{R}}$ (in fact, the arguments here should be treated with
care because the relation is given only almost everywhere, but both $\varphi$
and $\check{\varphi}$ analytically continue and in the domain of analyticity
it holds everywhere).
Such a $\varphi$ is a Longo-Witten unitary only for $\kappa\geq 0$. The
constant factor $\mathrm{e}^{\mathrm{i}\theta}$ corresponds to the factor
$V(\theta)$. ∎
###### Theorem 4.3.
Let $\varphi$ be an inner function as above. The endomorphism implemented by
the second quantization $V_{\varphi}$ of the operator constructed above
restricts to the ${\rm U(1)}$-current subnet. The restriction cannot be
implemented by any second quantization operator if
$\varphi(p)\neq\mathrm{e}^{\mathrm{i}(\kappa p+\theta)}$.
###### Proof.
The operator $V_{\varphi}$ restricts to the subnet ${{\mathcal{A}}^{(0)}}$ by
the general argument in Proposition 2.1. It cannot be implemented by a second
quantization operator, since any second quantization operator preserves the
particle number, while $V_{\varphi}$ does not for non-exponential $\varphi$ as
we saw above, and a Longo-Witten endomorphism is uniquely implemented up to
scalar (see Section 2.1). ∎
###### Remark 4.4.
By the construction in [LW11], each unitary
$V=V_{\varphi}|_{{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}}$ related to
an inner function $\varphi$ from above gives rise to a local, time-translation
covariant net of von Neumann algebras on the Minkowski half-space
$M_{+}=\\{(t,x)\in{\mathbb{R}}^{2}:x>0\\}$. This net is associated with the
${\rm U(1)}$-current net ${{\mathcal{A}}^{(0)}}$ and defined by
${{\mathcal{A}}^{(0)}}_{V}(O)={{\mathcal{A}}^{(0)}}(I_{1})\vee
V{{\mathcal{A}}^{(0)}}(I_{2})V^{\ast}$, where $O=I_{1}\times
I_{2}=\\{(t,x)\in{\mathbb{R}}^{2}:t-x\in I_{1},t+x\in I_{2}\\}$ is a double
cone with $\overline{O}\subset M_{+}$ corresponding uniquely to the two
intervals $I_{1}$ and $I_{2}$ with disjoint closures. In the case where
$\varphi$ is not exponential $V_{\varphi}$ does not come from second
quantization—in contrast to the unitaries constructed by Longo and Witten in
[LW11]—and therefore gives new examples.
## 5 Interacting wedge-local net with particle production
### 5.1 Construction of scattering operators
In the previous section we saw that, in the basis
$\\{e_{n}+e_{-n},e_{n}-e_{-n}\\}$ the matrix operator
$\left(\begin{matrix}a(P)&\mathrm{i}b(P)\\\
-\mathrm{i}b(P)&a(P)\end{matrix}\right)$ implements a Longo-Witten
endomorphism if $a$ is symmetric and $b$ is antisymmetric, and after the
simultaneous diagonalization it becomes $\left(\begin{matrix}\varphi(P)&0\\\
0&\check{\varphi}(P)\end{matrix}\right)$ where $\varphi$ is an inner function
and $\check{\varphi}(p)=\overline{\varphi(-p)}$ (note that if $\varphi$
extends to an analytic function $\varphi(z)$ on ${\mathbb{H}}$, then
$\check{\varphi}(z)=\overline{\varphi(-\overline{z})}$ also extends to
${\mathbb{H}}$, hence $\check{\varphi}$ is again an inner function). By the
same argument one sees that $\left(\begin{matrix}\check{\varphi}(P)&0\\\
0&\varphi(P)\end{matrix}\right)$ implements an endomorphism since
$\check{\check{\varphi}}=\varphi$.
With respect to the basis after diagonalization, we split the Hilbert space
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}=:{\mathcal{H}}_{+}\oplus{\mathcal{H}}_{-}$
and the generator of translation $P_{1}=:P_{+}\oplus P_{-}$. Then the tensor
product space can be written as follows:
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}=\left({\mathcal{H}}_{+}\otimes{\mathcal{H}}_{+}\right)\oplus\left({\mathcal{H}}_{+}\otimes{\mathcal{H}}_{-}\right)\oplus\left({\mathcal{H}}_{-}\otimes{\mathcal{H}}_{+}\right)\oplus\left({\mathcal{H}}_{-}\otimes{\mathcal{H}}_{-}\right)\,.$
According to this decomposition into a direct sum of four subspaces, we define
an operator
$M_{\varphi}:=\varphi(P_{+}\otimes P_{+})\oplus\check{\varphi}(P_{+}\otimes
P_{-})\oplus\check{\varphi}(P_{-}\otimes P_{+})\oplus\varphi(P_{-}\otimes
P_{-})\,.$
Then this restricts to the subspace
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}\otimes{\mathcal{H}}_{+}=\left({\mathcal{H}}_{+}\oplus{\mathcal{H}}_{-}\right)\otimes{\mathcal{H}}_{+}$
and it is $\varphi(P_{+}\otimes P_{+})\oplus\check{\varphi}(P_{-}\otimes
P_{+})$, or we can decompose it with respect to the spectral measure of
$P_{+}$:
$\int_{{\mathbb{R}}_{+}}\left(\begin{matrix}\varphi(pP_{+})&0\\\
0&\check{\varphi}(pP_{-})\end{matrix}\right)\otimes\,\mathrm{d}E_{+}(p).$
Similarly, the restriction to
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}\otimes{\mathcal{H}}_{-}$ is
written as
$\int_{{\mathbb{R}}_{+}}\left(\begin{matrix}\check{\varphi}(pP_{+})&0\\\
0&\varphi(pP_{-})\end{matrix}\right)\otimes\,\mathrm{d}E_{-}(p).$
Using the two-point set ${\mathbb{Z}}_{2}=\\{+,-\\}$ we define
$\varphi_{+}(p,+):=\varphi(p),\,\,\,\,\varphi_{+}(p,-)=\check{\varphi}(p),\,\,\,\,\varphi_{-}(p,+)=\check{\varphi}(p),\,\,\,\,\varphi_{-}(p,-)=\varphi(p).$
By defining the spectral measure $E_{1}=E_{+}\oplus E_{-}$ on
${\mathcal{H}}^{1}$, $M_{\varphi}$ can be simply written as
$M_{\varphi}=\int_{{\mathbb{R}}_{+}\times{\mathbb{Z}}_{2}}\left(\begin{matrix}\varphi_{+}(pP_{+},\iota)&0\\\
0&\varphi_{-}(pP_{-},\iota)\end{matrix}\right)\otimes\,\mathrm{d}E_{1}(p,\iota),$
where $\iota=\pm$.
As in [Tan11a], we construct the scattering matrix first on the unsymmetrized
Fock space, then restrict it to the antisymmetric space. For an operator $A$
on
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}$,
we denote by $A^{m,n}_{i,j}$ on
$({\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1})^{\otimes
m}\otimes({\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1})^{\otimes n}$ the
operator which acts only on the $i$-th factor in
$({\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1})^{\otimes m}$ and $j$-th
factor in $({\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1})^{\otimes n}$ as
$A$. As a convention, $A^{m,n}_{i,j}$ equals to the identity operator if $m$
or $n$ is $0$. Let us denote simply
$\widetilde{\varphi}(p,\iota):=\left(\begin{matrix}\varphi_{+}(p,\iota)&0\\\
0&\varphi_{-}(p,\iota)\end{matrix}\right)$ and
$\widetilde{\varphi}(P_{1},\iota):=\left(\begin{matrix}\varphi_{+}(P_{+},\iota)&0\\\
0&\varphi_{-}(P_{-},\iota)\end{matrix}\right)$. From the observation above, it
is straightforward to see that
$(M_{\varphi})^{m,n}_{i,j}=\int\left({\mathbbm{1}}\otimes\cdots\otimes\underset{i\mbox{-th}}{\widetilde{\varphi}(p_{j}P_{1},\iota_{j})}\otimes\cdots\otimes{\mathbbm{1}}\right)\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})$
(the case where $m$ or $n$ is $0$ is treated separately). Then we define, as
in [Tan11a],
$\displaystyle S_{\varphi}^{m,n}$ $\displaystyle:=$
$\displaystyle\prod_{i,j}(M_{\varphi})^{m,n}_{i,j}$ $\displaystyle
S_{\varphi}$ $\displaystyle:=$
$\displaystyle\bigoplus_{m,n}S_{\varphi}^{m,n}\,.$
Let ${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$ be the unsymmetrized
Fock space based on ${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1}$. Note that
$S_{\varphi}$ is defined on
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$,
and it naturally restricts to
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$,
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}$
and
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}$.
This $S_{\varphi}$ will be interpreted as the scattering matrix. In order to
confirm this, we have to take the spectral decomposition of $S_{\varphi}$ only
with respect to the right or left component. Namely,
$\displaystyle S_{\varphi}$ $\displaystyle:=$
$\displaystyle\bigoplus_{m,n}\prod_{i,j}(M_{\varphi})^{m,n}_{i,j}$
$\displaystyle=$
$\displaystyle\bigoplus_{m,n}\prod_{i,j}\int\left({\mathbbm{1}}\otimes\cdots\otimes\underset{i\mbox{-th}}{\widetilde{\varphi}(p_{j}P_{1},\iota_{j})}\otimes\cdots\otimes{\mathbbm{1}}\right)\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})$
$\displaystyle=$
$\displaystyle\bigoplus_{m,n}\int\prod_{i,j}\left({\mathbbm{1}}\otimes\cdots\otimes\underset{i\mbox{-th}}{\widetilde{\varphi}(p_{j}P_{1},\iota_{j})}\otimes\cdots\otimes{\mathbbm{1}}\right)\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})$
$\displaystyle=$
$\displaystyle\bigoplus_{n}\int\bigoplus_{m}\prod_{j}\left(\widetilde{\varphi}(p_{j}P_{1},\iota_{j})\right)^{\otimes
m}\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})$
$\displaystyle=$
$\displaystyle\bigoplus_{n}\int\prod_{j}\bigoplus_{m}\left(\widetilde{\varphi}(p_{j}P_{1},\iota_{j})\right)^{\otimes
m}\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})$
$\displaystyle=$
$\displaystyle\bigoplus_{n}\int\prod_{j}\Lambda(\widetilde{\varphi}(p_{j}P_{1},\iota_{j}))\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})\,,$
where the integral and the product commute in the third equality since the
spectral measure is disjoint for different values of $p$’s and $\iota$’s, and
the sum and the product commute in the fifth equality since the operators in
the integrand act on mutually disjoint spaces, namely on
$({\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{1})^{\otimes
m}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$ for different
$m$. In the final expression, all operators appearing in the integrand are the
second quantization operators, thus this formula naturally restricts to the
partially antisymmetrized space
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$.
Now we define
* •
${\mathcal{M}}_{\varphi}:=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S_{\varphi}({\mathbbm{1}}\otimes
y):x\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{-}),y\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{+})\\}^{\prime\prime}$,
* •
$T(t,x):=T_{0}(\frac{t-x}{\sqrt{2}})\otimes T_{0}(\frac{t+x}{\sqrt{2}})$,
* •
$\Omega:=\Omega_{0}\otimes\Omega_{0}$.
As the net $\mathrm{Fer}_{\mathbb{C}}$ is fermionic by nature, the
interpretation of the scattering theory of [Buc75] is not clear. Nevertheless,
we can show the following by an almost same proof as in [Tan11a, Lemma 5.2,
Theorem 5.3].
###### Lemma 5.1.
The triple $({\mathcal{M}}_{\varphi},T,\Omega)$ is a Borchers triple.
###### Proof.
To apply Proposition 2.2, it is immediate that $S_{\varphi}$ commutes with
translation since it is defined through the spectral measure as above. It
preserves ${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes\Omega_{0}$ and
$\Omega_{0}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}$ pointwise, since
these subspaces correspond to the case where $m$ or $n$ is $0$ in the above
decomposition and $S_{\varphi}$ acts as the identity operator by definition.
What remains to show is the commutation property.
As we saw above, the operator $S_{\varphi}$ can be written as
$S_{\varphi}=\bigoplus_{n}\int\prod_{j}\Lambda(\widetilde{\varphi}(p_{j}P_{1},\iota_{j}))\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})\,.$
The point is that the operators which appear in the integrand implement Longo-
Witten endomorphisms as we saw above since $p_{j}\geq 0$ in the support of the
integration.
Let $x^{\prime}\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{+})$ and consider
$x^{\prime}\otimes{\mathbbm{1}}$ as an operator on
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$.
We have
${\hbox{\rm
Ad\,}}S_{\varphi}(x^{\prime}\otimes{\mathbbm{1}})=\bigoplus_{n}\int{\hbox{\rm
Ad\,}}\left(\prod_{j}\Lambda(\widetilde{\varphi}(p_{j}P_{1},\iota_{j}))\right)(x^{\prime})\otimes\,\mathrm{d}E_{1}(p_{1},\iota_{1})\otimes\cdots\otimes\,\mathrm{d}E_{1}(p_{n},\iota_{n})\,.$
Although this formula is not closed on
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}$,
the left hand side obviously restricts there. One sees that the integrand
remains in $\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{+})$.
Recall the operator $Z_{0}$ which gives the graded locality of
$\mathrm{Fer}_{\mathbb{C}}$. One has to remind that
$Z_{0}=\frac{{\mathbbm{1}}-\mathrm{i}\Gamma_{0}}{1-\mathrm{i}}$ where
$\Gamma_{0}=\Lambda(-{\mathbbm{1}})$, hence $Z_{0}$ commutes with any second
quantization operator. Then by the disintegration above (and the corresponding
disintegration with respect to the left component), it is easy to see that
$Z_{0}\otimes{\mathbbm{1}}$ commutes with $S_{\varphi}$.
Let us check the commutation property of the assumptions in Proposition 2.2.
Note that ${\hbox{\rm Ad\,}}Z_{0}(x)\otimes{\mathbbm{1}}$ and ${\hbox{\rm
Ad\,}}Z_{0}(x)\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{+})^{\prime}$ for
$x\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{-})$. Since
$Z_{0}\otimes{\mathbbm{1}}$ and $S_{\varphi}$ commute as we saw above, to
prove the first commutation relation, it is enough to show that $[{\hbox{\rm
Ad\,}}Z_{0}(x)\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S_{\varphi}(x^{\prime}\otimes{\mathbbm{1}})]=0$ for
$x\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{-})$ and
$x^{\prime}\in\mathrm{Fer}_{\mathbb{C}}({\mathbb{R}}_{+})$. As operators
acting on
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}^{\Sigma}$,
this is done by the above disintegration of ${\hbox{\rm
Ad\,}}S_{\varphi}(x^{\prime}\otimes{\mathbbm{1}})$. Then both operators
naturally restrict to
${\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{\mathbb{C}}}$,
and we obtain the claim (cf. [Tan11a, Lemma 5.2, Theorem 5.3]). The second
commutation relation for Proposition 2.2 can be proven analogously. ∎
Finally we arrive at a new family of interacting Borchers triples with
asymptotic algebra ${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$.
###### Theorem 5.2.
Let us define
* •
${\mathcal{N}}_{\varphi}:=\\{x\otimes{\mathbbm{1}},{\hbox{\rm
Ad\,}}S_{\varphi}({\mathbbm{1}}\otimes y):x\in\mathrm{Fer}_{\mathbb{C}}^{{\rm
U(1)}}({\mathbb{R}}_{-}),y\in\mathrm{Fer}_{\mathbb{C}}^{{\rm
U(1)}}({\mathbb{R}}_{+})\\}^{\prime\prime}$,
* •
$T(t,x):=T_{0}(\frac{t-x}{\sqrt{2}})\otimes T_{0}(\frac{t+x}{\sqrt{2}})$,
* •
$\Omega:=\Omega_{0}\otimes\Omega_{0}$.
Then the triple $({\mathcal{N}}_{\varphi},T,\Omega)$, restricted to
$\overline{{\mathcal{N}}_{\varphi}\Omega}$, is an asymptotically complete,
interacting Borchers triple with the asymptotic algebra
${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ and scattering operator
$S_{\varphi}|_{\overline{{\mathcal{N}}_{\varphi}\Omega}}$. It also holds that
$\overline{{\mathcal{N}}_{\varphi}\Omega}=\overline{{{\mathcal{A}}^{(0)}}(I_{+})\otimes{{\mathcal{A}}^{(0)}}(I_{-})\Omega}$
for arbitrary intervals $I_{+},I_{-}$.
###### Proof.
Substantial arguments are already done: In Lemma 5.1 we constructed Borchers
triples with $\mathrm{Fer}_{\mathbb{C}}\otimes\mathrm{Fer}_{\mathbb{C}}$ as
the asymptotic algebra. We have seen in Section 3.3 the ${\rm U(1)}$-current
net ${{\mathcal{A}}^{(0)}}$ is the fixed point subnet of
$\mathrm{Fer}_{\mathbb{C}}$ with respect to the action of ${\rm U(1)}$. From
the construction in Section 5.1 and Theorem 4.3, it is easy to see that
$S_{\varphi}$ commutes with the product action of the inner symmetries. Then
all the statements of the Theorem follow from the general consideration of
Proposition 2.5. ∎
### 5.2 Action of the S-matrix on the 1+1 particle space
In this Section we want to analyze the action of the S-matrix of the models
constructed in Section 5.1 on the 1+1 particle space
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$,
i.e. one left and one right moving particle, where we use the word particle in
the sense of Fock space excitations. We note that on the $n$+0 and 0+$n$
particle spaces ${\mathcal{H}}_{n}\otimes{\mathbb{C}}\Omega_{0}$ and
${\mathbb{C}}\Omega_{0}\otimes{\mathcal{H}}_{n}$, respectively, the S-matrix
$S$ acts trivially. A typical vector in
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$
is of the form $\Psi:=J(f)\Omega_{0}\otimes J(g)\Omega_{0}$ which we express
as the function $\Psi(p,\bar{p})=\hat{f}(p)\hat{g}(\bar{p})$. The embedding
$\iota:L^{2}({\mathbb{R}}_{+},p\,\mathrm{d}p)\otimes
L^{2}({\mathbb{R}}_{+},\bar{p}\,\mathrm{d}\bar{p})\cong{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\hookrightarrow{\mathcal{H}}_{1,1}\otimes{\mathcal{H}}_{1,1}\subset{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}\otimes{\mathcal{H}}_{\mathrm{Fer}_{{\mathbb{C}}}}$
is given by
$\iota(\Psi)_{1,1;1,1}(p,q,\bar{p},\bar{q})=\frac{1}{(2\pi)^{2}}\Psi(p+q,\bar{p}+\bar{q})$.
We have an analogue of Lemma 4.2
###### Proposition 5.3.
Let $\varphi$ be some inner function. The unitary $S_{\varphi}$ satisfies
$S_{\varphi}({\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1})\subset{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$
if and only if $\varphi(p)=\mathrm{e}^{\mathrm{i}(kp+\theta)}$.
###### Proof.
The action of $S_{\varphi}$ on
$\Psi\in{\mathcal{H}}_{1}\otimes{\mathcal{H}}_{1}$ is given by
$S_{\varphi}\Psi(p+q,\bar{p}+\bar{q})=\varphi(p\cdot\bar{p})\check{\varphi}(q\cdot\bar{p})\check{\varphi}(p\cdot\bar{q})\varphi(q\cdot\bar{q})\Psi(p+q,\bar{p}+\bar{q})$
which is again in
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$
if it can be written as a function $\tilde{\Psi}(p+q,\bar{p}+\bar{q})$, in
particular if
$\varphi(p\cdot\bar{p})\check{\varphi}(q\cdot\bar{p})\check{\varphi}(p\cdot\bar{q})\varphi(q\cdot\bar{q})=\widetilde{\varphi}(p+q,\bar{p}+\bar{q})$.
Setting $\bar{p}=1$ and $\bar{q}=0$, we have
$\varphi(p)\check{\varphi}(q)=\widetilde{\varphi}(p+q,1)$. The rest follows as
Lemma 4.2.
∎
###### Remark 5.4.
In the case $\varphi(p)=\mathrm{e}^{\mathrm{i}\kappa p}$, one gets the models
obtained in [DT11] using warped convolution.
###### Proposition 5.5.
Let $e$ be the projection on
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$,
then $eS_{\varphi}e=\tilde{\varphi}(P\otimes P)$, where $\tilde{\varphi}$ is
boundary value of an analytic function in $\mathbb{H}$ with
$|\tilde{\varphi}(p)|\leq 1$ and $P$ is the generator of translation
restricted to the one-particle space (which gives rise the irreducible
standard pair).
###### Proof.
It can be checked that
$(e_{0}f)(p,q)=\frac{1}{p+q}\int_{0}^{p+q}f(p+q-x,x)\,\mathrm{d}x$
is the projection on
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\subset{\mathcal{H}}_{1,1}$.
Then the action of $eS_{\varphi}$ on a
$f\in{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$
can be calculated to be $\varphi^{\prime}(P\otimes 1,1\otimes P)$ with
$\varphi^{\prime}(p,q)=\frac{1}{p\cdot
q}\int_{0}^{p}\int_{0}^{q}\varphi((p-x)\cdot(q-y))\varphi(x\cdot
y)\check{\varphi}((p-x)\cdot
y)\check{\varphi}(x\cdot(q-y))\,\mathrm{d}y\,\mathrm{d}x$
and it is easy to check that with $\tilde{\varphi}(p):=\varphi^{\prime}(p,1)$
it holds $\varphi^{\prime}(p,q)=\tilde{\varphi}(p\cdot q)$ for all $p,q>0$.
That $|\tilde{\varphi}(p)|\leq 1$ can be checked directly or follows from the
fact that $S_{\varphi}$ is unitary. ∎
###### Remark 5.6.
It is a general feature of asymptotically complete Borchers triples with
asymptotic algebra ${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ that
the restriction of the scattering matrix $S$ to $eSe$ is a functional calculus
of $P\otimes P$. Indeed, both $e$ and $S$ commute with the translation $T$,
but $T$ is maximally abelian when restricted to
${\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}\otimes{\mathcal{H}}_{\mathrm{{{\mathcal{A}}^{(0)}}}}^{1}$,
hence there is a function $\varphi_{S}$ such that
$eSe=\varphi_{S}(P\otimes{\mathbbm{1}},{\mathbbm{1}}\otimes P)$. Furthermore,
both $e$ and $S$ commute with boosts, so does $\varphi_{S}$ and one obtains
the form $eSe=\varphi_{S}^{\prime}(P\otimes P)$.
We note that the proof above shows that $|\tilde{\varphi}(M^{2}/2)|$ is the
probability that an improper state in
${\mathcal{H}}^{1}_{{\mathcal{A}}^{(0)}}\otimes{\mathcal{H}}^{1}_{{\mathcal{A}}^{(0)}}$
with mass $M^{2}$ is scattered elastically in the sense of Fock space
particles, where
$\tilde{\varphi}(p)=\frac{1}{p}\int_{0}^{p}\int_{0}^{1}\varphi((p-x)(1-y))\varphi(xy)\check{\varphi}((p-x)y)\check{\varphi}(x(1-y))\,\mathrm{d}y\,\mathrm{d}x\,.$
As we discussed in 3.1, the Hilbert space of the ${\rm U(1)}$-current net, and
hence the tensor product of two copies of it, admit the bosonic Fock space
structure, hence we can consider the particle number. Although we admit that
this concept does not have an intrinsic meaning, we claim that it is possible
to interpret this as the number of massless particles.
An evidence comes from the comparison with massive cases. In [Lec08] Lechner
has constructed a family of massive interacting models parametrized by so-
called scattering functions, and later he reinterpreted them as deformations
of the massive free field [Lec11]. If one applies the same deformation
procedure to the derivative of the massless free field whose net is
${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ (with scattering functions
satisfying $S_{2}(0)=1$), he obtains the Borchers triples with
${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ as the asymptotic net
constructed in [Tan11a] 222Private communication with Gandalf Lechner and Jan
Schlemmer. This will be presented elsewhere.. Hence the models in [Tan11a]
should be considered as the massless versions of the models in [Lec08].
Likewise, it can be said that the models constructed in the present paper are
the deformed (in an appropriate sense) version of the massless free field.
In massive case, there is a mass gap in the spectrum of the spacetime
translation and the one-particle space of the Fock space has an intrinsic
meaning. In massless case, such an intrinsic interpretation is lost but there
is still the Fock space structure. Thus we think that, if the two-particle
space in the Fock structure is not preserved by the S-matrix, as in the case
where $\varphi$ is not exponential (see Proposition 5.3), then it represents
massless particle production.
## 6 Conclusion and outlook
In this paper we have constructed a new family of Longo-Witten endomorphisms
on ${{\mathcal{A}}^{(0)}}$ through the inclusion
${{\mathcal{A}}^{(0)}}=\mathrm{Fer}_{\mathbb{C}}^{\rm
U(1)}\subset\mathrm{Fer}_{\mathbb{C}}$. We combined them to construct
interacting wedge-local nets with
${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$ as the asymptotic algebra
and showed that their S-matrices do not preserve the $n$-particle space of the
bosonic Fock space. Particle production is a necessary feature of interacting
models in higher dimensions [Aks65], thus this result opens up some hope for
algebraic construction of higher dimensional interacting models.
However, there are at least two shortcomings with the present method. The
first is that we proved only wedge-locality of the models. As already shown in
[Tan11a], a wedge-local net can be dilation-covariant and at the same time
interacting. On the other hand, a strictly local dilation-covariant
(asymptotically complete) net is necessarily not interacting [Tan11b]. Hence,
interaction of wedge-local nets could be just a false-positive and strict
locality is desired. The second is the fact that the concept of particle in
massless case is not intrinsically defined. Although the Fock space structure
is easily understood, its interpretations should be treated with care.
These issues could be overcome by considering massive cases. As for strict
locality, it has been shown that the deformation of the massive free field by
a suitably regular function is again strictly local [Lec08, Lec11]. On the
other hand, in massless situation, even the simplest case $\varphi(p)=-1$
(where $\varphi$ is an inner symmetric function used in [Tan11a] to deform
directly ${{\mathcal{A}}^{(0)}}\otimes{{\mathcal{A}}^{(0)}}$) is already not
strictly local [Tan11a]. Hence we believe that strict locality should be
addressed in massive models. Furthermore, for a massive asymptotically
complete model, the notion of particle production is intrinsic. Fortunately,
it is known that the construction in [Tan11a] coincides with the deformation
of the massive free field as we remarked in the last section, hence a further
correspondence between massive and massless cases are expected. We hope to
investigate this problem in a future publication.
Of course, interacting models in higher dimensions are always one of the most
important issues. Although conformal nets themselves are not interacting
[BF77], some new constructions based on CFT could be possible and ideas from
the present article could be useful.
#### Acknowledgment.
We thank our supervisor Roberto Longo for his constant support and useful
suggestions. Y. T. thanks Gandalf Lechner and Jan Schlemmer for discussions
on the relation between the present construction and the deformation of
[Lec11].
## References
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* [BS08] D. Buchholz and S.J. Summers. Warped convolutions: a novel tool in the construction of quantum field theories. In Quantum field theory and beyond, pages 107–121. World Sci. Publ., Hackensack, NJ, 2008.
* [BSM90] D. Buchholz and H. Schulz-Mirbach. Haag duality in conformal quantum field theory. Rev. Math. Phys., 2(1):105–125, 1990.
* [Buc75] D. Buchholz. Collision theory for waves in two dimensions and a characterization of models with trivial $S$-matrix. Comm. Math. Phys., 45(1):1–8, 1975.
* [CKL08] S. Carpi, Y. Kawahigashi, and R. Longo. Structure and classification of superconformal nets. Ann. Henri Poincaré, 9(6):1069–1121, 2008.
* [DF77] W. Driessler and J. Fröhlich. The reconstruction of local observable algebras from the euclidean green’s functions of relativistic quantum field theory. Annales de L’Institut Henri Poincare Section Physique Theorique, 27:221–236, 1977.
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* [Lec08] G. Lechner. Construction of quantum field theories with factorizing $S$-matrices. Comm. Math. Phys., 277(3):821–860, 2008.
* [Lec11] G. Lechner. Deformations of quantum field theories and integrable models. Commun. Math. Phys., 212:265–302, 2012.
* [Lon08] R. Longo. Real Hilbert subspaces, modular theory, ${\rm SL}(2,{\bf R})$ and CFT. In Von Neumann algebas in Sibiu: Conference Proceedings, pages 33–91. Theta, Bucharest, 2008.
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http://www.theorie.physik.uni-goettingen.de/~rehren/ps/cqft.pdf
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* [Tan11a] Y. Tanimoto. Construction of wedge-local nets of observables through longo-witten endomorphisms. Commun. Math. Phys., 314(2):443–469, 2012.
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|
arxiv-papers
| 2011-11-07T18:35:11 |
2024-09-04T02:49:24.082565
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/4.0/",
"authors": "Marcel Bischoff and Yoh Tanimoto",
"submitter": "Yoh Tanimoto",
"url": "https://arxiv.org/abs/1111.1671"
}
|
1111.1805
|
# Repulsive and attractive Casimir interactions in liquids
Anh D. Phan University of South Florida, Tampa, USA anhphan@mail.usf.edu N.
A. Viet Institute of Physics, Hanoi, Vietnam
###### Abstract
The Casimir interactions in the solid-liquid-solid systems as a function of
separation distance have been studied by the Lifshitz theory. The dielectric
permittivity functions for a wide range of materials are described by Drude,
Drude-Lorentz and oscillator models. We find that the Casimir forces between
gold and silica or MgO materials are both the repulsive and attractive. We
also find the stable forms for the systems. Our studies would provide a good
guidance for the future experimental studies on the dispersion interactions.
###### pacs:
Valid PACS appear here
## I Introduction
The dispersion interactions, the Casimir force, between neutral objects have
brought attraction for many years. There are a lot of factors affecting on the
value of force, such as geometry and material properties. Each of them gives
rise to hot subjects of ongoing investigation. Some experiments have examined
the influence of the dielectric properties of objects on the Casimir force 1 ;
2 ; 3 ; 6 . A number of settings used to study the interaction in terms of
theory are ideal metals, real metals and semiconductors 3 ; 4 ; 6 ,
metamaterials, and two objects placed in liquids 1 ; 2 ; 5 . These studies
have significantly advanced our understanding of the subtle effect of geometry
and material on the Casimir-Lifshitz interactions, especially for designing
nanodevices and nanotechnologies.
In the Lifshitz theory, the dispersion interactions primarily depend on
dielectric permittivity functions of materials. Changing dielectric function
alters the Casimir interactions. There are some ways to modify dielectric
functions, including illuminating a light on the silicon 7 ; 8 , which make
drifting carriers on semiconductor materials. In principle, there are some
models to describe dielectric response functions of real materials, for
example, plasma and Drude models for metals 6 ; 9 ; 11 , Drude-Lorentz and
oscillator models for liquids 2 ; 11 , oxides and others 10 ; 11 ; 12 . Based
on these models, the Casimir forces were obtained by numerical integrations
and series expansion methods 19 .
It has been theoretically shown that the attractive Casimir interaction always
occurs between two (non-magnetic) dielectric bodies related by reflection.
Therefore, the repulsive force is a striking feature creating inspiration for
scientists to make accurate measurements of nano electromechanical machines
where the repulsive force plays an important role and might resolve the
stiction problems. The repulsive Casimir forces can be observed in systems
which have the presence of liquids 2 , metamaterials and metallic geometries
16 . Recent experiments have pointed out that there a repulsive force exists
between a gold sphere and a silica plate, separated by bromobenzen 2 . As a
matter of fact, the repulsive Casimir forces between solids arise when the
dielectric of material surfaces 1 and 2 and an intervening liquid obey the
relation
$\varepsilon_{1}(i\xi)>\varepsilon_{liquid}(i\xi)>\varepsilon_{2}(i\xi)$ over
a wide imaginary frequency range $\xi$.
A previous theoretical 14 study has noticed that it is difficult to establish
an equilibrium configuration of sytems in a vacuum medium. In the reference 15
, the authors showed that they were able to form some stable configurations of
Teflon-Si and Silica-Si immersed in ethanol. The equilibrium is explicitly
explained by dispersion properties. In the present work, our theoretical
studies have shown that the equilibra can be obtained by placing Au-MgO,
Silica-MgO and Au-Silica sytems in bromobenzen.
In this paper, the Casimir-Lifshitz forces between material plate systems made
in oxides and metals immersed in bromobenzen are calculated. The combination
between these results and the proximity force approximation (PFA) method allow
us to compute the Casimir interactions in different configurations. We find
that the magnitude of the Casimir force between two dielectric bodies depends
on the configuration and distance between two bodies. The shape usually used
in experiments is a combination of a sphere and a plate because one can avoid
the problem of alignment and easily control the distance between them. The
energy interactions between a plate-plate system per unit area can be obtained
by using the relationship between the Casimir energy of two plannar objects
and the dispersion force of a sphere-plate system.
The rest of the paper is organized as follows: In Sec. II the theoretical
formulations of Casimir-Lifshitz force interaction are introduced. In section
III, the numerical results for the Casimir force between two bodies are
presented. Important conclusions and discussions are finally given in section
IV.
## II Lifshitz theory for force calculations
For the force calculations, we used Lifshitz theory without considering effect
of temperature. The separations used here were less than 1 $\mu$m, therefore
thermal corrections at $T=300K$ are not significant. As previously noted in 2
; 3 ; 20 ; 21 , the Lifshitz formula at zero temperature for the Casimir force
acting between between two parallel flat bodies per unit area, separated by a
distance $d$ are given by
$\displaystyle F(d)$ $\displaystyle=$
$\displaystyle-\dfrac{\hbar}{2\pi^{2}}\int_{0}^{\infty}qk_{\perp}dk_{\perp}\int_{0}^{\infty}d\xi$
(1)
$\displaystyle\times\left(\dfrac{r_{TM}^{(1)}r_{TM}^{(2)}}{e^{2qd}-r_{TM}^{(1)}r_{TM}^{(2)}}+\dfrac{r_{TE}^{(1)}r_{TE}^{(2)}}{e^{2qd}-r_{TE}^{(1)}r_{TE}^{(2)}}\right).$
Here the reflection coefficients $r_{TM,TE}^{(1)}$ and $r_{TM,TE}^{(2)}$ for
two independent polarizations of the electromagnetic field (transverse
magnetic and transverse electric fields) are
$\displaystyle
r_{TM}^{(p)}=r_{TM}^{(p)}\left({\xi,k_{\bot}}\right)=\frac{{\varepsilon^{(p)}\left({i\xi}\right)q-\varepsilon^{(2)}\left({i\xi}\right)k^{(p)}}}{{\varepsilon^{(p)}\left({i\xi}\right)q+\varepsilon^{(2)}\left({i\xi}\right)k^{(p)}}},$
(2) $\displaystyle
r_{TE}^{(p)}=r_{TE}^{(p)}\left({\xi,k_{\bot}}\right)=\frac{{\mu^{(2)}(i\xi)k^{(p)}-\mu^{(p)}(i\xi)q}}{{\mu^{(2)}(i\xi)k^{(p)}+\mu^{(p)}(i\xi)q}},$
(3)
where
$\displaystyle
q=\sqrt{k_{\bot}^{2}+\varepsilon^{(2)}\left({i\xi}\right)\mu^{(2)}\left({i\xi}\right)\frac{{\xi^{2}}}{{c^{2}}}},$
(4) $\displaystyle
k^{(p)}=\sqrt{k_{\bot}^{2}+\varepsilon^{(p)}\left({i\xi}\right)\mu^{(p)}\left({i\xi}\right)\frac{{\xi^{2}}}{{c^{2}}}}.$
(5)
in which $\varepsilon^{(p)}\left({i\xi}\right)$ and
$\mu^{(p)}\left({i\xi}\right)$ are the dielectric permittivity and the
magnetic permeability of the first body (p = 1) and the second body (p=3),
respectively. $\varepsilon^{(2)}\left(\omega\right)$ and
$\mu^{(2)}\left({i\xi}\right)$ are the dielectric function and the
permeability of a liquid filled between two bodies. Here, medium ‘2’ selected
is a bromobenzen so $\mu^{(2)}\left({i\xi}\right)=1$. Moreover, in this paper,
the non-magnetic materials used such as germanium, gold and oxides have also
$\mu^{(p)}\left({i\xi}\right)=1$. $k_{\bot}$ magnitude of the wave vector
component perpendicular on the plate, is frequency variable along the
imaginary axis ($\omega=i\xi$).
We recall that Lifshitz formula, routinely used to interpret current
experiments, express the Casimir force between two parallel plates as an
integral over imaginary frequencies $i\xi$ of a quantity involving the
dielectric permittivities of the plates $\omega=i\xi$. It is important to note
that, in principle, recourse to imaginary frequencies is not mandatory because
it is possible to rewrite Lifshitz formula in a mathematically equivalent
form, involving an integral over the real frequency axis. In this case,
however, the integrand becomes a rapidly oscillating function of the
frequency, which hampers any possibility of numerical evaluation. Another
remarkable point is that occurrence of imaginary frequencies in the expression
of the Casimir force is a general feature of all recent formalisms hence
extending Lifshitz theory to non-planar geometries 17 ; 18 . The problem is
that the electric permittivity $\varepsilon(i\xi)$ at imaginary frequencies
cannot be measured directly by any experiment. The only way to determine it by
means of dispersion relations, which allow the expression of
$\varepsilon(i\xi)$ in terms of the observable real-frequency electric
permittivity $\varepsilon(i\xi)$. In the standard works on the Casimir effect,
$\varepsilon(i\xi)$ is expressed with the Kramers-Kronig relation in terms of
an integral of a quantity involving the imaginary part of the electric
permittivity 13
$\displaystyle\varepsilon(i\xi)=1+\frac{2}{\pi}\int\limits_{0}^{\infty}{d\omega\frac{{\omega{\mathop{\rm
Im}\nolimits}\varepsilon(\omega)}}{{\omega^{2}+\xi^{2}}}},$ (6)
where ${\mathop{\rm Im}\nolimits}\varepsilon(\omega)$ is calculated using the
tabulated optical data for the complex index of refraction.
The well-known dielectric function described for gold is the Drude model 13
$\displaystyle\varepsilon(i\xi)=1+\frac{{\omega_{p}^{2}}}{{\xi(\xi+\gamma)}},$
(7)
where $\omega_{p}=9.0$ eV, $\gamma=0.035$ eV are the plasma frequency and the
relaxation parameter of Au, respectively.
The imaginary part of the resulting dielectric function at 6 and 295 K of pure
MgO are shown in 12 . The optical features have been fitted to a classical
oscillator model using the complex dielectric function
$\displaystyle\varepsilon(\omega)=\varepsilon_{\infty}+\sum\limits_{j}{\frac{{\omega_{p,j}^{2}}}{{\omega_{TO,j}^{2}-\omega^{2}-i2\omega\gamma_{i}}}},$
(8)
where $\varepsilon_{\infty}$ is a high-frequency contribution, and
$\omega_{TO,j},{2\gamma_{i}}$ and $\omega_{p,j}$ are the frequency, full width
and effective plasma frequency of the $j$th vibration. The values of these
parameters can be found in 12 . Of course with such simple model for the
permittivity of MgO, there is no need to use dispersion relations to obtain
the expression of $\varepsilon(i\xi)$, for this can be simply done by the
substitution $\omega\to i\xi$ in the r.h.s of Eq.(8) 24
$\displaystyle\varepsilon(i\xi)=\varepsilon_{\infty}+\sum\limits_{j}{\frac{{\omega_{p,j}^{2}}}{{\omega_{TO,j}^{2}+\xi^{2}+2\xi\gamma_{j}}}}.$
(9)
In the case of bromobenzen and silica, it has recently been used for
measurement of repulsive forces between gold and silica surfaces. Extremely
weak repulsion was measured, indicating that the dielectric functions of
bromobenzen and silica are very similar in magnitude. In fact, oscillator
models are constructed to represent the dielectric function at imaginary
frequencies. The form of the oscillator model is given by
$\displaystyle\varepsilon(i\xi)=1+\sum_{i}{\frac{{C_{i}}}{{1+\xi^{2}/\omega_{i}^{2}}}},$
(10)
where the coefficients $C_{i}$ are the oscillator’s strengths corresponding to
(resonance) frequencies $\omega_{i}$ 1 ; 2 ; 25 . The dielectric data was
fitted in a wide frequency range 2 . They are more accurate in comparison with
other simple oscillator models. Moreover, many older references used limited
dielectric data, so the oscillator models with second or third order may lead
to the difference in Casimir force calculations. The parameters we used in the
present paper for bromobenzen and also silica come from 2 .
## III Numerical results and discussions
The Casimir attractive force usually occurs in experiments and theoretical
calculations. When bromobenzen is filled in the gap between two bodies, the
Casimir force is attractive if the dielectric functions do not satisfy one
condition
$\varepsilon_{1}(i\xi)>\varepsilon_{liquid}(i\xi)>\varepsilon_{2}(i\xi)$ for
all frequencies $\xi$. Therefore, by describing Fig. 1 as the dielectric
response function as a function of the frequency gives us some predictions of
repulsive and attractive forces.
Figure 1: (Color online) The dielectric function of various materials plotted
at imaginary frequencies $\xi$.
This graph shows that
$\varepsilon_{Au}(i\xi)>\varepsilon_{MgO}(i\xi)>\varepsilon_{liquid}(i\xi)$
and
$\varepsilon_{MgO}(i\xi)>\varepsilon_{Au}(i\xi)>\varepsilon_{liquid}(i\xi)$ at
$\xi<6.5$ eV, thus the interactions between Au and MgO body immersed in
bromobenzen liquid and in vacuum are attractive in this range. In the range of
$\xi>6.5$ eV,
$\varepsilon_{MgO}(i\xi)>\varepsilon_{liquid}(i\xi)>\varepsilon_{Au}(i\xi)$,
it causes the repulsive interaction. Similarly, in the gold-bromobenzen-silica
system, at extremely small frequencies , the forces are attractive. In the
larger frequency region, the Casimir forces are repulsive. Besides, the
similar explainations are applied to understand the interaction in the MgO-
bromobenzen-Au system. The numerical calculations of the normalized Casimir
force are provided in Fig. 2
Figure 2: (Color online) Relative Casimir force between two semi-infinite
plates normalized by the perfect metal force $F_{o}(d)=-\pi^{2}\hbar
c/240d^{4}$. The liquid used in this calculation is bromobenzen.
In the MgO-bromobenzen-Silica system, it can be clearly seen that there are
two positions in each curve where the Casimir force is equal to zero. The
first points are corresponding to unstable equilibria $d_{us}^{(1)}\approx 13$
nm because the interaction force changes from the attactive force to the
repulsive force, the second point $d_{s}^{(1)}\approx 110$ nm is a stable
position. There is only one position in the Au-bromobenzen-Silica system and
the Au-bromobenzen-MgO sytem, The interaction forces disappear at
$d_{s}^{(2)}\approx 275$ nm and $d_{s}^{(3)}\approx 5.5$ nm, stable position
of each sytem, respectively.
Figure 3: (Color online) Schematic picture of the setting considered in our
calculations. A sphere is located in bromobenzen at a distance $d$ away from a
material plate.
In order to consider the Casimir interactions between a spherical body and a
plate at a distance of close approach $d$ at a temperature T = 300 K, it is
very useful to utilize the PFA method to calculate. Experimental results for
the Casimir force in the plane-sphere geometry are usually compared with PFA-
based theoretical models. The spherical surface is assumed to be nearly flat
over the scale of $d$. Although the Casimir force is not additive, PFA is
often expected to provide an accurate description when $R\gg d$. Here, the
radius of Au sphere that is used in configurations is $R=40$ $\mu$m in order
to calculate Casimir interaction by the proximity force approximation (PFA)
method because the ratio of $d$ to $R$ is small enough to PFA results becoming
enormously accurate. It can be described by Fig. 3. In this approach, the
surfaces of the bodies are treated as a superposition of infinitesimal
parallel plates 22 .
$\displaystyle
F_{sp}^{PFA}(d)=\int\limits_{0}^{R}{F_{pp}(d+R-\sqrt{R^{2}-r^{2}})}2\pi rdr.$
(11)
here ${F_{pp}}$ is the Casimir force for two parallel plates of unit area.
When using the PFA method, one important point is that the interactions
between a gold sphere or a magnesium oxide sphere and a silica plate are equal
to the interactions between a magnesium oxide plate, which has the same
radius, and a gold plate or a silica plate. There is no difference in
calculations and results as well because the PFA method does not consider a
structure of bodies when their shape is modified or is spherical or
cylindrical shape. The equivalent situations occur in other materials. The
resulting Casimir forces are shown in Fig. 4.
Figure 4: (Color online) The Casimir forces of various sphere-bromobenzen-
plate systems are estimated as a function of separation, described in the text
with the spherical redius $R=40$ $\mu$m.
In the reference 2 , the authors experimentally measured and theoretically
calculated the Casimir interaction between a gold sphere and a silica plate
immersed in bromobenzen in the range from $20$ nm to $60$ nm. Our results in
this range are the same for this range. But when we extend the considered
range of distance, the attractive-repulsive transition occurs at approximately
$190$ nm. This position makes this system stable. Another consequence of Fig.
4 demonstrates that stable position of the Au-bromobenzen-MgO sytem moves to
$3.5$ nm to balance between the attractive and repulsive forces. It can be
explained that increasing the separation distance of infinitesimal parallel
plates causes the fast reduction of the dispersion interaction. At the same
minimal separation distance $d$, the attractive force acting on a sphere is
less than that of a plate in the same effective area. Finally, in the system
of a MgO sphere and a silica plate embeded in bromobenzen, there is only one
presence of non-interaction posion at nearly $10$ nm. It is an unstable
position.
Figure 5: (Color online) The Casimir energy is calculated as a function of
separation for different materials .
In addition, PFA formula and Eq.(11) can allow us to estimate the Casimir
energy per unit area between two plate bodies illustrated in Fig. 5. The
Casimir energy is approximated by 22
$\displaystyle F_{sp}^{PFA}(d)=2\pi RE(d),$ (12)
where $E(d)$ is the Casimir energy per a unit area for planar bodies.
We have also applied the PFA method to calculate the Casimir force in sphere-
sphere systems, we continue to calculate by the PFA method. The formula for
this calculation is given
$\displaystyle F_{ss}^{PFA}(d)=2\pi\int\limits_{0}^{R_{2}}rdr\times$
$\displaystyle
F_{pp}(d+R_{1}-\sqrt{R_{1}^{2}-r^{2}}+R_{2}-\sqrt{R_{2}^{2}-r^{2}}),$ (13)
where the radii of two spherical objects are ${R_{1}}$ and ${R_{2}}$,
respectively. It is assumed that $R_{2}<R_{1}$. In this study, we consider
$R_{1}=40$ $\mu$m and the case of $R_{2}=R_{1}$, $R_{1}=2R_{2}$ and
$R_{1}=2R_{2}$.
Here, having calculated $F_{ss}^{PFA}(d)$ in a sphere-sphere system using
Eq.(13) and $F_{sp}^{PFA}(d)$ in a sphere-plate system using Eq.(11). These
results obtained show that, when increasing $d$, the ratio
$F_{ss}^{PFA}(d)/F_{sp}^{PFA}(d)$ does not depend on the distance $d$. It is a
constant with its magnitude as a function of the radius of two spheres,
$F_{ss}^{PFA}(d)/F_{sp}^{PFA}(d)=1/2$ when $R_{1}=R_{2}$ and
$F_{ss}^{PFA}(d)/F_{sp}^{PFA}(d)=1/3$ when $R_{1}=2R_{2}$. Generalizing this
ratio, if $R_{1}=nR_{2}$, we have $F_{ss}^{PFA}(d)/F_{sp}^{PFA}(d)=1/(n+1)$.
This character is likely to be explained by the results in 22 . When the
second sphere is extremely small in comparison with the first one, the
interaction force goes to zero. In this case, the Lifshitz formula used to
calculate the Casimir force should be transformed to the Casimir-Polder
formula describing the interaction between an atom and a microscopic object.
Moreover, the PFA method is not accurate because this approach is useful if
the size of objects is much larger than the separation between them. On the
other hand, we have
$\displaystyle F_{ss}^{PFA}(d)=\frac{1}{n+1}F_{sp}^{PFA},$ (14)
where $R_{1}=nR_{2}$. If $F_{sp}^{PFA}=0$, $F_{ss}^{PFA}$ must be zero.
Therefore, the unstable and stable positions are constant and unchanged when a
radius of a second sphere varies.
One demonstrated that the Casimir force between two objects embedded in
liquids can be derived from the well-known Lifshitz formula at least if the
object is not made of nonabsorbing materials 5 . That explains why the
Lifshitz expressions is used in order to calculate the Casimir force and
compare with experimental data. Nothing changes in the dielectric functions of
bodies immersed in liquids. On the other hand, several experiments verified
that when metals are placed in liquids, there is a variation of Drude
parameters in the metal 23 . The discrepacy of the interaction between “dry”
and “wet” can reach to 15 $\%$ in this case. But they measured the Casimir
force between two metal plates and got the error. Besides, maybe the
dielectric of liquids and low index materials play much more important role in
Casimir force. In the reference 2 , the change of Drude parameters are not
taken into account but the theoretical calculations are close to the
experimental data curves when we have liquids and the low index materials.
## IV Conclusions
In this work, we have extended the Lifshitz theory to calculate the Casimir
force. Liquid, silica and magnesium oxide are represented by oscillator
models. Although further studies are required to determine the repulsive
Casimir force accurately, our results show that MgO and silica is a good
candidate for the demonstration of quantum levitation. The contribution of
bromobenzen is important because it is an important factor making the purely
repulsive force or the repulsive-attractive transition. After calculating the
Casimir force between two bodies per unit area and associating proximity force
approximation method, it is easy to compute the interaction of different
material plates with a material sphere. Based on the formula of the Casimir
force between a sphere and a plate, it is convenient to estimate the free
energy interaction of bodies. The result is a prediction for further
experimental studies.
###### Acknowledgements.
This work was supported by the Nafosted Grant No. 103.02.57.09. We thank Prof.
Lilia M. Woods and Prof. P. J. van Zwol for helpful discussions and comments.
## References
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* (18) O. Kenneth, and I. Klich, Phys. Rev. B 78, 014103 (2008).
* (19) R. S. Decca, E. Fischbach, G. L. Klimchitskaya, D. E. Krause, D. L ̵́ pez, and V. M. Mostepanenko, Phys. Rev. A 84, 042502 (2011).
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|
arxiv-papers
| 2011-11-08T05:11:39 |
2024-09-04T02:49:24.100820
|
{
"license": "Public Domain",
"authors": "Anh D. Phan, N. A. Viet",
"submitter": "Anh Phan",
"url": "https://arxiv.org/abs/1111.1805"
}
|
1111.1810
|
Harmonic analysis of the functions $\tilde{\Delta}(x)$ and $N(T)$
Jining Gao
Department of Mathematics, Shanghai Jiaotong University, Shanghai ,P. R. China
In this paper, under the Riemann hypothesis, we study the Fourier analysis
about the functions $\tilde{\Delta}(x)$ and $N(T)$ .
## 1 INTRODUCTION
Riemann hypothesis has been studied in many different ways, in this paper, we
will try to use somewhat new angles to study RH. Most results of this paper
are obtained under the RH. As we know, Guinand formula is a representation of
$N(T)$ which is the distribution function of Riemann zeros in term of series
of the prime number powers, although Guinand formula [2]is a result under the
assumption of RH,it provides an explicit method to figure out all non trivial
Riemann zeros. Actually, this fact is far from trivial because once we have
prime number representation of $N(T)$ (Guinand formula) at hand, we can
immediately restore a function via the distribution of it’s zeros,so Guinand
formula is equivalent to RH and we will prove it in the first section. Since
Guinand formula is very important in this paper and Guinand original proof is
complicated and full of the favor of harmonic analysis, we will first of all
give another simple and elementary proof based on the lemma, which is the
ground stone of this paper, besides ,our new proof gives out a stronger
conclusion than the original statement of Guinand formula. This stronger
result will help us to check the truth of RH much more efficiently. In the
second section we rewrite Guinand formula and Riemann-Mangoldt formula as two
integral equations of two ”functional variables” $\tilde{\Delta}(x)$ and
$S(T)$, which seems to imply Guinand formula and Riemann-Mangoldt formula are
reciprocal to each other and such integral representation will be used in the
4th section. In the third section ,first of all, we derive an elementary
formula based on functional equation of Riemann zeta function and lemma. This
formula provides infinitely many non trivial integral equations of $N(T)$,
also, we use the elementary formula to prove a theorem which claims
$\mid\tilde{\Delta}(x)\mid$ has a non-zero measurement of a positive lower
bound.
## 2 Guinand formula with an error term and it’s inverse theorem
In this section, we will give out a proof of Guinand formula with the
uniformly convergent error term, besides, we also give out an inverse theorem
of Guinand formula. First of all we need following notations and formulas
which will be used throughout this paper[1],[3].
Chebyshev function
$\displaystyle\psi(x)=\sum_{n<x}\Lambda(n)$
Where the Von Mangoldt function $\Lambda(n)=logp$ if $n=p^{k}$ for some $k$
and some prime number $p$ ,$\Lambda(n)=0$ otherwise.
###### Theorem 1
( Mangoldt and Riemann explicit formula)
$\displaystyle\psi(x)=x-lim_{T\rightarrow\infty}\sum_{\left|Im\rho\right|<T}\frac{x^{\rho}}{\rho}-log(2\pi)+\sum_{n=1}^{\infty}\frac{x^{-2n}}{2n}$
Where $\rho$ runs through all non-trivial Riemann zeros.
###### Theorem 2
[3]
$\displaystyle\psi(x)=x-\sum_{\left|Im\rho\right|<T}\frac{x^{\rho}}{\rho}+O(\frac{xlog^{2}x}{T})$
we set
$\displaystyle\tilde{\psi}(x)=\frac{x^{2}}{2}-\sum_{\rho}\frac{x^{\rho+1}}{\rho(\rho+1)}-xln(2\pi)-\sum_{n=1}^{\infty}\frac{x^{-2n+1}}{2n(2n-1)}$
(1)
It’s easy to prove that when $x$ is not equal to any integer,$\tilde{\psi}(x)$
is differentialble and it’s derivative is just $\psi(x)$ and it’s continuous
when $x>0$ [1]. We set $\Delta(x)=\psi(x)-x$ and
$\tilde{\Delta}(x)=\tilde{\psi}(x)-\frac{x^{2}}{2}$ .
The following lemma will is important for deducing the Guinand formula with
the error term.
###### Lemma 3
when $s\neq\rho$,
$\displaystyle\frac{\zeta^{\prime}}{\zeta}(s)=-\sum_{n<X}\frac{\Lambda(n)}{n^{s}}+\psi(X)X^{-s}+s\tilde{\psi}(X)X^{-s-1}-$
$\displaystyle\frac{s(s+1)}{2(s-1)}X^{1-s}+\sum_{\rho}\frac{s(s+1)X^{\rho-s}}{\rho(\rho+1)(s-\rho)}+\sum_{n\geq
1}\frac{s(s+1)X^{-2n-2s}}{2n(2n-1)(s+2n)}$ (2)
Proof. Let
$f_{X}(s)=\sum_{n<X}\frac{\Lambda(n)}{n^{s}}$
Using integration by parts twicely, we have that
$\displaystyle
f_{X}(s)=\int_{1}^{X}x^{-s}d\psi(x)=\psi(X)X^{-s}-\int_{1}^{X}\psi(x)dx^{-s}=\psi(X)X^{-s}+s\int_{1}^{X}\psi(x)x^{-s-1}dx$
$\displaystyle=\psi(X)X^{-s}+s\int_{1}^{X}x^{-s-1}d\tilde{\psi}(x)$
$\displaystyle=\psi(X)X^{-s}+s\tilde{\psi}(X)X^{-s-1}+s(s+1)\int_{1}^{X}\tilde{\psi}(x)x^{-s-2}dx$
(3)
and by formula 1,we can further get
$\displaystyle\int_{1}^{X}\tilde{\psi}(x)x^{-s-2}dx=\int_{1}^{X}\frac{\frac{x^{2}}{2}-\sum_{\rho}\frac{x^{\rho+1}}{\rho(\rho+1)}-ln(2\pi)x-\sum_{n=1}^{\infty}\frac{x^{-2n+1}}{2n(2n-1)}}{x^{s+2}}dx$
$\displaystyle=\frac{1}{2}(\frac{X^{1-s}}{1-s}-\frac{1}{1-s})-\sum_{\rho}\frac{1}{\rho(\rho+1)}(\frac{X^{\rho-s}}{\rho-s}-\frac{1}{\rho-s})$
$\displaystyle+ln(2\pi)(\frac{X^{-s}}{s}-\frac{1}{s})+\sum_{n\geq
1}\frac{1}{2n(2n-1)}(\frac{X^{-s-2n}}{s+2n}-\frac{1}{s+2n})$
We collect all above terms as two groups $J_{X}(s)$ and $I(s)$,obviously,
$I(s)=\frac{1}{2}\frac{1}{s-1}-\sum_{\rho}\frac{1}{\rho(\rho+1)}\frac{1}{s-\rho}-\frac{ln(2\pi)}{s}-\sum_{n\geq
1}\frac{1}{2n(2n-1)}\frac{1}{s+2n}$
By formula 3 and using the notations $J_{X}(s)$ and $I(s)$, we get
$\displaystyle
f_{X}(s)=\psi(X)X^{-s}+s\tilde{\psi}(X)X^{-s-1}+s(s+1)J_{X}(s)+s(s+1)I(s)$ (4)
and
$s(s+1)I(s)=\frac{s(s+1)}{2(s-1)}-\sum_{\rho}\frac{s(s+1)}{\rho(\rho+1)(s-\rho)}-(s+1)ln(2\pi)-\sum_{n\geq
1}\frac{s(s+1)}{2n(2n-1)(s+2n)}$
By the following identity,
$\displaystyle\frac{s(s+1)}{z(z+1)(s-z)}=\frac{s}{z(z+1)}+\frac{1}{s-z}+\frac{1}{z}$
we have that
$\displaystyle s(s+1)I(s)=-\frac{\zeta^{\prime}}{\zeta}(s)+as+b$ (5)
where $a,b$ are some constants which can be determined immediately. According
4 and 5 ,we get a new representation of $\frac{\zeta^{\prime}}{\zeta}(s)$ when
$s\neq\rho$ as follows:
$\displaystyle\frac{\zeta^{\prime}}{\zeta}(s)=-\sum_{n<X}\frac{\Lambda(n)}{n^{s}}+\psi(X)X^{-s}+s\tilde{\psi}(X)X^{-s-1}$
$\displaystyle-\frac{s(s+1)}{2(s-1)}X^{1-s}+\sum_{\rho}\frac{s(s+1)X^{\rho-s}}{\rho(\rho+1)(s-\rho)}+\sum_{n\geq
1}\frac{s(s+1)X^{-2n-s}}{2n(2n-1)(s+2n)}+as+b$ (6)
Since when $Res>1$,
$\frac{\zeta^{\prime}}{\zeta}(s)=-\sum_{n=1}^{\infty}\frac{\Lambda(n)}{n^{s}}$
Let $X\rightarrow\infty$ on the right side of 6 when $Res>1$, we immediately
get $a=b=0$, that follows our theorem.
As we know,
$log\zeta(s_{0})=\int_{2}^{s_{0}}\frac{\zeta^{\prime}}{\zeta}(s)ds$
,where the integral path is a positive orient half rectangle with vertices
$2,2+iT,\sigma+iT$ and $s_{0}=\sigma+iT$ $s_{0}\neq\rho$. Taking this complex
integral on both sides of 6, we directly get following theorem:
###### Theorem 4
When $s_{0}\neq\rho$, we have that
$\displaystyle
log\zeta(s_{0})=\sum_{n<X}\frac{\Lambda(n)}{(logn)n^{s_{0}}}-\frac{\Delta(X)}{(logX)X^{s_{0}}}-\frac{\tilde{\Delta}(X)}{(logX)X^{s_{0}+1}}(s_{0}+\frac{1}{lnX})$
$\displaystyle+\int_{2}^{s_{0}}\frac{X^{1-s}}{1-s}ds+\tilde{J}_{X}(s_{0})+C_{0}$
(7)
Where
$\displaystyle\tilde{J}_{X}(s_{0})=\int_{2}^{s_{0}}s(s+1)J(X)ds=-\frac{1}{lnX}\sum_{\rho}\frac{1}{\rho(\rho+1)}[\frac{s_{0}(s_{0}+1)X^{\rho-
s_{0}}}{s_{0}-\rho}$
$\displaystyle-\int_{2}^{s_{0}}X^{\rho-s}(\frac{2s+1}{s-\rho}-\frac{s^{2}+s}{(s-\rho)^{2}})ds]$
$\displaystyle-\frac{1}{lnX}\sum_{n\geq
1}\frac{1}{2n(2n-1)}[\frac{s_{0}(s_{0}+1)X^{-2n-s}}{s+2n}-\int_{2}^{s_{0}}X^{-2n-s}(\frac{2s+1}{s+2n}-\frac{s^{2}+s}{(s+2n)^{2}})ds]$
and $C_{0}$ is a real constant.
Proof. Using integration by parts and collecting all terms containing
$X^{1-s}$, we immediately get above results. setting $s_{0}=\frac{1}{2}+iT$ in
the formula 7 and taking imaginary parts on both sides, we have that
###### Theorem 5
If the Riemann hypothesis is true, and $\delta$ is the distance between $T$
and the coordinate of the nearest Riemann zero,we have
$\displaystyle\pi
S(T)=-\sum_{n<X}\frac{\Lambda(n)sin(Tlogn)}{\sqrt{n}}+\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}+Im(\int_{2}^{\frac{1}{2}+iT}\frac{X^{1-s}}{1-s}ds)+O(\frac{T^{3}}{\delta^{2}lnX})$
(8)
in the limit language, we have
$\displaystyle\pi
S(T)=-lim_{X\rightarrow\infty}[\sum_{n<X}\frac{\Lambda(n)sin(Tlogn)}{\sqrt{n}(logn)}$
$\displaystyle-\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}-Im(\int_{2}^{\frac{1}{2}+iT}\frac{X^{1-s}}{1-s}ds)]$
(9)
From now on ,we will prove formula 9 is the same as Guinand formula. To
achieve it, we need to make some simplification as follows:
Let’s first simplify the term
$Im(\int_{2}^{\frac{1}{2}+iT}\frac{X^{1-s}}{1-s}ds)$
,Let’s transform the original integral path which is half rectangle with
vertices $2,2+iT,\frac{1}{2}+iT$ to another half rectangle with vertices
$2,\frac{1}{2},\frac{1}{2}+iT$ and orient is clockwise, we get
$\displaystyle\int_{2}^{\frac{1}{2}+iT}\frac{X^{1-s}}{1-s}ds=i\int_{0}^{T}\frac{X^{\frac{1}{2}-it}}{\frac{1}{2}-it}dt+\int_{\Gamma_{r}}\frac{X^{1-s}}{1-s}ds$
(10)
Where $\gamma_{r}=[\frac{1}{2},1-r]\cup S_{r}\cup[1+r,2]$ and $S_{r}$ is upper
half semi-circle with radius $r$ and centered at $z=1$ Taking imaginary part
on both side of 10,we get the first term of right hand side is equal to
$\sqrt{X}\int_{0}^{T}\frac{2cos(tlogX)+4tsin(tlogX)}{1+4t^{2}}dt$
which is set to be $f_{1}(T,X)$ For the second term of right side of 10, we
have that
$\displaystyle
Im(\int_{\Gamma_{r}}\frac{X^{1-s}}{1-s}ds)=Im(\int_{J_{r}}\frac{X^{1-s}}{1-s}ds)$
$\displaystyle=lim_{r\rightarrow 0}Im(\int_{J_{r}}\frac{X^{1-s}}{1-s}ds)=-\pi$
(11)
Let’s pick up the second term of righ side of 47 i.e.
$\int_{1}^{X}\frac{sin(Tlogt)}{\sqrt{t}logt}dt$ and set it to be $f_{2}(T,X)$,
Since $\frac{sin(Tlogt)}{\sqrt{t}logt}$ is continously differentiable with
respect to $T$, we have that
$\displaystyle\frac{\partial f_{2}}{\partial
T}=\int_{1}^{X}\frac{cos(Tlogt)}{\sqrt{t}}dt$
$\displaystyle=\int_{0}^{lnX}e^{\frac{1}{2}u}cos(Tu)du=\frac{\sqrt{X}}{1+4T^{2}}[2cos(TlogX)+4Tsin(TlogX]-\frac{2}{1+4T^{2}}$
(12)
in which we have used the substituation $u=logt$
We can also notice that
$\frac{\partial f_{1}}{\partial
T}=\frac{\sqrt{X}}{1+4T^{2}}[2cos(TlogX)+4Tsin(TlogX]$
By 12, we get
$\displaystyle\frac{\partial f_{1}}{\partial T}-\frac{\partial f_{2}}{\partial
T}=\frac{2}{1+4T^{2}}$
Thus,
$\displaystyle f_{1}(T,X)-f_{2}(T,X)=\int_{0}^{T}\frac{2}{1+4t^{2}}dt$ (13)
$\displaystyle=arctan2T$
Consequently,by 10, 11,13
$\displaystyle
Im(\int^{\frac{1}{2}+iT}\frac{X^{1-s}}{1-s}ds)-\int_{1}^{X}\frac{sin(Tlogt)}{\sqrt{t}logt}dt=arctan2T-\pi$
(14)
Let’s single out the term
$\frac{sin(TlogX)}{logX}\left\\{\sum_{n<X}\Lambda(n)n^{-\frac{1}{2}}-2X^{\frac{1}{2}}\right\\}$
in right hand of 47 and get it simplified as follows:
$\displaystyle\frac{sin(TlogX)}{logX}\left\\{\sum_{n<X}\Lambda(n)n^{-\frac{1}{2}}-2X^{\frac{1}{2}}\right\\}$
$\displaystyle=\frac{sin(TlogX)}{logX}[\int_{1}^{X}x^{-\frac{1}{2}}d\psi(x)-2\sqrt{X}]$
$\displaystyle=\frac{sin(TlogX)}{logX}[\psi(X)X^{-\frac{1}{2}}+\frac{1}{2}\int_{1}^{X}\psi(x)x^{-\frac{3}{2}}dx-2\sqrt{X}]$
$\displaystyle=\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}+\frac{sin(TlogX)}{2logX}\int_{1}^{X}\Delta(x)x^{-\frac{3}{2}}dx$
(15)
and by the theorem 2 , we have that
$\displaystyle\int_{1}^{X}\Delta(x)x^{-\frac{3}{2}}dx=-\sum_{\rho}\frac{X^{\rho-\frac{1}{2}}-1}{\rho(\rho-\frac{1}{2})}+2ln(2\pi)(X^{-\frac{1}{2}}-1)$
$\displaystyle-\sum_{n\geq
1}\frac{X^{-2n-\frac{1}{2}}-1}{2n(2n+\frac{1}{2})}=O(1)$
With formula 15, we have
$\displaystyle\frac{sin(TlogX)}{logX}[\sum_{n\leq
X}\Lambda(n)n^{-\frac{1}{2}}-2X^{\frac{1}{2}}]=\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}+O(\frac{1}{logX})$
(16)
Since
$\displaystyle N(T)=\frac{1}{\pi}arg\xi(\frac{1}{2}+iT)$
$\displaystyle=\frac{1}{\pi}args(s-1)\pi^{-\frac{s}{2}}\Gamma(\frac{s}{2})\zeta(s)|_{s=\frac{1}{2}+iT}$
$\displaystyle=\frac{1}{\pi}arg(-\frac{1}{4}-T^{2})-\frac{Tln\pi}{2\pi}+\frac{1}{\pi}arg\Gamma(\frac{1}{4}+\frac{iT}{2})+S(T)$
$\displaystyle=1-\frac{Tln\pi}{2\pi}+\frac{1}{\pi}arg\Gamma(\frac{1}{4}+\frac{iT}{2})+S(T)$
(17)
Whenever $T$ is not equal to any cordinates of some Riemann zeros,we can
rewrite Guinand formula 47 as follows
$\displaystyle\pi
S(T)=F_{X}(T)+arctan(2T)-\pi+\frac{1}{2}arg\Gamma(\frac{1}{2}+iT)-arg\Gamma(\frac{1}{4}+\frac{iT}{2})$
$\displaystyle-\frac{Tln2}{2}-\frac{1}{4}arctan(sinh\pi T)$ (18)
Where
$\displaystyle F_{X}(T)=-lim_{X\rightarrow\infty}[\sum_{n\leq
X}\Lambda(n)\frac{sin(Tlogn)}{\sqrt{n}logn}-\int_{1}^{X}\frac{sin(Tlogt)}{\sqrt{t}logt}dt$
$\displaystyle-\frac{sin(TlogX)}{logX}[\sum_{n\leq
X}\Lambda(n)n^{-\frac{1}{2}}-2X^{\frac{1}{2}}]]$
Using equation 18(Guinand formula) minus equation 9 and notice 14 and 16, we
get that
$\displaystyle
0=\frac{1}{2}arg\Gamma(\frac{1}{2}+iT)-arg\Gamma(\frac{1}{4}+\frac{iT}{2})-\frac{Tln2}{2}-\frac{1}{4}arctan(sinh\pi
T)$ (19)
When $T$ is not cordinates of some Riemann zeros. We set right side of 19 to
be $d(T)$,then we just need to prove that $d(T)\equiv 0$ when $T>0$. Let’s
show it as follows: By rewriting
$\frac{Tln2}{2}=arg2^{\frac{iT}{2}}$
and
$arctan(sinh\pi T)=arg(1+isinh\pi T)$
, we have that
$\displaystyle
4d(T)=arg\frac{\Gamma^{2}(\frac{1}{2}+iT)}{\Gamma^{4}(\frac{1}{4}+\frac{iT}{2})4^{iT}(1+isinh\pi
T)}$
$\displaystyle=Imlog\frac{\Gamma^{2}(\frac{1}{2}+iT)}{\Gamma^{4}(\frac{1}{4}+\frac{iT}{2})4^{iT}(1+isinh\pi
T)}$
$\displaystyle=\frac{1}{2i}[log\frac{\Gamma^{2}(\frac{1}{2}+iT)}{\Gamma^{4}(\frac{1}{4}+\frac{iT}{2})4^{iT}(1+isinh\pi
T)}-log\frac{\Gamma^{2}(\frac{1}{2}-iT)}{\Gamma^{4}(\frac{1}{4}-\frac{iT}{2})4^{-iT}(1-isinh\pi
T)}]$
Let $s=iT$, then
$sinh\pi T=-isin\pi s$
and
$\displaystyle
4d(T)=\frac{1}{2i}log\frac{\Gamma^{2}(\frac{1}{2}+s)\Gamma^{4}(\frac{1}{4}-\frac{s}{2})(1-sin\pi
s)}{\Gamma^{2}(\frac{1}{2}-s)\Gamma^{4}(\frac{1}{4}+\frac{s}{2})4^{2s}(1+sin\pi
s)}$ (20)
Let
$g(s)=\frac{\Gamma^{2}(\frac{1}{2}+s)\Gamma^{4}(\frac{1}{4}-\frac{s}{2})(1-sin\pi
s)}{\Gamma^{2}(\frac{1}{2}-s)\Gamma^{4}(\frac{1}{4}+\frac{s}{2})4^{2s}(1+sin\pi
s)}$
then $g(s)$ is a meromorphic function in the whole complex number plane and by
the formula 20, $g(s)|_{s=iT}=1$ when $T>0$ . For the convenience of
factorizing $g(s)$, let’s set $s=\frac{1}{2}-z$ and reset $f(z)=g(s)$ we have
that
$\displaystyle
f(z)=\frac{\Gamma^{2}(1-z)\Gamma^{4}(\frac{1}{2}z)sin^{2}(\frac{\pi}{2}z)}{\Gamma^{2}(z)\Gamma^{4}(\frac{1}{2}-\frac{1}{2}z)4^{1-2z}cos^{2}(\frac{\pi}{2}z)}$
(21)
We just need to prove that $f(z)\equiv 1$ for any $z\in C$ , that can be
derived by the formula
$\Gamma(z)\Gamma(1-z)=\frac{\pi}{sin(\pi z)}$
With the formulas 16,18,19, we can rewrite the formula 8 as:
$\pi
S(T)=-\sum_{n<X}\frac{\Lambda(n)sin(Tlogn)}{\sqrt{n}}+\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}+\int_{1}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy$
$+arctan(2T)-\pi+O(\frac{T^{3}}{\delta^{2}lnX})$
Furthermore, we can rewrite above formula as an integral equation, first of
all,
$\sum_{n<X}\frac{\Lambda(n)sin(Tlogn)}{\sqrt{n}}=\int_{a}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}d\psi(y)$
$=\int_{a}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy+\int_{a}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}d\Delta(y)$
$=\int_{a}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy+\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}-\frac{\Delta(a)sin(Tloga)}{\sqrt{a}(loga)}$
$-\int_{a}^{X}\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}\Delta(y)dy$
Where $1<a<2$.
We substitute above formula into 2,we get that
$S(T)=-\frac{1}{\pi}\int_{a}^{X}\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}\Delta(y)dy$
$-\frac{1}{\pi}[\int_{a}^{1}\frac{sin(Tlogy)}{\sqrt{y}logy}dy-\frac{\Delta(a)sin(Tloga)}{\sqrt{a}(loga)}+arctan(2T)-\pi]$
and
$\frac{1}{\pi}\int_{a}^{X}\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}\Delta(y)dy$
$=\frac{1}{\pi}\int_{a}^{X}\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}d\tilde{\Delta}(y)$
$=\tilde{\Delta}(X)\frac{Tcos(TlnX)-sin(TlnX)(\frac{lnX}{2}+1)}{X\sqrt{X}ln^{2}X}-\tilde{\Delta}(a)\frac{Tcos(Tlna)-sin(Tlna)(\frac{lna}{2}+1)}{a\sqrt{a}ln^{2}a}$
$-\int_{a}^{X}\tilde{\Delta}(y)d\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}$
and
$\displaystyle\int_{a}^{X}\tilde{\Delta}(y)d\frac{Tcos(Tlny)-sin(Tlny)(\frac{lny}{2}+1)}{y\sqrt{y}ln^{2}y}=\int_{a}^{X}F(T,y)\tilde{\Delta}(y)dy$
(22)
Where
$F(T,y)=-\frac{T^{2}sin(Tlny)}{y^{\frac{5}{2}}lny}-\frac{2Tcos(Tlny)}{y^{\frac{5}{2}}lny}+\frac{3sin(Tlny)}{4y^{\frac{5}{2}}lny}$
$-\frac{2Tcos(Tlny)}{y^{\frac{5}{2}}ln^{2}y}+\frac{2sin(Tlny)}{y^{\frac{5}{2}}ln^{2}y}+\frac{2sin(Tlny)}{y^{\frac{5}{2}}ln^{3}y}$
By 2,2,22 and when t is not the cordinate of a Riemann zero, let
$X\rightarrow\infty$, we have following integral equation
$\displaystyle
S(t)=-\frac{1}{\pi}\int_{a}^{\infty}F(t,y)\tilde{\Delta}(y)dy+g(a,t)$ (23)
Where
$g(a,t)=-\frac{1}{\pi}[\int_{a}^{1}\frac{sin(tlogy)}{\sqrt{y}logy}dy-\frac{\Delta(a)sin(tloga)}{\sqrt{a}(loga)}$
$+\tilde{\Delta}(a)\frac{tcos(tlna)-sin(tlna)(\frac{lna}{2}+1)}{a\sqrt{a}ln^{2}a}+arctan(2t)-\pi]$
, $1<a<2$
## 3 Representing $\tilde{\Delta}(x)$ in term of $S(T)$
In this section, under the RH , we will represent $\tilde{\Delta}(x)$ as an
integral of $S(T)$ via Riemann-Von Mangoldt formula. Let $N(T)$ be a function
counting the number of non-trivial Riemann zeros whose imaginary is between
$0$ and $T$,under the RH,we can rewrite the formula 1in term of $N(T)$ as
follows,
$\displaystyle\tilde{\Delta}(x)=-\sum_{\rho}\frac{x^{\rho+1}}{\rho(\rho+1)}-xln(2\pi)-\sum_{n=1}^{\infty}\frac{x^{-2n+1}}{2n(2n-1)}$
(24)
$\displaystyle=-\int_{0}^{\infty}[\frac{x^{\frac{3}{2}+it}}{(\frac{1}{2}+it)(\frac{3}{2}+it)}+\frac{x^{\frac{3}{2}-it}}{(\frac{1}{2}-it)(\frac{3}{2}-it)}]dN(t)+f(x)$
(25)
$\displaystyle=-2x^{\frac{3}{2}}\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dN(t)+f(x)$
(26)
Where
$f(x)=-xln(2\pi)-\sum_{n=1}^{\infty}\frac{x^{-2n+1}}{2n(2n-1)}$
Noticing 17,set
$g(t)=1-\frac{tln\pi}{2\pi}+\frac{1}{\pi}arg\Gamma(\frac{1}{4}+\frac{it}{2})$
Thus we have
$\displaystyle\tilde{\Delta}(x)=-2x^{\frac{3}{2}}\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dS(t)$
(27)
$\displaystyle-2x^{\frac{3}{2}}\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dg(t)+f(x)$
(28)
Putting the last two terms together and setting it to be $\tilde{f}(x)$, we
get that
$\displaystyle\tilde{\Delta}(x)=-2x^{\frac{3}{2}}\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dS(t)+\tilde{f}(x)$
(29)
Using integration by parts and noticing $S(t)=O(logt)$, we have
$\displaystyle\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dS(t)$
$\displaystyle=-\frac{4S(0)}{3}-\int_{0}^{\infty}S(t)d\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}$
(30)
From now on, we are going to evaluate the second term of 30 in detail for the
convenience of checking.
$\int_{0}^{\infty}S(t)d\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}$
$=\int_{0}^{\infty}S(t)\frac{[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)]^{\prime}[4t^{2}+(\frac{3}{4}-t^{2})^{2}]}{[4t^{2}+(\frac{3}{4}-t^{2})^{2}]^{2}}dt$
$-\int_{0}^{\infty}S(t)\frac{[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)][4t^{2}+(\frac{3}{4}-t^{2})^{2}]^{\prime}}{[4t^{2}+(\frac{3}{4}-t^{2})^{2}]^{2}}dt$
$=\int_{0}^{\infty}S(t)\frac{[-2tcos(tlnx)-(\frac{3}{4}lnx)sin(tlnx)+t^{2}lnxsin(tlnx)+2sin(tlnx)+2tlnxcos(tlnx)][t^{4}+\frac{5}{2}+\frac{9}{16}]}{[4t^{2}+(\frac{3}{4}-t^{2})^{2}]^{2}}dt$
$-\int_{0}^{\infty}S(t)\frac{[\frac{3}{4}cos(tlnx)-t^{2}cos(tlnx)+2tsin(tlnx)][4t^{3}+5t]}{[4t^{2}+(\frac{3}{4}-t^{2})^{2}]^{2}}dt$
$=\int_{0}^{\infty}S(t)\frac{-2t^{5}cos(tlnx)-\frac{3}{4}t^{4}sin(tlnx)+t^{6}lnxsin(tlnx)+2t^{4}sin(tlnx)+2t^{5}lnxcos(tlnx)}{(t^{4}+\frac{5}{2}t^{2}+\frac{9}{16})^{2}}dt\\\
$
$+\int_{0}^{\infty}S(t)\frac{-5t^{3}cos(tlnx)-\frac{15}{8}t^{2}sin(tlnx)+\frac{5}{2}t^{4}lnxsin(tlnx)+5t^{2}sin(tlnx)+5t^{3}lnxcos(tlnx)}{(t^{4}+\frac{5}{2}t^{2}+\frac{9}{16})^{2}}dt$
$+\int_{0}^{\infty}S(t)\frac{-\frac{9}{8}tcos(tlnx)-(\frac{27}{64}lnx)sin(lnx)+\frac{9}{16}t^{2}lnxsin(tlnx)+\frac{9}{8}sin(tlnx)+\frac{9}{8}tlnxcos(tlnx)}{(t^{4}+\frac{5}{2}t^{2}+\frac{9}{16})^{2}}dt$
$-\int_{0}^{\infty}S(t)\frac{3t^{3}-4t^{5}cos(tlnx)+8t^{4}sin(tlnx)+\frac{15}{4}tcos(tlnx)-5t^{3}cos(tlnx)+10t^{2}sin(tlnx)}{(t^{4}+\frac{5}{2}t^{2}+\frac{9}{16})^{2}}dt$
To summarize, we get
$\displaystyle\tilde{\Delta}(x)=-2x^{\frac{3}{2}}\int_{0}^{\infty}K(x,t)S(t)dt+f(x)$
(31)
Where we still let
$f(x)=-2x^{\frac{3}{2}}\int_{0}^{\infty}\frac{(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)}{4t^{2}+(\frac{3}{4}-t^{2})^{2}}dg(t)-xln(2\pi)-\sum_{n=1}^{\infty}\frac{x^{-2n+1}}{2n(2n-1)}+\frac{8S(0)}{3}x^{\frac{3}{2}}$
and where
$K(x,t)=lnx[t^{6}sin(tlnx)+2t^{5}cos(tlnx)+\frac{7}{4}t^{4}sin(tlnx)+5t^{3}cos(tlnx)-\frac{21}{16}t^{2}sin(tlnx)$
$+\frac{9}{8}tcos(tlnx)-\frac{27}{64}sin(tlnx)]+2t^{5}cos(tlnx)-6t^{4}sin(tlnx)-3t^{3}cos(tlnx)-5t^{2}sin(tlnx)$
$-\frac{39}{8}tcos(tlnx)+\frac{9}{8}sin(tlnx)$
By the equations 23,31, we can get a system of integral equations
$\displaystyle\left\\{\begin{array}[]{ccc}\tilde{\Delta}(x)&=&-2x^{\frac{3}{2}}\int_{0}^{\infty}K(x,t)S(t)dt+f(x)\\\
\\\
S(t)&=&-\frac{1}{\pi}\int_{a}^{\infty}F(t,y)\tilde{\Delta}(y)dy+g(a,t)\end{array}\right.$
(35)
We shall prove an inverse theorem of the Guinnand formula as follows.
###### Theorem 6
Let $f(t)$ be a function which is continous on the interval $[0,+\infty]$
except some discrete points,which forms a set $E$ ,and
$f(t)=lim_{X\rightarrow\infty}f_{X}(t)$,
$f_{X}(t)=-\sum_{n<X}\frac{\Lambda(n)sin(Tlogn)}{\sqrt{n}}+\frac{\Delta(X)sin(TlogX)}{\sqrt{X}(logX)}+\int_{1}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy$
and $f_{X}(t)=O(logt)$ on $[0,+\infty]\setminus E$,then the RH holds
Proof. Considering the function
$F(s)=\int_{0}^{+\infty}\frac{2(s-1)}{(s-\frac{1}{2})^{2}+t^{2}}df(t)$
and
$f_{X}(s)=\int_{0}^{+\infty}\frac{2(s-1)}{(s-\frac{1}{2})^{2}+t^{2}}df_{X}(t)$
, where $Res>\frac{1}{2}$ Using integration by parts,
$F(s)=-\frac{2f(0)}{s-\frac{1}{2}}+\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}f(t)dt$
and
$f_{X}(s)=-\frac{2f_{X}(0)}{s-\frac{1}{2}}+\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}f_{X}(t)dt$
Since $f(t)=lim_{X\rightarrow\infty}f_{X}(t)$ and $f_{X}(t)=O(logt)$ on
$[0,+\infty]\setminus E$, by the Lebesque CCL,
$lim_{X\rightarrow\infty}f_{X}(s)=F(s)$ and
$\displaystyle\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}f_{X}(t)dt=\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}[-\sum_{n<X}\frac{\Lambda(n)sin(tlogn)}{\sqrt{n}}+\frac{\Delta(X)sin(tlogX)}{\sqrt{X}(logX)}+\int_{1}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy]dt$
(36)
Using residue theorem, we have that
$\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}\frac{\Lambda(n)sin(tlogn)}{\sqrt{n}}dt=\frac{\Lambda(n)}{n^{s}}$
Similarly,
$\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}[\frac{\Delta(X)sin(tlogX)}{\sqrt{X}(logX)}]dt=\frac{\Delta(X)}{X^{s+\frac{1}{2}}logX}$
and
$\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}[\int_{1}^{X}\frac{sin(Tlogy)}{\sqrt{y}logy}dy]dt=\frac{1}{1-s}[X^{1-s}-1]$
To summarize,
$\displaystyle\int_{0}^{+\infty}\frac{4(s-1)t}{[(s-\frac{1}{2})^{2}+t^{2}]^{2}}f_{X}(t)dt=-\sum_{n<X}\frac{\Lambda(n)}{n^{s}}+\frac{\Delta(X)}{X^{s+\frac{1}{2}}logX}+\frac{1}{1-s}[X^{1-s}-1]$
(37)
Since for any $Res>\frac{1}{2}$,$lim_{X\rightarrow\infty}f_{X}(s)=F(s)$, and
$F(s)$ is analytical in the half plane $Res>\frac{1}{2}$, we notice that when
$Res>1$, we have
$\displaystyle
lim_{X\rightarrow\infty}f_{X}(s)=-\sum_{n=2}^{\infty}\frac{\Lambda(n)}{n^{s}}+\frac{1}{s-1}$
(38) $\displaystyle=-\frac{\zeta^{\prime}(s)}{\zeta(s)}+\frac{1}{s-1}=F(s)$
(39)
Which means that RH is true.
## 4 Lower bound of $\tilde{\Delta}(x)$
First of all , Let’s derive a formula based on the functional equation and
formula, since
$\displaystyle\frac{\zeta^{\prime}(s)}{\zeta(s)}=-s(s+1)\int_{1}^{X}\tilde{\Delta}(x)x^{-s-2}dx$
$\displaystyle-\frac{s(s+1)}{2(s-1)}X^{1-s}+\sum_{\rho}\frac{s(s+1)X^{\rho-s}}{\rho(\rho+1)(s-\rho)}+\sum_{n\geq
1}\frac{s(s+1)X^{-2n-s}}{2n(2n-1)(s+2n)}$ (40)
As we know, by the functional equation and 40, we have
$\displaystyle
Re\frac{\zeta^{\prime}(s)}{\zeta(s)}|_{s=\frac{1}{2}+it}=\frac{ln\pi}{2}-\frac{1}{2}Re\frac{\Gamma^{\prime}}{\Gamma}(\frac{1}{2}+it)$
(41)
Evaluating real part at$s=\frac{1}{2}+it$ on both sides of 42, we get that
$\displaystyle\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)]dx-\sum_{t_{\rho}}\frac{(\frac{3}{4}-t^{2})(\frac{3}{4}-t_{\rho}^{2})+4tt_{\rho}}{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}\frac{sin(t_{\rho}-t)lnX}{t_{\rho}-t}$
$\displaystyle-\sum_{t_{\rho}}\frac{(\frac{3}{4}-t^{2})(\frac{3}{4}-t_{\rho}^{2})-4tt_{\rho}}{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}\frac{sin(t_{\rho}+t)lnX}{t_{\rho}+t}$
$\displaystyle+\sum_{t_{\rho}}\frac{\frac{3}{2}+2tt_{\rho}}{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}cos(t_{\rho}-t)lnX$
$\displaystyle+\sum_{t_{\rho}}\frac{\frac{3}{2}-2tt_{\rho}}{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}cos(t_{\rho}+t)lnX=g_{X}(t)$
(42)
Where
$g_{X}(t)=Re[-\frac{\zeta^{\prime}(s)}{\zeta(s)}+\sum_{n\geq
1}\frac{s(s+1)X^{-2n-s}}{2n(2n-1)(s+2n)}]_{s=\frac{1}{2}+it}$
By above formula, when $t\neq t_{\rho}$,
$\displaystyle\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)]dx=O(1)$
(43)
otherwise,
$\displaystyle\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}[(\frac{3}{4}-t_{\rho}^{2})cos(t_{\rho}lnx)+2tsin(t_{\rho}lnx)]dx=lnX+O(1)$
(44)
Where $\rho=\frac{1}{2}+it_{\rho}$ are Riemann zeros, to simplify LHS of 44,
set $\theta_{\rho}=arctan\frac{\frac{3}{4}-t_{\rho}^{2}}{2t_{\rho}}$,then we
have
$\displaystyle\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}sin(t_{\rho}lnx+\theta_{\rho})dx=\frac{1}{\sqrt{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}}lnX+O(1)$
(45)
Let $u=lnx$, above formula can be reduced to
$\int_{1}^{u}\frac{\tilde{\Delta}(e^{u})}{e^{\frac{3}{2}u}}sin(t_{\rho}y+\theta_{\rho})dy=\frac{1}{\sqrt{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}}u+R_{\rho}(u)$
Set
$f_{\rho}(u)=\frac{\tilde{\Delta}(e^{u})}{e^{\frac{3}{2}u}}sin(t_{\rho}u+\theta_{\rho})$,
and $A_{\rho}=\frac{1}{\sqrt{(\frac{3}{4}-t_{\rho}^{2})^{2}+4t_{\rho}^{2}}}$
Thus
We have simplified form
$\int_{0}^{X}f_{\rho}(t)dt=A_{\rho}X+R_{\rho}(X)$
Let $g(t)=\frac{\tilde{\Delta}(e^{t})}{e^{\frac{3}{2}t}}$ and $max_{0\leq
t<+\infty}\mid g(t)\mid=C_{1}$ , $\mu(x)=m\\{t|\mid g(t)\mid\leq x\\}$,
$\overline{\mu}(x)=X-\mu(x)$, and $E_{x}=\\{t|\mid g(t)\mid\leq x\\}$
By choosing any $x<C_{1}$ , we have following estimate
$A_{\rho}X+R_{\rho}(X)=\int_{0}^{X}f_{\rho}(t)dt\leq\int_{0}^{X}\mid g(t)\mid
dt=\int_{E_{x}}\mid g(t)\mid dt+\int_{[0,X]\setminus E_{x}}\mid g(t)\mid dt$
$\leq x\mu(x)+C_{1}(X-\mu(x))$
When $X$ is big enough, we have
$\mu(x)\leq\frac{C_{1}-A_{\rho}}{C_{1}-x}X+\frac{C_{0\rho}}{C_{1}-x}$
or
$\overline{\mu}(x)\geq\frac{A_{\rho}-x}{C_{1}-x}X-\frac{C_{0\rho}}{C_{1}-x}$
Where $C_{0\rho}=max_{0<t<\infty}\mid R_{\rho}(t)\mid$ and set
$F_{x}^{X}=\\{u\mid|\frac{\tilde{\Delta}(u)}{u^{\frac{3}{2}}}|>x,0<u<X\\}$
Therefore
$m(F_{x}^{X})\geq\frac{A_{\rho}-x}{C_{1}-x}X-\frac{C_{0\rho}}{C_{1}-x}$
Where $\tilde{\Delta}(u)=\sum_{n\leq u}(n-\psi(n))\Lambda(n)-\frac{u^{2}}{2}$
Choosing $\rho=\rho_{0}=\frac{1}{2}+14.134....$ which is the first non-trivial
Riemann zero, we have following theorem
###### Theorem 7
When $X$ is big enough,
$m(F_{x}^{X})\geq\frac{A_{\rho_{0}}-x}{C_{1}-x}X-\frac{C_{0\rho_{0}}}{C_{1}-x}$
Set
$F_{X}(t)=\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)]dx$
, and $G_{X}(t)=\int_{0}^{t}F_{X}(y)dy$, by the formula 43,45,and Guinand
formula i.e $lim_{X\rightarrow\infty}G_{X}(t)=N(t)$, we can conjecture that
when $X\rightarrow\infty$, $F_{X}(t)$ will behave like a distribution more
than an ordinary function. Let’s verify it as follows:
First of all, we notice that
$\displaystyle\int_{0}^{+\infty}e^{-\epsilon
u}[(\frac{3}{4}-k^{2})cos(ku)+2ksin(ku)]\frac{\tilde{\Delta}(e^{u})}{e^{\frac{3}{2}u}}du$
$\displaystyle=Re[s(s+1)\int_{1}^{+\infty}\tilde{\Delta}(x)x^{-s-\epsilon-2}dx]_{s=\frac{1}{2}+ik}$
Using integration by parts couple of times, we get
$\displaystyle
s(s+1)\int_{1}^{X}\tilde{\Delta}(x)x^{-s-\epsilon-2}dx=s\tilde{\Delta}(1)-\Delta(1)+\int_{1}^{X}x^{-s-\epsilon}d\Delta(x)+\epsilon\int_{1}^{X}x^{-s-\epsilon}\Delta(x)dx$
$\displaystyle-s\epsilon\int_{1}^{X}x^{-s-\epsilon-2}\tilde{\Delta}(x)dx-
sx^{-s-\epsilon-1}\tilde{\Delta}(x)$
and
$\displaystyle\int_{1}^{X}x^{-s-\epsilon}d\Delta(x)=\int_{1}^{X}x^{-s-\epsilon}d\psi(x)-\int_{1}^{X}x^{-s-\epsilon}dx=\sum_{n<X}\frac{\Lambda(n)}{n^{s+\epsilon}}-\frac{X^{1-s-\epsilon}}{1-s-\epsilon}+1$
By the formula 2, when $Res\geq\frac{1}{2}$we get that
$\displaystyle\int_{1}^{+\infty}x^{-s-\epsilon}d\Delta(x)=-\frac{\zeta^{\prime}(s+\epsilon)}{\zeta(s+\epsilon)}+1$
(46)
Therefore
$\displaystyle\int_{0}^{+\infty}e^{-\epsilon
u}[(\frac{3}{4}-k^{2})cos(ku)+2ksin(ku)]\frac{\tilde{\Delta}(e^{u})}{e^{\frac{3}{2}u}}du=Re[\frac{\zeta^{\prime}(s+\epsilon)}{\zeta(s+\epsilon)}]\mid_{s=\frac{1}{2}+ik}+\varphi_{\epsilon}(k)$
(47)
Where
$\varphi_{\epsilon}(k)=Re[-s\epsilon\int_{1}^{\infty}x^{-s-\epsilon-2}\tilde{\Delta}(x)dx+\epsilon\int_{1}^{\infty}x^{-s-\epsilon-1}\Delta(x)dx]\mid_{s=\frac{1}{2}+ik}$
For the simplicity, we denote LHS of 47 by $J_{\epsilon}(k)$. Choosing any
test function $g(k)\in C_{0}^{\infty}(R^{+})$ ,where $C_{0}^{\infty}(R^{+})$
is the set of all smooth functions which have compact supports on $R^{+}$ ,
let’s compute following inner product,by 47
$\displaystyle lim_{\epsilon\rightarrow 0}\langle
J_{\epsilon}(k),g(k)\rangle=Re[\frac{\zeta^{\prime}(s+\epsilon)}{\zeta(s+\epsilon)}]\mid_{s=\frac{1}{2}+ik}+\varphi_{\epsilon}(k),g(k)\rangle$
$\displaystyle=lim_{\epsilon\rightarrow 0}\langle
Re[\frac{\zeta^{\prime}(s+\epsilon)}{\zeta(s+\epsilon)}]\mid_{s=\frac{1}{2}+ik},g(k)\rangle+lim_{\epsilon\rightarrow
0}\langle\varphi_{\epsilon}(k),g(k)\rangle$ (48)
and it’s not difficult to verify that $lim_{\epsilon\rightarrow
0}\langle\varphi_{\epsilon}(k),g(k)\rangle=0$
Using integration by parts twice, we have that
$\displaystyle\langle
Re[\frac{\zeta^{\prime}(s+\epsilon)}{\zeta(s+\epsilon)}]\mid_{s=\frac{1}{2}+ik},g(k)\rangle=\langle
h(\frac{1}{2}+\epsilon+ik),g^{\prime\prime}(k)\rangle$ (49)
Where $h(z)=\int_{1}^{z}ln\zeta(s)ds$ and $Rez>\frac{1}{2}$,the integration
path is the conventional contour from $1$ to $z$
Since $h(\frac{1}{2}+\epsilon+ik)=O(klogk)$ uniformly for any small $\epsilon$
and $lim_{\epsilon\rightarrow 0}h(\frac{1}{2}+\epsilon+ik)=S_{1}(k)$
Where $S_{1}(k)=\int_{0}^{k}S(t)dt$
Finally, we have
$\displaystyle lim_{\epsilon\rightarrow 0}\langle
J_{\epsilon}(k),g(k)\rangle=\langle S_{1}(k),g^{\prime\prime}(k)\rangle$ (50)
Before ending this section, let’s take look at the equation 42 again, we can
rewrite this equation in term of integral equation as follows:
$\displaystyle\int_{0}^{\infty}K_{X}(t,t^{\prime})dN(t^{\prime})=H_{X}(t)$
(51)
Where
$K_{X}(t)=-\frac{(\frac{3}{4}-t^{2})(\frac{3}{4}-t^{\prime
2})+4tt^{\prime}}{(\frac{3}{4}-t^{\prime 2})^{2}+4t^{\prime
2}}\frac{sin(t^{\prime}-t)lnX}{t^{\prime}-t}-\frac{(\frac{3}{4}-t^{2})(\frac{3}{4}-t^{\prime
2})-4tt^{\prime}}{(\frac{3}{4}-t^{\prime 2})^{2}+4t^{\prime
2}}\frac{sin(t^{\prime}+t)lnX}{t^{\prime}+t}$
$+\frac{\frac{3}{2}+2tt^{\prime}}{(\frac{3}{4}-t^{\prime 2})^{2}+4t^{\prime
2}}cos(t^{\prime}-t)lnX+\frac{\frac{3}{2}-2tt^{\prime}}{(\frac{3}{4}-t^{\prime
2})^{2}+4t^{\prime 2}}cos(t^{\prime}+t)lnX$
and
$H_{X}(t)=\int_{1}^{X}\frac{\tilde{\Delta}(x)}{x^{\frac{5}{2}}}[(\frac{3}{4}-t^{2})cos(tlnx)+2tsin(tlnx)]dx+\frac{ln\pi}{2}-\frac{1}{2}Re\frac{\Gamma^{\prime}}{\Gamma}(\frac{1}{2}+it)$
Actually, above equation depends on the parameter $X$,for every fixed $X$,we
get a non trivial integral equation of $N(t)$, so we obtain a family of
integral equations, noticing that the integral kernel $K_{X}(t,t^{\prime})$ is
an explicit function,it’s expected that exploring these integral equations
will help us to understand RH further, besides we can consider similar results
for $L$ function which satisfies functional equation.
## References
* [1] H. M. Edwards Riemann’s zeta function 74-74
* [2] A. P. Guinand A summation formula in the theory of prime numbers Pro. London Math. Soc. (2) 50 (1948) 107-119
* [3] A. A. Karatsuba S. M. Voronin The Riemann zeta function De Gruyter Exposition in Mathmatics. 5 43-56
|
arxiv-papers
| 2011-11-08T06:18:07 |
2024-09-04T02:49:24.107243
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jining Gao",
"submitter": "Jining Gao",
"url": "https://arxiv.org/abs/1111.1810"
}
|
1111.1860
|
# Superconducting proximity effect to the block antiferromagnetism in
KyFe2-xSe2
Hong-Min Jiang Department of Physics and Center of Theoretical and
Computational Physics, The University of Hong Kong, Hong Kong, China
Department of Physics, Hangzhou Normal University, Hangzhou, China Wei-Qiang
Chen Department of Physics, South University of Science and Technology of
China, Shenzhen, China Department of Physics and Center of Theoretical and
Computational Physics, The University of Hong Kong, Hong Kong, China Zi-Jian
Yao Department of Physics and Center of Theoretical and Computational
Physics, The University of Hong Kong, Hong Kong, China Fu-Chun Zhang
Department of Physics and Center of Theoretical and Computational Physics, The
University of Hong Kong, Hong Kong, China Department of Physics, Zhejiang
University, Hangzhou, China
###### Abstract
Recent discovery of superconducting (SC) ternary iron selenides has block
antiferromagentic (AFM) long range order. Many experiments show possible
mesoscopic phase separation of the superconductivity and antiferromagnetism,
while the neutron experiment reveals a sizable suppression of magnetic moment
due to the superconductivity indicating a possible phase coexistence. Here we
propose that the observed suppression of the magnetic moment may be explained
due to the proximity effect within a phase separation scenario. We use a two-
orbital model to study the proximity effect on a layer of block AFM state
induced by neighboring SC layers via an interlayer tunneling mechanism. We
argue that the proximity effect in ternary Fe-selenides should be large
because of the large interlayer coupling and weak electron correlation. The
result of our mean field theory is compared with the neutron experiments semi-
quantitatively. The suppression of the magnetic moment due to the SC proximity
effect is found to be more pronounced in the $d$-wave superconductivity and
may be enhanced by the frustrated structure of the block AFM state.
###### pacs:
74.20.Mn, 74.25.Ha, 74.62.En, 74.25.nj
## I introduction
The recent discovery of high-$T_{c}$ superconductivity in the ternary iron
selenides AyFe2-xSe2 (A=K; Rb; Cs;…) JGuo ; AKrzton ; MHFang has triggered a
new surge of interest in study of iron-based superconductors (Fe-SC). The
fascinating aspect of these material lies in the tunable Fe-vacancies in these
materials, which substantially modifies the normal-state metallic behavior and
enhances the transition temperature $T_{c}$ to above 30K from 9K for the
binary system FeSe at ambient pressure. JGuo ; MHFang ; YZhang1 Particular
attention has been focused on the vacancy ordered 245 system, K0.8Fe1.6Se2, as
it introduces a novel magnetic structure into the already rich magnetism of
Fe-SC. Unlike the collinear Cruz1 ; FMa1 ; XWYan1 or bi-collinear FMa2 ;
WBao2 ; SLi1 AFM order observed in the parent compounds of other Fe-SC, the
neutron diffraction experiment has clearly shown that these materials have a
block AFM (BAFM) order. WBao1 Meanwhile, the AFM order with an
unprecedentedly large magnetic moment of $3.31\mu_{B}$/Fe below the Néel
temperature is the largest one among all the known parent compounds of Fe-SC.
ZShermadini ; WBao1 Moreover, the carrier concentration is extremely low,
indicating the parent compound to be a magnetic insulator/semiconductor, yzhou
; RYu in comparison with a metallic spin-density-wave (SDW) state of the
parent compound in other Fe-SC. RHYuan ; MHFang
The relation between the novel magnetism and superconductivity in ternary Fe
selenides is currently an interesting issue under debate. The question is
whether the superconductivity and the BAFM order are phase separated or co-
exist in certain region of the phase diagram. The neutron experiment shows the
suppression of the AFM ordering below SC transition point, WBao1 suggesting
the coexistence. Some other experiments, such as two-magnon Raman-scattering,
AMZhang and muon-spin rotation and relaxation Shermadini1 are consistent
with this picture. On the other hand, the ARPES, FChen1 NMR Torchetti1 and
TEM ZWang1 experiments indicate a mesoscopic phase separation between the
superconductivity and the insulating AFM state. Most recently, Li et al.
showed the superconductivity and the BAFM orders to occur at different layers
of the Fe-selenide planes in the STM measurement. wli1
The vacancy in Fe-selenides is an interesting but complicated issue. The
vacancy in the Fe-selenide carries a negative charge since the Fe-ion has a
valence of 2+. In the equilibrium, we expect the vacancies to repel each other
at short distance for the Coulomb interaction and to attract to each other at
a long distance for the elastic strain. Such a scenario would be in favor of
the phase separation to form a vacancy rich and vacancy poor regions in the
compound. The challenge is then to explain the observed suppression of
magnetic moment of the BAFM due to the superconductivity. At the
phenomenological level, the suppression of magnetism due to superconductivity
has been reported previously, aeppli and such phenomenon may be explained by
Ginzburg-Landau theory. varma
In this paper, we propose that the proximity effect of superconductivity to
the BAFM in a mesoscopically phase separated Fe-selenides may be large to
account the suppression of the AFM moments observed in neutron experiment.
More specifically, we use a microscopic model to study the proximity effect on
a layer of the BAFM state induced by adjacent SC layer. The proximity effect
in Fe-selenides is expected to be important for the two reasons. One is the
weaker correlation effect, and the other is the larger interlayer hopping
amplitude, compared with those in cuprates. Both of them may enhance proximity
effect on the magnetism from the neighboring SC layer. Our model calculations
show the proximity effect in a mesoscopically phase separated state of Fe-
selenides may explain various seemingly conflicted experiments.
## II MODEL AND MEAN FIELD THEORY
AyFe2-xSe2 is a layered material with FeSe layers separated by alkali atoms,
similar to the 122 material in iron pnictides family. To investigate the
proximity effect to the BAFM layer, we consider a single BAFM layer next to a
SC layer as shown schematically in Fig. 1. The electronic Hamiltonian
describing the BAFM layer is given by
$\displaystyle H$ $\displaystyle=H_{0}+H_{inter},$ (1)
where $H_{0}$ describes the electron motion and spin couplings in the BAFM
layer and $H_{inter}$ describes the coupling to the neighboring SC layer. We
consider a two-orbital model to describe $H_{0}$,
$\displaystyle H_{0}=$
$\displaystyle-\sum_{ij,\alpha\beta,\sigma}t_{ij,\alpha\beta}C^{{\dagger}}_{i,\alpha\sigma}C_{j,\beta\sigma}-\mu\sum_{i,\alpha\sigma}C^{{\dagger}}_{i,\alpha\sigma}C_{i,\alpha\sigma}$
(2)
$\displaystyle+J_{1}\sum_{<ij>,\alpha\beta}\textbf{S}_{i,\alpha}\cdot\textbf{S}_{j,\beta}+J_{2}\sum_{<<ij>>,\alpha\beta}\textbf{S}_{i,\alpha}\cdot\textbf{S}_{j,\beta}$
$\displaystyle+J^{\prime}_{1}\sum_{<ij>^{\prime},\alpha\beta}\textbf{S}_{i,\alpha}\cdot\textbf{S}_{j,\beta}+J^{\prime}_{2}\sum_{<<ij>>^{\prime},\alpha\beta}\textbf{S}_{i,\alpha}\cdot\textbf{S}_{j,\beta},$
where $C_{i,\alpha\sigma}$ annihilates an electron at site $i$ with orbital
$\alpha$ ($d_{xz}$ and $d_{yz}$) and spin $\sigma$, $\mu$ is the chemical
potential. $t_{ij,\alpha\beta}$ are the hopping integrals, and $<ij>$
($<ij>^{\prime}$) and $<<ij>>$ ($<<ij>>^{\prime}$) denote the intra-block
(inter-block) nearest (NN) and next nearest neighbor (NNN) bonds, respectively
[see the upper layer in Fig. 1]. $J_{1}$ ($J_{2}$) are the exchange coupling
constants for NN (NNN) spins in the same block, and $J^{\prime}_{1}$
($J^{\prime}_{2}$) are for the two NN (NNN) spins in different blocks. The
two-orbital model is a crude approximation for electronic structure. However,
it may be a minimal model to capture some of basic physics in examining the
proximity effect. The band structure around the obtained Fermi energy from the
two-orbital model is shown in Fig. 2, which is very similar to the result
obtained in density functional theory. The main shortcoming in using the two-
orbital model is that the magnetic moment is $2\mu$B at largest, smaller than
the experimentally measured $3.31\mu$B. We consider this to be a quantitative
issue, and will not qualitatively change our results.
Figure 1: (color online) Schematic diagram of the system in our model to study
proximity effect to the block AFM state (upper layer) induced by
superconductivity at the lower layer via a pair tunneling process $H_{inter}$
in Eq. (3).
We now consider $H_{inter}$, the coupling between SC layer and BAFM layer.
Because of the semiconducting gap of the BAFM layer and the SC gap in the SC
layer, the leading order of the interlayer coupling are the pairing hopping
between the SC and BAFM layers. According to the crystal structures, wli1
such a coupling term can be expressed as mcmillan1
$\displaystyle H_{inter}=$
$\displaystyle\frac{t^{2}_{\tau}}{\omega_{c}}\sum_{ij,\alpha\beta,\alpha^{\prime}\beta^{\prime},\sigma}(C^{{\dagger}}_{i,\alpha,\sigma}C^{B}_{i,\beta,\sigma}C^{{\dagger}}_{j,\alpha^{\prime},\bar{\sigma}}C^{B}_{j,\beta^{\prime},\bar{\sigma}}$
(3) $\displaystyle+H.C.),$
where $\omega_{c}$ is a characteristic energy, and $t_{\tau}$ is the
interlayer hopping integral, the superscript $B$ represents the SC layer. With
the mean field approximation
$\Delta_{ij,\alpha\alpha^{\prime},\sigma\bar{\sigma}}=<C^{B}_{i,\alpha,\sigma}C^{B}_{j,\alpha^{\prime},\bar{\sigma}}>$,
we have
$\displaystyle
H_{inter}=\sum_{ij,\beta\beta^{\prime},\sigma}(V_{\tau,ij}C^{{\dagger}}_{i,\beta,\sigma}C^{{\dagger}}_{j,\beta^{\prime},\bar{\sigma}}+H.C.),$
(4)
where
$V_{\tau,ij}=\frac{t^{2}_{\tau}}{\omega_{c}}\sum_{\alpha\alpha^{\prime}}\Delta_{ij,\alpha\alpha^{\prime},\sigma\bar{\sigma}}$.
The $\sqrt{5}\times\sqrt{5}$ vacancy order and the BAFM order lead to an
enlarged unit cell with eight sites per unit cell. We use a mean field theory
for the Ising spins in Eq. (2) and obtain the Bogoliubov-de Gennes equations
in the enlarged unit cell
$\displaystyle\sum_{k}{{}^{\prime}}\sum_{j,\beta}\left(\begin{array}[]{lr}H_{ij,\alpha\beta,\sigma}+\tilde{H}_{ij,\alpha\beta,\sigma}&H_{c,ij,\alpha\beta}\\\
H^{\ast}_{c,ij,\alpha\beta}&-H^{\ast}_{ij,\alpha\beta,\bar{\sigma}}+\tilde{H}_{ij,\alpha\beta,\sigma}\end{array}\right)$
(7)
$\displaystyle\times\exp[i\textbf{k}\cdot(\textbf{r}_{j}-\textbf{r}_{i})]\left(\begin{array}[]{lr}u^{k}_{n,j,\beta,\sigma}\\\
v^{k}_{n,j,\beta,\bar{\sigma}}\end{array}\right)=E^{k}_{n}\left(\begin{array}[]{lr}u^{k}_{n,i,\alpha,\sigma}\\\
v^{k}_{n,i,\alpha,\bar{\sigma}}\end{array}\right),$ (12)
where, the summation of $k$ are over the reduced Brillouin zone, and
$\displaystyle H_{ij,\alpha\beta,\sigma}=$ $\displaystyle-
t_{ij,\alpha\beta}-\mu,$ $\displaystyle\tilde{H}_{ij,\alpha\beta,\sigma}=$
$\displaystyle\sum_{\tau}(J_{\tau,intra}+J_{\tau,inter})<S_{i+\tau,\beta}>\delta_{ij}$
$\displaystyle H_{c,ij,\alpha\beta}=$ $\displaystyle
V_{\tau,ij}\sum_{\alpha^{\prime}\beta^{\prime}}\Delta_{ij,\alpha^{\prime}\beta^{\prime},\sigma\bar{\sigma}}.$
(13)
Here, $J_{\tau,intra}=J_{1}$ ($J_{\tau,inter}=J^{\prime}_{1}$) if
$\tau=\pm\hat{x},\pm\hat{y}$ and $J_{\tau,intra}=J_{2}$
($J_{\tau,inter}=J^{\prime}_{2}$) if $\tau=\pm\hat{x}\pm\hat{y}$. $\hat{x}$
and $\hat{y}$ denote the unit vectors along $x$ and $y$ directions,
respectively. $<S_{i+\tau,\beta}>$ is defined as
$(n_{i+\tau,\beta,\uparrow}-n_{i+\tau,\beta,\downarrow})/2$.
$u^{k}_{n,j,\alpha,\sigma}$ ($u^{k}_{n,j,\beta,\bar{\sigma}}$),
$v^{k}_{n,j,\alpha,\sigma}$ ($v^{k}_{n,j,\beta,\bar{\sigma}}$) are the
Bogoliubov quasiparticle amplitudes on the $j$-th site with corresponding
eigenvalues $E^{k}_{n}$. The self-consistent equations of the mean fields are
$\displaystyle n_{i,\beta,\uparrow}=$
$\displaystyle\sum_{k,n}|u^{k}_{n,i,\beta,\uparrow}|^{2}f(E^{k}_{n})$
$\displaystyle n_{i,\beta,\downarrow}=$
$\displaystyle\sum_{k,n}|v^{k}_{n,i,\beta,\downarrow}|^{2}[1-f(E^{k}_{n})].$
(14)
The magnitude of the magnetic order on the $i$-th site and the induced SC
pairing correlation in the BAFM layer are defined as,
$\displaystyle M(i)=$
$\displaystyle\frac{1}{2}\sum_{\beta}(n_{i,\beta,\uparrow}-n_{i,\beta,\downarrow})$
$\displaystyle\Delta^{A}_{ij,\alpha\beta}=$
$\displaystyle\frac{1}{4}\sum_{k,n}{{}^{\prime}}(u^{k}_{n,i,\alpha,\sigma}v^{k\ast}_{n,j,\beta,\bar{\sigma}}e^{-i\textbf{k}\cdot(\textbf{r}_{j}-\textbf{r}_{i})}$
(15)
$\displaystyle+v^{k\ast}_{n,i,\alpha,\bar{\sigma}}u^{k}_{n,j,\beta,\sigma}e^{i\textbf{k}\cdot(\textbf{r}_{j}-\textbf{r}_{i})})\tanh(\frac{E^{k}_{n}}{2k_{B}T}).$
In the calculations, we choose the hopping integrals as follows: wgyin1 Along
the $y$ direction, the $d_{xz}-d_{xz}$ NN hopping integral $t_{1}=0.4$ eV and
the $d_{yz}-d_{yz}$ NN hopping integral $t_{2}=0.13$ eV; they are exchanged to
the $x$ direction; the NN interorbital hoppings are zero; the NNN intraorbital
hopping integral $t_{3}=-0.25$ eV for both $d_{xz}$ and $d_{yz}$ orbitals, and
the NNN interorbital hopping is $t_{4}=0.07$ eV. The hopping integral $t_{1}$
is taken as the energy unit. We keep
$J_{1}:J^{\prime}_{1}:J_{2}:J^{\prime}_{2}=-4:-1:1:2$. CCao1 ; YZYou1 The
doping level is given by $\delta=n-2.0$.
## III results
To begin with, we present the energy band structure at half filling with $n=2$
in Fig. 2(a), where $J_{1}=2.0$ is so chosen to get a band gap $\sim 500$meV
being in agreement with the first principle calculations. CCao1 ; XWYan2 For
the electron doping with $n=2.1$, the Fermi level crosses an energy band
around the center of the Brillouin zone [$\Gamma$ point in Fig. 2(b)], while
it intersects with an energy band around the zone corner at the hole doping
with $n=1.9$ [$M$ point in Fig. 2(c)]. Although a simple two-orbital model is
adopted here, both the electron and hole doping cases with $\delta=0.1$ are
qualitatively consistent with the first principle calculations. XWYan2 In the
presence of the ordered vacancies and BAFM order, the original two-band
structures are splitting to sixteen subbands as a result of the enlarged unit
cell with $8$ sites. At half filling, $8$ lower bands are occupied, i.e.,
$1/4$ electron per one subband, while another $8$ bands above the Fermi energy
are unoccupied, resulting in a band gap in Fig. 2. For the electron and hole
doping with $\delta=0.1$, the chemical potential crosses one subband which
produce the characteristic features of the Fermi surface and the metallic BAFM
state.
Figure 2: (color online) Electronic band structures of $H_{0}$ given by Eq.
(2). The parameters are given at the end of section II of the text, and
$J_{1}=2.0$. (a): at half filling or $n=2.0$; (b): at electron doping $n=2.1$;
and (c): at hole doping $n=1.9$. The color scale indicates the relative
spectra weight.
Motivated by the agreement of the self-consistent mean-field solutions with
the mentioned first principle calculations, we consider now the proximity
effect in BAFM layer induced by the SC in SC layer. For explicit reason, we
choose two possible singlet pairing symmetries in the SC layer, i.e., the NNN
$s_{\pm}$-wave and the NN $d$-wave symmetries with their respective gap
functions $\Delta_{s_{\pm}}=\Delta_{0}\cos(k_{x})\cos(k_{y})$ and
$\Delta_{d}=\Delta_{0}[\cos(k_{x})-\cos(k_{y})]$, where the former results in
the NNN bond and the latter the NN bond couplings in the BAFM layer. The
interlayer hopping constant $t_{\tau}$ is assumed to be site independent. Fig.
3 displays the moment of the BAFM order as a function of the effective
tunneling strength $V_{\tau,ij}$. At the half filling, both symmetries of the
SC order in the SC layer introduce the decrease of the BAFM order and
simultaneously induce the SC correlation with the same symmetries in the BAFM
layer as the tunneling strength increases. A main difference between the
$s_{\pm}$\- and the $d$-wave symmetries is the more pronounced proximity
effect in reducing the moment of the BAFM order produced by the $d$-wave
symmetry as the tunneling strength increase, as shown in Figs. 3(a) and 3(b).
In the case of electron doping with $n=2.1$, where the metallic BAFM state
results, although the proximity effect is more pronounced, the magnetic and
the induced SC correlation remain the qualitatively unchanged, due possibly to
the very low total carrier concentration.
Figure 3: (color online) Block AFM moment and the SC pairing correlation as
functions of the effective tunneling strength $V_{\tau,ij}$. Black curves are
for next nearest neighbor $s_{\pm}$-wave pairing, and red for nearest neighbor
$d$-wave pairing. Upper panel (a) and (b): $n=2.0$ and lower panel (c) and
(d): $n=2.1$.
The unique feature of the the effective tunneling in the second order is it’s
temperature dependence via the SC pairing
$\Delta_{ij,\alpha\alpha^{\prime},\sigma,\sigma^{\prime}}$, which differs from
that in one particle tunneling process. mmori1 ; jxzhu1 ; andersen1 The
temperature dependence of the SC pairing parameter is modeled by a
phenomenological form with $\Delta=\Delta_{0}\sqrt{1-T/T_{c}}$. We present the
temperature dependence of the magnetic moment in Fig. 4(a) for the typical
choice of the coupling constants $g_{\tau}=2V_{\tau,ij}\Delta_{0}=0.25$. As
temperature decrease, the magnetic order increases when temperature is above
$T_{c}$, while it decreases when temperature is below $T_{c}$, resulting in a
broad peak around $T_{c}$. We note that the temperature dependence of the AFM
moment is reminiscent of the neutron diffraction and the two-magnon
experiments [Fig. 4(b)]. WBao1 ; AMZhang There is another scenario that the
competition between the AFM and the SC orders in the microscopic coexistence
of them may also produce the decrease of the AFM moment below $T_{c}$. The
study of such possibility is currently under way and the results will be
published elsewhere. It is worthwhile to notice that the sizable proximity
effect relies on the substantial interlayer hoping constant $t_{\tau}$. Based
on the first principle calculation, the interlayer hopping $t_{\tau}$ was
estimated to have a comparable magnitude with $t_{1}$ possibly due to the high
values of electron mobility from the intercalated alkaline atoms, CCao2 and
leads to the highly three dimensional Fermi surface. CCao1 ; XWYan2
Figure 4: (color online) (a) Temperature dependence of the block AFM moment at
$n=2.0$. Black and red curves are for next nearest neighbor $s_{\pm}$-wave and
nearest neighbor $d$-wave pair couplings, respectively. (b): the re-plotted
curve of the neutron data from Ref. WBao1, .
## IV SUMMARY AND DISCUSSIONS
In summary, we have proposed that various seemingly conflict experiments on
the phase separation or coexistence of superconductivity and BAFM may be
explained within a phase separation scenario by taking into account of the
proximity effect of superconductivity to the neighboring layer of BAFM. We
have theoretically studied the proximity effect to a BAFM layer induced by
adjacent SC layers in a simplified two-orbital model for Fe-selenides. The
proximity effect in reducing the moment of the BAFM order highly depends on
the coupling constant $V_{\tau,ij}$. For realistic parameters of the
interlayer tunneling, our calculation shows that the superconductivity
proximity effect may result in substantial suppression of the magnetic moment.
This is in contrary to that in the cuprate superconductor, where the coupling
constant $V_{\tau,ij}$ is very small because of small c-axis hopping integral
due to the large anisotropy, and because of the renormalization of
$V_{\tau,ij}$ by a factor proportional to hole concentrations due to the no
double-occupation condition. maekawa In iron-based superconductor, the
anisotropy of iron-based material is much smaller than in the cuprate, which
lead to a relative larger $t_{\tau}$. And the moderate correlation effect in
iron-based superconductor leads to a moderate renormalization factors. As a
consequence, the coupling constant $V_{\tau,ij}$ in iron chalcogenide
superconductor should be moderate.
We remark that we’d be careful in drawing a concrete conclusion to compare
with the experiments. The approximation that only $d_{xz}$ and $d_{yz}$
orbitals are important in the bands close to Fermi energy is good in terms of
the band structures. CCao1 ; XWYan2 But the maximum magnetic moment in two-
orbital model is only $2\mu$B, smaller than the moment of $3.31\mu$B measured
in experiments. ZShermadini ; WBao1 The other effect is that we only
calculated the suppression of the BAFM order of the surface layer of the BAFM
domain. According to TEM experiment, ZWang1 each BAFM domain has around ten
layers. And the suppression of BAFM order of the layers in the middle of
domain may be more complicated. In brief, the suppression of the BAFM moment
is sizable because of the moderate coupling constant $V_{\tau,ij}$, and our
calculation may be viewed as a semi-quantitative result.
We also investigated proximity effect for various pairing symmetry of the SC
phase. It has shown that the SC pairing with NN $d$-wave symmetry resulted a
more pronounced proximity effect in reducing the moment of the BAFM order than
the NNN $s_{\pm}$-wave pairing. The second order process induced proximity
effect has a temperature dependent as the SC pairing, which may be relevant to
the experimental observations. More remarkable proximity effect was found in
the BAFM state by comparison with the conventional AFM state, which was the
consequence of the frustrated structure and the associated anisotropic
exchange interactions.
## V acknowledgement
We thank W. Bao, G. Aeppli, Y. Zhou, and T. M. Rice for helpful discussions.
This work is supported in part by Hong Kong’s RGC GRF HKU706809 and
HKUST3/CRF/09. HMJ is grateful to the NSFC (Grant No. 10904062), Hangzhou
Normal University (HSKQ0043, HNUEYT).
## VI appendix
In the following, we compare the above proximity effect with that in the
single band conventional AFM (CAFM) system. In order to make the comparison
more convincing, we choose the dispersion
$\varepsilon_{k}=-2t[\cos(k_{x})+\cos(k_{y})]-4t^{\prime}\cos(k_{x})\cos(k_{y})-\mu$
with $t=t_{1}$ and $t^{\prime}=t_{3}$, which gives rise to the similar energy
band width with that in the above two-orbital model and is close to the case
of the cuprates. The AFM order is introduced by the AFM exchange interaction
$J\sum_{\langle ij\rangle}\textbf{S}_{i}\cdot\textbf{S}_{j}$ between the NN
sites. At the half filling $n=1$, we find that $J=1.6$ produces the comparable
band gap and the electron polarization as in the above BAFM state. In Fig. 5,
we present the magnitude of the magnetic order and the induced pairing
correlation as a function of the effective tunneling $V_{\tau,ij}$. The upper
panel shows the results for the NNN $s_{\pm}$-wave pairing and the lower panel
the results for the NN $d$-wave pairing. In the figure, the magnitude of the
magnetic order in both cases is renormalized. The proximity effect in reducing
the AFM order is more pronounced for the BAFM state as shown in Figs. 5(a) and
5(c). As for the induced pairing correlation, the larger correlation is found
in the BAFM state for the $s_{\pm}$-wave paring and in the CAFM state for the
$d$-wave pairing, as displayed in Figs. 5(b) and 5(d), respectively.
Figure 5: (color online) Comparison of the proximity effect between the block
and conventional AFM states. Left column shows the moment of the AFM order,
and right column the SC pairing correlation as functions of the effective
tunneling strength $V_{\tau,ij}$. Upper panel: for the next nearest neighbor
$s_{\pm}$-wave pairing and lower panel for the nearest neighbor $d$-wave
pairing. Figure 6: (color online) Comparison of the spin structures and their
respective NN and NNN bonds. (a): block AFM state; (b): conventional AFM
state.
We can understand the above results by considering the different spin
configurations of the BAFM and CAFM orders, as shown in Fig. 6. In the BAFM
state, when two electrons transfer from the BAFM layer to the SC one, the
energy changes due to the bonds breaking for the NN bond coupling are $\Delta
E^{\uparrow\uparrow}_{NN}=|J_{1}|$ [A1 and A2 bonds in Fig. 6(a)] and $\Delta
E^{\uparrow\downarrow}_{NN}=7|J_{1}|/4$ [A3 bond in Fig. 6(a)] for the
electron pairs with ferromagnetic and antiferromagnetic alignments,
respectively. On the other hand, the energy changes due the bonds breaking for
the NNN bond are $\Delta E^{\uparrow\uparrow}_{NNN}=9|J_{1}|/4$ [B3 bond in
Fig. 6(a)] and $\Delta E^{\uparrow\downarrow}_{NNN}=5|J_{1}|/2$ [B1 and B2
bonds in Fig. 6(a)]. As a result, the proximity effect in reducing the moment
of the AFM order by the $d$-wave pairing is more remarkable than that by the
$s_{\pm}$ one. However, the energy changes due to the bonds breaking in the
CAFM state are $\Delta E_{NN}=7|J|$ [C bond in Fig. 6(b)] and $\Delta
E_{NNN}=8|J|$ [D bond in Fig. 6(b)] for the NN and NNN band couplings.
Therefore, the proximity effect in reducing the moment of the AFM in the CAFM
state is rather weak for both the $s_{\pm}$\- and $d$-wave pairing couplings.
As for the induced pairing correlation, the extent of the match between the
AFM order configuration and the singlet SC pairing largely determines the
magnitude of the induced pairing correlation. For example, the CAFM matches
well with the NN $d$-wave pairing, so that one can expect a large induced
pairing correlation without the severe decrease of the AFM order as displayed
in Figs. 5(c) and 5(d).
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|
arxiv-papers
| 2011-11-08T10:33:14 |
2024-09-04T02:49:24.115189
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hong-Min Jiang, Wei-Qiang Chen, Zi-Jian Yao, and Fu-Chun Zhang",
"submitter": "Hong-Min Jiang",
"url": "https://arxiv.org/abs/1111.1860"
}
|
1111.1937
|
2011 Vol. 0 No. XX, 000–000
11institutetext: School of Physics, Damghan University, Damghan, Iran;
kfaghei@du.ac.ir
Received [year] [month] [day]; accepted [year] [month] [day]
# Numerical study of self-gravitating protoplanetary discs
Kazem Faghei
###### Abstract
In this paper, the effect of self-gravity on the protoplanetary discs is
investigated. The mechanisms of angular momentum transport and energy
dissipation are assumed to be the viscosity due to turbulence in the accretion
disc. The energy equation is considered in situation that the released energy
by viscosity dissipation is balanced with cooling processes. The viscosity is
obtained by equality of dissipation and cooling functions, and is used for
angular momentum equation. The cooling rate of the flow is calculated by a
prescription, $du/dt=-u/\tau_{cool}$, that $u$ and $\tau_{cool}$ are internal
energy and cooling timescale, respectively. The ratio of local cooling to
dynamical timescales $\Omega\tau_{cool}$ is assumed as a constant and also as
a function of local temperature. The solutions for protoplanetary discs show
that in situation of $\Omega\tau_{cool}=constant$, the disc does not show any
gravitational instability in small radii for a typically mass accretion rate,
$\dot{M}=10^{-6}M_{\odot}yr^{-1}$, while by choosing $\Omega\tau_{cool}$ as a
function of temperature, the gravitational instability for this amount of mass
accretion rate or even less can occur in small radii. Also, by study of the
viscous parameter $\alpha$, we find that the strength of turbulence in the
inner part of self-gravitating protoplanetary discs is very low. These results
are qualitatively consistent with direct numerical simulations of
protoplanetary discs. Also, in the case of cooling with temperature
dependence, the effect of physical parameters on the structure of the disc is
investigated. The solutions represent that disc thickness and Toomre parameter
decrease by adding the ratio of disc mass to central object mass. While, the
disc thickness and Toomre parameter increase by adding mass accretion rate.
Furthermore, for typically input parameters such as mass accretion rate
$10^{-6}M_{\odot}yr^{-1}$, the ratio of the specific heats $\gamma=5/3$, and
the ratio of disc mass to central object mass $q=0.1$, the gravitational
instability can occur in whole radii of the discs excluding very near to the
central object.
###### keywords:
accretion, accretion discs — planetary systems: protoplanetary discs —
planetary systems: formation
## 1 Introduction
Accretion discs are important for many astrophysical phenomena, including
protoplanetary systems, different types of binary stars, binary X-ray sources,
quasars, and Active Galactic Nuclei (AGN). Historically, theory of accretion
discs had concentrated in the case of non self-gravitating and occasionally
the effect of self-gravity had studied (Paczyński 1978; Kolykhalov & Sunyaev
1979; Lin & Pringle 1987, 1990). On the other hand, in recent years, the
importance of study of disc self-gravity has increased, especially in the
protostellar discs and Active Galactic Nuclei (AGN) discs. It can be due to
increase of computational resources in simulation of self-gravitating
accretion discs and the observational evidences that have confirmed the
existence of self-gravity on all scale discs, from AGN to protostars (Lodato
2008 and references therein). Also, it appears the development of
gravitational instability is important for cool regions of accreting gas that
angular momentum transport by magneto-rotational instability (MRI) becomes
weak (Fleming et al. 2000; Masada & Sano 2008; Faghei 2011) and angular
momentum can transport by gravitational instability.
The structure of self-gravitating discs has been studied both through self-
similar solutions assuming steady and unsteady state (Mineshige & Umemura
1996, 1997; Tsuribe 1999; Bertin & Lodato 1999, 2001; Shadmehri & Khajenabi
2006; Abbassi et. al. 2006; Shadmehri 2009) and through direct numerical
simulations (Gammie 2001; Rice et al. 2003, 2005, 2010; Rice & Armitage 2009;
Cossins et al. 2010; Meru & Bate 2011a).
Mineshige & Umemura (1996) investigated the role of self-gravity on the
classical self-similar solution of advection dominated accretion flows (ADAF,
Narayan & Yi 1994) and found global one-dimensional solutions influenced by
self-gravity in both the radial and the perpendicular directions of the disc.
They extended the previous steady state solutions to the time-dependent case
while the effect of the self-gravity of the disc was taken into account. They
used an isothermal equation, and so their solutions describe a viscous
accretion discs in the slow accretion limit. Tsuribe (1999) studied unsteady
viscous accretion in self-gravitating discs. Taking into account the growth of
the central point mass, Tsuribe (1999) derived a series of self-similar
solutions for rotating isothermal discs. The solutions showed, as a core mass
increases, the rotation law changes from flat rotation to Keplerian rotation
in the inner disc and in addition to the central point mass, the inner disc
grows by mass accumulation due to the differing mass accretion rates in the
inner and outer radii. Bertin & Lodato (1999) considered a class of steady-
state self-gravitating accretion discs for which efficient cooling mechanisms
are assumed to operate so that the disc is self-regulated at a condition of
approximate marginal Jeans stability. They investigated the entire parameter
space available for such self-regulated accretion discs. In another study,
Bertin & Lodato (2001) followed the model that, when the disc is sufficiently
cold, the stirring due to Jeans-related instabilities acts as a source of
effective heating. The corresponding reformulation of the energy equations,
they demonstrated how self-regulation can be established, so that the
stability parameter $Q$ is maintained close to a threshold value, with weak
dependence on radius. Abbassi et al. (2006) studied the effect of viscosity on
the time evolution of axisymmetric, polytropic self-gravitating discs around a
new born central object. Thus, they ignored from the gravitational effect of
central object and only self-gravity of the disc played an important role.
They compared effects of $\alpha$-viscosity prescription (Shakura & Sunyaev)
and $\beta$-viscosity prescription (Duschel et al. 2000) on disc structure.
They found that accretion rate onto the central object for $\beta$-discs more
than $\alpha$-discs at least in the outer regions where $\beta$-discs are more
efficient. Also, their results showed gravitational instability can occur
everywhere on the $\beta$-discs and thus they suggested that $\beta$-discs can
be a good candidate for the origin of planetary systems. Shadmehri & Khejenabi
(2006) examined steady self-similar solutions of isothermal self-gravitating
discs in the presence of a global magnetic field. Similar to Abbassi et. al.
(2006) they neglected from the mass of the central object to the disc mass. By
study of Toomre parameter they showed that magnetic field can be important in
gravitational stability of the disc.
An accretion discs can become gravitationally unstable if Toomre parameter
becomes smaller than its critical value, $Q<Q_{crit}$ (Toomre 1964). For
axisymmetric instabilities $Q_{crit}\sim 1$, while for non-axisymmetric
instabilities $Q_{crit}$ values as high as $1.5\,-\,1.7$ (Durisen et al.
2007). One possible outcome is that unstable discs fragment to produce bound
objects and has been suggested as a possible mechanism for forming giant
planets (Boss 1998, 2002). However, recently it has been realized that above
condition is not sufficient to guarantee fragmentation. Gammie (2001) showed
that in addition to the above instability criterion, the disc must cool at a
fast enough rate. Let the cooling timescale $\tau_{cool}$ be defined as the
gas internal energy divided by the volumetric cooling rate. For power-law
equation of state and with $\tau_{cool}$ prescribed to be some value over a
annulus of the disc, the thin shearing box simulations of Gammie (2001) show
that fragmentation occurs if and only if
$\Omega\tau_{cool}\lesssim\beta_{crit}$, where $\beta_{crit}\sim 3$ and
$\Omega$ is angular velocity of the disc or inverse of dynamical timescale
$\tau_{dyn}=\Omega^{-1}$. The critical value of $\Omega\tau_{cool}$ can be
somewhat larger than three for more massive and physically thicker discs (Rice
et al. 2003), larger adiabatic index (Rice et al. 2005), and more resolution
of simulations (Meru & Bate 2011b). Cossins et al. (2010) by SPH simulation
studied the effects of opacity regimes on the stability of self-gravitating
protoplanetary discs to fragmentation into bound objects. They showed that
$\Omega\tau_{cool}$ has a strong dependence on the local temperature. As, they
found that without temperature dependence, for radii $\lesssim 10AU$ a very
large accretion rate $~{}10^{-3}M_{\odot}yr^{-1}$ is required for
fragmentation, but that this is reduced to $10^{-4}$ with cooling of dependent
on temperature.
As mentioned, typically semi-analytical studies of self-gravitating discs are
regarding polytropic discs (Abbassi et al. 2006), isothermal discs (Mineshige
& Umemura 1996, 1997; Tsuribe 1999; Shadmehri & Khajenabi 2006), ADAFs in the
extreme limit of no radiative cooling (Shadmehri 2004), and discs without
central object (Mineshige & Umemura 1996, 1997; Tsuribe 1999; Shadmehri &
Khajenabi 2006; Abbassi et al. 2006). In this paper, it will be interesting to
understand under which conditions gravitational instability can occur in
accretion discs by a suitable energy equation and assuming a Newtonian
potential of a mass point that stands in the disc centre. Thus, to obtain
these conditions, we will use a prescription for cooling rate that is
introduced by Gammie (2001), $du/dt=-u/\tau_{cool}$, that $u$ and
$\tau_{cool}$ are internal energy and cooling timescale, respectively. The
ratio of local cooling to dynamical timescales $\Omega\tau_{cool}$ is assumed
a power-law function of temperature in adapting Cossins et al. (2010),
$\Omega\tau_{cool}=\beta_{0}(T/T_{0})^{\delta}$, where $T_{0}$ and $\delta$
are free parameters, and $\beta_{0}$ is free parameter in Gammie (2001). In
$\delta=0$, $\Omega\tau_{cool}$ reduces to Gammie (2001) model that
$\Omega\tau_{cool}$ is a constant, while non-zero $\delta$ is qualitatively
consistent with results of Cossins et al. (2010). We will examine the effects
of $\delta$ parameter on gravitational stability of disc. We will show that
the present model is qualitatively consistent with direct numerical
simulations (Rice & Armitage 2009; Cossin et al. 2010; Rice et al. 2010) and
can provide conditions that gravitational instability occur in whole radii
excluding very near to the central object.
In section 2, the basic equations of constructing a model for steady self-
gravitating disc will be defined. In section 3, we will find asymptotic
solutions for outer edge of the disc. In section 4, by exploit of asymptotic
solutions as boundary conditions for system equations, we will investigate
numerically the effects of physical parameters on structure and stability of
the disc. The summary and discussion of the model will appear in section 5.
## 2 Basic Equations
We use cylindrical coordinate $(r,\varphi,z)$ centered on the accreting object
and make the following standard assumptions:
* (i)
The flow is assumed to be steady and axisymmetric
$\partial_{t}=\partial_{\varphi}=0$, so all flow variables are a function of
$r$ and $z$ ;
* (ii)
The gravitational force of central object on a fluid element is characterized
by the Newtonian potential of a point mass, $\Psi=-{GM_{*}}/{r}$, with $G$
representing the gravitational constant and $M_{*}$ standing for the mass of
the central star;
* (iii)
The equations written in cylindrical coordinates are integrated in the
vertical direction, hence all quantities of the flow variables will be
expressed in terms of cylindrical radius $r$;
The governing equations on the self-gravitating accretion disc for such
assumptions is as follows. The continuity equation is
$\frac{1}{r}\frac{d}{dr}(r\Sigma v_{r})=0,$ (1)
where $v_{r}$ is the radial infall velocity and $\Sigma$ is the surface
density, which is defined as $\Sigma=2\rho h$, and $\rho$ and $h$ are density
and the disk half-thickness, respectively. The half-thickness of the disc with
assume of hydrostatic equilibrium in vertical direction is $h=c_{s}/\Omega$,
where $c_{s}$ is sound speed, which is defined as $c_{s}^{2}=p/\rho$, $p$
being the gas pressure and $\Omega$ represents angular velocity of the flow.
The equation (1) implies that
$\displaystyle\dot{M}=-2\pi r\Sigma v_{r}=constant$
where $\dot{M}$ is the mass accretion rate and is a constant in the present
model. The simulation results of protoplanetary discs show that the disc
reaches a quasi-steady state in 20000 years or less and might imply that these
systems are rarely out of equilibrium. Also, the simulations show that the
mass of the disc redistribute itself to produce a state in which the accretion
rate, $\dot{M}$, is largely independent of $r$ (Rice & Armitage 2009, Rice et
al. 2010). Thus, we can use the mass accretion as a constant and it can not be
a limitation for the present model. The momentum equations are
$v_{r}\frac{dv_{r}}{dr}=-\frac{1}{\Sigma}\frac{d}{dr}(\Sigma
c_{s}^{2})-G\left[\frac{M_{*}+M(r)}{r^{2}}\right]+r\Omega^{2},$ (2)
$\Sigma v_{r}\frac{d}{dr}(r^{2}\Omega)=\frac{1}{r}\frac{d}{dr}\left[\nu\Sigma
r^{3}\frac{d\Omega}{\partial r}\right],$ (3)
where $\nu$ is the kinematic viscosity coefficient, and $\gamma$ is the
adiabatic index, and $M(r)$ is the mass of a disc within a radius $r$. As in
Mineshige & Umemura (1997), we adopt the monopole approximation for the radial
gravitational force due to the self-gravity of the disc, which considerably
simplifies the calculations and is not expected to introduce any significant
error as long as the surface density profile is steeper than $1/r$ (e.g. Li &
Shu 1997; Saigo & Hanawa 1998; Tsuribe 1999; Krasnopolsky & Konigl 2002;
Shadmehri 2009). Now, we can write
$\frac{dM(r)}{dr}=2\pi r\Sigma.$ (4)
The energy equation is
$\frac{\Sigma v_{r}}{\gamma-1}\frac{dc_{s}^{2}}{dr}+\frac{\Sigma
c_{s}^{2}}{r}\frac{d}{dr}\left(rv_{r}\right)=\Gamma-\Lambda,$ (5)
where $\Gamma$ is the heating rate of the gas by dissipation processes such as
turbulent viscosity and $\Lambda$ represents the energy loss through radiative
cooling processes. The forms of the dissipation and cooling functions can be
written as
$\Gamma=r^{2}\Sigma\nu|\frac{d\Omega}{dr}|^{2}$ (6)
$\Lambda=\frac{1}{\gamma(\gamma-1)}\frac{\Sigma c_{s}^{2}}{\tau_{cool}}$ (7)
where $\tau_{cool}$ is cooling timescale. As noted in the introduction, we are
interest to consider the effect of cooling function on the structure of self-
gravitating discs. Thus, similar to Rice & Armitage (2009) we will study the
effects of it in the case of the heating rate in disc is equal to cooling
rate, $\Gamma=\Lambda$.
Since fragmentation requires fast cooling, Gammie (2001) suggested the cooling
timescale can be parameterized as $\beta=\Omega\tau_{cool}$ , where $\beta$ is
a free parameter. Gammie (2001) showed fragmentation requires
$\beta~{}\lesssim~{}\beta_{crit}$, where $\beta_{crit}\approx 3$ for the
adiabatic index of $\gamma=2$. Rice et al. (2005) performed 3D simulations to
show the dependence of $\beta_{crit}$ on $\gamma$: for discs with $\gamma=5/3$
and $7/5$, $\beta_{crit}\approx 6-7$ and $\approx 12-13$, respectively.
Recently, Cossins et al. (2010) studied $\beta$ as a function of temperature.
They showed that $\beta$ has a strong dependence on the local temperature.
They found that without temperature dependence, for radii $\lesssim~{}10au$ a
very large accretion rate $~{}10^{-3}~{}M_{\odot}~{}yr^{-1}$ is required for
fragmentation, but that this is reduced to $10^{-4}~{}M_{\odot}~{}yr^{-1}$
with cooling of dependent on temperature. So, for simplicity in this paper we
will use a cooling timescale with a power-law dependence on temperature for
study of the equations (1)-(5)
$\displaystyle\tau_{cool}=\frac{\beta_{0}}{\Omega}(\frac{T}{T_{0}})^{\delta}$
$\displaystyle=\frac{\beta_{0}}{\Omega}(\frac{c_{s}}{c_{s_{0}}})^{2\delta}$
(8)
that $\delta$ and $\beta_{0}$ are free parameters, and if we select $T_{0}$ as
a temperature of the outer part of the disc, then $c_{s_{0}}$ will be sound
speed in there. From equation (8) and $\delta=0$, we expect that
$\Omega\tau_{cool}$ becomes a constant that is same with Gammie (2001) model.
While non-zero $\delta$ is qualitatively consistent with Cossins et al. (2010)
model. It is important to stress that the above description for cooling rate
is not meant to reproduce any specific cooling law, but is just a convenient
way of exploring the role of the cooling timescale in the outcome of the
gravitational instability.
Here, the kinematic coefficient of viscosity can be obtained by equating of
the heating and cooling rates
$\nu=\frac{1}{\gamma(\gamma-1)}\frac{\left|\frac{d\Omega}{dr}\right|^{-2}}{r^{2}}\frac{c_{s}^{2}}{\tau_{cool}}.$
(9)
Thus, by exploit of equation (9) we do not need to use of viscosity
descriptions, such as $\alpha$ and $\beta$ prescriptions that are introduced
by Shakura & Sunyaev (1973) and Duschel et al. (2000), respectively. Equation
(9) implies that the kinematic coefficient of viscosity in the present model
depends on physical quantities of the system, specially cooling timescale. The
kinematic coefficient of viscosity in $\alpha$-prescription is $\nu=\alpha
c_{s}h$, where $\alpha$ is a free parameter and is less than unity (Shakura &
Sunyaev 1973). By using equation (9) for $\alpha$ parameter we can write
$\displaystyle\alpha=\frac{\nu}{c_{s}h}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}$
$\displaystyle=\frac{1}{\gamma(\gamma-1)}\frac{\left|\frac{d\Omega}{dr}\right|^{-2}}{r^{2}h}\frac{c_{s}}{\tau_{cool}}.$
(10)
The above equation implies that the $\alpha$ parameter is not a constant and
varies by position and strongly depends on cooling timescale. We will study
the $\alpha$ parameter in section 4 and will show that in the present model it
increases by radii.
As mentioned in the introduction, the gravitational stability of the disc can
be investigated by Toomre parameter (Toomre 1964). The Toomre parameter for an
epicyclic motion can be written as
$Q=\frac{c_{s}k}{\pi G\Sigma}$ (11)
where
$k=\Omega\sqrt{4+2\frac{d\log\Omega}{d\log r}}$ (12)
is the epicyclic frequency which can be replaced by the angular frequency,
$\Omega$.
The equations of (1)-(5) and (9) provide a set of ordinary differential
equations that describes physical properties of the self-gravitating disc.
Since, these equations are nonlinear, we will need suitable boundary
conditions to solve it numerically. Thus, in next section we will try to
obtain asymptotic solution in outer edge of the disc and then by exploit of
this asymptotic solution as a boundary condition, we will able to integrate
system equations inward from a point very near to the outer edge of the disc.
Before next sections and the numerical study of the model, we shall express
all quantities in units with values typical protostellar disc. We will choose
astronomical unit ($au$) and the sun mass ($M_{\odot}$) as the units of length
and mass, respectively. Thus, the time unit is given by
$\sqrt{au^{3}/GM_{\odot}}$ that is equal to a year divided to $2\pi$.
## 3 Outer Limit
Here, the asymptotic behavior of the system equations as $r\rightarrow R$ is
investigated that $R$ is the outer radius of the disc. The asymptotic
solutions are given by
$\Sigma(r)\sim\frac{\Sigma_{0}}{R^{1/2}}~{}\left(1+a_{1}\frac{s}{R}+\cdot\cdot\cdot\right)$
(13)
$v_{r}(r)\sim-
c_{1}\sqrt{\frac{M_{*}+M_{disc}}{R}}~{}\left(1+a_{2}\frac{s}{R}+\cdot\cdot\cdot\right)$
(14)
$\Omega(r)\sim
c_{2}\sqrt{\frac{M_{*}+M_{disc}}{R^{3}}}~{}\left(1+a_{3}\frac{s}{R}+\cdot\cdot\cdot\right)$
(15)
$c_{s}^{2}(r)\sim
c_{3}\frac{M_{*}+M_{disc}}{R}~{}\left(1+a_{4}\frac{s}{R}+\cdot\cdot\cdot\right)$
(16)
$M(r)\sim M_{disc}-\int^{R}_{r}2\pi r^{\prime}\Sigma(r^{\prime})dr^{\prime}$
(17)
where $s=R-r$, $M_{disc}$ is the disc mass, and the coefficients of $c_{i}$,
$a_{i}$, and $\Sigma_{0}$ must be determined. Using these solutions, from the
continuity, momentum, angular momentum, energy, and viscosity equations
[(1)-(5), and (9)], we can obtain the coefficients of $c_{i}$ that have the
following forms:
$c_{1}=\frac{\dot{M}}{2\pi\Sigma_{0}\sqrt{M_{*}+M_{disc}}}$ (18)
$\displaystyle
c_{2}^{2}+\left[\frac{a_{3}\gamma(\gamma-1)\beta_{0}\dot{M}(a_{3}-2)(a_{1}+a_{4})}{2\pi\Sigma_{0}\sqrt{M_{*}+M_{disc}}(a_{1}+a_{3}+a_{4}-1)}\right]c_{2}$
$\displaystyle+\left[\frac{a_{2}\dot{M}^{2}}{4\pi^{2}\Sigma_{0}^{2}(M_{*}+M_{disc})}-1\right]=0$
(19)
$c_{3}=\left(\frac{a_{3}\gamma\beta_{0}(a_{3}-2)(\gamma-1)\dot{M}}{2\pi\Sigma_{0}(a_{1}+a_{3}+a_{4}-1)\sqrt{M_{*}+M_{disc}}}\right)c_{2}$
(20)
where
$a_{4}=(1+a_{2})(1-\gamma).$ (21)
The amount of mass accretion rate can be determined by observational evidences
of the protoplanetary discs. Also, $\Sigma_{0}$ approximately can be
determined by disc mass, $M_{disc}\sim\pi R^{2}\Sigma$. Thus, knowing the
amounts of $\Sigma_{0}$ and $\dot{M}$ from the observations, the value of
$c_{3}$ coefficient is only depended on value of $c_{2}$. On the other hand,
the value of $c_{2}$ can be obtained by equation (19). Since, we have only one
equation for coefficients of $a_{i}$ (equation 21), we will select below
values for them in duration of numerical integration of system equations to
obtain physical results
$a_{1}<-2+\frac{3}{2}\,\gamma,~{}~{}3\,a_{2}=a_{3}=\frac{3}{2},~{}~{}a_{4}=(1+a_{2})(1-\gamma).$
(22)
Figure 1: Surface density, thickness, temperature, and Toomre parameter of the
disc as a function of radius, for several values of $\delta$. The surface
density and the temperature are in $cgs$ system, and the thickness and the
distance are in $au$ unit. The solid lines represent $\delta=0$, the dashed
lines represent $\delta=0.75$, and the dotted lines represent $\delta=1.5$.
The input parameters are set to the disc mass $M_{disc}=0.1M_{\odot}$, the
star mass $M_{*}=M_{\odot}$, the mass accretion rate
$\dot{M}=10^{-6}M_{\odot}yr^{-1}$, the ratio of the specific heats is set to
be $\gamma=5/3$, and $\beta_{0}=2$.
## 4 Numerical Results
If the value of $R$ is guessed, the equations by Fehlberg-Runge-Kutta fourth-
fifth order method can be integrated inwards from a point very near to the
outer edge of the disc, using the above expansions. Examples of such solutions
for surface density, half-thickness of the disc, temperature, Toomre
parameter, and the viscous parameter of $\alpha$ as a function of radius are
presented in Figs 1-5. The delineated quantities of $T$ in Figs 1-4 is the
mid-plane temperature and can then be determined using
$\displaystyle T=\left(\frac{\mu m_{p}}{k_{B}}\right)c_{s}^{2}$
where $\mu=2$ is the mean molecular weight, $m_{p}$ is the proton mass, and
$k_{B}$ is Boltzmann’s constant.
### 4.1 The influences of physical parameters on the results
The free parameters in the present model are the importance degree of
temperature in cooling timescale, $\delta$, the mass accretion rate,
$\dot{M}$, the parameter of $\beta_{0}$, the ratio of disc mass to star mass,
$q=M_{disc}/M_{*}$.
#### 4.1.1 $\delta$ parameter
The effects of $\delta$ parameter on the physical quantities are presented in
Fig. 1. The profiles of surface density and temperature show that they
increase by adding $\delta$. But, the increase of surface density is more than
temperature. Thus, the Toomre parameter ($Q\propto
c_{s}/\Sigma\propto\sqrt{T}/\Sigma$) decreases by adding $\delta$ parameter.
The profiles of Toomre parameter represent that for small $\delta$, only outer
part of the disc gravitationally is unstable, and the gravitational
instability can extend to inner radii by adding $\delta$ parameter. For
$\delta_{crit}\sim 1.5$, the Toomre parameter in radii $\gtrsim 5\,au$ becomes
smaller than critical Toomre parameter ($Q_{cri}\sim 1$) and the disc becomes
gravitationally unstable. In the other words, the profiles of Toomre parameter
represent the gravitational instability of the flow strongly depends on
cooling timescale with temperature dependence. This result is qualitatively
consistent with direct numerical simulations (e. g. Cossins et al. 2010). The
disc thickness increases by adding $\delta$ parameter. It can be due to the
increase of the temperature ($h\propto c_{s}\propto\sqrt{T}$).
Equations 8 and 9 imply that
$\frac{\nu_{(\delta\neq
0)}}{\nu_{(\delta=0)}}=\left(\frac{c_{s}}{c_{s_{0}}}\right)^{-2\delta}.$ (23)
Since $c_{s}\geq c_{s_{0}}$ the right-hand side of above equation is equal or
less than one. On the other hand, non-zero $\delta$ constrains lower viscosity
for hotter regions of the disc. The study of gravitational instability shows
that it enhances in lower viscosity (Abbassi et al. 2006; Shadmehri &
Khajenabi 2006; Khajenabi & Shadmehri 2007). Thus, the gravitational
instability can be enhanced by adding the $\delta$ parameter for hotter
regions. But there is a limitation for the amount of $\delta$ parameter that
we discuss it in next section.
Figure 2: Surface density, thickness, temperature, and Toomre parameter of the
disc as a function of radius, for several values of $\beta_{0}$. The surface
density and the temperature are in $cgs$ system, and the thickness and the
distance are in $au$ unit. The solid lines represent $\beta_{0}=1$, the dashed
lines represent $\beta_{0}=5.0$, and the dotted lines represent
$\beta_{0}=10$. The input parameters are set to the disc mass
$M_{disc}=0.1M_{\odot}$, the star mass $M_{*}=M_{\odot}$, the mass accretion
rate $\dot{M}=10^{-6}M_{\odot}yr^{-1}$, the ratio of the specific heats is set
to be $\gamma=5/3$, and $\delta=1.0$.
#### 4.1.2 $\beta_{0}$ parameter
The influences of parameter of $\beta_{0}$ are shown in Fig. 2. As, we know
from the simulations of self-gravitating disc (Gammie 2001; Rice et al. 2003),
the reduce of this parameter provides conditions that the disc places on
gravitational instability and consequently fragmentation. The profiles of
surface density show that it does not change by adding the $\beta_{0}$
parameter and only it shows small deviations in large radii. The disc
temperature increases by adding the $\beta_{0}$ parameter. Because, the
increase of this parameter reduces the rate of cooling. In large amount of
$\beta_{0}$ ($\sim 10$), the disc is gravitationally stable, while by reduce
of its value to $5$, the gravitational instability can occur in large radii,
and for the small value of it ($\beta_{0}\sim 1$), we can expect gravitational
instability in whole of the disc excluding near to the star. These results are
qualitatively consistent with direct numerical simulations of protoplanetary
disc (Gammie 2001; Rice et al. 2003; Cossins et al. 2010). Also, the solutions
show that the disc thickness increases by adding the $\beta_{0}$ parameter.
Figure 3: Surface density, thickness, temperature, and Toomre parameter of the
disc as a function of radius, for several values of $\dot{M}$. The surface
density and the temperature are in $cgs$ system, and the thickness and the
distance are in $au$ unit. The solid lines represent
$\dot{M}=10^{-7}M_{\odot}yr^{-1}$, the dashed lines represent $\dot{M}=5\times
10^{-7}M_{\odot}yr^{-1}$, and the dotted lines represent
$\dot{M}=10^{-6}M_{\odot}yr^{-1}$. The input parameters are set to the disc
mass $M_{disc}=0.1M_{\odot}$, the star mass $M_{*}=M_{\odot}$, the ratio of
the specific heats is set to be $\gamma=5/3$, $\beta_{0}=10$ and $\delta=1.0$.
#### 4.1.3 The mass accretion rate
Rice & Armitage (2009) showed that beyond of $1\,au$ the disc reaches a quasi-
steady state in $20000$ years and mass is redistributing itself to produce a
state in which the accretion rate is largely independent of $r$. The mass
accretion rate in their simulations finally reached to
$10^{-6}-10^{-7}M_{\odot}/yr$ (see Fig 4 in their paper). We will study the
behavior of the present model in Fig 3 for several values of the mass
accretion rate ($10^{-7}$, $5\times 10^{-7}$, and $10^{-6}M_{\odot}/yr$). The
solutions imply that the disc temperature is sensitive to the amount of mass
accretion rate and increases by adding the mass accretion rate. While, the
surface density is not sensitive to the mass accretion rate and only shows
small variations in large radii. Thus, the behavior of the temperature only
specifies the behavior of the Toomre parameter ($Q\propto\sqrt{T}/\Sigma$).
The profiles of Toomre parameter represent that it increases by adding the
mass accretion rate. Also, the solutions show the disc thickness increases by
adding mass accretion rate, that is due to increase of the disc temperature.
The solutions show that for a low mass accretion rate ($\sim
10^{-7}M_{\odot}/yr$), but cooling timescale with temperature dependence
($\delta\sim 1$), the gravitational instability can occur for radii $\gtrsim
10\,au$.
Figure 4: Surface density, thickness, temperature, and Toomre parameter of the
disc as a function of radius, for several values of $q=M_{disc}/M_{*}$. The
surface density and the temperature are in $cgs$ system, and the thickness and
the distance are in $au$ unit. The solid lines represent $q=0.05$, the dashed
lines represent $q=0.1$, and the dotted lines represent $q=0.15$. The input
parameters are set to the star mass $M_{*}=M_{\odot}$, the mass accretion rate
$\dot{M}=10^{-6}M_{\odot}yr^{-1}$, the ratio of the specific heats is set to
be $\gamma=5/3$, $\beta_{0}=2$ and $\delta=1.5$.
#### 4.1.4 Mass ratio
As noted in the introduction, semi-analytical studies of self-gravitating
discs are regarding discs without central object. This simplification is
relevant to protostellar discs at the beginning of the accretion phase, during
which the mass of the central object is small and only self-gravity of the
disk plays an important role. Also, this simplification can correspond to
discs at large radii because the effects of the central mass become
unimportant in the outer regions of the disc. As, the central object attends
in the present model and its effects are not ignored. Thus, the present model
does not have limitations of previous studies of semi-analytical self-
gravitating discs and can be applied for all region of the disc. Fig. 4
represents the effects of the ratio of the disc mass to the star mass
$q=M_{*}/M_{disc}$ on the present model. The solutions show the surface
density increases and the temperature decreases. Each of the surface density
increasing and the temperature decreasing individually can reduce the Toomre
parameter. Thus, we expect that Toomre parameter decreases by adding $q$
parameter that the profiles of Toommre parameter confirm it. The disc
thickness profiles represent the disc thickness decreases by adding the disc
mass. This property is qualitatively consistent with two-dimensional study of
self-gravitating disc (e. g. Ghanbari & Abbassi 2004).
Figure 5: The viscous parameter of $\alpha$ as a function of radius ($au$).
The input parameters are set to the star mass $M_{*}=M_{\odot}$, the mass
accretion rate $\dot{M}=10^{-6}M_{\odot}yr^{-1}$, the ratio of the specific
heats is set to be $\gamma=5/3$. Left panel is for several values of Gammie’s
parameter $\beta_{0}$, the solid line represents $\beta_{0}=1$, the dashed
line represents $\beta_{0}=5$, and the dotted line represents $\beta_{0}=10$,
and $\delta=1.5$. Right panel is for several values of parameter of $\delta$,
the solid line represents $\delta=0.5$, the dashed line represents
$\delta=1.0$, and the dotted line represents $\delta=1.5$, and
$\beta_{0}=1.0$.
### 4.2 The viscous parameter of $\alpha$
In the present model, the viscous parameter of $\alpha$ depends on the
physical quantities of the disc (Equation 10), especially the local cooling
rate which depends on the local temperature. The profiles of the viscous
parameter of $\alpha$ show that it increases by radii that this property is
agree with simulation results of Rice & Armitage (2009) and Rice et al.
(2010). As, mentioned in the introduction, the minimum cooling timescale
depends on the equation of state (Rice et al. 2005) with fragmentation
occurring for $\tau_{cool}\leq 3\Omega^{-1}$ when the specific heat ratio
$\gamma=5/3$ (Gammie 2001). Rice et al. (2005) showed that fragmentation
occurs for $\alpha>0.06$ and this boundary is independent of the specific heat
ratio $\gamma$. Left panel of Fig 5 represents the viscous parameter of
$\alpha$ as a function of radius for several values of the $\beta_{0}$
parameter. The solutions show the viscous $\alpha$ strongly depends on the
$\beta_{0}$ parameter. As, the $\alpha$ parameter decreases by factor of
$\beta_{0}$. Also, the solutions for small values of $\beta_{0}$ show the
viscous $\alpha$ can reach to its critical value for fragmentation. Right
panel of Fig 5 represents the viscous parameter of $\alpha$ as a function of
radius for several values of the $\delta$ parameter. The solutions represent
the $\alpha$ parameter excluding the outer region of the disc strongly depends
on the $\delta$ parameter. In $\delta=0.5$, the value of viscous $\alpha$ in
whole of the disc is in the region for fragmentation. However, Rafikov (2005)
suggested that it is extremely difficult to see how fragmentation can occur
within $10\,au$ even for the relatively massive discs. In $\delta=1.0$ and
$\delta=1.5$, the viscous $\alpha$ in the inner disc ($r\lesssim 10$ and
$40\,au$, respectively) is well below that required for fragmentation.
The requirements for fragmentation are $Q\lesssim 1$ and $\alpha>0.06$ (Rice
et al. 2005, 2010; Rice & Armitage 2009). In the present model, apparently the
increase of the $\delta$ parameter reduces possibility of fragmentation (Right
panel of Fig 5). On the other hand, the increase of $\delta$ parameter can
place the disc in gravitational instability (Fig 1). Thus, by a suitable value
for the $\delta$ parameter, the disc can obtain two requirements for
fragmentation. The Figs 1 and 5 imply that this value for small $\beta_{0}$
can be between $0.5$ and $1.0$.
## 5 Summary and Discussion
In this paper, we have studied self-gravitating accretion discs in presence of
a Newtonian potential of a point mass. We have used a prescription for cooling
that is introduced by Gammie (2001). But, due to recent results of Cossins et
al. (2010), we have assumed that cooling timescale in unit of dynamical
timescale is a power-law function of temperature. Because of, the system
equations are non-linear and there is not self-similar solution for it. First,
we have obtained asymptotic solutions for system equations and then by them as
boundary conditions, we integrated system equations numerically.
The solutions showed that the structure of the disc strongly depends on the
present cooling function. As, by adding importance degree of temperature in
cooling timescale, gravitational instability extends from outer to inner
radii. The solutions showed that in the case of cooling with temperature
dependence, the disc thickness increases. But, this change of thickness is
important in region with smaller Toomre parameter. In the present model, the
effect of physical parameters studied such as mass accretion rate, $\beta_{0}$
parameter, and the ratio of the disc mass to central object mass. The results
showed the structure of the disc is sensitive to these parameters. For
example, the disc becomes gravitationally stable in larger mass accretion
rate. While, the gravitational instability can occur in the larger disc mass.
Also, the disc thickness increases by adding the mass accretion rate and
decreases by adding the ratio of the disc mass to the star mass. The study of
the viscous parameter $\alpha$ in the present model shows that it increases by
radii that this result is consistent with direct numerical simulations (e. g.
Rice & Armitage 2009; Rice et al. 2010). Also, the solution implies that the
viscous $\alpha$ in the outer part of the disc becomes larger than its
critical value ($\sim 0.06$) that might mean condition for fragmentation.
Here, the solutions represented that the disc thickness is very sensitive to
input parameters. Thus, study of the present in a two dimensional approach may
be interesting subject for future works. Also, it will be interesting to
obtain a suitable $\delta$ value for fragmentation by direct numerical
simulations.
## Acknowledgments
I would like to acknowledge useful discussions with Alireza Khesali.
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arxiv-papers
| 2011-11-08T15:27:44 |
2024-09-04T02:49:24.123875
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kazem Faghei",
"submitter": "Kazem Faghei",
"url": "https://arxiv.org/abs/1111.1937"
}
|
1111.1962
|
# Quantum solution to a three player Kolkata restaurant problem using
entangled qutrits
Puya Sharif† and Hoshang Heydari
Department of physics, Stockholm university 10691 Stockholm Sweden
$\dagger$ Email: ps@puyasharif.net
###### Abstract
Three player quantum Kolkata restaurant problem is modeled using three
entangled qutrits. This first use of three level quantum states in this
context is a step towards a $N$-choice generalization of the $N$-player
quantum minority game. It is shown that a better than classical payoff is
achieved by a Nash equilibrium solution where the space of available
strategies is spanned by subsets of SU(3) and the players share a tripartite
entangled initial state.
Keywords: Quantum information theory, Quantum game theory, Quantum minority
games, Qutrits, Three level systems, Multipartite entangled states.
## 1 Introduction
Quantum game theory is a fairly recent extension of game theoretical analysis
to situations formulated in the framework of quantum information theory. The
first papers appeared in 1999. Meyer showed with a model of a penny-flip game
that a player making a _quantum move_ always comes out as a winner against a
player making a _classical_ move regardless of the classical players choice
[1]. The same year Eisert et al. published a quantum protocol in which they
overcame the dilemma in Prisoners dilemma [2]. In 2003 Benjamin and Hayden
generalized Eisert’s protocol to handle multiplayer quantum games and
introduced the quantum minority game together with a solution for the four
player case which outperformed the classical randomization strategy [3]. This
result was later generalized to the $n$-players by Chen et al. in 2004 [4].
Multiplayer minority games has since then been extensively investigated by
Flitney et al. [5, 6, 7].
We will here extend quantum minority games to situations where there are not
only multiple players, but also multiple choices. A quantum version of the
Kolkata restaurant problem, which is a generalized minority game will be
presented. The players uses maximally entangled qutrits as a quantum resource
and selects their strategy by locally acting with a general SU(3) operator on
the qutrit in their possession.
### 1.1 Kolkata restaurant problem
The Kolkata restaurant problem is a minority-type game [8, 9, 10, 11, 12]. In
its most general form $N$ non-communicating agents (players), have to choose
between $n$ choices. The agents receive a gain in their utility if their
choice is not too crowded, i.e the number of agents that made the same choice
is under some threshold limit. The choices can also have different values of
utility associated with them, accounting for a preference profile over the set
of choices. The original formulation comes with a story of workers in Kolkata
that during lunch hours has to choose between a fixed number of cheap
restaurants. Each restaurant can only serve a finite number of customers, so
workers arriving to a crowded restaurant will simply miss the opportunity of
having lunch. Often is the number of agents taken to be equal to the number of
restaurants, and the maximum number of costumers per restaurant limited to
one. The problem is usually modeled as an iterative game where agents ought to
base their decision on information about the distribution of agents over
choices in the previous iterations. The Kolkata restaurant problem offers
therefore a method for modeling heard behavior and market dynamics, where
visiting a restaurant translates to buying a security, in which case an agent
wishes to be the only bidder.
### 1.2 The model
In our simplified model there are just three agents, Alice, Bob and Charlie.
They have three possible choices: security 0, security 1 and security 2. They
receive a payoff $\$$ of one unit if their choice is unique, i.e that nobody
else has made the same choice, otherwise they receive $\$=0$. The game is so
called _one shoot_ , which means that it is non-iterative, and the agents have
no information from previous rounds to base their decisions on. Under the
constraint that they cannot communicate, there is nothing left to do other
than randomizing between the choices. Given the symmetric nature of the
problem, any deterministic strategy would lead all three agents to the same
strategy, which in turn would mean that all three would leave empty handed.
There are $27$ different strategy profiles possible, i.e combinations of
choices. $12$ of which gives a payoff of $\$=1$ to each one of them.
Randomization gives therefore agent $i$ an expected payoff of
$E^{c}(\$)=\frac{4}{9}$, where the superscript denotes that the result is due
to the best _classical_ strategy (as opposed to _quantum_ strategy).
In the framework of quantum game theory [13, 14, 15, 16], Alice, Bob and
Charlie shares a quantum resource. Each has a part of a multipartite quantum
state. They play their strategy by manipulating their own part of the combined
system, before measuring their subsystems and choosing accordingly. Whereas
classically the players would be allowed randomizing over a discrete set of
choices, in the quantum version each subsystem is allowed to be transformed
with the full machinery of quantum operations. A strategy, or choice therefore
translates to choosing a unitary operator $U$. In the absence of entanglement,
quantum games of this type usually yield the same payoffs as their classical
counterparts, whereas the combination of unitary operators (or a subset
therein) and entanglement, sometimes strongly outperforms classical games and
decision theoretic models. We will here present such a case.
## 2 Qutrits and parametrization of SU(3)
A qutrit is a 3-level quantum system on 3-dimensional Hilbert space
$\mathcal{H}=\mathbb{C}^{3}$ , written in the computational basis as:
$|\psi\rangle=a_{0}|0\rangle+a_{1}|1\rangle+a_{2}|2\rangle\in\mathbb{C}^{3},$
(2.0.1)
with $a_{0},a_{1},a_{2}\in\mathbb{C}$ and
$|a_{0}|^{2}+|a_{1}|^{2}+|a_{2}|^{2}=1$. A general $N$-qutrit system
$\left|\Psi\right\rangle$ is a vector on $3^{N}$-dimensional Hilbert space,
and is written as a linear combination of $3^{N}$ orthonormal basis vectors.
$\left|\Psi\right\rangle=\sum_{x_{N},..,x_{1}=0}^{2}a_{x_{N}...x_{1}}\left|x_{N}\cdots
x_{1}\right\rangle,$ (2.0.2)
where
$\left|x_{N}\cdots
x_{1}\right\rangle=\left|x_{N}\right\rangle\otimes\left|x_{N-1}\right\rangle\otimes\cdots\otimes\left|x_{1}\right\rangle\in\mathcal{H}=\overbrace{\mathbb{C}^{3}\otimes...\otimes\mathbb{C}^{3}}^{\text{$N$-times}},$
(2.0.3)
with $x_{i}\in\\{0,1,2\\}$ and complex coefficients $a_{x_{i}}$, obeying
$\sum|a_{x_{N}...x_{1}}|^{2}=1$.
Single qutrits are transformed with unitary operators $U\in$ SU(3), i.e
operators from the special unitary group of dimension 3, acting on
$\mathcal{H}$ as $U:\mathcal{H}\rightarrow\mathcal{H}$. In a multi-qutrit
system, operations on single qutrits are said to be local. They affect the
state-space of the corresponding qutrit only. The transformation of a multi-
qutrit state vector under local operations is given by the tensor products of
the individual operators:
$\left|\Psi_{fin}\right\rangle=U_{N}\otimes U_{N-1}\otimes\cdots\otimes
U_{1}\left|\Psi_{in}\right\rangle,$ (2.0.4)
where $\left|\Psi_{in}\right\rangle$ and $\left|\Psi_{fin}\right\rangle$
denotes the initial and final state of the system respectively.
There are a number of ways you can parameterize SU(3) [18, 19]. One common
approach is through the Lie algebra of the group, the eight traceless $3\times
3$ Gell-Mann matrices. We are using a different and maybe slightly more
intuitive parametrization [17]. Let $\bar{x},\bar{y},\bar{z}$ be three
general, mutually orthogonal complex unit vectors, such that
$\bar{x}\cdot\bar{y}=0$ and $\bar{x}^{*}\times\bar{y}=\bar{z}$. We construct a
SU(3) matrix by placing $\bar{x},\bar{y}^{*}$ and $\bar{z}$ as its columns.
Now a general complex unit vector is given by:
$\bar{x}=\left(\begin{array}[]{c}\sin\theta\cos\phi e^{i\alpha_{1}}\\\
\sin\theta\sin\phi e^{i\alpha_{2}}\\\ \cos\theta
e^{i\alpha_{3}}\end{array}\right),$ (2.0.5)
and one complex unit vector orthogonal to $\bar{x}$ is given by:
$\bar{y}=\left(\begin{array}[]{c}\cos\chi\cos\theta\cos\phi
e^{i(\beta_{1}-\alpha_{1})}+\sin\chi\sin\phi e^{i(\beta_{2}-\alpha_{1})}\\\
\cos\chi\cos\theta\sin\phi e^{i(\beta_{1}-\alpha_{2})}-\sin\chi\cos\phi
e^{i(\beta_{2}-\alpha_{2})}\\\ -\cos\chi\sin\theta
e^{i(\beta_{1}-\alpha_{3})}\end{array}\right),$ (2.0.6)
where $0\leq\phi,\theta,\chi,\leq\pi/2$ and
$0\leq\alpha_{1},\alpha_{2},\alpha_{3},\beta_{1},\beta_{2}\leq 2\pi$. We have
a general SU(3) matrix $U$, given by:
$U=\left(\begin{array}[]{ccc}x_{1}&y_{1}^{*}&x_{2}^{*}y_{3}-y_{3}^{*}x_{2}\\\
x_{2}&y_{2}^{*}&x_{3}^{*}y_{1}-y_{1}^{*}x_{3}\\\
x_{3}&y_{3}^{*}&x_{1}^{*}y_{2}-y_{2}^{*}x_{1}\end{array}\right),$ (2.0.7)
and it is controlled by eight real parameters
${\phi,\theta,\chi,\alpha_{1},\alpha_{2},\alpha_{3},\beta_{1},\beta_{2}}$.
## 3 The scheme
The scheme under study is a development of one first introduced by Eisert et
al. [2], and later generalized by Benjamin and Hayden [3]. It starts out with
Alice, Bob and Charlie, $A,B$ and $C$ respectively, sharing a quantum
resource, an entangled tripartite 3-level quantum state. We need to allow the
quantum states to have a common origin, since creating entanglement is a
global operation, and can t be done by acting locally on the subsystems. We
assume that there exists an unbiased referee that prepares the state and
distributes the subsystems among the players. From that point on, no
communication is allowed between the players and the referee. Each qutrit is
due to be measured by the player owning it, at the end of the protocol in the
$\\{|0\rangle,|1\rangle,|2\rangle\\}$ -basis, where basis vector corresponds
to one of the three choices: security 0, security 1, and security 2. The
players plays their strategy by applying an operator from the set of allowed
strategies $S$, followed by a local measurement which determines their final
choice. The unitary operations done by $A,B,C$ are done locally, which means
that the operator is applied on the subsystem held by the player. As
mentioned, this translates to the transformation of $|\psi_{in}\rangle$ by the
tensor product of the unitary operators applied by the players.
We want to create a quantization of the classical game in which we expand the
set of available strategies to include quantum moves. While we are proposing a
quantum game which in some sense is fundamentally different from the classical
version, we require it to be an extension, not an addition to the classical
Kolkata restaurant problem. Tracing the steps of the predecessors of this
protocol, we restrict our formulation to have the classical game fully present
at all times, accessible in the form of restrictions on the set of allowed
local operations. We simply require that there exists a set of operators that
when applied locally on an entangled initial state gives the same outcomes as
in the classical non-quantum version. Lets first look at the classical game
presented with quantum formalism. Note that there is nothing quantum
mechanical happening at this point. The initial state
$|\psi_{in}\rangle=|000\rangle=|0\rangle_{C}\otimes|0\rangle_{B}\otimes|0\rangle_{A}$
corresponds to the case where the three players chooses security 0, by
default. The individual choices are made by applying operators
$s_{i},s_{j},s_{k}\in S=\\{s_{0},s_{1},s_{2}\\}$ to each subsystem. The exact
form of these operators can be left to discuss later. The only restriction at
the moment is that they obey:
$s_{0}|0\rangle=|0\rangle,\,s_{1}|0\rangle=|1\rangle,\,s_{2}|0\rangle=|2\rangle$,
resulting in fully deterministic outcomes:
$s_{i}\otimes s_{j}\otimes s_{k}|000\rangle=|i\,j\,k\rangle.$ (3.0.1)
As mentioned earlier there are 27 different such outcomes, each linked to
different combinations of the operators $\\{s_{0},s_{1},s_{2}\\}$, 12 of which
gives a player a payoff $\$=1$, and the rest $\$=0$. Clearly there are no
operators available corresponding to mixed strategies, so randomization
processes leading to classical mixed strategies are here lifted outside the
protocol and is done by the players before the choices are made. Having
finished the first step in the quantization process our task is now to keep
the classical game present throughout the coming steps while we add quantum
structure by choosing an entangled initial state and expanding the set of
strategy operators to include any $U\in$ $S$ = SU(3). Since the game is
symmetric and unbiased in regards to permutation of player positions, then
this is a property that has to be true of the initial state
$|\psi_{in}\rangle$, to assure that the payoff functions
$\$_{i}(|\psi_{in}\rangle,U_{A},U_{B},U_{C})$ of all three players $i$ are
identical up to some permutation of $U_{A},U_{B},U_{C}$. Note that when
dealing with mixed classical and quantum strategies the payoff function
becomes an expectation value $E(\$)$ of a probability distribution over the
different outcomes. We summarize the criteria for choosing an initial state:
1. 1.
The state ought to be entangled, to accommodate for correlated randomization
among the players.
2. 2.
It should be symmetric and unbiased in regards to player positions.
3. 3.
It must allow for classical game to be accessed by restrictions on the space
of available strategy operators.
The three qutrit GHZ-type-state:
$\mid\psi_{in}\rangle=\frac{1}{\sqrt{3}}\left(|000\rangle+|111\rangle+|222\rangle\right),$
(3.0.2)
not only fulfills the above criteria, it is also a _maximally_ entangled state
on $\mathcal{H}=\mathbb{C}^{3}\otimes\mathbb{C}^{3}\otimes\mathbb{C}^{3}$. It
has the additional property of initializing the game in an maximally undesired
state. i.e. one in which none of the players receives any payoff. In order to
change their situation, they will have to make an active choice. It is left to
show that we can define a set of operators corresponding to classical pure
strategies that gives raise to deterministic classical payoffs when applied to
the entangled initial state. This problem was addressed by Eisert et. al. [2],
and further developed by Benjamin et.al. [3] for cases of $n$ players and two
choices, by defining an entangling operator $J$ and its inverse $J^{\dagger}$,
acting on a $n$-_qubit_ product state $|00\cdots 0\rangle$ with Hermitian
strategy operators $\hat{s_{0}},\hat{s_{1}}$, sandwiched in between. By
showing that any combination of the classical strategies
$\hat{s}_{x_{1}}\otimes\hat{s}_{x_{2}}\otimes\cdots\otimes\hat{s}_{x_{n}},x_{i}\in\\{0,1\\}$
commutes with $J$, one guarantees that the classical game is embedded in the
quantum version. That route is not possible when formulating a game with aid
of higher dimensional quantum states like qutrits, since at least _two
different_ Lie-algebra elements of su(3) must appear in the Hamiltonian of $J$
(For the GHZ-type-state), whereby commutation is no longer a fact in the
general case. We need a set of operators that replicates the classical payoffs
when applied directly on our entangled initial state $\mid\psi_{in}\rangle$.
The cyclic group of order three, $C_{3}$, generated by the matrix:
$s=\ \left(\begin{array}[]{ccc}0&0&1\\\ 1&0&0\\\ 0&1&0\end{array}\right)\ ,$
(3.0.3)
where $s^{3}=s^{0}=I$ and $s^{2}=s^{T}$, has the properties we are after. The
set of classical strategies $S=\\{s^{0},s^{1},s^{2}\\}$ with $s^{i}\otimes
s^{j}\otimes s^{k}|000\rangle=|i\,j\,k\rangle$ acts on the GHZ-state as:
$s^{i}\otimes s^{j}\otimes
s^{k}\frac{1}{\sqrt{3}}\left(|000\rangle+|111\rangle+|222\rangle\right)=\\\
=\frac{1}{\sqrt{3}}\left(|0+i\;0+j\;0+k\rangle+|1+i\;1+j\;1+k\rangle+|2+i\;2+j\;2+k\rangle\right).$
(3.0.4)
Note that the superscripts denotes powers of the generator and that the
addition is modulo 3. In the case under study, where there is no preference
profile over the different choices, any combination of the operators in
$S=\\{s^{0},s^{1},s^{2}\\}$ leads to the same payoffs when applied to the GHZ-
state as to $|000\rangle$.
Now that an entangled initial state $\mid\psi_{in}\rangle$ is chosen, the
scheme for the quantum game proceeds as follows. We form a density matrix
$\rho_{in}$ out of the initial state $\mid\psi_{in}\rangle$ and add noise that
can be controlled by the parameter $f$ [7]. We get:
$\rho_{in}=f\mid\psi_{in}\rangle\langle\psi_{in}\mid+\frac{1-f}{27}\mathbb{I_{\mathrm{27}}},$
(3.0.5)
where $\mathbb{I_{\mathrm{27}}}$ is the $27\times 27$ identity matrix. Alice,
Bob and Charlie now applies a unitary operator $U$ that maximizes their
chances of receiving a payoff $\$=1$, and thereby the initial state
$\rho_{in}$ is transformed into the final state $\rho_{fin}$.
$\rho_{fin}=U^{\dagger}\otimes U^{\dagger}\otimes U^{\dagger}\rho_{in}U\otimes
U\otimes U.$ (3.0.6)
Note that they are all applying the same operator $U$ since in the absence of
communication, coordination of which operator to be applied by whom, would be
impossible. We define for each player $i$ a payoff-operator $P_{i}$ , which
contains the sum of orthogonal projectors associated with the states for which
player $i$ receives a payoff $\$=1$ . For Alice this would correspond to
$P_{A}=\left(\sum_{x_{3},x_{2},x_{1}=0}^{2}|x_{3}x_{2}x_{1}\rangle\langle
x_{3}x_{2}x_{1}|,\,x_{3}\neq x_{2},x_{3}\neq x_{1},x_{2}\neq x_{1}\right)+\\\
+\left(\sum_{x_{3},x_{2},x_{1}=0}^{2}|x_{3}x_{2}x_{1}\rangle\langle
x_{3}x_{2}x_{1}|,\,x_{3}=x_{2}\neq x_{1}\right).$ (3.0.7)
The expected payoff $E_{i}(\$)$ of player $i$ is calculated by taking the
trace of the product of the final state $\rho_{fin}$ and the payoff-operator
$P_{i}$:
$E_{i}(\$)=\mathrm{Tr\left(\mathit{P}_{i}\rho_{fin}\right)}.$ (3.0.8)
## 4 Optimal strategy
The problem now is to find the unitary operator
$U(\phi,\theta,\chi,\alpha_{1},\alpha_{2},\alpha_{3},\beta_{1},\beta_{2})$
that maximizes the expected payoff. Due to the symmetry of the problem,
optimization can be done with respect to the $P_{i}$ of any of the three
players. Doing so one arrives at a maximum expected payoff of
$E(\$)=\frac{6}{9}$, assuming ($f=1$), compared to the classical
$E^{c}(\$)=\frac{4}{9}$. Which is an $50\%$ increase. This occurs when the
players applies the optimal unitary operator $U^{opt}$, whose parameters are
listed in table 1.
Parameter | $\phi$ | $\theta$ | $\chi$ | $\alpha_{1}$ | $\alpha_{2}$ | $\alpha_{3}$ | $\beta_{1}$ | $\beta_{2}$
---|---|---|---|---|---|---|---|---
Value | $\frac{\pi}{4}$ | $\cos^{-1}\left(\frac{1}{\sqrt{3}}\right)$ | $\frac{\pi}{4}$ | $\frac{5\pi}{18}$ | $\frac{5\pi}{18}$ | $\frac{5\pi}{18}$ | $\frac{\pi}{3}$ | $\frac{11\pi}{6}$
Table 1: $U^{opt}$ in the given parametrization.
Because of the periodic nature of the solution, there could be more than one
unique choice for some of the parameters within the allowed ranges
$0\leq\phi,\theta,\chi\leq\pi/2$ and
$0\leq,\alpha_{1},\alpha_{2},\alpha_{3},\beta_{1},\beta_{2}\leq 2\pi$. This is
the case for $\alpha_{1},\alpha_{2},\alpha_{3}$, where maximum expected payoff
is achieved for $\left(\frac{5+12n}{18}\right)\pi$, $n\in\\{0,1,2\\}$. Noting
that the the center of SU(3), Z(3) =$\\{I,e^{\pm\frac{i2\pi}{3}}I\\}$ only
adds a global phase and leaves the density matrix invariant, one concludes
that the transformation belongs to SU(3)/Z(3). This removes the above
ambiguity, ending up with $\alpha_{1},\alpha_{2},\alpha_{3}=\frac{5\pi}{18}$.
The final state arrived at by playing $U^{opt}$ is given by:
$\mid\psi_{fin}\rangle=\frac{1}{3}\left(|000\rangle+|012\rangle+|021\rangle+|102\rangle\right.+\\\
\left.|111\rangle+|120\rangle+|201\rangle+|210\rangle+|222\rangle\right).$
(4.0.1)
This is an even distribution of all the states that leads to payoff to all
three players and the states which gives payoff to none and shows that the
$U^{opt}\otimes U^{opt}\otimes U^{opt}$-operation fails to make the state
fully depart from the space spanned by $|000\rangle,|111\rangle,|222\rangle$.
This failure accounts for the expected payoff not reaching unity.
Now by setting $\alpha_{1},\alpha_{2}=0$ and $\alpha_{3}=\alpha$, in the
parametrization, one arrives at a six parameter subset of SU(3), given by
operators $U(\phi,\theta,\chi,\alpha,\beta_{1},\beta_{2})$. The optimum is at
the same value as with the transformation belonging to its domain. There is
thereby a $V^{opt}$ in this subset, given in table 2 below, that gives each
player an expected payoff of $E(\$)=\frac{6}{9}=\frac{2}{3}$.
Parameter | $\phi$ | $\theta$ | $\chi$ | $\alpha$ | $\beta_{1}$ | $\beta_{2}$
---|---|---|---|---|---|---
Value | $\frac{\pi}{4}$ | $\cos^{-1}\left(\frac{1}{\sqrt{3}}\right)$ | $\frac{\pi}{4}$ | $\frac{\pi}{2}$ | $\frac{\pi}{3}$ | $\frac{5\pi}{6}$
Table 2: $V^{opt}$ in the reduced parametrization.
$U^{opt}$ and $V^{opt}$ differs only by a constant phase factor, so for our
purposes, what’s true of one is true of the other. We will therefore regard
the reduced parametrization when showing that the solution is a Nash
equilibrium in the next section.
If we further reduce the parametrization by letting all phase parameters
$\alpha_{1},\alpha_{2},\alpha_{3},\beta_{1},\beta_{2}=0$, we end up with an
operator $O(\phi,\theta,\chi)\in$ SO(3), i.e. the elements of the special
orthogonal group of dimension 3. These operators corresponds to rotations in
$\mathbb{R}^{3}$. In quantum games with two choices, like quantum prisoners
dilemma and in minority games, local _orthogonal_ operations merely achieves
to replicate the results of classical mixed strategies and offers no
improvement in the expected payoff, even with a maximally entangled initial
state. In this case though, there exists an $O^{opt}\in$ SO(3), given in table
3 that outperforms the classical expected payoff by a small margin. Each
player would in this case receive a payoff of $E(\$)=\frac{40}{81}$, compared
to the classical $E(\$)=\frac{4}{9}=\frac{36}{81}$. This result might open up
the possibility of a new classification of quantum games, where there could
exist a category of quantum games with classical strategies that are
fundamentally different than classical games with classical strategies.
Parameter | $\phi$ | $\theta$ | $\chi$
---|---|---|---
Value | $\frac{\pi}{6}$ | $\cos^{-1}\left(\frac{1}{3}\right)$ | $\frac{\pi}{6}$
Table 3: $O^{opt}$ in the given parametrization.
### 4.1 Nash equilibrium
To show that this solution is valid from a game-theoretical point of view, we
need to show that $V_{opt}$ is a Nash equilibrium, i.e. that none of the
players gains by unilaterally changing strategy from $V^{opt}$ to any other
strategy $U(\phi,\theta,\chi,\alpha,\beta_{1},\beta_{2})$. Without loss of
generality, we show for the expected payoff $E_{A}(\$)$ of Alice that the
following inequality holds for any $V$:
$E_{A}(\$)(V_{C}^{opt}\otimes V_{B}^{opt}\otimes V_{A}^{opt})\geq
E_{A}(\$)(V_{C}^{opt}\otimes V_{B}^{opt}\otimes V_{A}).$ (4.1.1)
We show that this is the case by letting Alice act with a general unitary
operator $V(\phi,\theta,\chi,\alpha,\beta_{1},\beta_{2})\in$ SU(3), while Bob
and Charlie acts with $V^{opt}\;(U^{opt})$. Then we take the partial
derivatives of $E_{A}(\$)$ with respect of each of the parameters while
keeping the rest at optimal value. Vanishing partial derivatives together with
a negative definite Hessian matrix at the values of $V^{opt}$ proves that
$V^{opt}$ is Alice’s dominant strategy and because of the symmetry of the
protocol, thereby a Nash equilibrium.
$\left.\frac{\partial
E_{A}(\$)}{\partial\phi}\right|_{\phi=\phi^{{}^{\prime}}}=\left.\frac{2}{9}\cos(2\phi)\right|_{\phi=\phi^{{}^{\prime}}}=0,\qquad\left.\frac{\partial^{2}E_{A}(\$)}{\partial\phi^{2}}\right|_{\phi=\phi^{{}^{\prime}}}<0,$
(4.1.2)
$\left.\frac{\partial
E_{A}(\$)}{\partial\theta}\right|_{\theta=\theta^{{}^{\prime}}}=\left.\frac{1}{27}\left(-\sqrt{3}\sin(\theta)+3\sqrt{2}\cos(2\theta)+\right.\right.\\\
\left.\left.\left(3\sin(\theta)+\sqrt{6}\right)\cos(\theta)\right)\right|_{\theta=\theta^{{}^{\prime}}}=0,\qquad\left.\frac{\partial^{2}E_{A}(\$)}{\partial\theta^{2}}\right|_{\theta=\theta^{{}^{\prime}}}<0,$
(4.1.3)
$\left.\frac{\partial
E_{A}(\$)}{\partial\chi}\right|_{\chi=\chi^{{}^{\prime}}}=\left.\frac{2}{27}(\cos(\chi)-\sin(\chi))\left(\sin(\chi)+\cos(\chi)+\sqrt{2}\right)\right|_{\chi=\chi^{{}^{\prime}}}=0,\\\
\left.\frac{\partial^{2}E_{A}(\$)}{\partial\chi^{2}}\right|_{\chi=\chi^{{}^{\prime}}}<0,$
(4.1.4)
$\left.\frac{\partial
E_{A}(\$)}{\partial\alpha}\right|_{\alpha=\alpha^{{}^{\prime}}}=\left.\frac{4\cos(\text{$\alpha$})}{27}\right|_{\alpha=\alpha^{{}^{\prime}}}=0,\qquad\left.\frac{\partial^{2}E_{A}(\$)}{\partial\alpha^{2}}\right|_{\alpha=\alpha^{{}^{\prime}}}<0,$
(4.1.5)
$\left.\frac{\partial
E_{A}(\$)}{\partial\beta_{1}}\right|_{\beta_{1}=\beta_{1}^{{}^{\prime}}}=\left.\frac{1}{54}\left(-3\sin(\text{$\beta_{1}$})+\sin(2\text{$\beta_{1}$})+\right.\right.\\\
\left.\left.+3\sqrt{3}\cos(\text{$\beta_{1}$})+\sqrt{3}\cos(2\text{$\beta_{1}$})\right)\right|_{\beta_{1}=\beta_{1}^{{}^{\prime}}}=0,\qquad\left.\frac{\partial^{2}E_{A}(\$)}{\partial\beta_{1}^{2}}\right|_{\beta_{1}=\beta_{1}^{{}^{\prime}}}<0,$
(4.1.6)
$\left.\frac{\partial
E_{A}(\$)}{\partial\beta_{2}}\right|_{\beta_{2}=\beta_{2}^{{}^{\prime}}}=\left.\frac{1}{54}\left((2\sin(\text{$\beta_{2}$})+3)(-\cos(\text{$\beta_{2}$}))-\right.\right.\\\
\left.\left.-\sqrt{3}(3\sin(\text{$\beta_{2}$})+\cos(2\text{$\beta_{2}$}))\right)\right|_{\beta_{2}=\beta_{2}^{{}^{\prime}}}=0,\qquad\left.\frac{\partial^{2}E_{A}(\$)}{\partial\beta_{2}^{2}}\right|_{\beta_{2}=\beta_{2}^{{}^{\prime}}}<0.$
(4.1.7)
By calculating the Hessian $H$ with
$\left.H_{ij}=\frac{\partial^{2}}{\partial a_{i}\partial
a_{j}}E_{A}(\$)\right|_{a_{i}=a_{i}^{opt},a_{j}=a_{j}^{opt}},$ (4.1.8)
where $a_{i},a_{j}\in\\{\phi,\theta,\chi,\alpha,\beta_{1},\beta_{2}\\}$, and
confirming that all eigenvalues are negative, we conclude that
$V^{opt}\;(U^{opt})$ is indeed a Nash equilibrium.
### 4.2 Adjusting entanglement and fidelity
We have included a simple model of noise, to show the behavior of the expected
payoff, when the initial state was adjusted towards a completely mixed state.
This was done by controlling the fidelity $f$ of the initial state, by mixing
it with an even distribution of all basis states in
$\mathcal{H}=\mathbb{C}^{3}\otimes\mathbb{C}^{3}\otimes\mathbb{C}^{3}$.
Clearly as $f\rightarrow 0$, we should expect the entanglement as a resource
in the initial state to vanish. This is of course the case and we have
$E(\$)(U^{opt},f)=(2(2+f))/9$. For $f=0$ we simply end up with the classical
result.
A way of directly adjusting the strength of entanglement in the initial state,
while keeping the state pure is to start with
$\mid\psi_{in}\rangle=\sin\vartheta\cos\varphi|000\rangle+\sin\vartheta\sin\varphi|111\rangle+\cos\vartheta|222\rangle,$
(4.2.1)
where $0\leq\vartheta\leq\pi$ and $0\leq\varphi\leq 2\pi$. We retrieve the
maximally entangled state (3.0.2) for $\varphi=\frac{\pi}{4},\frac{3\pi}{4}$
and $\vartheta=\pm\cos^{-1}(1/\sqrt{3})$. The expected payoff is given by:
$E(\$)(U^{opt},\vartheta,\varphi)=\frac{1}{9}\left(\sin(\varphi)\sin(2\vartheta)+\cos(\varphi)\left(2\sin(\varphi)\sin^{2}(\vartheta)+\sin(2\vartheta)\right)+4\right),$
(4.2.2)
which shows that any deviation from maximal entanglement reduces the expected
payoff towards the classical $E^{c}(\$)$, graphically shown in figure 1.
A point to note here is that the maximum expected payoff strongly depends on
the choice of initial state $|\psi_{in}\rangle$, and that there can exist more
or less suitable initial states depending on the task. We chose the GHZ-state
for this protocol because it is an unbiased maximally entangled state, which
lets the classical game be present and accessible trough restrictions on $S$.
Would our preferences been different and we had chosen for example the
antisymmetric Aharonov state instead:
$|\mathcal{A_{-}}\rangle=\frac{1}{\sqrt{6}}\sum_{x_{3},x_{2},x_{1}=0}^{2}\epsilon_{x_{3}x_{2}x_{1}}|x_{3}x_{2}x_{1}\rangle,$
(4.2.3)
where $\epsilon_{x_{3}x_{2}x_{1}}$ is the completely antisymmetric tensor,
then the expected payoff would have been $E(\$)=1$, just by letting the
players apply the identity operator. This state would guarantee that everybody
ends up with a unique choice every time. But that wouldn’t be of any interest
from a game theoretical perspective since outcomes would have resembled a
classical game with unrestricted communication. However, due to the the
invariance of $|\mathcal{A_{-}}\rangle$ under local unitary transformations of
the form $U\otimes U\otimes U$, superpositions of $|\mathcal{A_{-}}\rangle$
and $|000\rangle$ under some restricted set of operators resembling the set of
mixed classical strategies, could model a classical game under different
amounts of communication.
Figure 1: Expected payoff $E(\$)(U^{opt},\vartheta,\varphi)$ as a function of
$\vartheta$ and $\varphi$ at the Nash equilibrium strategy.
## 5 Conclusions
We have created the first quantum model for a three player, three restaurant
Kolkata restaurant problem. We have shown that when the players share an
initial tripartite entangled state, there exists a local unitary operation for
which the players can increase their expected payoff $E(\$)$ by 50% compared
with classical randomization. This solution is a Nash equilibrium and
therefore a natural attractor in the space of available strategies. The
achievement of this performance is highly dependent on the strength of
entanglement and the fidelity of the initial state.
Acknowledgments: We wish to thank Ole Andersson for valuable inputs and
fruitful discussions. This study was supported by the Swedish Research Council
(VR).
## References
* [1] D. Meyer, ”Quantum strategies”, Physical Review Letters 82, (1999), 1052 1055.
* [2] J. Eisert, M. Wilkens, M. Lewenstein, ”Quantum games and quantum strategies” Physical Review Letters 83,(1999) 3077 3080.
* [3] S. Benjamin, P. Hayden, ”Multiplayer quantum games”, Physical Review A 64,(2001) 030301.
* [4] Q. Chen, Y. Wang, ”N-player quantum minority game”, Physics Letters A, 327 (2004), 98, 102.
* [5] A. Flitney, L.C.L. Hollenberg, ”Multiplayer quantum minority game with decoherence”, Quant. Inform. Comput. 7 (2007) 111-126.
* [6] A. Flitney, A. Greentree, ”Coalitions in the quantum minority game: classical cheats and quantum bullies”, Physics Letters A 362 (2007) 132 137.
* [7] C. Schmid, A. P. Flitney, W. Wieczorek, N. Kiesel, H. Weinfurter, L. C. L. Hollenberg, ”Experimental implementation of a four-player quantum game”, New J. Phys. 12 (2010) 063031.
* [8] A. Chakrabarti, New Economic Windows Series, Springer, Milan, (2007), pp. 220-227.
* [9] A. S. Chakrabarti, B. K. Chakrabarti, A. Chatterjee, M. Mitra, ”The Kolkata Paise Restaurant problem and resource utilization”, Physica A 388,(2009) 2420-2426.
* [10] A. Ghosh, A. S. Chakrabarti, B. K. Chakrabarti, 2010, ”Kolkata Paise Restaurant problem in some uniform learning strategy limits”, in Econophysics & Economis of Games, Social Choices & Quantitative Techniques, New Economic Windows, Eds. B. Basu, B. K. Chakrabarti, S. R. Chakravarty, K. Gangopadhyay, Springer, Milan, pages 3-9.
* [11] A. Ghosh, A. Chatterjee, M. Mitra, B. K. Chakrabarti, New J Phys 12 (2010) 075033.
* [12] Arthur, W. B., 1994, ”Inductive reasoning and bounded rationality: El Farol problem”, Am. Eco. Assoc. Papers & Proc. 84, 406.
* [13] A. P. Flitney, ”Review of quantum games”, in Game Theory Strategies, Equilibria, and Theorems, I. N. Haugen, A. S. Nilsen, eds.,(2008) 1 40.
* [14] E. W. Piotrowski, J. Sladkowski, ”An invitation to quantum game theory”, International Journal of Theoretical Physics 42 (2003) 1089.
* [15] F. S. Khan, S. J.D. Phoenix, ”Nash equilibrium in quantum superpositions”, Proceedings of SPIE, Vol. 8057 80570K-1.
* [16] S. E. Landsburg, ”Quantum Game Theory”, Wiley Encyclopedia of Operations Research and Management Science, (2011).
* [17] M. Mathur and D. Sen, ”Coherent states for SU(3)”, J.Math.Phys. 42 (2001) 4181-4196.
* [18] A. T. Bolukbasi, T. Dereli, ”On the SU(3) parametrization of qutrits”,12th Central European Workshop on Quantum Optics, 6-9 (June 2005), Bilkent Univ., Ankara, Turkey.
* [19] D. E. Burlankov,”The SU3 space and its quotient spaces”, Theoretical and Mathematical Physics, 138(1): 78 87 (2004).
|
arxiv-papers
| 2011-11-08T16:34:40 |
2024-09-04T02:49:24.132141
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Puya Sharif and Hoshang Heydari",
"submitter": "Puya Sharif",
"url": "https://arxiv.org/abs/1111.1962"
}
|
1111.2083
|
# Energy spectrum and effective mass using a non-local 3-body interaction
Alexandros Gezerlis1 and G. F. Bertsch1,2 1Department of Physics, University
of Washington, Seattle, WA 98195–1560 USA 2Institute for Nuclear Theory,
University of Washington, Seattle, WA 98195–1560 USA
###### Abstract
We recently proposed a nonlocal form for the 3-body induced interaction that
is consistent with the Fock space representation of interaction operators but
leads to a fractional power dependence on the density. Here we examine the
implications of the nonlocality for the excitation spectrum. In the two-
component weakly interacting Fermi gas, we find that it gives an effective
mass that is comparable to the one in many-body perturbation theory. Applying
the interaction to nuclear matter, it predicts a huge enhancement to the
effective mass. Since the saturation of nuclear matter is partly due to the
induced 3-body interaction, fitted functionals should treat the effective mass
as a free parameter, unless the two- and three-body contributions are
determined from basic theory.
Zero- and finite-range nuclear energy-density functionals have a long history
and a successful track record, allowing the description of heavy nuclei
without region-specific parametrizations.Bender:2003 The most popular
functionals use interactions that depend on fractional powers of density,
which causes serious problems when one tries to extend the theory to include
correlations Duguet:2003 ; Robledo:2007 ; Duguet:2009 . Ideally, to avoid
these problems the effective theory should be based on a Fock-space
Hamiltonian operator. As a partial solution, one can consider energy
functionals of integral powers of the density; there have been a number of
attempts to construct functionals of this kind Baldo:2010 ; Erler:2010 .
With this in mind, we recently proposed a nonlocal effective three-body
interaction that achieves a fractional dependence on density using only
integral powers of the density matrix Gezerlis:2010a . This was derived using
the many-body perturbation theory of the dilute, weakly interacting Fermi gas.
By construction, the interaction gives the correct Lee-Yang contribution
Lee:1957 to the Fermi-gas energy to order $\rho^{7/3}$. The interaction was
validated for finite systems in a harmonic trap by comparing with numerically
accurate calculations performed by the Green’s Function Monte Carlo method. At
very weak coupling, the new operator led to results that are identical with
the Lee-Yang dependence, while for stronger coupling the contribution of the
new 3-body operator turned out to be more repulsive than in Lee-Yang (though
with the same power-law behavior), thus providing a more accurate description
of the microscopic simulation.
Using the new interaction, the internal energy of the dilute Fermi gas can be
expressed in terms of the one-body density matrix as:
$\begin{split}E&=\frac{\hbar^{2}k_{F}^{2}}{m}\int
d^{3}r_{1}\,\left(\frac{\nabla_{r_{1}}\cdot\nabla_{r_{2}}}{2}\rho({\bf
r_{1},r_{2}})|_{{\bf r_{1}=r_{2}}}+4\pi a\rho_{\downarrow}({\bf r_{1}},{\bf
r_{1}})\rho_{\uparrow}({\bf r_{1}},{\bf r_{1}})\right)\\\ &+C\int
d^{3}r_{1}d^{3}r_{2}\,\frac{\rho_{\uparrow}({\bf
r_{1},r_{2}})\rho_{\downarrow}({\bf r_{1},r_{2}})\rho({\bf
r_{1},r_{2}})}{|{\mathbf{r}_{1}}-{\mathbf{r}_{2}}|}.\end{split}$ (1)
The subscript $i$ on $\rho_{i}$ denotes the spin state, $\rho$ without a
subscript is the total density. Also, if $a$ is the scattering length
associated with the two-body interaction, $C$ is a constant proportional to
$a^{2}$. The value of $C$ was derived in Ref. Gezerlis:2010a by demanding
that the formula reproduce the Lee-Yang energy in the uniform Fermi gas. The
energy $E$ or energy density $\cal E$ is given by
$\frac{E}{A}=\frac{\cal
E}{\cal\rho}=\frac{\hbar^{2}k_{F}^{2}}{2m}\left(\frac{3}{5}+\frac{2}{3\pi}ak_{F}+\frac{4}{35\pi^{2}}\left(11-2\ln
2\right)\left(ak_{F}\right)^{2}\right)~{}.$ (2)
It is convenient for later use to rederive from Eq. (1) the formula for $C$,
which was originally derived from the perturbation theory in a momentum space
representation. We insert in Eq. (1) the free Fermi gas density matrix and
drop one of the integrals to get the energy density. The Fermi gas density
matrix only depends on the relative coordinate ${\bf r}={\bf r_{1}}-{\bf
r_{2}}$ and can be written
$\rho_{0i}({\bf r})=\int_{0}^{k_{F}}\frac{d^{3}k}{(2\pi)^{3}}e^{i{\bf
k}\cdot{\bf r}}=\rho_{0i}(0)F(k_{F}r)$ (3)
where $F(x)=3j_{1}(x)/x$. The integral to be evaluated may be expressed
$\frac{E_{3}}{A}=C\frac{24\pi^{3}}{k_{F}^{5}}\rho_{0i}^{3}(0)\int_{0}^{\infty}xdxF^{3}(x)$
(4)
The integration can be performed analytically; the final result for the
strength parameter $C$ is
$C={\hbar^{2}a^{2}\over m}{64\pi(11-\ln 2)\over 3(92-27\ln 3)}$ (5)
We now calculate the single-particle energy with functional Eq. (1) and the
value of $C$ fixed by Eq. (5). The density matrix with a particle added to the
Fermi sea is
$\rho_{i}({\bf r)}=\rho_{0i}({\bf r})+\rho_{ki}e^{i{\bf k}\cdot{\bf r}}$ (6)
The second term represents a particle of momentum $k$ in spin state $i$; the
coefficient $\rho_{ki}$ has dimensions of density. With this definition the
single-particle energy may be computed as
$\varepsilon_{i}(k)=\frac{d{\cal E}}{d\rho_{ki}}\Big{|}_{\rho_{ki}=0}.$ (7)
Carrying out the differentiation on the energy expression Eq. (1), the first
term gives the usual kinetic energy and the second term is independent of $k$.
The third term is rather complicated. Assuming equal populations of spin up
and spin down in $\rho_{0}$, the derivative is given by the integral:
$\varepsilon_{3}(k)=3C\int d^{3}r\,e^{i{\bf k}\cdot({\bf
r})}\frac{\rho_{0i}^{2}({\bf r})}{r}=3C{4\pi\over
k_{F}^{2}}\rho_{0i}^{2}(0)\int_{0}^{\infty}dx~{}x\,\,j_{0}\Big{(}{k\over
k_{F}}x\Big{)}F^{2}(x).$ (8)
The factor of 3 is a direct consequence of the spin structure of the numerator
in the third term of Eq. (1). The integral can also be expressed analytically:
$\begin{split}\int_{0}^{\infty}dx~{}x\,j_{0}(yx)F^{2}(x)&=\frac{3}{160}\Big{[}4(22+y^{2})+\frac{(y-2)^{3}}{y}(y^{2}+6y+4)\log(2-y)\\\
&-2y^{2}(y^{2}-20)\log
y+\frac{(y+2)^{3}}{y}(y^{2}-6y+4)\log(2+y)\Big{]},\end{split}$ (9)
where $y=k/k_{F}$. The 3-body contribution to the single-particle energy
$\varepsilon_{3}$ is plotted in Fig. 1, with the dimensionful factors divided
out.
Figure 1: Single-particle energy in the dilute Fermi gas, normalized to
$\hbar^{2}a^{2}\rho_{0i}^{2}(0)/mk_{F}^{2}$. The solid line shows result using
the effective 3-body interaction, Eq. (8). The dashed line shows the
contribution to the quasiparticle energy obtained by Galitskii. The dotted
line is the slope of the Galitskii expression.
Galitskii’s expression for the real part of the quasiparticle energy
(Galitskii:1958, , Eq. (34)) is plotted with the same normalization on the
graph. The perfect agreement of the two at the Fermi momentum is not
accidental: the single-particle energy at the Fermi surface is identical to
the chemical potential $\mu$, which can be extracted from the interaction
energy by the formula $\mu=\partial{\cal E}/\partial\rho$. Since we fit the
total 3-body interaction energy to the dilute Fermi gas, the chemical
potential must agree as well.
The momentum-dependence of the single-particle energy gives rise to an
effective mass $m^{*}$ for the quasiparticle spectrum,
$\frac{m^{*}}{m}=\left(1+{m\over\hbar^{2}k_{F}}{\partial\epsilon_{3}(k)\over\partial
k}\Big{|}_{k=k_{F}}\right)^{-1}$ (10)
The derivative in this expression is negative, implying that the effective
mass will be larger than $m$. Fig. 1 also shows the derivative for Galitskii’s
quasiparticle energy, as the straight line (see also Ref. Fetter:1971 ). We
note that the slope for the 3-body single-particle energy is smaller, implying
less of an effective-mass enhancement. Even so, the two results are close
enough in magnitude to motivate the application of the new operator to a
nuclear energy functional.
As stated in the introduction, our main interest is to find an improved
effective Hamiltonian for nuclear structure theory. There is no reliable low-
density expansion in the nuclear many-body problem, and in fact one must
impose some length scale in the interactions to avoid collapse. Nevertheless,
in some formulations there will be a contribution to saturation coming from
the Pauli effects that we are concerned with here. To assess the importance of
the nonlocality, we take $C$ as an adjustable parameter to be fitted in the
functional, similar to the parameter $t_{3}$ of the Skyrme interaction. The
counting of the contributing graphs is different in the four-component Fermi
system than in the two-component case treated by Galitskii, but the scaling
between the total energy and the single-particle energy remains the same under
plausible assumptions about the spin-isospin character of the interaction.
Thus we may use the same formulas, only remembering that in the nuclear
context $\rho_{0i}$ is the density associated with a specific spin-isospin
projection, e.g. neutrons with spin up.
Table 1: Contributions to the energy of 208Pb in density functional theory. The numbers for the Skyrme Ska and Gogny D1S functionals were obtained with the ev8 code Bonche:2005 and the HFBaxial code robledo , respectively. | Ska | D1S
---|---|---
Kinetic | 3863 | 3920
Coulomb direct/exchange | 831/-31 | 832/-31
Spin-orbit | -97 | -105
Central 2B | -12480 | -12783
$t_{3}$ | 6274 | 6530
Total | -1640 | -1637
While we cannot calculate $C$, we can at least put a bound on its value using
the magnitude of the 3-body interaction energy that is obtained from
phenomenological energy functionals. With our form for the interaction, the
relation between the 3-body energy and the effective mass is
$\frac{m^{*}}{m}=\left(1+d\frac{E_{3}/A}{\hbar^{2}k_{F}^{2}/(2m)}\right)^{-1}$
(11)
where $d\approx-1.32$, and the two-body contribution has been omitted.
To see what the scale of the effect would be, we show in Table 1 the various
contributions to the energy of 208Pb found using the Ska Skyrme functional and
the D1S Gogny functional. Both these functionals have the same $\rho^{1/3}$
density-dependent interaction as in the Lee-Yang expansion. One sees that the
decomposition into the two-body and three-body contributions is quite similar,
although the two-body interactions have a very different construction. Eq.
(11) gives a negative effective mass for both functionals, which is of course
unphysical. The two-body nonlocality gives a contribution of the opposite
sign, but not enough to produce an effective mass in the physical range
($m^{*}/m\sim 1$). As mentioned earlier, there must be other 3-body
contributions containing intrinsic length scales in order to achieve nuclear
saturation. However, unless the nonlocalities can be calculated in detail, it
does not seem feasible to derive a theoretical effective mass to be used with
an effective Hamiltonian. The extreme sensitivity to the induced 3-body
interaction suggests that the effective mass may need to be an unconstrained
free parameter when constructing an effective Hamiltonian for mean-field
theory and its extensions.
In summary, we have applied our newly proposed non-local effective 3-body
operator to the study of the single-particle excitation spectrum, both at weak
coupling and at strong coupling. At weak coupling we see that the new operator
has similar behavior to that found by Galitskii. We also applied the new
operator to the nuclear case. The effects pointed to are very large, implying
that the effective mass cannot be simply taken to be reduced from the bare
mass based on mean-field theory: as long as no dependable ab initio results
are available, the effective mass should also be treated like an undetermined
parameter.
We would like to thank L. Robledo for providing us with the energies of 208Pb
for the Gogny D1S interaction. This work was supported by DOE Grant Nos. DE-
FG02-97ER41014 and DE-FG02-00ER41132.
## References
* (1) M. Bender, P.-H. Heenen, and P.-G. Reinhard, Rev. Mod. Phys. 75, 121 (2003).
* (2) T. Duguet and P. Bonche, Phys. Rev. C 67, 054308 (2003).
* (3) L. M. Robledo, Int. J. Mod. Phys. E 16, 337 (2007).
* (4) T. Duguet, M. Bender, K. Bennaceur, D. Lacroix, and T. Lesinski, Phys. Rev. C 79, 044320 (2009).
* (5) M. Baldo, L. M. Robledo, P. Schuck, X. Viñas, J. Phys. G 37, 064015 (2010).
* (6) J. Erler, P. Klüpfel, P.-G. Reinhard, Phys. Rev. C 82, 044307 (2010).
* (7) A. Gezerlis and G. F. Bertsch, Phys. Rev. Lett. 105, 212501 (2010).
* (8) T. D. Lee and C. N. Yang, Phys. Rev. 105, 1119 (1957).
* (9) V. M. Galitskii, Sov. Phys. (JETP) 34, 151 (1958).
* (10) A. L. Fetter and J. D. Walecka, Quantum Theory of Many-Particle Systems (McGraw-Hill, New York, 1971).
* (11) P. Bonche, H. Flocard, and P.H. Heenen, Comp. Phys. Comm. 171, 49 (2005).
* (12) ??
|
arxiv-papers
| 2011-11-09T00:20:35 |
2024-09-04T02:49:24.141614
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Alexandros Gezerlis, G. F. Bertsch",
"submitter": "George F. Bertsch",
"url": "https://arxiv.org/abs/1111.2083"
}
|
1111.2157
|
\+ daix@aphy.iphy.ac.cn
# Implementation of LDA+DMFT with pseudo-potential-plane-wave method
Jian-Zhou Zhao1,2, Jia-Ning Zhuang1, Xiao-Yu Deng1, Yan Bi2, Ling-Cang Cai2,
Zhong Fang1, and Xi Dai1 1Institute of Physics, Chinese Academy of Sciences,
Beijing 100190, People’s Republic of China 2National key Laboratory of Shock
Wave and Detonation Physics, Institute of Fluid Physics, China Academy of
Engineering Physics, Mianyang, 621900, China
###### Abstract
In this paper, we propose an efficient implementation of combining Dynamical
Mean field theory (DMFT) with electronic structure calculation based on the
local density approximation (LDA). The pseudo-potential-plane-wave method is
used in the LDA part, which makes it possible to be applied to large systems.
The full loop self consistency of the charge density has been reached in our
implementation which allows us to compute the total energy related properties.
The procedure of LDA+DMFT is introduced in detail with a complete flow chart.
We have also applied our code to study the electronic structure of several
typical strong correlated materials, including Cerium, Americium and NiO. Our
results fit quite well with both the experimental data and previous studies.
###### pacs:
71.27.+a, 71.15.Mb, 71.15.Nc, 71.20.-b
## 1 Introduction
The first principle calculation based on the density functional theory (DFT)
with the local density approximation (LDA) and its generalization generalized
gradient approximation (GGA) is very successful in predicting ground-state
properties and band structures of a wide range of real materials. However it
is known for a long time that it is not sufficient to calculate the electronic
structure of those strongly correlated materials (i.e. transition metal
oxides, actinide and lanthanide-based materials, high $T_{c}$ superconductors)
by these methods alone even on the qualitative level. In order to overcome
these shortcomings of the traditional DFT-LDA scheme, remedies such as LDA+$U$
have been proposed, which can describe some of the strong correlated materials
with long range order in their ground states. While even LDA+$U$ can not be
applied to many of the strongly correlated materials, for example, the
paramagnetic Mott insulator phase and the materials containing unfilled
f-shell with strong multiplet effects.
Therefore it is important to have a first principle method which can be
applied to the strongly correlated materials with the ability to capture the
dynamical properties, finite temperature properties and multiplet effects. In
the last two decades, the dynamical mean field theory (DMFT) has been
developed very fast to be the standard tool to deal with the on-site
correlation effect in the limit of large dimension [1]. After being
successfully applied to many model systems, DMFT has been considered as a
powerful tool to capture the on-site correlation effects based on the Hubbard
like Hamiltonians containing both the local interaction terms and the single
particle Hamiltonians extracted from LDA. Therefore the LDA+DMFT method which
combines the DFT-based band structure techniques with DMFT has been proposed
and developed quickly in the last decade. By DMFT the local correlation effect
can be well described by the self energy, which has frequency dependence and
in general takes the matrix form within the subspace spanned by the
correlation orbitals. Using the Green’s function containing self energy, many
of the physical properties of strongly correlated materials can be calculated,
i.e. the electronic spectral function, the total energy, the optical
conductivity and the local spin susceptibility.
Most of the LDA+DMFT calculation till now have been performed by the partially
self-consistent scheme, where the local self-energy is obtained by the DMFT
calculation with a fixed LDA Hamiltonian generated from the fixed LDA charge
density. Therefore, in this simplified scheme one neglects the inference of
the strong on-site Coulomb interaction on the charge density. The above
mentioned “one-shot” DMFT calculation works quite well for the electronic
structure. While in order to obtain reliable total energy related properties,
i.e. the equilibrium volume, the elastic constants and the phonon frequencies,
the LDA+DMFT scheme with full charge density self consistency is needed. Up to
date, there have been several fully self-consistent LDA+DMFT schemes [2, 3] as
well as actual implementations [4, 5, 6, 7]. As mentioned in [7], there are
three major issues that have to be addressed in DFT+DMFT implementations: i)
quality of the basis set, ii) quality of the impurity solvers, iii) choice of
correlated orbitals onto which the full Green’s function is projected. In the
implementations of LDA+DMFT mentioned above with full charge self-consistency
, the basis set of the linear muffin-tin orbitals (LMTO) are used, and the
local correlated orbitals are naturally chosen based on the muffin-tin
orbitals . While on the other hand, there are very few implementations of
LDA+DMFT with full charge self-consistency are based on the plane wave
methods.
In the present paper, we implement a full self-consistent LDA+DMFT scheme
using pseudo-potential-plane-wave as the basis set, and we use either atomic
orbitals or Wannier orbitals depending on different physics systems. The
Hubbard-I approximation is used as the impurity solver in the present paper,
but it can be replace by more precise solver like continuous time quantum
Monte Carlo (CTQMC) for metallic systems. The present paper is organized as
follows. In section 2, we describe the way to choose and construct the
correlated orbitals onto which the local interactions exert. In section 3, we
show the full-loop flow and LDA+DMFT(Hubbard-I) formalism in detail. We apply
this LDA+DMFT approach to several correlated materials such as Ce, Am, and NiO
in section 4. Finally we conclude our work in section 5.
## 2 Projection onto localized orbitals
In the LDA+DMFT calculation, all the energy bands are divided into two groups.
And only the on-site interactions among those local orbitals are treated more
precisely by DMFT. Therefore the choice of localized orbitals does affect the
results obtained by LDA+DMFT method and becomes one of the important issues in
LDA+DMFT. A natural choice for the local orbitals is a set of atomic like wave
functions with corresponding $d$ or $f$ characters. Typical atomic like local
orbitals is the linear muffin-tin orbitals (LMTOs)[8] adopted in early
LDA+DMFT implementations[9, 10, 11]. However a more physical choice should be
wannier functions, since the shape of the local orbitals will be altered in
crystals especially when there is strong hybridization between localized
orbital and delocalized $p$-type or $s$-type orbitals. The wannier functions
are not uniquely determined by the Bloch wave functions, so different choice
can be made, for example, the Nth-order muffin-tin orbitals[12], the projected
wannier functions[13] and the Maximally localized Wannier functions
(MLWFs)[14, 15]. Comparisons of these choices have been made in previous
literature. In this paper, two kinds of local orbitals, atomic like functions
and the projected wannier functions are adopted according to the different
systems.
This implementation of LDA+DMFT method reported in this paper is developed on
an existing DFT package BSTATE (Beijing Simulation Tool of Atomic TEchnique),
which is based on the pseudo-potentials and plane waves. Unlike the LMTO
methods, the local orbitals do not enter the basis set of plane waves and
should be constructed in a suitable way. Projected wannier functions or MLWFs
as local orbitals have been used by previous reported LDA+DMFT implementations
based on pseudo-potentials (or projector augmented wave (PAW)) method) and
plane wave scheme[16, 17, 18]. In this report, two kinds of orbitals are used
according to the need, one is directly derived from the atomic wave function
of an isolated atom and the other is projected wannier functions constructed
from these atomic wave functions. In order to be self-contained the
construction procedure is presented below in detail. Since all our
calculations are based on plane wave set, the local orbitals will be projected
onto these plane waves.
First we consider the atomic like local orbitals. The localized nature of the
correlated bands with $d$ or $f$ characters in crystals ensures that these
local orbitals are very similar to the corresponding atomic wave functions of
an isolated atom. The atomic wave functions can be picked as the local
orbitals directly if the hybridization in crystal could be neglected. All
atomic wave functions for the isolated atom could be obtained by the solving
the all-electron radical Schrödinger equation
$(T+V(r))\psi_{nl}^{ALL}(r)=E_{nl}\psi_{nl}^{ALL}(r)$ (1)
In this manner, all orbitals are well defined by the primer quantum number $n$
and angular quantum number $l$. In realistic systems, the local orbitals
usually come from the partially filled $d$ or $f$ shell and could be picked
according to its character. Often there is only one type of local orbitals, so
the local subspace can be labeled by $l$. Of course, considering the spherical
symmetry of the isolated atom, the spherical harmonics should be multiplied
the radical wave functions to form a complete set. The local orbitals are
$|\phi_{lm,I}\rangle=|\psi_{l,I}^{ALL}Y_{lm}\rangle,$ (2)
and in which $m=1\dots 2l+1$ denote different angular components.
In methods based on plane waves, it is not desirable to use the all-electron
wave function directly since it requires a large amount of plane waves to
expand in momentum space. To avoid this problem, in PAW method all-electron
atomic partial waves can be used, while in pseudo-potentials method, pseudo
atomic wave functions can be chosen. The latter is used in this paper. During
the generation of pseudo-potentials, the Schrödinger equation obeyed by the
pseudo wave functions is as below
$(T+V^{PS}(r))\psi_{nl}^{PS}(r)=E_{nl}^{PS}\psi_{nl}^{PS}(r)$ (3)
The details on the generation of pseudo-potentials can be referred to previous
literature[19, 20]. The spirit of pseudo-potentials is quite straight forward.
Beyond a core radius $r_{c}$, the pseudo-potential $V^{PS}$ play the role of a
scattering center just as a real atomic potential $V^{ALL}$ do. Thus, the
pseudo eigen energy $E_{nl}^{PS}$ should be the same as the realistic one
$E_{nl}^{ALL}$, and both the pseudo wave functions $|\psi_{n,I}^{PS}\rangle$
and the pseudo-potential $V^{PS}$ coincide with the exact wave functions
$|\psi_{n,I}^{ALL}\rangle$ and the realistic potential $V^{ALL}$ of an
isolated atoms beyond this core radius $r_{c}$, respectively.
$\eqalign{E_{nl}^{PS}=E_{nl};\cr
V^{PS}(r)=V(r),\qquad\psi^{PS}(r)=\psi^{ALL}(r)\qquad r>r_{c}}$ (4)
The quality of pseudo wave functions is the same as the quality of the pseudo-
potentials, which could be justified by comparing simple DFT calculations with
accepted results. The pseudo wave functions bear the same atomic features as
the exact atomic wave functions, and can be picked as the local orbitals as
above (angular quantum number $l$ ignored since usually only one kind of local
orbital are considered).
$|\phi_{m,I}\rangle=|\psi_{l,I}^{PS}Y_{lm}\rangle,$ (5)
It is convenient to transform the local orbitals from real space to momentum
space in crystal calculations
$|\phi_{m,{\mathbf{k}}}\rangle=\sum_{I}e^{i{\mathbf{k}}\cdot{\mathbf{R}}_{I}}|\phi_{m,I}\rangle.$
(6)
Since the atomic wave functions at different site $I$ are not orthogonal, an
orthogonalization procedure is needed for the above local orbitals.
The DFT calculations give a set of Bloch waves $|\Psi_{n\mathbf{k}}\rangle$
spanning the total Hilbert space in which lies the local subspace spanning by
the local orbitals. Thus, the physical choice of local orbitals is Wannier
functions derived from the Bloch waves. The Wannier functions are localized in
real space and could be constructed from the Bloch waves via a unitary
transformation. However, this unitary transformation is not unique since the
phase factors of Bloch waves are uncertain, which results in that there are
many kinds of Wannier functions in use, e.g., the projected Wannier functions,
NMTOs and MLWFs. Also it is not necessary to include all the Bloch waves but
the relevant bands lying in a certain energy range to construct the Wannier
functions. The local orbitals could be picked as a subset of these Wannier
functions with specified $d$ or $f$ localized character.
The projecting Wannier functions method is a simple way to construct wannier
functions, in which trial wave functions are projecting onto physical relevant
Bloch bands. The atomic like local orbitals above are used as trial functions
here. The Bloch bands can be selected by band indices ($N_{i},...N_{j}$) or by
an energy interval ($E_{i},E_{j}$) enclosing them.
$\displaystyle|W_{m,{\mathbf{k}}}\rangle$ $\displaystyle=$
$\displaystyle\sum_{n=N_{i}}^{N_{j}}|\psi_{n{\mathbf{k}}}\rangle\langle\psi_{n{\mathbf{k}}}|\phi_{lm,{\mathbf{k}}}\rangle$
(7) $\displaystyle=$
$\displaystyle\sum_{n(E_{i}<\epsilon_{n{\mathbf{k}}}<E_{j})}|\psi_{n{\mathbf{k}}}\rangle\langle\psi_{n{\mathbf{k}}}|\phi_{m,{\mathbf{k}}}\rangle.$
These orbitals need normalization and orthogonalization to be true Wannier
functions.
$|\tilde{W}_{m,{\mathbf{k}}}\rangle=\sum_{m^{\prime}}O_{m,m^{\prime}}({\mathbf{k}})^{-1/2}|{W}_{m,{\mathbf{k}}}\rangle$
(8)
in which $O$ is the overlap matrix between different
$|W_{m,{\mathbf{k}}}\rangle$s
$O_{m,m^{\prime}}({\mathbf{k}})=\langle
W_{m,{\mathbf{k}}}|W_{m^{\prime},{\mathbf{k}}}\rangle$ (9)
Although the construction method can not give the most localized Wannier
orbitals, the basic features of the localized bands are captured to a very
good extend as indicated by a few reports[13, 16, 17, 18], and also proved by
our calculations in this paper.
## 3 LDA+DMFT(Hubbard-I) Formalism
We will introduce the detail procedge of LDA+DMFT with full loop charge
density self consistency in this section. First we plot the general flow chart
in figure 1, which contains a inner DMFT loop and a outer charge density loop.
However, if we use Hubbard-I approximation as the impurity solver, the inner
DMFT loop is thus neglected.
Figure 1: The most general flow chart for LDA+DMFT scheme. The first LDA
calculation gives out the band structure on LDA level, and constructs local
orbitals, then the full Hamiltonian is established and solved by a DMFT loop,
after which the charge density can be recalculated by Fourier transforming the
k-space density matrix to real space. The “non-interacting” Hamiltonian is
thus regenerated based on the new charge density profile and the calculation
will be completed when the self consistency has been reached for the full
loop.
In the following subsections, we are going to describe the whole process in
detail.
### 3.1 First LDA Calculation
The first step of the full loop LDA+DMFT is an LDA self-consistent
calculation, whose main purpose is to generate an effective single particle
Hamiltonian $\hat{H}_{LDA}$ and construct the correlated orbitals.
Generally the effective LDA Hamiltonian can be expressed as the following:
$\displaystyle\hat{H}_{LDA}=\sum_{n\mathbf{k}}E_{n\mathbf{k}}\hat{c}^{\dagger}_{n\mathbf{k}}\hat{c}_{n\mathbf{k}}$
(10)
In the above equation, $n=1\sim N_{band}$ are the joint indices of band and
spin; $E_{n\mathbf{k}}$ is the eigen energy of the Kohn-Sham equation
determined by
$\hat{H}_{LDA}|\varphi_{n\mathbf{k}}\rangle=E_{n\mathbf{k}}|\varphi_{n\mathbf{k}}\rangle$,
where $|\varphi_{n\mathbf{k}}\rangle$ represents a set of orthonormal Bloch
functions.
The LDA calculation also complement the construction of local correlated
orbitals. In the present implementation the atomic orbitals and the Wannier
functions are two types of commonly used local basis. In LDA+DMFT, the local
interactions have been considered on the DMFT level only within the local
orbitals and in order to set up the DMFT self consistent equation we need to
obtain the overlap matrix between local orbitals and Bloch wave functions,
which takes the form of
$S^{\mathbf{k}}_{\alpha,n}=\langle\alpha_{\mathbf{k}}|\varphi_{n\mathbf{k}}\rangle$,
with Greek letter $\alpha$ denote the index of local orbitals. The
completeness of the Bloch basis set gives
$\displaystyle\sum_{n}|S^{\mathbf{k}}_{\alpha,n}|^{2}=1$ (11)
for any local orbital $\alpha$ and any k-point.
### 3.2 DMFT Loop
The purpose of the DMFT loop is to calculate the self-energy caused by the
local interactions through the DMFT self consistent loop. As we have mentioned
before, the DMFT loop is not necessity if we use Hubbard-I approximation as
the impurity solver. However, here we first introduce the algorithm with more
general impurity solvers.
By adding the local interaction terms to LDA, we get the total Hamiltonian of
the system which is to be solved by DMFT,
$\displaystyle\hat{H}_{LDA+DMFT}=\hat{H}_{LDA}+\sum_{i}(\hat{H}^{i}_{U}-\hat{V}^{i}_{DC})-\mu\hat{N}$
(12)
where $i$ is the site index. For a specified site, we remove the $i$ index and
$\displaystyle\hat{H}_{U}=\frac{1}{2}\sum_{\alpha\beta\gamma\delta}U_{\alpha\beta\gamma\delta}\hat{f}^{\dagger}_{\alpha}\hat{f}^{\dagger}_{\beta}\hat{f}_{\gamma}\hat{f}_{\delta}$
(13)
is written as a general form of two-body local interactions, in which the
creation and annihilation operators of correlated orbitals
$\hat{f}^{\dagger}_{\alpha}$ and $\hat{f}_{\alpha}$ are associated with Bloch
operators in term of
$\displaystyle\hat{f}^{\dagger}_{\alpha}=\sum_{n\mathbf{k}}\hat{c}^{\dagger}_{n\mathbf{k}}\langle\varphi_{n\mathbf{k}}|\alpha_{\mathbf{k}}\rangle$
(14)
$\displaystyle\hat{f}_{\alpha}=\sum_{n\mathbf{k}}\hat{c}_{n\mathbf{k}}\langle\alpha_{\mathbf{k}}|\varphi_{n\mathbf{k}}\rangle$
(15)
The third term of (12) is the double-counting term correspond to the
correlated energy that has already been considered in LDA calculation at a
Hartree-Fock mean field level,
$\hat{V}_{DC}=\sum_{\alpha}V_{DC}^{\alpha}\hat{f}^{\dagger}_{\alpha}\hat{f}_{\alpha}$
(16)
and this term will be discussed in section 3.2.3. The last term of (12)
$\hat{N}=\sum_{n\mathbf{k}}\hat{c}^{\dagger}_{n\mathbf{k}}\hat{c}_{n\mathbf{k}}$
(17)
is the total number of particle operator, while the chemical potential $\mu$
controls the occupation number of the unit cell.
#### 3.2.1 Quantum Impurity Hamiltonian
In DMFT, the correlation problem on the lattice can be mapped to a quantum
impurity model, which contains the same on-site interaction and reads,
$\displaystyle\hat{H}_{imp}=\sum_{q\alpha}\epsilon_{q\alpha}\hat{c}^{\dagger}_{q\alpha}\hat{c}_{q\alpha}+\sum_{q\alpha}V_{q\alpha}(\hat{c}^{\dagger}_{q\alpha}\hat{f}_{\alpha}+h.c.)+\sum_{\alpha}E^{imp}_{\alpha}\hat{f}^{\dagger}_{\alpha}\hat{f}_{\alpha}+\hat{H}_{U}$
(18)
The inference from the rest of the lattice site besides the one considered in
the impurity model is simulated by a non-interacting “heat bath”, which is
described by the first term in (18). The second term describes the coupling
between the impurity site and the heat bath and the rest two terms describe
the local interactions.
#### 3.2.2 Hybridization function and Weiss field
The above quantum impurity model can be solved by the impurity solver, i.e.
Hubbard-I, and the self energy $\hat{\Sigma}(i\omega)$ is then obtained, with
which we can construct the lattice Green’s function by applying the same self
energy term as,
$[\hat{G}^{\mathbf{k}}_{lattice}]^{-1}=\rmi\omega-\hat{H}_{\mathbf{k}}-\hat{\Sigma}_{\mathbf{k}}(\rmi\omega)+\mu$
(19)
The hybridization function, which characterizes the dynamics of the “heat
bath” is defined as
$\Delta(\rmi\omega)_{\alpha\beta}=\delta_{\alpha\beta}\sum\limits_{q}\frac{|V_{q\alpha}|^{2}}{\rmi\omega-\epsilon_{q\alpha}}$
and can be obtained iteratively by the following DMFT self consistent
equation.
$\hat{\Delta}(\rmi\omega)=\rmi\omega-\hat{E}_{imp}-\hat{\Sigma}+\mu-\left[\sum_{k}G^{\mathbf{k}}_{lattice}(\rmi\omega)\right]^{-1}$
(20)
The above equation is obtained by requiring that the local Green’s function on
the lattice should equal to the Green’s function of the quantum impurity
problem. The equations (19) and (20) form a closed self consistent loop, which
determines both the self energy and the hybridization function iteratively.
#### 3.2.3 Double Counting Term
In this section we discuss the explicit expression of $E^{imp}_{\alpha}$ as
well as the double-counting term in (12). For 3d system under cubic symmetry,
the local interaction can be simplified as coulomb interaction $U$ and Hund’s
rule coupling $J$, which can be written as,
$\hat{H}_{U}=U\sum_{b}\hat{n}_{b,\uparrow}\hat{n}_{b,\downarrow}+(U-2J)\sum_{b<b^{\prime}\atop\sigma\sigma}\hat{n}_{b,\sigma}\hat{n}_{b^{\prime},\sigma}-J\sum_{b<b^{\prime}\atop\sigma}\hat{n}_{b\sigma}\hat{n}_{b^{\prime}\sigma}$
(21)
And the double counting energy in this case has been studied by Held et
al[21]. and can be approximately chosen as
$\displaystyle E_{DC}=\frac{1}{2}\bar{U}n_{\mbox{\tiny
LDA+DMFT}}(n_{\mbox{\tiny LDA+DMFT}}-1)$ (22)
where
$\displaystyle\bar{U}$ $\displaystyle=$
$\displaystyle\frac{U+2(M-1)(U-2J)-(M-1)J}{2M-1}$ (23) $\displaystyle M$
$\displaystyle=$ $\displaystyle 2l+1$ (24)
and $n_{\mbox{\tiny LDA+DMFT}}$ is the total number of electrons on correlated
orbitals for a specific atom, hence in (10) becomes
$\displaystyle\hat{V}_{DC}$ $\displaystyle=$
$\displaystyle\sum_{\alpha}\frac{\partial E_{DC}}{\partial
n_{\alpha}}\hat{n}_{\alpha}$ (25) $\displaystyle=$
$\displaystyle\sum_{\alpha}\bar{U}(n_{\mbox{\tiny
LDA+DMFT}}-\frac{1}{2})\hat{n}_{\alpha}$
The final expression of $E^{imp}$ is
$\displaystyle E^{imp}_{\alpha}$ $\displaystyle=$
$\displaystyle\sum_{k}\langle\alpha_{k}|\hat{H}_{LDA}|\alpha_{k}\rangle-
V_{DC}^{\alpha}$ (26)
where
$\displaystyle V_{DC}^{\alpha}=\bar{U}(n_{\mbox{\tiny LDA+DMFT}}-\frac{1}{2})$
(27)
The way to remove the form of double-counting energy is not unique, and in
fact this process needs intuition of physics. The double-counting discussed in
earlier sections is a very usual way which is also used in LDA+$U$ method.
Besides, the following two ways are also commonly used
* •
to remove self-energy at zero-frequency
$\hat{\Sigma}(\rmi\omega)\rightarrow\hat{\Sigma}(\rmi\omega)-\hat{\Sigma}(0)$
* •
to remove self-energy at infinity-frequency
$\hat{\Sigma}(\rmi\omega)\rightarrow\hat{\Sigma}(\rmi\omega)-\hat{\Sigma}(\infty)$
Furthermore, the double-counting energy is also regarded as “impurity solver
dependent”. For example, it is reasonable to use an integer to replace
$n_{\mbox{\tiny LDA+DMFT}}$ in (22) for the Hubbard-I solver[6], and in this
paper we also follow their suggestion.
#### 3.2.4 Hubbard-I solver
The essential point of Hubbard-I solver is to neglect the effect of the heat
bath and take the atomic self energy as the zeroth order approximation for the
quantum impurity problem, which can be written as
$\displaystyle\hat{\Sigma}^{atom}=[\rmi\omega_{n}+\mu_{atom}-\hat{E}^{imp}]^{-1}-[\hat{G}^{atom}]^{-1}$
(28)
In the above equation $\hat{G}^{atom}$ is the Green’s function for a single
atom, which can be expressed as
$\displaystyle
G^{atom}_{\alpha\alpha^{\prime}}(\rmi\omega)=\sum_{\Gamma\Gamma^{\prime}}\frac{(F^{\alpha})_{\Gamma\Gamma^{\prime}}(F^{\alpha^{\prime}\dagger})_{\Gamma^{\prime}\Gamma}(X_{\Gamma}+X_{\Gamma^{\prime}})}{\rmi\omega-
E_{\Gamma^{\prime}}+E_{\Gamma}}$ (29)
where $|\Gamma\rangle$ and $|\Gamma^{\prime}\rangle$ are the atomic
eigenstates obtained by exact diagonalization of a single atom problem,
$X_{\Gamma^{\prime}}={e^{-\beta E_{\Gamma^{\prime}}}}/({\sum_{\Gamma}e^{-\beta
E_{\Gamma}}})$ represents the occupation probability of the local
configuration $|\Gamma\rangle$, and
$F^{\alpha}_{\Gamma\Gamma^{\prime}}=\langle\Gamma|f_{\alpha}|\Gamma^{\prime}\rangle$.
Hubbard-I approximation is
$\displaystyle\hat{\Sigma}\approx\hat{\Sigma}^{atom}$ (30)
In reference [22] , Savrasov et alpropose that the above self energy can be
written by the summation of a set of poles, which greatly simplifies the DMFT
calculation, because it is not necessary to handle the full frequency
dependence of the Green’s function. In Hubbard-I approximation, it can be
proved that both the atomic Green’s function and self energy are diagonal and
can be written in the form of pole expansion if the following two necessary
conditions are satisfied. 1) The single particle basis used here should be the
one which diagonalize the local density matrix; 2) The atomic Hamiltonian only
contains the two-body interaction terms. Obviously by chosen the proper single
particle basis the atomic Hamiltonian (18) in general will satisfy the above
two condition and thus can be always written in terms of pole expansion.
Unlike the reference[22], where they only use a few poles to capture the main
features of the self energy, in the present implementation, we keep all the
poles in the self energy, which makes it more accurate. From (29), it is very
clear that the atomic green’s function has already been expressed in terms of
poles and the pole expansion of the corresponding self energy can be obtained
using (28), which is introduced in details in the Appendix. Therefore in
general the above atomic self energy can be written as
$\hat{\Sigma}(\rmi\omega)=\hat{\Sigma}({\infty})+\sum^{np}_{i=1}\hat{V}^{\dagger}_{i}(\rmi\omega-
p_{i})^{-1}\hat{V}_{i}$ (31)
where $i$ labels the number of poles, $np$ is the total number of poles in the
self energy and $\hat{V}_{i}$ is a vector defined in the orbital space
describing the distribution of the $ith$ “pole states” among the local
orbitals.
### 3.3 Correction of the Density Matrices and Pole Expansion of the Self-
energy
Once we obtain the converged local self-energy, we are at the point to correct
the k-dependent lattice Green’s functions as well as the density matrices. In
general we need to do the summation over the Matsubara frequencies, which is
time consuming for realistic materials contains lots of bands. But it can be
greatly simplified when the self energy can be expressed as the summation of
poles, which can be written as
$\displaystyle\langle\varphi_{n_{1}\mathbf{k}}|\hat{\rho}_{LDA+DMFT}|\varphi_{n_{2}\mathbf{k}}\rangle$
(32) $\displaystyle=$
$\displaystyle\frac{1}{\beta}\sum_{i\omega_{n}}\langle\varphi_{n_{1}\mathbf{k}}|\frac{1}{\rmi\omega_{n}+\mu-\hat{H}_{LDA}+\hat{V}_{DC}-\hat{\Sigma}(i\omega_{n})}|\varphi_{n_{2}\mathbf{k}}\rangle$
As pointed out firstly in reference [22], when the self energy can be written
in terms of poles, the full green’s function can be expressed as the
“physical” part of an enlarged “pseudo Hamiltonian”, defined as
$\displaystyle\hat{H}_{ps}(\mathbf{k})=\left(\begin{array}[]{cc}\hat{H}_{LDA}(\mathbf{k})-\hat{V}_{DC}+\hat{\Sigma}(\infty)&\qquad-\hat{V}^{\dagger}\\\
-\hat{V}&\qquad\hat{P}\end{array}\right)$ (35)
where
$\hat{H}_{LDA}(\mathbf{k})=E_{n\mathbf{k}}\hat{c}^{\dagger}_{n\mathbf{k}}\hat{c}_{n\mathbf{k}}$
is the physical Hilbert space, $\hat{V}_{i}$ is defined in (31) and
$\hat{P}=Diag(p_{1},p_{2},....,p_{np})$. Therefore the physical green’s
function can be expressed in terms of eigenstate of the pseudo Hamiltonian
simply as
$\displaystyle\langle\varphi_{n_{1}\mathbf{k}}\bigg{|}\frac{1}{\rmi\omega_{n}-\hat{H}_{LDA}(\mathbf{k})+\hat{V}_{DC}-\hat{\Sigma}(i\omega_{n})}\bigg{|}\varphi_{n_{2}\mathbf{k}}\rangle$
(36) $\displaystyle=$
$\displaystyle\langle\varphi_{n_{1}\mathbf{k}}|\psi^{ps}_{l\mathbf{k}}\rangle\frac{1}{\rmi\omega_{n}-E^{ps}_{l}(\mathbf{k})}\langle\psi^{ps}_{l\mathbf{k}}|\varphi_{n_{2}\mathbf{k}}\rangle$
where $|\psi^{ps}_{l\mathbf{k}}\rangle$ and $E^{ps}_{l}(\mathbf{k})$ are the
eigenstate and eigenvalue of pseudo Hamiltonian respectively.
Then the sum of frequencies in (32) can be performed directly
$\displaystyle\langle\varphi_{n_{1}\mathbf{k}}|\hat{\rho}_{LDA+DMFT}|\varphi_{n_{2}\mathbf{k}}\rangle$
(37) $\displaystyle=$
$\displaystyle\frac{1}{\beta}\sum_{i\omega_{n}}\sum_{l\mathbf{k}}\langle\varphi_{n_{1}\mathbf{k}}|\psi^{ps}_{l\mathbf{k}}\rangle\frac{1}{\rmi\omega_{n}-E_{l}^{ps}(\mathbf{k})}\langle\psi^{ps}_{l\mathbf{k}}|\varphi_{n_{2}\mathbf{k}}\rangle$
$\displaystyle=$
$\displaystyle\sum_{l\mathbf{k}}\langle\varphi_{n_{1}\mathbf{k}}|\psi^{ps}_{l\mathbf{k}}\rangle\langle\psi^{ps}_{l\mathbf{k}}|\varphi_{n_{2}\mathbf{k}}\rangle
n_{F}[E_{l}^{ps}(\mathbf{k})-\mu]$
in which $\mu$ is exactly the chemical potential in (12), determined by the
occupation number of electrons $N_{tot}$ in the unit cell
$\displaystyle N_{tot}$ $\displaystyle=$
$\displaystyle\sum_{n\mathbf{k}}\langle\varphi_{n\mathbf{k}}|\hat{\rho}|\varphi_{n\mathbf{k}}\rangle$
(38) $\displaystyle=$
$\displaystyle\sum_{ln\mathbf{k}}|\langle\varphi_{n\mathbf{k}}|\psi^{ps}_{l\mathbf{k}}\rangle|^{2}n_{F}[E_{l}^{ps}(\mathbf{k})-\mu]$
### 3.4 Evaluation of Real Space Charge Density and Complete of the Full Loop
After DMFT obtains the corrected density matrix $\hat{\rho}$, the real space
charge density will be generated based on it by Fourier transformation, and
then, the new LDA Hamiltonian as well as the new overlap matrices between
Bloch states and local orbitals will be recalculated, which closes the full
iteration loop for the charge density self consistency. The occupation number
of the Bloch basis can also be obtained by,
$\displaystyle n_{\mbox{\tiny
LDA+DMFT}}=\sum_{\mathbf{k},n_{1}n_{2},\alpha}\langle\alpha_{\mathbf{k}}|\varphi_{n_{1}\mathbf{k}}\rangle\langle\varphi_{n_{1}\mathbf{k}}|\hat{\rho}|\varphi_{n_{2}\mathbf{k}}\rangle\langle\varphi_{n_{2}\mathbf{k}}|\alpha_{\mathbf{k}}\rangle$
(39)
### 3.5 Calculation of Physical Quantities
#### 3.5.1 Energy Functional
Most of the physical quantities we are interested in can be calculated once we
have reached the charge density self consistency. First we present the energy
functional of the full loop LDA+DMFT, which can be used to calculate many
quantities, i.e. the total energy, force and so on. It can be written as, [6]
$\displaystyle E=E_{LDA}[\rho]-\langle H_{KS}\rangle_{LDA}+\langle
H_{KS}\rangle+\langle H_{U}\rangle-E_{DC}$ (40)
in which $E_{LDA}[\rho]$ is the expression of the energy within density-
functional theory; $\langle
H_{KS}\rangle_{LDA}=\tr[\hat{H}_{LDA}\hat{G}_{LDA}]$ is the non-interacting
energy at LDA level; $\langle H_{KS}\rangle=\tr[\hat{H}_{LDA}\hat{G}]$ is the
non-interacting energy at DMFT level; $\langle
H_{U}\rangle=\frac{1}{2}\tr[\hat{\Sigma}\hat{G}]$ is the interaction energy
caused by local correlation interactions; the last term $E_{DC}$ is the double
counting energy given before in (22). In above the meaning for “trace” is
defined as
$\displaystyle\tr[\mathcal{A}]=\frac{1}{\beta}\sum_{n\mathbf{k}}\sum_{\rmi\omega_{n}}\langle\varphi_{n\mathbf{k}}|\mathcal{A}(\rmi\omega_{n})|\varphi_{n\mathbf{k}}\rangle$
(41)
and the integral path surrounds all the energies of occupation states.
The first and the second term is easily calculated in the LDA framework, and
the three terms remaining must be evaluated in the DMFT process, explicitly
$\displaystyle\langle
H_{KS}\rangle=\sum_{n\mathbf{k}}E_{n\mathbf{k}}\langle\varphi_{n\mathbf{k}}|\hat{\rho}|\varphi_{n\mathbf{k}}\rangle$
(42)
and
$\displaystyle\langle H_{U}\rangle$ $\displaystyle=$
$\displaystyle\frac{1}{2}\tr[\hat{\Sigma}\hat{G}]$ (43) $\displaystyle=$
$\displaystyle\frac{1}{2}\tr[(\hat{G}^{-1}_{0}-\hat{G}^{-1})\hat{G}]$
$\displaystyle=$
$\displaystyle\frac{1}{2}\left[\sum_{m,n\mathbf{k}}(E^{ps}_{m\mathbf{k}}-\mu)|\langle\varphi_{n\mathbf{k}}|\psi_{m\mathbf{k}}\rangle|^{2}n_{F}(E^{ps}_{m\mathbf{k}}-\mu)\right.$
$\displaystyle\left.-\langle H_{KS}\rangle+\langle\hat{V}_{DC}\rangle\right]$
while $\langle H_{KS}\rangle$ is expressed in (42), and
$\langle\hat{V}_{DC}\rangle$ is the expected value of $\hat{V}_{DC}$ at DMFT
level
$\displaystyle\langle\hat{V}_{DC}\rangle$ $\displaystyle=$
$\displaystyle\sum_{i}\bar{U}_{i}(n^{i}_{\mbox{\tiny
LDA+DMFT}}-\frac{1}{2})\sum_{\alpha}\langle\hat{n}^{i}_{\alpha}\rangle$ (44)
$\displaystyle=$ $\displaystyle\bar{U}_{i}(n^{i}_{\mbox{\tiny
LDA+DMFT}}-\frac{1}{2})n^{i}_{\mbox{\tiny LDA+DMFT}}$ (45)
Notice that here emerges the site index $i$. We should be very careful in the
case that the number of correlated atoms is larger than one. The final
expression of the total energy is
$\displaystyle E=E_{LDA}[\rho]-\langle H_{KS}\rangle_{LDA}+\langle
H_{KS}\rangle$
$\displaystyle+\frac{1}{2}\left[\sum_{m,nk}(E^{ps}_{mk}-\mu)|\langle\varphi_{nk}|\psi_{mk}\rangle|^{2}n_{F}(E^{ps}_{mk}-\mu)-\langle
H_{KS}\rangle\right]$
$\displaystyle+\frac{1}{4}\sum_{i}\bar{U}_{i}n^{i}_{\mbox{\tiny LDA+DMFT}}$
(46)
However, if we use Hubbard-I as the impurity solver, and remove double-
counting term by using an integer, the energy functional is given by Haule[7]
$\displaystyle E=E_{LDA}[\rho]-\langle H_{KS}\rangle_{LDA}+\langle
H_{KS}\rangle$
$\displaystyle+\frac{1}{2}\left[\sum_{m,n\mathbf{k}}(E^{ps}_{m\mathbf{k}}-\mu)|\langle\varphi_{n\mathbf{k}}|\psi_{m\mathbf{k}}\rangle|^{2}n_{F}(E^{ps}_{m\mathbf{k}}-\mu)-\langle
H_{KS}\rangle\right]$ $\displaystyle+\frac{1}{4}\sum_{i}\bar{U}_{i}n^{i}_{0}$
(47)
where $n^{i}_{0}$ refers to the integer used to remove double counting for the
$i$th correlated atom.
#### 3.5.2 Expressions of DOS and PDOS
In order to calculate the density of state (DOS) or the partial density of
state (PDOS) on correlated orbitals, it is necessary to calculate the retarded
form of the Green’s function. By using the virtue of pole expansion method, we
write directly
$\displaystyle
D(\epsilon)=(-\frac{1}{\pi})\mathrm{Im}\bigg{[}\sum_{\mathbf{k},m,n}\frac{\langle\varphi_{n\mathbf{k}}|\psi_{m\mathbf{k}}\rangle\langle\psi_{m\mathbf{k}}|\varphi_{n\mathbf{k}}\rangle}{\epsilon+\mu-E^{ps}_{m\mathbf{k}}+\rmi\eta}\bigg{]}\
(\eta\rightarrow 0^{+})$ (48)
and the PDOS of orbital $\alpha$
$\displaystyle D_{\alpha}(\epsilon)$ $\displaystyle=$
$\displaystyle(-\frac{1}{\pi})\mathrm{Im}\left[\sum_{\mathbf{k},m}\frac{\langle\alpha_{\mathbf{k}}|\psi_{m\mathbf{k}}\rangle\langle\psi_{m\mathbf{k}}|\alpha_{\mathbf{k}}\rangle}{\epsilon+\mu-E^{ps}_{m\mathbf{k}}+\rmi\eta}\right]\
(\eta\rightarrow 0^{+})$ (49)
where
$\displaystyle\langle\alpha_{\mathbf{k}}|\psi_{m\mathbf{k}}\rangle=\sum_{n}\langle\alpha_{\mathbf{k}}|\varphi_{n\mathbf{k}}\rangle\langle\varphi_{n\mathbf{k}}|\psi_{m\mathbf{k}}\rangle$
(50)
## 4 Benchmark
To validate the LDA+DMFT framework based on pseudo-potentials-planewave
package and pole expansion of self energy, we benchmark our implementation by
applying to $\gamma$-cerium, americium and paramagnetic NiO. These three are
typical strongly correlated systems in which the valence electrons are
believed to be on the localized side as reported in previous literatures. We
use the Hubbard-I method as impurity solver, which is good enough to capture
the atomic-like features in the Mott insulators, so we expect that these
systems could be well described by our method. We would emphasize here that in
our implementation we can replace the Hubbard-I solver by any of the solvers
as long as the self energy can be written in pole expansion form.
### 4.1 Cerium
The cerium metal attracts lots of research interests for its isostructural
volume-collapse transition from $\gamma$ phase to $\alpha$ phase. The volume
change is about $15\%$ during the transition, which is possibly driven by the
entropy change[23]. Normal LDA calculations could not give a correct
description of cerium, as show in table 1, the equillibrium volume of cerium
given by LDA calculations is smaller than the experimental volume of $\alpha$
phase. This is due to the fact that in LDA the $f$-electrons are treated as
itinerant and the strong correlation effects among them can not be well
captured. When $f$-electron is treated as core electrons, we see that the
equilibrium volume is larger but close to $\gamma$ phase. These simple LDA
results implies that in $\gamma$ phase, the $f$-electrons are more localized
than in $\alpha$-phase. Therefore $f$-electrons in $\gamma$ phase is quite
close to the localized picture like the situation in Mott insulators, which
makes it suitable to applying the Hubbard-I solver. The $\alpha$-Ce which is
stable in low temperature is a correlated metal, which is confirmed by many
other LDA+DMFT calculations[23], and also by our newly development
LDA+Gutzwiller method[24].
In our calculation for $\gamma$-Ce, we used ultrasoft pseudo-potentional in
which 4 _f_ , 5 _s_ , 5 _d_ , 6 _s_ orbitals are taken as valence orbitals.
The energy cutoff for plane-wave expansion is 12.5 Ha for convergence.
Calculations are performed with $10\times 10\times 10$ k-points. For DMFT
calculation, the on-site Coulomb interaction $U$ is chosen to be 6.0 eV as
suggested in previous reports.[25, 23, 26]
In figure 2 the total energy versus volume for $\gamma$-cerium obtained by
LDA, LDA+$U$ and LDA+DMFT are plotted. The equilibrium volume and bulk modulus
are also shown in table 1. We can find in figure 2 that the equilibrium volume
obtained by our LDA+DMFT calculation is very close to the experimental data,
which shows great improvements over LDA . The bulk modulus is also improved
but still larger than the experimental data, which may imply the contribution
for lattice vibration is inneglectable as discussed in reference [6].
Table 1: Lattice parameter and bulk modulus of $\gamma$-cerium according to
both experiment and calculations
* | Volume($\mathrm{\AA}^{3}$) | Bulk modulus(GPa)
---|---|---
Experiments | 34.35[27] | 21[28]
LDA | 22.36 | 62.09
LDA(fcore)[29] | 36.39 | 30.2
LDA+$U$/PAW[30] | 32.0 | 34
LDA+DMFT[6] | 30.10 | 48.49
LDA+DMFT | 33.29 | 38.27
Figure 2: Total energy versus volume curves obtained by LDA, LDA+$U$[6], and
LDA+DMFT. The dashed lines are fitted by Birch EOS. The energy were shifted
for different curves. The experimental equilibrium volume is from reference
[27]
Figure 3 shows the 4 _f_ partial DOS for $\gamma$ cerium at the equilibrium
volume, the position of the two Hubbard bands agree well with the experimental
result.
Figure 3: 4 _f_ pdos for $\gamma$ cerium in LDA+DMFT (Hubbard I), XPS and BIS
experimental data are from reference [31] and [32]
### 4.2 Americium
Americium, which is widely used in smoke detectors, can be viewed as the
”changing point” of the actinide series, where the 5 _f_ electrons changes
from delocalized to localized states[33]. Valence-band ultraviolet
photoemission experiment by Naegele et al[34] shows that the 5 _f_ states are
strongly localized. Am undergoes a series of structure phase transitions with
the increment of pressure, from double hexagonal close packed(dhcp, P63/mmc),
to face centered cubic(fcc, Fm3m, 6.1 GPa), face centered orthorhombic(Fddd,
10.0 GPa), and to primitive orthorhombic structure(Pnma, 16 PGa)[35].
Therefore Am provides an interesting playfield to investigate the relation
between $f$-electron localization and structure transition. Here we focus on
the fcc phase and dhcp phase.
In the calculation, we used ultrasoft pseudo-potentional contains 5 _f_ , 6
_p_ , 6 _d_ , 7 _s_ orbitals as valence orbitals, the plane-wave expansion
kinetic energy cutoff of 12.5 Ha. The LDA self-consistent calculations were
performed with $10\times 10\times 10$ k-points grid for fcc phase, $9\times
9\times 3$ k-points grid for dhcp phase. For the volume calculation, we only
considered density-density interaction, and on-site Coulomb interaction
$F^{(0)}=4.5$ eV, for the partial DOS calculation, we considered general
interaction, and take the atomic values $F^{(2)}=7.2$ eV, $F^{(4)}=4.8$ eV,
$F^{(6)}=3.6$ eV[36].
We present the optimized volume and bulk modulus for both dhcp and fcc phase
in table 2 and figure 4. Our LDA+DMFT results are quite close to the
experimental data and show great improvement over LDA. Also our results are
quite consistent with the reference [22], in which the full potential LMTO
methods are used in the LDA part. The equation of state of Am obtained by our
LDA+DMFT calculation is plotted in figure 5, which again shows very good
agreement with both the experimental data and previous LDA+DMFT calculation
based on LMTO method[22].
Table 2: Am equilibrium volume and bulk modulus for dhcp and fcc phase by
different methods, together with the experimental results[35, 37].
* | Volume($\mathrm{\AA}^{3}$) | Bulk modulus(GPa)
---|---|---
experiment | 29.250 | 29.9
GGA(dhcp)[38] | 19.916 | 70.0
LDA(dhcp) | 18.55 | 118.99
LDA(fcc) | 17.73 | 169.40
LDA+DMFT(dhcp) | 27.04 | 52.70
LDA+DMFT(fcc) | 26.99 | 50.75
LDA+DMFT[22] | 27.4 | 45
Figure 4: Total energy per atom of Am obtained by LDA and LDA+DMFT. The
LDA+DMFT results of Am for both phase are very close, and the experimental
equilibrium volume is denoted by dot line. Figure 5: Calculated P-V relation
of dhcp and fcc Am. Here, $V_{0}$ is experimental equilibrium volume, The
experimental P-V results[37] are shown by black dots.
Figure 6 shows the 5 _f_ partial DOS of fcc phase Am obtained both by our
LDA+DMFT calculation and LDA+DMFT(OCA) by Savrasov[22],together with the
photoemission data from ref.[34]. Our results is quite consistent with the
previous calculations and the experimental results.
Figure 6: The spectral function of Am evaluated by LDA+DMFT. The experimental
photoemission results is from reference [34] are denoted by solid dots.
### 4.3 Paramagnetic phase of NiO
NiO has been heavily studied as a prototype of Mott insulators. It has been
studied within the frame of LDA+DMFT by many groups[39, 40]. Therefore NiO can
be used as a good benchmark material for the implementation of LDA+DMFT. In
this paper, we apply our code to study the paramagnetic phase of NiO. First we
plot the partial density of states (PDOS) obtained by LDA in figure 7, which
clearly shows metallic behavior and is not consistent with the experiments. In
our LDA+DMFT(Hub1) calculation, we use the wannier functions as the local
basis and choose on-site interaction $U=8.0eV$, Hund’s coupling $J=1.0eV$. The
Mott insulating features of NiO then can be well captured by our LDA+DMFT
method, as shown by the PDOS in figure 8).
Figure 7: The partial DOS of NiO obtained by LDA. The main contribution to the
Fermi surface is attributed to eg-like Wannier orbitals, and the t2g-like
orbitals are almost fully occupied. Figure 8: The spectral function of NiO
obtained by LDA+DMFT calculation. As shown in the figure, compared with LDA
results, the t2g-like orbitals are still fully occupied while the eg-like
orbitals split into upper and lower Hubbard bands.
The energy gap obtained by our calculation is around 4.0eV, which is also
quite consistent with the experiments. The photo emission and BIS data
obtained by Sawatzky and Allen[41] are also plotted in figure 9. We can find
that our results fit the photo emission and BIS data very well.
Figure 9: The spectral function of NiO evaluated by LDA+DMFT, with the
comparison of experimental data[41]. The behavior of the density of state near
the fermi surface fits well with the experimental data.
## 5 Conclusions
The new implementation of LDA+DMFT based on the pseudo-potential plane-wave
method is introduced in detail in this paper. We choose the Hubbard-I method
as the impurity solver to solve the quantum impurity problem generated by
DMFT, which is quite suitable for the Mott insulator materials. The most
important advantage of Hubbard-I method is that it is simple enough, which
makes the full loop charge self consistent calculation accessible. We also
point out in the paper that the difficulty of handling frequency dependent
Green’s function can be completely removed by expressing the self energy in
terms of pole expansion, which greatly raise the efficiency of the method.
Finally we benchmark our implementation of LDA+DMFT on several important
correlation materials including $\gamma$-Ce,Am and NiO. Our results for all
these materials fit very well with both the experimental data and the previous
LDA+DMFT results.
## Appendix A Self-energy in pole expansion form
In this appendix, we will introduce how to evaluate self-energy from Green’s
function in the pole expansion form.
Green’s function can be expressed as
$G(i\omega)=\sum^{N_{G}}_{i=1}\frac{V_{i}}{i\omega-P_{i}}$ (51)
where $P_{i}$ is the $i$th pole of the Green’s function, $V_{i}$ is the weight
of the pole, and $N_{G}$ is the number of Green’s function poles. Besides the
self-energy can be expressed as
$\Sigma(i\omega)=\sum^{N_{S}}_{i=1}\frac{W_{i}}{i\omega-Q_{i}}+\Sigma(\infty)$
(52)
where $Q_{i}$ is the $i$th pole of the self-energy, $W_{i}$ is the weight of
the pole, and $N_{S}$ is the number of self-energy poles.
Then the Dyson’s equation (28) becomes
$\sum^{N_{S}}_{i=1}\frac{W_{i}}{i\omega-
Q_{i}}+\Sigma(\infty)=i\omega+\mu-\epsilon_{imp}-\left[\sum^{N_{G}}_{i=1}\frac{V_{i}}{i\omega-
P_{i}}\right]^{-1}$ (53)
First consider the $\Sigma(\infty)$, when $\omega\rightarrow\infty$
$\displaystyle\Sigma(\omega\rightarrow\infty)$
$\displaystyle=i\omega+\mu-\epsilon_{imp}-\left[\sum^{N_{G}}_{i=1}\frac{V_{i}}{i\omega-
P_{i}}\right]^{-1}$ (54)
$\displaystyle=i\omega+\mu-\epsilon_{imp}-\left[\frac{1}{i\omega}\sum^{N_{G}}_{i=1}\frac{V_{i}}{1-\frac{P_{i}}{i\omega}}\right]^{-1}$
$\displaystyle=i\omega+\mu-\epsilon_{imp}-i\omega\left[\sum^{N_{G}}_{i=1}V_{i}-\frac{1}{i\omega}\sum^{N_{G}}_{i=1}V_{i}P_{i}\right]^{-1}$
$\displaystyle=i\omega+\mu-\epsilon_{imp}-i\omega\left(1-\frac{\sum^{N_{G}}_{i=1}V_{i}P_{i}}{i\omega}\right)$
$\displaystyle=\mu-\epsilon_{imp}+\sum^{N_{G}}_{i=1}V_{i}P_{i}$
For the the self energy poles in (53), it corresponding to the zero in the
square bracket of the right side of (53), that is for $Q\in\\{Q_{i}\\}$, we
have
$\sum^{N_{G}}_{i=1}\frac{V_{i}}{Q-P_{i}}=0$ (55)
The left hand side is monotonically decreasing between two adjacent Green’s
function poles, so we could use bisection method to find the self energy
poles.
We rewrite (53) as
$\sum^{N_{s}}_{i=1}\frac{W_{i}}{i\omega-
Q_{i}}=i\omega+\mu-\epsilon-\Sigma(\infty)-\left[\sum^{N_{g}}_{i=1}\frac{V_{i}}{i\omega-
P_{i}}\right]^{-1}$ (56)
For the left hand side
$\displaystyle W_{i}$ $\displaystyle=\lim_{i\omega\rightarrow
Q_{i}}\left(\sum^{N_{S}}_{j=1}\frac{W_{j}}{i\omega-Q_{j}}\right)(i\omega-
Q_{i})$ (57) $\displaystyle=\lim_{i\omega\rightarrow
Q_{i}}\left[i\omega+\mu-\epsilon_{imp}-\Sigma(\infty)-\left(\sum^{N_{g}}_{j=1}\frac{V_{j}}{i\omega-
P_{j}}\right)^{-1}\right](i\omega-Q_{i})$
$\displaystyle=\lim_{i\omega\rightarrow Q_{i}}-(i\omega-
Q_{i})\left(\sum^{N_{G}}_{j=1}\frac{V_{j}}{i\omega-P_{j}}\right)^{-1}$
here we define $\delta=i\omega-Q_{i}$ and
$h(\delta)=\sum^{N_{G}}_{j=1}\frac{V_{j}}{Q_{i}+\delta-P_{j}}$, expand
$h(\delta)$ as
$h(\delta)\approx\delta\cdot
h^{\prime}(0)=-\sum^{N_{G}}_{j=1}\frac{V_{j}}{(Q_{i}-P_{j})^{2}}\cdot\delta$
(58)
so
$W_{i}=\left[\sum^{N_{G}}_{j=1}\frac{V_{j}}{(Q_{i}-P_{j})^{2}}\right]^{-1}$
(59)
here we have got both self energy poles and its weight.
The authors thank L.Wang, G. Xu, and F. Lu for their helpful discussions. We
acknowledge the support from the 973 Program of China (Grant No.
2011CBA00108), from the NSF of China (Grants No. NSFC 10876042 and No. NSFC
10874158), and that from the Fund of Laboratory of Shock Wave and Detonation
Physics, Institute of Fluid Physics, CAEP, Contract No. 2011-056000-0833F.
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|
arxiv-papers
| 2011-11-09T10:02:43 |
2024-09-04T02:49:24.149333
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jian-Zhou Zhao, Jia-Ning Zhuang, Xiao-Yu Deng, Yan Bi, Ling-Cang Cai,\n Zhong Fang, Xi Dai",
"submitter": "Feng Lu",
"url": "https://arxiv.org/abs/1111.2157"
}
|
1111.2183
|
See pages 1-3 of StartDoc.pdf
See pages 1-6 of icrc0726.pdf
See pages 1-4 of icrc0684.pdf
See pages 1-4 of icrc0935.pdf See pages 1-4 of icrc2011_0666_v2.pdf See pages
1-4 of icrc0936.pdf See pages 1-4 of icrc0706_v2.pdf See pages 1-4 of
icrc0580.pdf
See pages 1-4 of icrc1119_Kubo.pdf See pages 1-4 of icrc1091_final.pdf See
pages 1-4 of icrc0941_v2.pdf See pages 1-4 of icrc0692_submitted.pdf See pages
1-4 of 1065_v3.pdf
See pages 1-4 of icrc1021_v2.pdf See pages 1-4 of icrc0688.pdf See pages 1-4
of icrc0668_v5.pdf
See pages 1-4 of icrc0408.pdf
See pages 1-4 of Orr-Krennrich-2011-icrc1156.pdf See pages 1-4 of
icrc1086_v2.pdf See pages 1-4 of icrc0883_v2.pdf
|
arxiv-papers
| 2011-11-09T11:56:41 |
2024-09-04T02:49:24.158214
|
{
"license": "Public Domain",
"authors": "The CTA Consortium",
"submitter": "Bruno Kh\\'elifi",
"url": "https://arxiv.org/abs/1111.2183"
}
|
1111.2272
|
# Numerical and variational solutions of the dipolar Gross-Pitaevskii equation
in reduced dimensions
P. Muruganandam Instituto de Física Teórica, UNESP - Universidade Estadual
Paulista, 01.140-070 São Paulo, São Paulo, Brazil School of Physics,
Bharathidasan University, Palkalaiperur Campus, Tiruchirappalli 620024,
Tamilnadu, India S. K. Adhikari Instituto de Física Teórica, UNESP -
Universidade Estadual Paulista, 01.140-070 São Paulo, São Paulo, Brazil
###### Abstract
We suggest a simple Gaussian Lagrangian variational scheme for the reduced
time-dependent quasi-one- and quasi-two-dimensional Gross-Pitaevskii (GP)
equations of a dipolar Bose-Einstein condensate (BEC) in cigar and disk
configurations, respectively. The variational approximation for stationary
states and breathing oscillation dynamics in reduced dimensions agrees well
with the numerical solution of the GP equation even for moderately large
short-range and dipolar nonlinearities. The Lagrangian variational scheme also
provides much physical insight about soliton formation in dipolar BEC.
###### pacs:
03.75.Hh,03.75.Kk
## I Introduction
The time-dependent mean-field Gross-Pitaevskii (GP) equation can accurately
describe many static and dynamic properties of a harmonically trapped Bose-
Einstein condensate (BEC) rev ; rev-2 ; lpl-1 ; lpl-2 ; lpl-3 ; lpl-4 ;
ref-a01 ; ref-a02 ; ref-a03 ; ref-a04 ; ref-b01 ; ref-b02 ; ref-b03 ; ref-b04
; ref-b05 ; ref-b06 ; balaz2011 . However, the numerical solution of the
three-dimensional (3D) GP equation could often be a difficult task due to a
large nonlinear term CPC ; CPC-1 . Fortunately, in many experimental
situations the 3D axially symmetric harmonic trap has extreme symmetry so that
the BEC has either a cigar or a disk shape geom . In these cases the essential
statics and dynamics of a BEC take place in reduced dimensions. By integrating
out the unimportant dimensional variable(s), reduced GP equations have been
derived in lower dimensions luca ; other ; other-1 , which give a faithful
description of the BEC in disk and cigar shapes. For disk and cigar shapes the
reduced GP equation is written in two (2D) and one dimensions (1D),
respectively. The numerical solution of such 2D, or 1D equation, although
simpler than that of the original 3D equation, remain complex due to the
nonlinear nature of the GP equation. Hence, for small values of the
nonlinearity parameter, a Gaussian variational approximation is much useful
for the solution of these equations var .
The alkali metal atoms used in early BEC experiments have negligible dipole
moment. However, most bosonic atoms and molecules have large dipole moments
and a 52Cr lahaye ; pfau ; pfau-1 ; pfau-4 , and 164Dy dy ; dy-2 BEC with a
larger long-range dipolar interaction superposed on the short-range atomic
interaction, has been realized. Other atoms, like 166Er otherdi ; otherdi-1 ,
and molecules, such as 7Li-133Cs becmol , with much larger dipole moment are
being considered for BEC experiments. A 3D GP equation for a dipolar BEC with
a nonlocal nonlinear interaction has been suggested lahaye and successfully
used to describe many properties of these condensates dipsol ; jb ; YY ;
Yi2000 ; Dutta2007-1 ; Dutta2007-2 ; Dutta2007-3 ; Dutta2007-5 ; Dutta2007-6 .
The applicability of the nonlocal GP equation to the case of dipolar BEC has
been a subject of intensive study lahaye . After a detailed analysis, You and
Yi YY ; Yi2000 concluded that the GP equation is valid for the dipolar BEC.
Further support on the validity of this equation came from the study of
Bortolotti et al. BB ; BB-1 . They compared the solution of the dipolar GP
equation with the results of diffusive Monte Carlo calculations and found good
agreement between the two. However, the 3D GP equation for a dipolar BEC with
the nonlocal dipolar interaction has a complex structure and its numerical
solution, involving the Fourier transformation of the dipolar nonlinear term
to momentum space jb ; YY ; Yi2000 , is even more challenging than that of the
GP equation of a non-dipolar BEC.
Here we reconsider the dimensional reduction SS ; deu of the GP equation to
1D form for cigar-shaped dipolar BEC and obtain the precise 1D potential with
a dipolar contact-interaction term. Previous derivations SS ; deu of the 1D
reduced equation for dipolar BEC did not include the proper contact-
interaction term, lacking which the 1D model will not provide a correct
description of the full 3D system. We also consider the reduced 2D GP equation
fisch ; PS for a disk-shaped dipolar BEC. Though these reduced GP equations
for dipolar BEC are computationally less expensive than their 3D counterparts,
the numerical solution procedure remains complicated due to repeated forward
and backward Fourier transformations of the non-local dipolar term. As an
alternative, here we suggest time-dependent Gaussian Lagrangian variational
approximation of the 1D and 2D reduced equations. A direct attempt to derive
the variational Lagrangian density of the reduced equations is not
straightforward due to nonlocal integrals with error functions. We present an
indirect evaluation of the Lagrangian density avoiding the above complex
procedure. Thus, the present variational approximation involves algebraical
quantities without requiring any Fourier transformation to momentum space.
In case of dipolar BEC of Cr and Dy atoms we consider the numerical solution
of the 3D and the reduced 1D and 2D GP equations for cigar and disk shapes to
demonstrate the appropriateness of the solution of the reduced equations. The
variational approximation of the reduced equations provided results for
density, root-mean-square (rms) size, chemical potential, and breathing
oscillation dynamics in good agreement with the numerical solution of the
reduced and full 3D GP equations.
## II Analytical formulation
### II.1 3D GP Equation
We study a dipolar BEC of $N$ atoms, each of mass $m$, using the dimensionless
GP equation pfau ; pfau-1 ; pfau-4
$\displaystyle i\frac{\partial\phi({\bf r},t)}{\partial t}$
$\displaystyle\,=\biggr{[}-\frac{1}{2}\nabla^{2}+V({\bf r})+4\pi aN|\phi({\bf
r},t)|^{2}$ $\displaystyle\,+N\int U_{dd}({\bf r-r^{\prime}})|\phi({\bf
r^{\prime}},t)|^{2}d^{3}{r^{\prime}}\biggr{]}\phi({\bf r},t),$ (1)
with dipolar interaction $U_{dd}({\bf R})=3a_{dd}(1-3\cos^{2}\theta)$
$/R^{3},$ $\quad{\bf R=r-r^{\prime}}.$ Here $V({\bf r})$ is the confining
axially symmetric harmonic potential, $\phi({\bf r},t)$ the wave function at
time $t$ with normalization $\int|\phi({\bf r},t)|^{2}d{\bf r}=1$, $a$ the
atomic scattering length, $\theta$ the angle between $\bf R$ and the
polarization direction $z$. The constant
$a_{dd}=\mu_{0}\bar{\mu}^{2}m/(12\pi\hbar^{2})$ is a length characterizing the
strength of dipolar interaction and its experimental value for 52Cr is
$15a_{0}$ pfau ; pfau-1 ; pfau-4 , with $a_{0}$ the Bohr radius, $\bar{\mu}$
the (magnetic) dipole moment of a single atom, and $\mu_{0}$ the permeability
of free space. In equation (II.1) length is measured in units of
characteristic harmonic oscillator length $l\equiv\sqrt{\hbar/m\omega}$,
angular frequency of trap in units of $\omega$, time $t$ in units of
$\omega^{-1}$, and energy in units of $\hbar\omega$. The axial and radial
angular frequencies of the trap are $\Omega_{z}\omega$ and
$\Omega_{\rho}\omega$, respectively. The dimensionless 3D harmonic trap is
$V({\bf
r})=\frac{1}{2}\Omega_{\rho}^{2}\rho^{2}+\frac{1}{2}\Omega_{z}^{2}z^{2},$ (2)
where ${\bf r}\equiv(\vec{\rho},z)$, with $\vec{\rho}$ the radial coordinate
and $z$ the axial coordinate.
The Lagrangian density of equation (II.1) is given by
$\displaystyle{\mathcal{L}}$
$\displaystyle\,=\frac{i}{2}(\phi\phi^{\star}_{t}-\phi^{\star}\phi_{t})+\frac{|\nabla\phi|^{2}}{2}+V({\bf
r})|\phi|^{2}$ $\displaystyle\,+2\pi aN|\phi|^{4}+\frac{N}{2}|\phi|^{2}\int
U_{dd}({\mathbf{r}}-{\mathbf{r}^{\prime}})|\phi({\mathbf{r}^{\prime}})|^{2}d^{3}{r}^{\prime}.$
(3)
We use the Gaussian ansatz var ; jb ; YY ; Yi2000
$\phi({\bf
r},t)=\frac{\pi^{-3/4}}{w_{\rho}\sqrt{w_{z}}}\exp\left(-\frac{\rho^{2}}{2w_{\rho}^{2}}-\frac{z^{2}}{2w_{z}^{2}}+i\alpha\rho^{2}+i\beta
z^{2}\right)$ (4)
for a variational calculation, where $w_{\rho}$ and $w_{z}$ are time-dependent
radial and axial widths, and $\alpha$ and $\beta$ time-dependent phases. The
effective Lagrangian $L\equiv\int{\mathcal{L}}d^{3}{r}$ (per particle) becomes
$\displaystyle L$ $\displaystyle=$
$\displaystyle\,\left(w_{\rho}^{2}\dot{\alpha}+\frac{w_{z}^{2}\dot{\beta}}{2}\right)+\frac{\Omega_{\rho}^{2}w_{\rho}^{2}}{2}+\frac{\Omega_{z}^{2}w_{z}^{2}}{4}+\frac{1}{2{w_{\rho}^{2}}}+\frac{1}{4w_{z}^{2}}$
(5)
$\displaystyle\,+2w_{\rho}^{2}\alpha^{2}+w_{z}^{2}\beta^{2}+\frac{N}{(\sqrt{2\pi}w_{\rho}^{2}w_{z})}\left[{a}-{a_{dd}}f(\kappa)\right],$
with
$\displaystyle
f(\kappa)=\frac{1+2\kappa^{2}-3\kappa^{2}d(\kappa)}{(1-\kappa^{2})},$ (6)
$\displaystyle
d(\kappa)=\frac{\mbox{atanh}\sqrt{1-\kappa^{2}}}{\sqrt{1-\kappa^{2}}},\;\;\kappa=\frac{w_{\rho}}{w_{z}}.$
(7)
The Euler-Lagrange equations for variational parameters
$w_{\rho},w_{z},\alpha$ and $\beta$ yield the following equations for widths
$w_{\rho}$ and $w_{z}$
$\displaystyle\ddot{w}_{\rho}+\Omega_{\rho}^{2}w_{\rho}=\frac{1}{w_{\rho}^{3}}+\frac{N}{\sqrt{2\pi}}\frac{\left[2{a}-a_{dd}{g(\kappa)}\right]}{w_{\rho}^{3}w_{z}},$
(8)
$\displaystyle\ddot{w}_{z}+\Omega_{z}^{2}w_{z}=\frac{1}{w_{z}^{3}}+\frac{2N}{\sqrt{2\pi}}\frac{\left[{a}-a_{dd}c(\kappa)\right]}{w_{\rho}^{2}w_{z}^{2}},$
(9)
with
$\displaystyle
g(\kappa)=\frac{2-7\kappa^{2}-4\kappa^{4}+9\kappa^{4}d(\kappa)}{(1-\kappa^{2})^{2}},$
(10) $\displaystyle
c(\kappa)=\frac{1+10\kappa^{2}-2\kappa^{4}-9\kappa^{2}d(\kappa)}{(1-\kappa^{2})^{2}}.$
(11)
The chemical potential $\mu$ for a stationary state is
$\displaystyle\mu=$
$\displaystyle\,\frac{1}{2w_{\rho}^{2}}+\frac{1}{4w_{z}^{2}}+\frac{2N[a-a_{dd}f(\kappa)]}{\sqrt{2\pi}w_{z}w_{\rho}^{2}}+\frac{\Omega_{\rho}^{2}w_{\rho}^{2}}{2}+\frac{\Omega_{z}^{2}w_{z}^{2}}{4}.$
(12)
### II.2 1D reduction
For a cigar-shaped dipolar BEC with a strong radial trap
$(\Omega_{\rho}>\Omega_{z})$ one can write the following effective 1D equation
(details given in Appendix)
$\displaystyle i\frac{\partial\phi_{1D}(z,t)}{\partial
t}=\biggr{[}-\frac{\partial_{z}^{2}}{2}+\frac{\Omega_{z}^{2}z^{2}}{2}+\frac{2aN}{d_{\rho}^{2}}|\phi_{1D}|^{2}+\frac{2a_{dd}N}{d_{\rho}^{2}}$
$\displaystyle\times\int_{-\infty}^{\infty}\frac{dk_{z}}{2\pi}e^{ik_{z}z}\tilde{n}(k_{z})s_{1D}\left(\frac{k_{z}d_{\rho}}{\sqrt{2}}\right)\biggr{]}\phi_{1D}(z,t),$
(13)
where $s_{1D}$ is defined by equation (35) and $d_{\rho}\equiv
1/\sqrt{\Omega_{\rho}}$ is the radial harmonic oscillator length.
To solve equation (II.2), we use the Gaussian variational ansatz
$\phi_{1D}(z)=\frac{\pi^{-1/4}}{\sqrt{w_{z}}}\exp\left[-\frac{z^{2}}{2w_{z}^{2}}+i\beta
z^{2}\right].$ (14)
From equation (28) we see that the variational 1D ansatz (14) corresponds to
the following 3D wave function
$\phi({\bf
r},t)=\frac{\pi^{-3/4}}{d_{\rho}\sqrt{w_{z}}}\exp\left(-\frac{\rho^{2}}{2d_{\rho}^{2}}-\frac{z^{2}}{2w_{z}^{2}}+i\beta
z^{2}\right).$ (15)
The present variational wave function (15) is a special case of the 3D
variational wave function (4) with $w_{\rho}=d_{\rho}$ and $\alpha=0$. Hence,
the 1D variational Lagrangian can be written from the 3D Lagrangian (5),
(using $w_{\rho}=d_{\rho}$ and $\alpha=0$,) as
$\displaystyle L_{1D}$
$\displaystyle=\frac{w_{z}^{2}\dot{\beta}}{2}+\frac{1}{4w_{z}^{2}}+w_{z}^{2}\beta^{2}+\frac{\Omega_{z}^{2}w_{z}^{2}}{4}$
$\displaystyle+\frac{N}{\sqrt{2\pi}d_{\rho}^{2}w_{z}}\left[{a}-{a_{dd}}f(\kappa_{0})\right];\quad\kappa_{0}=\frac{d_{\rho}}{w_{z}},$
(16)
where we have removed the constant terms. This inductive derivation of the 1D
Lagrangian (16) avoids the construction of Lagrangian density involving error
functions in the 1D potential (36) and subsequent integration to obtain the
Lagrangian. The Euler-Lagrange equation for the variational parameter $w_{z}$
of Lagrangian (16) is
$\displaystyle\ddot{w}_{z}+\Omega_{z}^{2}w_{z}=\frac{1}{w_{z}^{3}}+\frac{2N[a-a_{dd}c(\kappa_{0})]}{\sqrt{2\pi}w_{z}^{2}d_{\rho}^{2}}.$
(17)
The variational chemical potential is given by
$\displaystyle\mu=\frac{1}{4w_{z}^{2}}+\frac{2N[a-a_{dd}f(\kappa_{0})]}{\sqrt{2\pi}w_{z}d_{\rho}^{2}}+\frac{\Omega_{z}^{2}w_{z}^{2}}{4}.$
(18)
Not only are the above variational results simple and yield a good
approximation to the 1D GP equation, much physical insight about the system
can be obtained from the variational Lagrangian (16). In a quasi-1D system,
the axial width is much larger than the transverse oscillator length:
$w_{z}\gg d_{\rho}$. Consequently, $\kappa_{0}\to 0$ and $f(\kappa_{0})\to 1$.
From equation (16), we see that the interaction term becomes in this limit
$N(a-a_{dd})/(\sqrt{2\pi}d_{\rho}^{2}w_{z})$. In equation (II.2), the dipolar
term involves a nonlocal integral. However, the variational approximation
suggests that the effect of the dipolar interaction integral is to reduce the
contact interaction term in equation (II.2) replacing the scattering length
$a$ by $(a-a_{dd})$. Immediately, one can conclude that the system effectively
becomes attractive for $a_{dd}>a$. So one can have the formation of bright
soliton even for positive (repulsive) scattering length $a$, provided that
$a_{dd}>a$.
### II.3 2D reduction
In the disk-shape, with a strong axial trap ($\Omega_{z}>\Omega_{\rho}$), the
dipolar BEC is assumed to be in the ground state
$\phi(z)=\exp(-z^{2}/2d_{z}^{2})/{(\pi d_{z}^{2})}^{1/4}$ of the axial trap
and the wave function $\phi({\bf r})$ can be written as fisch ; PS
$\phi({\bf r})=\frac{1}{{(\pi
d_{z}^{2})}^{1/4}}\exp\left(-\frac{z^{2}}{2d_{z}^{2}}\right)\phi_{2D}(x,y),$
(19)
where $\phi_{2D}(x,y)$ is the 2D wave function and
$d_{z}=\sqrt{1/\Omega_{z}}$. Using ansatz (19) in equation (II.1), the $z$
dependence can be integrated out to obtain the following effective 2D equation
fisch ; PS
$\displaystyle\,i\frac{\partial\phi_{2D}(\vec{\rho},t)}{\partial
t}=\biggr{[}-\frac{\nabla_{\rho}^{2}}{2}+\frac{\Omega_{\rho}^{2}\rho^{2}}{2}+\frac{4\pi
aN}{\sqrt{2\pi}d_{z}}|\phi_{2D}|^{2}+\frac{4\pi a_{dd}N}{\sqrt{2\pi}d_{z}}$
$\displaystyle\,\quad\times\int\frac{d^{2}k_{\rho}}{(2\pi)^{2}}\exp({i{\bf
k}_{\rho}\cdot{\vec{\rho}}}){\tilde{n}}({\bf
k_{\rho}})h_{2D}(\frac{k_{\rho}d_{z}}{\sqrt{2}})\biggr{]}\phi_{2D}(\vec{\rho},t),$
(20) $\displaystyle\,\tilde{n}({\bf k}_{\rho})=\int\exp\left(i{\bf
k}_{\rho}\cdot\vec{\rho}\right)|\phi_{2D}(\vec{\rho})|^{2}d\vec{\rho},$ (21)
where $h_{2D}(\xi)=2-3\sqrt{\pi}\xi e^{\xi^{2}}{\mbox{erfc}}(\xi)$ PS , ${\bf
k}_{\rho}\equiv(k_{x},k_{y})$, and the dipolar term is written in Fourier
space.
To solve equation (II.3), we use the Gaussian ansatz
$\phi_{2D}(\rho)=\frac{1}{w_{\rho}\sqrt{\pi}}\exp\left(-\frac{\rho^{2}}{2w_{\rho}^{2}}+i\alpha\rho^{2}\right).$
(22)
From equation (19) we see that the 2D wave function (22) corresponds to the
following 3D wave function
$\phi({\bf
r},t)=\frac{\pi^{-3/4}}{w_{\rho}\sqrt{d_{z}}}\exp\left(-\frac{\rho^{2}}{2w_{\rho}^{2}}-\frac{z^{2}}{2d_{z}^{2}}+i\alpha\rho^{2}\right).$
(23)
The present variational wave function (23) is a special case of the 3D
variational wave function (4) with $w_{z}=d_{z}$ and $\beta=0$. Hence, the 2D
variational Lagrangian can be written from the 3D Lagrangian (5) as
$\displaystyle L_{2D}$
$\displaystyle\,={w_{\rho}^{2}\dot{\alpha}}+\frac{w_{\rho}^{2}\Omega_{\rho}^{2}}{2}+\frac{1}{2w_{\rho}^{2}}+2w_{\rho}^{2}\alpha^{2}$
$\displaystyle\,+\frac{N}{\sqrt{2\pi}w_{\rho}^{2}d_{z}}\left[{a}-{a_{dd}}f(\bar{\kappa})\right];\quad\bar{\kappa}=\frac{w_{\rho}}{d_{z}},$
(24)
where we have removed the constant terms. The Euler-Lagrange variational
equation for width $w_{\rho}$ becomes
$\displaystyle\ddot{w}_{\rho}+{w_{\rho}\Omega_{\rho}^{2}}=\frac{1}{w_{\rho}^{3}}+\frac{N}{\sqrt{2\pi}}\frac{\left[2{a}-a_{dd}{g(\bar{\kappa})}\right]}{w_{\rho}^{3}d_{z}}.$
(25)
The chemical potential $\mu$ for a stationary state is
$\displaystyle\mu$ $\displaystyle=$
$\displaystyle\frac{1}{2w_{\rho}^{2}}+\frac{2N[a-a_{dd}f(\bar{\kappa})]}{\sqrt{2\pi}d_{z}w_{\rho}^{2}}+\frac{w_{\rho}^{2}\Omega_{\rho}^{2}}{2}.$
(26)
In a quasi-2D system, the radial width is much larger than the axial
oscillator length: $w_{\rho}\gg d_{z}$. Consequently, $\bar{\kappa}\to\infty$
and $f(\bar{\kappa})\to-2$. From equation (II.3), we see that the interaction
term becomes in this limit $N(a+2a_{dd})/(\sqrt{2\pi}w_{\rho}^{2}d_{z})$. The
variational approximation suggests that the effect of the dipolar interaction
in equation (II.3) is to increase the contact interaction term replacing $a$
by $(a+2a_{dd})$. Hence, for positive $a$, there cannot be any bright soliton
in 2D, which was found from a solution of the 2D GP equation (II.3) and
Bogoliubov theory tick . However, effectively the sign of the dipolar term in
the GP equation can be changed by rotating the external field that orients the
dipoles much faster than any other relevant time scale in the system nath . In
this fashion Nath et al. nath suggest changing the dipole interaction term by
a factor of $-1/2$, which changes the effective scattering length in the
Lagrange variational approximation to $(a-a_{dd})$, (as discussed in the 1D
case above,) leading to the formation of bright 2D solitons for $a_{dd}>a$.
These solitons were obtained by Nath et al. from a solution of the 2D GP
equation (II.3).
## III Numerical results
We solve the 1D, 2D, and 3D GP equations employing imaginary- and real-time
propagation with Crank-Nicolson method CPC ; CPC-1 . The dipolar interaction
is evaluated by fast Fourier transform jb ; YY .
We present results for 52Cr and 164Dy atoms. The 52Cr has a moderate dipole
moment with $a_{dd}=15a_{0}$ pfau ; pfau-1 ; pfau-4 , while the 164Dy atom has
a large dipole moment with $a_{dd}=130a_{0}$ dy ; dy-2 . In both cases we
present results for dipolar BEC of up to 10,000 atoms for $0<a<10$ nm and
choose the frequency $\omega$ such that the oscillator length $l=1\mu$m.
First, we present the results for density profiles obtained from a solution of
the reduced 1D and 2D equation and compare with the full 3D results. It is
known that the densities obtained from the reduced equations agree well with
the full 3D density, as the nonlinearity tends to zero and/or the trap
asymmetry is extreme luca . Hence in this study we consider a moderately small
trap asymmetry and a relatively large nonlinearity of experimental interest.
In the cigar (1D) case we consider 52Cr atoms with $a=6$ nm, and in the disk
(2D) case we consider 164Dy atoms with $a=6$ nm.
Figure 1: Linear density of a cigar-shaped 52Cr dipolar BEC of 1,000 atoms of
$a=6$ nm, with trap parameters $\Omega_{z}=1$ and (a) $\Omega_{\rho}=4,$ and
(b) $\Omega_{\rho}=9$ from a numerical (N) solution of the 3D equation (II.1)
and 1D equation (II.2), and its variational (V) result. Radial density of a
disk-shaped 164Dy dipolar BEC of 1,000 atoms of $a=6$ nm, with trap parameters
$\Omega_{\rho}=1$ and (c) $\Omega_{z}=4,$ and (d) $\Omega_{z}=9$ from a
numerical solution of the 3D equation (II.1) and 2D equation (II.3), and its
variational result.
In Figs. 1 (a) and (b), we plot results for linear density of a cigar-shaped
52Cr dipolar BEC of 1,000 atoms as calculated from the numerical solution of
the 3D equation (II.1) and the 1D equation (II.2) and its variational result
(17) for $\Omega_{z}=1$ and $\Omega_{\rho}=4$ and 9. We find, as the trap
asymmetry increases by changing $\Omega_{\rho}$ from 4 to 9, the agreement
between 3D and 1D models improves. In Figs. 1 (c) and (d), we plot results for
radial density of a disk-shaped 164Dy dipolar BEC of 1,000 atoms as calculated
from the numerical solution of the 3D equation (II.1) and the 2D equation
(II.3) and its variational approximation (25) for $\Omega_{\rho}=1$ and
$\Omega_{z}=4$ and 9. We find that, with the increase of the trap asymmetry
from $\Omega_{z}=4$ to 9, the agreement between the 3D and 2D models enhances.
In all cases the variational results of the reduced 1D and 2D equations are in
good agreement with those of the full 3D model.
After having established the appropriateness of the reduced 1D and 2D
equations in the cigar and disk shapes, it is realized that although the
numerical solution of these reduced GP equations are simpler than that of the
full 3D GP equation, they are still complicated due to the presence of the
nonlocal dipolar interaction. The variational approximation of these equations
presented here is relatively simple and could be used for approximate solution
of these equations. Now we test the variational results of the reduced 1D and
2D equations by comparing with the numericalsolution of these equations.
Figure 2: The numerical (N) and variational (V) rms size $\langle\rho\rangle$
versus scattering length $a$ of a disk-shaped dipolar BEC of 10,000 (a) 52Cr
and (b) 164Dy atoms for trap parameters $\Omega_{\rho}=1$ and $\Omega_{z}=4$
and 9 from a solution of the reduced 2D GP equation (II.3). The corresponding
chemical potential $\mu$ in these cases for (c) 52Cr and (d) 164Dy atoms.
In Figs. 2 we present the results for rms size $\langle\rho\rangle$ and
chemical potential $\mu$ of a disk-shaped 52Cr and 164Dy dipolar BEC of 10,000
atoms with the trap parameters $\Omega_{\rho}=1$ and $\Omega_{z}=4$ and 9 for
$0<a<10$ nm as calculated from numerical and variational approaches of the
reduced 2D equation (II.3). In Figs. 3 we exhibit the results for rms size
$\langle z\rangle$ and chemical potential $\mu$ of a cigar-shaped 52Cr and
164Dy dipolar BEC of 10,000 atoms with the trap parameters $\Omega_{z}=1$ and
$\Omega_{\rho}=4$ and 9 for $0<a<20$ nm as calculated from numerical and
variational approaches of the reduced 1D equation (II.2).
Figure 3: The numerical (N) and variational (V) rms length $\langle z\rangle$
versus scattering length $a$ of a cigar-shaped dipolar BEC of 10,000 (a) 52Cr
and (b) 164Dy atoms for trap parameters $\Omega_{z}=1$ and $\Omega_{\rho}=4$
and 9 from a solution of the reduced 1D GP equation (II.2). The corresponding
chemical potential $\mu$ in these cases for (c) 52Cr and (d) 164Dy atoms.
The dipolar interaction changes from strongly attractive in the extreme cigar
shape ($\Omega_{\rho}\gg\Omega_{z}$) to strongly repulsive in the extreme disk
shape ($\Omega_{\rho}\ll\Omega_{z}$) and its effect is minimum (nearly zero)
for $\Omega_{\rho}$ slightly less than $\Omega_{z}$. In Fig. 2 the dipolar
interaction is slightly attractive for $\Omega_{z}=4$ and $\Omega_{\rho}=1$.
Hence in the absence of any short-range interaction ($a=0$), the system will
collapse and no stable solution of the GP equation can be obtained. For Cr
atoms the dipolar interaction is weak, and for $a\geq 1$ nm, the short-range
repulsion for 10,000 atoms surpluses the dipolar attraction and a stable state
can be obtained for $\Omega_{z}=4$. For Dy atoms the dipolar interaction is
stronger, and a stable state can be obtained only for $a\geq 2$ nm for
$\Omega_{z}=4$. For $\Omega_{z}=9$, the dipolar interaction for both Cr and Dy
atoms are repulsive and a stable state is obtained in this case for $a>0$. In
Fig. 3 the dipolar interaction is attractive for both $\Omega_{\rho}=4$ and 9.
Hence the dipolar BEC can be stable only for scattering length $a$ greater
than a critical value. This is why the curves in this figure start above this
critical value. This critical value is larger for Dy atoms and
$\Omega_{\rho}=9$ compared to that of Cr atoms and $\Omega_{\rho}=4$ as can be
found in Fig. 3. As there is no real collapse in 1D models with cubic
nonlinearity; for confirming the collapse correctly one must solve the full 3D
GP equation.
Figure 4: The rms sizes $\langle z\rangle$ of a cigar-shaped Cr dipolar BEC
of 1,000 atoms versus time $t$ for (a) $\Omega_{\rho}=4$ and (b) 9 from a
numerical (N) and variational (V) results of the reduced 1D equation. The rms
sizes $\langle\rho\rangle$ of a disk-shaped Dy dipolar BEC of 1,000 atoms
versus time $t$ for (c) $\Omega_{z}=4$ and (d) 9. The oscillation was
initiated by jumping the scattering length $a$ from 6 nm to $6.15$ nm for Cr
($6.3$ nm for Dy) at $t=0$ from a solution of the reduced 2D equation.
Next we study, by numerical and variational solutions of the reduced 1D and 2D
equations, the dynamics of breathing oscillation of the four dipolar BEC of
cigar- and disk-shaped Cr and Dy atoms shown in Figs. 1 started by a small
change of the scattering length. This can be implemented experimentally by a
Feshbach resonance fesh . In Fig. 4 this dynamics is shown for a cigar-shaped
Cr dipolar BEC of 1,000 atoms for (a) $\Omega_{\rho}=4$ and (b)
$\Omega_{\rho}=9$ from a solution of the reduced 1D equation, and for a disk-
shaped Dy dipolar BEC of 1,000 atoms for (c) $\Omega_{z}=4$ and (d)
$\Omega_{z}=9$ from a solution of the reduced 2D equation. In these figures we
also show the results from a numerical solution of the 3D Eq. (II.1). The
agreement between the numerical and variational results is good in all cases.
We also calculated the angular frequencies of these oscillations. In case of
Cr in Figs. 4 (a) and (b), the axial frequencies are 1.75 (variational, 1D),
1.76 (numerical, 1D) and 1.63 (numerical, 3D), and in case of Dy in Figs. 4
(c) and (d), the radial frequencies are 1.93 (variational, 2D), 1.89
(numerical, 2D) and 1.76 (numerical, 3D). For quasi-linear systems, these
angular frequencies are expected to be 2 stringari . The deviation from this
value is due to the large nonlinearity of the dipolar BECs considered here.
## IV Conclusion
The usual GP equation provides a good description of statics and dynamics of a
normal nondipolar BEC. For a dipolar BEC the numerical solution of the GP
equation is a difficult task due to the nonlocal dipolar interaction. For a
cigar- and disk-shaped dipolar BEC, the reduced 1D and 2D equations provide an
alternative to the full 3D equation. Nevertheless, the solution of these
reduced equations is also challenging involving Fourier and inverse Fourier
transformations. As an alternative, we suggest a time-dependent variational
scheme for these reduced equations, not requiring any Fourier transformation.
The variational approximation of these reduced equations provides results for
stationary cigar- and disk-shaped dipolar BEC as well as for breathing
oscillation of the same in good agreement with the numerical solution of the
respective GP equations. This is illustrated for large Cr and Dy dipolar BECs
of 10,000 atoms and large atomic scattering lengths $a$ up to 20 nm. We also
study the breathing oscillation of a bright soliton of 1,000 Cr atoms using
the numerical solution of the 3D equation as well as the numerical and
variational approaches to the 1D equation. A typical dipolar BEC considered
here corresponds to a large short-range cubic nonlinearity of about $4\pi
aN\approx 1250$ for $a=10$ nm and $N=10,000$ and a large dipolar nonlinearity
of $4\pi a_{dd}\approx 865$ for Dy atoms for $a_{dd}=130a_{0}$ and $N=10,000$.
The variational approximations considered here provided good results for such
large nonlinearities and should be useful for analyzing the statics and
dynamics of realistic dipolar BECs under appropriate experimental conditions.
###### Acknowledgements.
We thank FAPESP (Brazil), CNPq (Brazil), DST (India), and CSIR (India) for
partial support.
## Appendix A 1D reduction
For a cigar-shaped dipolar BEC with a strong radial trap
($\Omega_{\rho}>\Omega_{z}$), we assume that in the radial direction the
dipolar BEC is confined in the ground state
$\displaystyle\phi({\bf\rho})=\exp(-(\rho^{2}/2d_{\rho}^{2})/(d_{\rho}\sqrt{\pi})$
(27)
of the transverse trap and the wave function $\phi({\bf r})=\phi_{1D}(z)$
$\times\phi({\bf\rho})$ can be written as SS ; deu
$\displaystyle\phi({\bf r})=\frac{1}{\sqrt{\pi
d_{\rho}^{2}}}\exp\left[-\frac{\rho^{2}}{2d_{\rho}^{2}}\right]\phi_{1D}(z);\quad\Omega_{\rho}d_{\rho}^{2}=1,$
(28)
where $d_{\rho}$ is the radial harmonic oscillator length.
The contribution of the dipole potential to energy is
$\displaystyle H_{dd}$ $\displaystyle=$ $\displaystyle\frac{N}{2}\int
d^{3}r\int d^{3}r^{\prime}n({\bf r})U_{dd}({\bf r-r^{\prime}})n({\bf
r^{\prime}}),$ (29) $\displaystyle=$
$\displaystyle\frac{1}{2}\frac{N}{(2\pi)^{3}}\int d^{3}k\tilde{n}({\bf
k})\tilde{U}_{dd}({\bf k})\tilde{n}(-{\bf k}),$ (30)
where $n({\bf r})\equiv|\phi({\bf r})|^{2}$ is the density and in Eq. (30) we
used a convolution of the respective variables to Fourier space and where
tilde denotes Fourier transformations: jb ; YY ; fisch ; PS
$\displaystyle\tilde{U}_{dd}({\bf k})$
$\displaystyle=\frac{4\pi}{3}3a_{dd}\biggr{[}\frac{3k_{z}^{2}}{k^{2}}-1\biggr{]},$
(31) $\displaystyle\tilde{n}({\bf k})$
$\displaystyle=\exp\left[-\frac{k_{\rho}^{2}d_{\rho}^{2}}{4}\right]\tilde{n}_{1D}(k_{z}).$
(32)
The $k_{x},k_{y}$ integrals in (30) can now be done and
$\displaystyle H_{dd}$ $\displaystyle=$ $\displaystyle\frac{4\pi
N}{3}\frac{3a_{dd}}{2}\frac{1}{2\pi}\int_{-\infty}^{\infty}dk_{z}\tilde{n}_{1D}(k_{z})\tilde{n}_{1D}(-k_{z})\frac{1}{(2\pi)^{2}}$
(33) $\displaystyle\times$
$\displaystyle\int_{-\infty}^{\infty}\int_{-\infty}^{\infty}dk_{x}dk_{y}\left[\frac{3k_{z}^{2}}{k_{\rho}^{2}+k_{z}^{2}}-1\right]\exp\left[-\frac{k_{\rho}^{2}d_{\rho}^{2}}{2}\right],$
$\displaystyle\equiv$
$\displaystyle\frac{N}{2}\frac{1}{2\pi}\int_{-\infty}^{\infty}dk_{z}\tilde{n}_{1D}(k_{z})\tilde{n}_{1D}(-k_{z})V_{1D}(k_{z}),$
where the 1D potential in Fourier space is
$\displaystyle V_{1D}(k_{z})=$
$\displaystyle\,{2a_{dd}}\int_{0}^{\infty}dk_{\rho}k_{\rho}\left[\frac{3k_{z}^{2}}{k_{\rho}^{2}+k_{z}^{2}}-1\right]\exp\left[-\frac{k_{\rho}^{2}d_{\rho}^{2}}{2}\right],$
$\displaystyle\equiv$
$\displaystyle\,\frac{2a_{dd}}{d_{\rho}^{2}}s_{1D}(\frac{k_{z}d_{\rho}}{\sqrt{2}}),$
(34) $\displaystyle
s_{1D}(\zeta)=\int_{0}^{\infty}du\left[\frac{3\zeta^{2}}{u+\zeta^{2}}-1\right]e^{-u}.$
(35)
The 1D potential in configuration space is
$\displaystyle
U_{dd}^{1D}(Z)=\frac{1}{2\pi}\int_{-\infty}^{\infty}dk_{z}e^{ik_{z}z}V_{1D}(k_{z})$
$\displaystyle=\frac{6a_{dd}}{(\sqrt{2}d_{\rho})^{3}}\left[\frac{4}{3}\delta(\sqrt{t})+2\sqrt{t}-\sqrt{\pi}(1+2t)e^{t}{\mbox{erfc}}(\sqrt{t})\right],$
(36)
where $t=[Z/(\sqrt{2}d_{\rho})]^{2},Z=|z-z^{\prime}|$. Similar, but not
identical, 1D reduced potential was derived in SS ; deu , where the
$\delta$-function term was absent. To derive the effective 1D equation for the
cigar-shaped dipolar BEC, we substitute the ansatz (28) in Eq. (II.1),
multiply by the ground-state wave function $\phi(\rho)$ and integrate in
$\rho$ to get the 1D equation
$\displaystyle i\frac{\partial\phi_{1D}(z,t)}{\partial
t}=\biggr{[}-\frac{\partial_{z}^{2}}{2}+\frac{\Omega_{z}^{2}z^{2}}{2}+\frac{2aN}{d_{\rho}^{2}}|\phi_{1D}|^{2}$
$\displaystyle+\frac{2a_{dd}N}{d_{\rho}^{2}}\int_{-\infty}^{\infty}\frac{dk_{z}}{2\pi}e^{ik_{z}z}\tilde{n}(k_{z})s_{1D}(\frac{k_{z}d_{\rho}}{\sqrt{2}})\biggr{]}\phi_{1D}(z,t),$
(37)
$\displaystyle\equiv\biggr{[}-\frac{\partial_{z}^{2}}{2}+\frac{\Omega_{z}^{2}z^{2}}{2}+\frac{2aN|\phi_{1D}|^{2}}{d_{\rho}^{2}}$
$\displaystyle+N\int_{-\infty}^{\infty}U_{dd}^{1D}(Z)|\phi_{1D}({z^{\prime}},t)|^{2}d{z^{\prime}}\biggr{]}\phi_{1D}(z,t).$
(38)
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|
arxiv-papers
| 2011-11-09T16:56:46 |
2024-09-04T02:49:24.163958
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "P. Muruganandam and S. K. Adhikari",
"submitter": "Paulsamy Muruganandam",
"url": "https://arxiv.org/abs/1111.2272"
}
|
1111.2357
|
EUROPEAN ORGANIZATION FOR NUCLEAR RESEARCH (CERN)
CERN-PH-EP-2011-172 LHCb-PAPER-2011-018
Measurement of $b$ hadron production fractions in 7 TeV $pp$ collisions
The LHCb Collaboration 111Authors are listed on the following pages.
Measurements of $b$ hadron production ratios in proton-proton collisions at a
centre-of-mass energy of 7 TeV with an integrated luminosity of 3 pb-1 are
presented. We study the ratios of strange $B$ meson to light $B$ meson
production $f_{s}/(f_{u}+f_{d})$ and $\mathchar 28931\relax_{b}^{0}$ baryon to
light $B$ meson production $f_{\Lambda_{b}}/(f_{u}+f_{d})$ as a function of
the charmed hadron-muon pair transverse momentum $p_{\rm T}$ and the $b$
hadron pseudorapidity $\eta$, for $p_{\rm T}$ between 0 and 14 GeV and $\eta$
between 2 and 5. We find that $f_{s}/(f_{u}+f_{d})$ is consistent with being
independent of $p_{\rm T}$ and $\eta$, and we determine $f_{s}/(f_{u}+f_{d})$
= 0.134$\pm$0.004${}^{+0.011}_{-0.010}$, where the first error is statistical
and the second systematic. The corresponding ratio
$f_{\Lambda_{b}}/(f_{u}+f_{d})$ is found to be dependent upon the transverse
momentum of the charmed hadron-muon pair,
$f_{\Lambda_{b}}/(f_{u}+f_{d})=(0.404\pm 0.017{\rm(stat)}\pm
0.027{\rm(syst)}\pm 0.105{\rm(Br)})\times[1-(0.031\pm 0.004{\rm(stat)}\pm
0.003{\rm(syst)})\times p_{\rm T}{\rm(GeV)}]$, where Br reflects an absolute
scale uncertainty due to the poorly known branching fraction ${\cal
B}(\mathchar 28931\relax_{c}^{+}\rightarrow pK^{-}\pi^{+})$. We extract the
ratio of strange $B$ meson to light neutral $B$ meson production $f_{s}/f_{d}$
by averaging the result reported here with two previous measurements derived
from the relative abundances of $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{+}_{s}\pi^{-}$ to
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow D^{+}K^{-}$ and
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow D^{+}\pi^{-}$. We
obtain $f_{s}/f_{d}=0.267^{+0.021}_{-0.020}$.
R. Aaij23, C. Abellan Beteta35,n, B. Adeva36, M. Adinolfi42, C. Adrover6, A.
Affolder48, Z. Ajaltouni5, J. Albrecht37, F. Alessio37, M. Alexander47, G.
Alkhazov29, P. Alvarez Cartelle36, A.A. Alves Jr22, S. Amato2, Y. Amhis38, J.
Anderson39, R.B. Appleby50, O. Aquines Gutierrez10, F. Archilli18,37, L.
Arrabito53, A. Artamonov 34, M. Artuso52,37, E. Aslanides6, G. Auriemma22,m,
S. Bachmann11, J.J. Back44, D.S. Bailey50, V. Balagura30,37, W. Baldini16,
R.J. Barlow50, C. Barschel37, S. Barsuk7, W. Barter43, A. Bates47, C. Bauer10,
Th. Bauer23, A. Bay38, I. Bediaga1, K. Belous34, I. Belyaev30,37, E. Ben-
Haim8, M. Benayoun8, G. Bencivenni18, S. Benson46, J. Benton42, R. Bernet39,
M.-O. Bettler17, M. van Beuzekom23, A. Bien11, S. Bifani12, A. Bizzeti17,h,
P.M. Bjørnstad50, T. Blake37, F. Blanc38, C. Blanks49, J. Blouw11, S. Blusk52,
A. Bobrov33, V. Bocci22, A. Bondar33, N. Bondar29, W. Bonivento15, S.
Borghi47, A. Borgia52, T.J.V. Bowcock48, C. Bozzi16, T. Brambach9, J. van den
Brand24, J. Bressieux38, D. Brett50, S. Brisbane51, M. Britsch10, T.
Britton52, N.H. Brook42, H. Brown48, A. Büchler-Germann39, I. Burducea28, A.
Bursche39, J. Buytaert37, S. Cadeddu15, J.M. Caicedo Carvajal37, O. Callot7,
M. Calvi20,j, M. Calvo Gomez35,n, A. Camboni35, P. Campana18,37, A. Carbone14,
G. Carboni21,k, R. Cardinale19,i,37, A. Cardini15, L. Carson36, K. Carvalho
Akiba2, G. Casse48, M. Cattaneo37, M. Charles51, Ph. Charpentier37, N.
Chiapolini39, K. Ciba37, X. Cid Vidal36, G. Ciezarek49, P.E.L. Clarke46,37, M.
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Dzhelyadin34, S. Easo45, U. Egede49, V. Egorychev30, S. Eidelman33, D. van
Eijk23, F. Eisele11, S. Eisenhardt46, R. Ekelhof9, L. Eklund47, Ch.
Elsasser39, D. Esperante Pereira36, L. Estève43, A. Falabella16,e, E.
Fanchini20,j, C. Färber11, G. Fardell46, C. Farinelli23, S. Farry12, V.
Fave38, V. Fernandez Albor36, M. Ferro-Luzzi37, S. Filippov32, C.
Fitzpatrick46, M. Fontana10, F. Fontanelli19,i, R. Forty37, M. Frank37, C.
Frei37, M. Frosini17,f,37, S. Furcas20, A. Gallas Torreira36, D. Galli14,c, M.
Gandelman2, P. Gandini51, Y. Gao3, J-C. Garnier37, J. Garofoli52, J. Garra
Tico43, L. Garrido35, D. Gascon35, C. Gaspar37, N. Gauvin38, M. Gersabeck37,
T. Gershon44,37, Ph. Ghez4, V. Gibson43, V.V. Gligorov37, C. Göbel54, D.
Golubkov30, A. Golutvin49,30,37, A. Gomes2, H. Gordon51, M. Grabalosa
Gándara35, R. Graciani Diaz35, L.A. Granado Cardoso37, E. Graugés35, G.
Graziani17, A. Grecu28, E. Greening51, S. Gregson43, B. Gui52, E. Gushchin32,
Yu. Guz34, T. Gys37, G. Haefeli38, C. Haen37, S.C. Haines43, T. Hampson42, S.
Hansmann-Menzemer11, R. Harji49, N. Harnew51, J. Harrison50, P.F. Harrison44,
J. He7, V. Heijne23, K. Hennessy48, P. Henrard5, J.A. Hernando Morata36, E.
van Herwijnen37, E. Hicks48, K. Holubyev11, P. Hopchev4, W. Hulsbergen23, P.
Hunt51, T. Huse48, R.S. Huston12, D. Hutchcroft48, D. Hynds47, V. Iakovenko41,
P. Ilten12, J. Imong42, R. Jacobsson37, A. Jaeger11, M. Jahjah Hussein5, E.
Jans23, F. Jansen23, P. Jaton38, B. Jean-Marie7, F. Jing3, M. John51, D.
Johnson51, C.R. Jones43, B. Jost37, M. Kaballo9, S. Kandybei40, M. Karacson37,
T.M. Karbach9, J. Keaveney12, U. Kerzel37, T. Ketel24, A. Keune38, B. Khanji6,
Y.M. Kim46, M. Knecht38, P. Koppenburg23, A. Kozlinskiy23, L. Kravchuk32, K.
Kreplin11, M. Kreps44, G. Krocker11, P. Krokovny11, F. Kruse9, K.
Kruzelecki37, M. Kucharczyk20,25,37, S. Kukulak25, R. Kumar14,37, T.
Kvaratskheliya30,37, V.N. La Thi38, D. Lacarrere37, G. Lafferty50, A. Lai15,
D. Lambert46, R.W. Lambert37, E. Lanciotti37, G. Lanfranchi18, C.
Langenbruch11, T. Latham44, R. Le Gac6, J. van Leerdam23, J.-P. Lees4, R.
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Storaci23, M. Straticiuc28, U. Straumann39, N. Styles46, V.K. Subbiah37, S.
Swientek9, M. Szczekowski27, P. Szczypka38, T. Szumlak26, S. T’Jampens4, E.
Teodorescu28, F. Teubert37, E. Thomas37, J. van Tilburg11, V. Tisserand4, M.
Tobin39, S. Topp-Joergensen51, N. Torr51, E. Tournefier4,49, M.T. Tran38, A.
Tsaregorodtsev6, N. Tuning23, M. Ubeda Garcia37, A. Ukleja27, P. Urquijo52, U.
Uwer11, V. Vagnoni14, G. Valenti14, R. Vazquez Gomez35, P. Vazquez Regueiro36,
S. Vecchi16, J.J. Velthuis42, M. Veltri17,g, K. Vervink37, B. Viaud7, I.
Videau7, X. Vilasis-Cardona35,n, J. Visniakov36, A. Vollhardt39, D. Voong42,
A. Vorobyev29, H. Voss10, S. Wandernoth11, J. Wang52, D.R. Ward43, A.D.
Webber50, D. Websdale49, M. Whitehead44, D. Wiedner11, L. Wiggers23, G.
Wilkinson51, M.P. Williams44,45, M. Williams49, F.F. Wilson45, J. Wishahi9, M.
Witek25, W. Witzeling37, S.A. Wotton43, K. Wyllie37, Y. Xie46, F. Xing51, Z.
Xing52, Z. Yang3, R. Young46, O. Yushchenko34, M. Zavertyaev10,a, F. Zhang3,
L. Zhang52, W.C. Zhang12, Y. Zhang3, A. Zhelezov11, L. Zhong3, E. Zverev31, A.
Zvyagin 37.
1Centro Brasileiro de Pesquisas Físicas (CBPF), Rio de Janeiro, Brazil
2Universidade Federal do Rio de Janeiro (UFRJ), Rio de Janeiro, Brazil
3Center for High Energy Physics, Tsinghua University, Beijing, China
4LAPP, Université de Savoie, CNRS/IN2P3, Annecy-Le-Vieux, France
5Clermont Université, Université Blaise Pascal, CNRS/IN2P3, LPC, Clermont-
Ferrand, France
6CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
7LAL, Université Paris-Sud, CNRS/IN2P3, Orsay, France
8LPNHE, Université Pierre et Marie Curie, Université Paris Diderot,
CNRS/IN2P3, Paris, France
9Fakultät Physik, Technische Universität Dortmund, Dortmund, Germany
10Max-Planck-Institut für Kernphysik (MPIK), Heidelberg, Germany
11Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg, Heidelberg,
Germany
12School of Physics, University College Dublin, Dublin, Ireland
13Sezione INFN di Bari, Bari, Italy
14Sezione INFN di Bologna, Bologna, Italy
15Sezione INFN di Cagliari, Cagliari, Italy
16Sezione INFN di Ferrara, Ferrara, Italy
17Sezione INFN di Firenze, Firenze, Italy
18Laboratori Nazionali dell’INFN di Frascati, Frascati, Italy
19Sezione INFN di Genova, Genova, Italy
20Sezione INFN di Milano Bicocca, Milano, Italy
21Sezione INFN di Roma Tor Vergata, Roma, Italy
22Sezione INFN di Roma La Sapienza, Roma, Italy
23Nikhef National Institute for Subatomic Physics, Amsterdam, Netherlands
24Nikhef National Institute for Subatomic Physics and Vrije Universiteit,
Amsterdam, Netherlands
25Henryk Niewodniczanski Institute of Nuclear Physics Polish Academy of
Sciences, Cracow, Poland
26Faculty of Physics & Applied Computer Science, Cracow, Poland
27Soltan Institute for Nuclear Studies, Warsaw, Poland
28Horia Hulubei National Institute of Physics and Nuclear Engineering,
Bucharest-Magurele, Romania
29Petersburg Nuclear Physics Institute (PNPI), Gatchina, Russia
30Institute of Theoretical and Experimental Physics (ITEP), Moscow, Russia
31Institute of Nuclear Physics, Moscow State University (SINP MSU), Moscow,
Russia
32Institute for Nuclear Research of the Russian Academy of Sciences (INR RAN),
Moscow, Russia
33Budker Institute of Nuclear Physics (SB RAS) and Novosibirsk State
University, Novosibirsk, Russia
34Institute for High Energy Physics (IHEP), Protvino, Russia
35Universitat de Barcelona, Barcelona, Spain
36Universidad de Santiago de Compostela, Santiago de Compostela, Spain
37European Organization for Nuclear Research (CERN), Geneva, Switzerland
38Ecole Polytechnique Fédérale de Lausanne (EPFL), Lausanne, Switzerland
39Physik-Institut, Universität Zürich, Zürich, Switzerland
40NSC Kharkiv Institute of Physics and Technology (NSC KIPT), Kharkiv, Ukraine
41Institute for Nuclear Research of the National Academy of Sciences (KINR),
Kyiv, Ukraine
42H.H. Wills Physics Laboratory, University of Bristol, Bristol, United
Kingdom
43Cavendish Laboratory, University of Cambridge, Cambridge, United Kingdom
44Department of Physics, University of Warwick, Coventry, United Kingdom
45STFC Rutherford Appleton Laboratory, Didcot, United Kingdom
46School of Physics and Astronomy, University of Edinburgh, Edinburgh, United
Kingdom
47School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom
48Oliver Lodge Laboratory, University of Liverpool, Liverpool, United Kingdom
49Imperial College London, London, United Kingdom
50School of Physics and Astronomy, University of Manchester, Manchester,
United Kingdom
51Department of Physics, University of Oxford, Oxford, United Kingdom
52Syracuse University, Syracuse, NY, United States
53CC-IN2P3, CNRS/IN2P3, Lyon-Villeurbanne, France, associated member
54Pontifícia Universidade Católica do Rio de Janeiro (PUC-Rio), Rio de
Janeiro, Brazil, associated to 2
aP.N. Lebedev Physical Institute, Russian Academy of Science (LPI RAS),
Moscow, Russia
bUniversità di Bari, Bari, Italy
cUniversità di Bologna, Bologna, Italy
dUniversità di Cagliari, Cagliari, Italy
eUniversità di Ferrara, Ferrara, Italy
fUniversità di Firenze, Firenze, Italy
gUniversità di Urbino, Urbino, Italy
hUniversità di Modena e Reggio Emilia, Modena, Italy
iUniversità di Genova, Genova, Italy
jUniversità di Milano Bicocca, Milano, Italy
kUniversità di Roma Tor Vergata, Roma, Italy
lUniversità di Roma La Sapienza, Roma, Italy
mUniversità della Basilicata, Potenza, Italy
nLIFAELS, La Salle, Universitat Ramon Llull, Barcelona, Spain
oHanoi University of Science, Hanoi, Viet Nam
## 1 Introduction
The fragmentation process, in which a primary $b$ quark forms either a
$b\bar{q}$ meson or a $bq_{1}q_{2}$ baryon, cannot be reliably predicted
because it is driven by strong dynamics in the non-perturbative regime. Thus
fragmentation functions for the various hadron species must be determined
experimentally. The LHCb experiment at the LHC explores a unique kinematic
region: it detects $b$ hadrons produced in a cone centered around the beam
axis covering a region of pseudorapidity $\eta$, defined in terms of the polar
angle $\theta$ with respect to the beam direction as $-\ln(\tan{\theta/2})$,
ranging approximately between 2 and 5. Knowledge of the fragmentation
functions allows us to relate theoretical predictions of the $b\bar{b}$ quark
production cross-section, derived from perturbative QCD, to the observed
hadrons. In addition, since many absolute branching fractions of $B^{-}$ and
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ decays have been well
measured at $e^{+}e^{-}$ colliders [1], it suffices to measure the ratio of
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ production to either
$B^{-}$ or $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ production to
perform precise absolute $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ branching fraction
measurements. In this paper we describe measurements of two ratios of
fragmentation functions: $f_{s}/(f_{u}+f_{d})$ and
$f_{\Lambda_{b}}/(f_{u}+f_{d})$, where $f_{q}\equiv{\cal B}(b\rightarrow
B_{q})$ and $f_{\Lambda_{b}}\equiv{\cal B}(b\rightarrow\Lambda_{b})$. The
inclusion of charged conjugate modes is implied throughout the paper, and we
measure the average production ratios.
Previous measurements of these fractions have been made at LEP [2] and at CDF
[3]. More recently, LHCb measured the ratio $f_{s}/f_{d}$ using the decay
modes $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow
D^{+}\pi^{-}$, $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow
D^{+}K^{-}$, and $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{+}_{s}\pi^{-}$
[4] and theoretical input from QCD factorization [5, 6]. Here we measure this
ratio using semileptonic decays without any significant model dependence. A
commonly adopted assumption is that the fractions of these different species
should be the same in high energy $b$ jets originating from $Z^{0}$ decays and
high $p_{\rm T}$ $b$ jets originating from $p\bar{p}$ collisions at the
Tevatron or $pp$ collisions at LHC, based on the notion that hadronization is
a non-perturbative process occurring at the scale of $\Lambda_{\rm QCD}$.
Nonetheless, the results from different experiments are discrepant in the case
of the $b$ baryon fraction [2].
The measurements reported in this paper are performed using the LHCb detector
[7], a forward spectrometer designed to study production and decays of hadrons
containing $b$ or $c$ quarks. LHCb includes a vertex detector (VELO),
providing precise locations of primary $pp$ interaction vertices, and of
detached vertices of long lived hadrons. The momenta of charged particles are
determined using information from the VELO together with the rest of the
tracking system, composed of a large area silicon tracker located before a 4
Tm dipole magnet, and a combination of silicon strip and straw drift chamber
detectors located after the magnet. Two Ring Imaging Cherenkov (RICH)
detectors are used for charged hadron identification. Photon detection and
electron identification are implemented through an electromagnetic calorimeter
followed by a hadron calorimeter. A system of alternating layers of iron and
chambers provides muon identification. The two calorimeters and the muon
system provide the energy and momentum information to implement a first level
(L0) hardware trigger. An additional trigger level is software based, and its
algorithms are tuned to the experiment’s operating condition.
In this analysis we use a data sample of 3 pb-1 collected from 7 TeV centre-
of-mass energy $pp$ collisions at the LHC during 2010. The trigger selects
events where a single muon is detected without biasing the impact parameter
distribution of the decay products of the $b$ hadron, nor any kinematic
variable relevant to semileptonic decays. These features reduce the systematic
uncertainty in the efficiency. Our goal is to measure two specific production
ratios: that of $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$
relative to the sum of $B^{-}$ and $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$, and that of $\mathchar
28931\relax_{b}^{0}$, relative to the sum of $B^{-}$ and $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$. The sum of the $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$, $B^{-}$, $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ and $\mathchar
28931\relax_{b}^{0}$ fractions does not equal one, as there is other $b$
production, namely a very small rate for $B_{c}^{-}$ mesons, bottomonia, and
other $b$ baryons that do not decay strongly into $\mathchar
28931\relax_{b}^{0}$, such as the $\Xi_{b}$. We measure relative fractions by
studying the final states $D^{0}\mu^{-}\overline{\nu}X$,
$D^{+}\mu^{-}\overline{\nu}X$, $D^{+}_{s}\mu^{-}\overline{\nu}X$, $\mathchar
28931\relax_{c}^{+}\mu^{-}\overline{\nu}X$,
$D^{0}K^{+}\mu^{-}\overline{\nu}X$, and $D^{0}p\mu^{-}\overline{\nu}X$. We do
not attempt to separate $f_{u}$ and $f_{d}$, but we measure the sum of $D^{0}$
and $D^{+}$ channels and correct for cross-feeds from $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ and $\mathchar
28931\relax_{b}^{0}$ decays. We assume near equality of the semileptonic decay
width of all $b$ hadrons, as discussed below. Charmed hadrons are
reconstructed through the modes listed in Table 1, together with their
branching fractions. We use all $D^{+}_{s}\rightarrow K^{-}K^{+}\pi^{+}$
decays rather than a combination of the resonant $\phi\pi^{+}$ and
$\overline{K}^{*0}K^{+}$ contributions, because these $D^{+}_{s}$ decays
cannot be cleanly isolated due to interference effects of different
amplitudes.
Table 1: Charmed hadron decay modes and branching fractions. Particle | Final state | Branching fraction (%)
---|---|---
$D^{0}$ | $K^{-}\pi^{+}$ | 3.89$\pm$0.05 [1]
$D^{+}$ | $K^{-}\pi^{+}\pi^{+}$ | 9.14$\pm$0.20 [19]
$D_{s}^{+}$ | $K^{-}K^{+}\pi^{+}$ | 5.50$\pm$0.27 [20]
$\Lambda_{c}^{+}$ | $pK^{-}\pi^{+}$ | 5.0$\pm$1.3 [1]
Each of these different charmed hadron plus muon final states can be populated
by a combination of initial $b$ hadron states. $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ mesons decay semileptonically into
a mixture of $D^{0}$ and $D^{+}$ mesons, while $B^{-}$ mesons decay
predominantly into $D^{0}$ mesons with a smaller admixture of $D^{+}$ mesons.
Both include a tiny component of $D^{+}_{s}K$ meson pairs. $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ mesons decay predominantly
into $D_{s}^{+}$ mesons, but can also decay into $D^{0}K^{+}$ and
$D^{+}K_{S}^{0}$ mesons; this is expected if the $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ decays into a $D_{s}^{**}$
state that is heavy enough to decay into a $DK$ pair. In this paper we measure
this contribution using $D^{0}K^{+}X\mu^{-}\overline{\nu}$ events. Finally,
$\mathchar 28931\relax_{b}^{0}$ baryons decay mostly into $\mathchar
28931\relax_{c}^{+}$ final states. We determine other contributions using
$D^{0}pX\mu^{-}\overline{\nu}$ events. We ignore the contributions of
$b\rightarrow u$ decays that comprise approximately 1% of semileptonic $b$
hadron decays [8], and constitute a roughly equal portion of each $b$ species
in any case.
The corrected yields for $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ or
$B^{-}$ decaying into $D^{0}\mu^{-}\overline{\nu}X$ or
$D^{+}\mu^{-}\overline{\nu}X$, $n_{\rm corr}$, can be expressed in terms of
the measured yields, $n$, as
$\displaystyle n_{\rm corr}(B\rightarrow D^{0}\mu)$ $\displaystyle=$
$\displaystyle\frac{1}{{\cal{B}}(D^{0}\rightarrow
K^{-}\pi^{+})\epsilon(B\rightarrow D^{0})}\times$
$\displaystyle\left[n(D^{0}\mu)-n(D^{0}K^{+}\mu)\frac{\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{0})}{\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{0}K^{+})}-n(D^{0}p\mu)\frac{\epsilon(\mathchar
28931\relax_{b}^{0}\rightarrow D^{0})}{\epsilon(\mathchar
28931\relax_{b}^{0}\rightarrow D^{0}p)}\right],$
where we use the shorthand $n(D\mu)\equiv n(DX\mu^{-}\overline{\nu})$. An
analogous abbreviation $\epsilon$ is used for the total trigger and detection
efficiencies. For example, the ratio $\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{0})/\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+})$ gives
the relative efficiency to reconstruct a charged $K$ in semi-muonic $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ decays producing a $D^{0}$
meson. Similarly
$\displaystyle n_{\rm corr}(B\rightarrow D^{+}\mu)$ $\displaystyle=$
$\displaystyle\frac{1}{\epsilon(B\rightarrow
D^{+})}\left[\frac{n(D^{+}\mu^{-})}{{\cal{B}}(D^{+}\rightarrow
K^{-}\pi^{+}\pi^{+})}-\right.$ (2)
$\displaystyle\left.\frac{n(D^{0}K^{+}\mu^{-})}{{\cal{B}}(D^{0}\rightarrow
K^{-}\pi^{+})}\frac{\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{+})}{\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+})}\right.$
$\displaystyle\left.-\frac{n(D^{0}p\mu^{-})}{{\cal{B}}(D^{0}\rightarrow
K^{-}\pi^{+})}\frac{\epsilon(\mathchar 28931\relax_{b}\rightarrow
D^{+})}{\epsilon(\mathchar 28931\relax_{b}\rightarrow D^{0}p)}\right].$
Both the $D^{0}X\mu^{-}\overline{\nu}$ and the $D^{+}X\mu^{-}\overline{\nu}$
final states contain small components of cross-feed from $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ decays to
$D^{0}K^{+}X\mu^{-}\overline{\nu}$ and to $D^{+}K^{0}X\mu^{-}\overline{\nu}$.
These components are accounted for by the two decays $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D_{s1}^{+}X\mu^{-}\overline{\nu}$ and $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D_{s2}^{*+}X\mu^{-}\overline{\nu}$ as reported in a recent LHCb publication
[9]. The third terms in Eqs. 1 and 2 are due to a similar small cross-feed
from $\mathchar 28931\relax_{b}^{0}$ decays.
The number of $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ resulting
in $D^{+}_{s}X\mu^{-}\overline{\nu}$ in the final state is given by
$\displaystyle n_{\rm corr}(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{+}_{s}\mu)$
$\displaystyle=$ $\displaystyle\frac{1}{\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{+}_{s})}\left[\frac{n(D^{+}_{s}\mu)}{{\cal B}(D^{+}_{s}\rightarrow
K^{+}K^{-}\pi^{+})}-\right.$ $\displaystyle\left.N(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}+B^{-}){\cal B}(B\rightarrow
D^{+}_{s}K\mu)\epsilon(\bar{B}\rightarrow D^{+}_{s}K\mu)\right],$
where the last term subtracts yields of $D^{+}_{s}KX\mu^{-}\overline{\nu}$
final states originating from $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ or $B^{-}$ semileptonic decays,
and $N(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}+B^{-})$ indicates the
total number of $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ and $B^{-}$
produced. We derive this correction using the branching fraction ${\cal
B}(B\rightarrow D_{s}^{(*)+}K\mu\nu)=(6.1\pm 1.2)\times 10^{-4}$ [10] measured
by the BaBar experiment. In addition, $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ decays semileptonically into
$DKX\mu^{-}\overline{\nu}$, and thus we need to add to Eq. 1
$n_{\rm corr}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
DK\mu)=2\frac{n(D^{0}K^{+}\mu)}{{\cal B}(D^{0}\rightarrow
K^{-}\pi^{+})\epsilon(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+}\mu)},$
(4)
where, using isospin symmetry, the factor of 2 accounts for $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D\kern
1.99997pt\overline{\kern-1.99997ptK}{}^{0}X\mu^{-}\overline{\nu}$ semileptonic
decays.
The equation for the ratio $f_{s}/(f_{u}+f_{d})$ is
$\frac{f_{s}}{f_{u}+f_{d}}=\frac{n_{\rm corr}(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D\mu)}{n_{\rm
corr}(B\rightarrow D^{0}\mu)+n_{\rm corr}(B\rightarrow
D^{+}\mu)}\frac{\tau_{B^{-}}+\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}}}{2\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}_{s}}}.$ (5)
where $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D\mu$
represents $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ semileptonic
decays to a final charmed hadron, given by the sum of the contributions shown
in Eqs. 1 and 4, and the symbols $\tau_{B_{i}}$ indicate the $B_{i}$ hadron
lifetimes, that are all well measured [1]. We use the average $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ lifetime, 1.472$\pm$0.025 ps
[1]. This equation assumes equality of the semileptonic widths of all the $b$
meson species. This is a reliable assumption, as corrections in HQET arise
only to order 1/$m_{b}^{2}$ and the SU(3) breaking correction is quite small,
of the order of 1% [11, 12, 13].
The $\mathchar 28931\relax_{b}^{0}$ corrected yield is derived in an analogous
manner. We determine
$n_{\rm corr}(\mathchar 28931\relax_{b}^{0}\rightarrow
D\mu)=\frac{n(\Lambda_{c}^{+}\mu^{-})}{{\cal{B}}(\Lambda_{c}^{+}\rightarrow
pK^{-}\pi^{+})\epsilon(\mathchar 28931\relax_{b}^{0}\rightarrow\mathchar
28931\relax_{c}^{+})}+2\frac{n(D^{0}p\mu^{-})}{{\cal{B}}(D^{0}\rightarrow
K^{-}\pi^{+})\epsilon(\mathchar 28931\relax_{b}^{0}\rightarrow D^{0}p)},$ (6)
where $D$ represents a generic charmed hadron, and extract the
$\Lambda_{b}^{0}$ fraction using
$\frac{f_{\Lambda_{b}}}{f_{u}+f_{d}}=\frac{n_{\rm
corr}(\Lambda_{b}^{0}\rightarrow D\mu)}{n_{\rm corr}(B\rightarrow
D^{0}\mu)+n_{\rm corr}(B\rightarrow D^{+}\mu)}\frac{\tau_{B^{-}}+\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}}}{2\tau_{\Lambda_{b}^{0}}}(1-\xi).$
(7)
Again, we assume near equality of the semileptonic widths of different $b$
hadrons, but we apply a small adjustment $\xi=4\pm$2%, to account for the
chromomagnetic correction, affecting $b$-flavoured mesons but not $b$ baryons
[11, 12, 13]. The uncertainty is evaluated with very conservative assumptions
for all the parameters of the heavy quark expansion.
## 2 Analysis method
To isolate a sample of $b$ flavoured hadrons with low backgrounds, we match
charmed hadron candidates with tracks identified as muons. Right-sign (RS)
combinations have the sign of the charge of the muon being the same as the
charge of the kaon in $D^{0}$, $D^{+}$, or $\mathchar 28931\relax_{c}^{+}$
decays, or the opposite charge of the pion in $D^{+}_{s}$ decays, while wrong-
sign (WS) combinations comprise combinations with opposite charge
correlations. WS events are useful to estimate certain backgrounds. This
analysis follows our previous investigation of $b\rightarrow
D^{0}X\mu^{-}\overline{\nu}$ [14]. We consider events where a well-identified
muon with momentum greater than 3 GeV and transverse momentum greater than 1.2
GeV is found. Charmed hadron candidates are formed from hadrons with momenta
greater than 2 GeV and transverse momenta greater than 0.3 GeV, and we require
that the average transverse momentum of the hadrons forming the candidate be
greater than 0.7 GeV. Kaons, pions, and protons are identified using the RICH
system. The impact parameter (IP), defined as the minimum distance of approach
of the track with respect to the primary vertex, is used to select tracks
coming from charm decays. We require that the $\chi^{2}$, formed by using the
hypothesis that each track’s IP is equal to 0, is greater than 9. Moreover,
the selected tracks must be consistent with coming from a common vertex: the
$\chi^{2}$ per number of degrees of freedom of the vertex fit must be smaller
than 6. In order to ensure that the charm vertex is distinct from the primary
$pp$ interaction vertex, we require that the $\chi^{2}$, based on the
hypothesis that the decay flight distance from the primary vertex is zero, is
greater than 100.
Charmed hadrons and muons are combined to form a partially reconstructed $b$
hadron by requiring that they come from a common vertex, and that the cosine
of the angle between the momentum of the charmed hadron and muon pair and the
line from the $D\mu$ vertex to the primary vertex be greater than 0.999. As
the charmed hadron is a decay product of the $b$ hadron, we require that the
difference in $z$ component of the decay vertex of the charmed hadron
candidate and that of the beauty candidate be greater than 0. We explicitly
require that the $\eta$ of the $b$ hadron candidate be between 2 and 5. We
measure $\eta$ using the line defined by connecting the primary event vertex
and the vertex formed by the $D$ and the $\mu$. Finally, the invariant mass of
the charmed hadron and muon system must be between 3 and 5 GeV for
$D^{0}\mu^{-}$ and $D^{+}\mu^{-}$ candidates, between 3.1 and 5.1 GeV for
$D^{+}_{s}\mu^{-}$ candidates, and between 3.3 and 5.3 GeV for $\mathchar
28931\relax_{c}^{+}\mu^{-}$ candidates.
We perform our analysis in a grid of 3 $\eta$ and 5 $p_{\rm T}$ bins, covering
the range $2<\eta<5$ and $p_{\rm T}\leq 14$ GeV. The $b$ hadron signal is
separated from various sources of background by studying the two dimensional
distribution of charmed hadron candidate invariant mass and ln(IP/mm). This
approach allows us to determine the background coming from false combinations
under the charmed hadron signal mass peak directly. The study of the ln(IP/mm)
distribution allows the separation of prompt charm decay candidates from
charmed hadron daughters of $b$ hadrons [14]. We refer to these samples as
Prompt and Dfb respectively.
### 2.1 Signal extraction
We describe the method used to extract the charmed hadron-$\mu$ signal by
using the $D^{0}X\mu^{-}\overline{\nu}$ final state as an example; the same
procedure is applied to the final states $D^{+}X\mu^{-}\overline{\nu}$,
$D^{+}_{s}X\mu^{-}\overline{\nu}$, and $\mathchar
28931\relax_{c}^{+}X\mu^{-}\overline{\nu}$. We perform unbinned extended
maximum likelihood fits to the two-dimensional distributions in $K^{-}\pi^{+}$
invariant mass over a region extending $\pm$80 MeV from the $D^{0}$ mass peak,
and ln(IP/mm). The parameters of the IP distribution of the Prompt sample are
found by examining directly produced charm [14] whereas a shape derived from
simulation is used for the Dfb component.
Figure 1: The logarithm of the IP distributions for (a) RS and (c) WS $D^{0}$
candidate combinations with a muon. The dotted curves show the false $D^{0}$
background, the small red-solid curves the Prompt yields, the dashed curves
the Dfb signal, and the larger green-solid curves the total yields. The
invariant $K^{-}\pi^{+}$ mass spectra for (b) RS combinations and (d) WS
combinations are also shown.
An example fit for $D^{0}\mu^{-}\overline{\nu}X$, using the whole $p_{\rm T}$
and $\eta$ range, is shown in Fig. 1. The fitted yields for RS are
27666$\pm$187 Dfb, 695$\pm$43 Prompt, and 1492$\pm$30 false $D^{0}$
combinations, inferred from the fitted yields in the sideband mass regions,
spanning the intervals between 35 and 75 MeV from the signal peak on both
sides. For WS we find 362$\pm$39 Dfb, 187$\pm$18 Prompt, and 1134$\pm$19 false
$D^{0}$ combinations. The RS yield includes a background of around 0.5% from
incorrectly identified $\mu$ candidates. As this paper focuses on ratios of
yields, we do not subtract this component. Figure 2 shows the corresponding
fits for the $D^{+}X\mu^{-}\overline{\nu}$ final state. The fitted yields
consist of 9257$\pm$110 Dfb events, 362$\pm$34 Prompt, and 1150$\pm$22 false
$D^{+}$ combinations. For WS we find 77$\pm$22 Dfb, 139$\pm$14 Prompt and
307$\pm$10 false $D^{+}$ combinations.
Figure 2: The logarithm of the IP distributions for (a) RS and (c) WS $D^{+}$
candidate combinations with a muon. The grey-dotted curves show the false
$D^{+}$ background, the small red-solid curves the Prompt yields, the blue-
dashed curves the Dfb signal, and the larger green-solid curves the total
yields. The invariant $K^{-}\pi^{+}\pi^{+}$ mass spectra for (b) RS
combinations and (d) WS combinations are also shown.
The analysis for the $D_{s}^{+}X\mu^{-}\overline{\nu}$ mode follows in the
same manner. Here, however, we are concerned about the reflection from
$\Lambda_{c}^{+}\rightarrow pK^{-}\pi^{+}$ where the proton is taken to be a
kaon, since we do not impose an explicit proton veto. Using such a veto would
lose 30% of the signal and also introduce a systematic error. We choose to
model separately this particular background. We add a probability density
function (PDF) determined from simulation to model this, and the level is
allowed to float within the estimated error on the size of the background. The
small peak near 2010 MeV in Fig. 3(b) is due to
$D^{*+}\rightarrow\pi^{+}D^{0},D^{0}\rightarrow K^{+}K^{-}$. We explicitly
include this term in the fit, assuming the shape to be the same as for the
$D^{+}_{s}$ signal, and we obtain 4$\pm$1 events in the RS signal region and
no events in the WS signal region. The measured yields in the RS sample are
2192$\pm$64 Dfb, 63$\pm$16 Prompt, 985$\pm$145 false $D^{+}_{s}$ background,
and 387$\pm$132 $\Lambda_{c}^{+}$ reflection background. The corresponding
yields in the WS sample are 13$\pm$19, 20$\pm$7, 499$\pm$16, and 3$\pm$3
respectively. Figure 3 shows the fit results.
Figure 3: The logarithm of the IP distributions for (a) RS and (c) WS
$D^{+}_{s}$ candidate combinations with a muon. The grey-dotted curves show
the false $D^{+}_{s}$ background, the small red-solid curves the Prompt
yields, the blue-dashed curves the Dfb signal, the purple dash-dotted curves
represent the background originating from $\mathchar 28931\relax_{c}^{+}$
reflection, and the larger green-solid curves the total yields. The invariant
$K^{-}K^{+}\pi^{+}$ mass spectra for RS combinations (b) and WS combinations
(d) are also shown. Figure 4: The logarithm of the IP distributions for (a)
RS and (c) WS $\Lambda_{c}^{+}$ candidate combinations with a muon. The grey-
dotted curves show the false $\Lambda_{c}^{+}$ background, the small red-solid
curves the Prompt yields, the blue-dashed curves the Dfb signal, and the
larger green-solid curves the total yields. The invariant $pK^{-}\pi^{+}$ mass
spectra for RS combinations (b) and WS combinations (d) are also shown.
The last final state considered is $\mathchar
28931\relax_{c}^{+}X\mu^{-}\overline{\nu}$. Figure 4 shows the data and fit
components to the ln(IP/mm) and $pK^{-}\pi^{+}$ invariant mass combinations
for events with $2<\eta<5$. This fit gives 3028$\pm$112 RS Dfb events,
43$\pm$17 RS Prompt events, 589$\pm$27 RS false $\mathchar
28931\relax_{c}^{+}$ combinations, 9$\pm$16 WS Dfb events, 0.5$\pm$4 WS Prompt
events, and 177$\pm$10 WS false $\mathchar 28931\relax_{c}^{+}$ combinations.
Figure 5: (a) Invariant mass of $D^{0}p$ candidates that vertex with each
other and together with a RS muon (black closed points) and for a
$\overline{p}$ (red open points) instead of a $p$; (b) fit to $D^{0}$
invariant mass for RS events with the invariant mass of $D^{0}p$ candidate in
the signal mass difference window; (c) fit to $D^{0}$ invariant mass for WS
events with the invariant mass of $D^{0}p$ candidate in the signal mass
difference window.
The $\mathchar 28931\relax_{b}^{0}$ may also decay into
$D^{0}pX\mu^{-}\overline{\nu}$. We search for these decays by requiring the
presence of a track well identified as a proton and detached from any primary
vertex. The resulting $D^{0}p$ invariant mass distribution is shown in Fig. 5.
We also show the combinations that cannot arise from $\Lambda_{b}^{0}$ decay,
namely those with $D^{0}\overline{p}$ combinations. There is a clear excess of
RS over WS combinations especially near threshold. Fits to the $K^{-}\pi^{+}$
invariant mass in the $[m(K^{-}\pi^{+}p)-m(K^{-}\pi^{+})+m(D^{0})_{\rm PDG}]$
region shown in Fig. 5(a) give 154$\pm$13 RS events and 55$\pm$8 WS events. In
this case, we use the WS yield for background subtraction, scaled by the RS/WS
background ratio determined with a MC simulation including $(B^{-}+\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow
D^{0}X\mu^{-}\overline{\nu})$ and generic $b\overline{b}$ events. This ratio
is found to be 1.4$\pm$0.2. Thus, the net signal is 76$\pm$17$\pm$11, where
the last error reflects the uncertainty in the ratio between RS and WS
background.
### 2.2 Background studies
Apart from false $D$ combinations, separated from the signal by the two-
dimensional fit described above, there are also physical background sources
that affect the RS Dfb samples, and originate from $b\overline{b}$ events,
which are studied with a MC simulation. In the meson case, the background
mainly comes from $b\rightarrow DDX$ with one of the $D$ mesons decaying semi-
muonically, and from combinations of tracks from the $pp\rightarrow b\bar{b}X$
events, where one $b$ hadron decays into a $D$ meson and the other $b$ hadron
decays semi-muonically. The background fractions are (1.9$\pm$0.3)% for
$D^{0}X\mu^{-}\overline{\nu}$, (2.5$\pm$0.6)% for
$D^{+}X\mu^{-}\overline{\nu}$, and (5.1$\pm$1.7)% for
$D_{s}^{+}X\mu^{-}\overline{\nu}$. The main background component for
$\mathchar 28931\relax_{b}^{0}$ semileptonic decays is $\mathchar
28931\relax_{b}^{0}$ decaying into $D^{-}_{s}\mathchar 28931\relax_{c}^{+}$,
and the $D^{-}_{s}$ decaying semi-muonically. Overall, we find a very small
background rate of (1.0$\pm$0.2)%, where the error reflects only the
statistical uncertainty in the simulation. We correct the candidate $b$ hadron
yields in the signal region with the predicted background fractions. A
conservative 3% systematic uncertainty in the background subtraction is
assigned to reflect modelling uncertainties.
Figure 6: Projections of the two-dimensional fit to the $q^{2}$ and
$m(D^{+}_{s}\mu)$ distributions of semileptonic decays including a $D^{+}_{s}$
meson. The $D_{s}^{*}/D_{s}$ ratio has been fixed to the measured $D^{*}/D$
ratio in light $B$ decays (2.42$\pm$0.10), and the background contribution is
obtained using the sidebands in the $K^{+}K^{-}\pi^{+}$ mass spectrum. The
different components are stacked: the background is represented by a black
dot-dashed line, $D^{+}_{s}$ by a red dashed line, $D_{s}^{*+}$ by a blue
dash-double dotted line and $D_{s}^{**+}$ by a green dash-dotted line.
### 2.3 Monte Carlo simulation and efficiency determination
In order to estimate the detection efficiency, we need some knowledge of the
different final states which contribute to the Cabibbo favoured semileptonic
width, as some of the selection criteria affect final states with distinct
masses and quantum numbers differently. Although much is known about the
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ and $B^{-}$ semileptonic
decays, information on the corresponding $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ and $\mathchar
28931\relax_{b}^{0}$ semileptonic decays is rather sparse. In particular, the
hadronic composition of the final states in $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ decays is poorly known [9],
and only a study from CDF provides some constraints on the branching ratios of
final states dominant in the corresponding $\mathchar 28931\relax_{b}^{0}$
decays [15].
In the case of the $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{+}_{s}$
semileptonic decays, we assume that the final states are $D^{+}_{s}$,
$D^{*+}_{s}$, $D_{s0}^{*}(2317)^{+}$, $D_{s1}(2460)^{+}$, and
$D_{s1}(2536)^{+}$. States above $DK$ threshold decay predominantly into
$D^{(*)}K$ final states. We model the decays to the final states
$D^{+}_{s}\mu^{-}\overline{\nu}$ and $D^{*+}_{s}\mu^{-}\overline{\nu}$ with
HQET form factors using normalization coefficients derived from studies of the
corresponding $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ and $B^{-}$
semileptonic decays [1], while we use the ISGW2 form factor model [16] to
describe final states including higher mass resonances.
In order to determine the ratio between the different hadron species in the
final state, we use the measured kinematic distributions of the quasi-
exclusive process $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{+}_{s}\mu^{-}\overline{\nu}X$. To reconstruct the squared invariant mass of
the $\mu^{-}\overline{\nu}$ pair ($q^{2}$), we exploit the measured direction
of the $b$ hadron momentum, which, together with energy and momentum
conservation, assuming no missing particles other than the neutrino, allow the
reconstruction of the $\nu$ 4-vector, up to a two-fold ambiguity, due to its
unknown orientation with respect to the $B$ flight path in its rest frame. We
choose the solution corresponding to the lowest $b$ hadron momentum. This
method works well when there are no missing particles, or when the missing
particles are soft, as in the case when the charmed system is a $D^{*}$ meson.
We then perform a two-dimensional fit to the $q^{2}$ versus $m(\mu D^{+}_{s})$
distribution. Figure 6 shows stacked histograms of the $D^{+}_{s}$,
$D^{*+}_{s}$, and $D_{s}^{**+}$ components. In the fit we constrain the ratio
${\cal B}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{*+}_{s}\mu^{-}\overline{\nu})/{\cal B}(\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{+}_{s}\mu^{-}\overline{\nu})$ to be equal to the average
$D^{*}\mu^{-}\overline{\nu}/D\mu^{-}\overline{\nu}$ ratio in semileptonic
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ and $B^{-}$ decays
(2.42$\pm$0.10) [1]. This constraint reduces the uncertainty of one $D^{**}$
fraction. We have also performed fits removing this assumption, and the
variation between the different components is used to assess the modelling
systematic uncertainty.
A similar procedure is applied to the $\mathchar 28931\relax_{c}^{+}\mu^{-}$
sample and the results are shown in Fig. 7. In this case we consider three
final states, $\mathchar 28931\relax_{c}^{+}\mu^{-}\overline{\nu}$,
${\it\Lambda_{c}}(2595)^{+}\mu^{-}\overline{\nu}$, and
${\it\Lambda_{c}}(2625)^{+}\mu^{-}\overline{\nu}$, with form factors from the
model of Ref. [17]. We constrain the two highest mass hadrons to be produced
in the ratio predicted by this theory.
Figure 7: Projections of the two-dimensional fit to the $q^{2}$ and
$m(\mathchar 28931\relax_{c}^{+}\mu^{-})$ distributions of semileptonic decays
including a $\mathchar 28931\relax_{c}^{+}$ baryon. The different components
are stacked: the dotted line represents the combinatoric background, the
bigger dashed line (red) represents the $\mathchar
28931\relax_{c}^{+}\mu^{-}\overline{\nu}$ component, the smaller dashed line
(blue) the ${\it\Lambda_{c}}(2595)^{+}$, and the solid line represents the
${\it\Lambda_{c}}(2625)^{+}$ component. The
${\it\Lambda_{c}}(2595)^{+}/{\it\Lambda_{c}}(2625)^{+}$ ratio is fixed to its
predicted value, as described in the text. Figure 8: Measured proton
identification efficiency as a function of the $\mathchar
28931\relax_{c}^{+}\mu^{-}$ $p_{\rm T}$ for $2<\eta<3$, $3<\eta<4$, $4<\eta<5$
respectively, and for the selection criteria used in the $\mathchar
28931\relax_{c}^{+}\rightarrow pK^{-}\pi^{+}$ reconstruction.
The measured pion, kaon and proton identification efficiencies are determined
using $K_{\rm S}^{0}$, $D^{*+}$, and $\Lambda^{0}$ calibration samples where
$p$, $K$, and $\pi$ are selected without utilizing the particle identification
criteria. The efficiency is obtained by fitting simultaneously the invariant
mass distributions of events either passing or failing the identification
requirements. Values are obtained in bins of the particle $\eta$ and $p_{\rm
T}$, and these efficiency matrices are applied to the MC simulation.
Alternatively, the particle identification efficiency can be determined by
using the measured efficiencies and combining them with weights proportional
to the fraction of particle types with a given $\eta$ and $p_{\rm T}$ for each
$\mu$ charmed hadron pair $\eta$ and $p_{\rm T}$ bin. The overall efficiencies
obtained with these two methods are consistent. An example of the resulting
particle identification efficiency as a function of the $\eta$ and $p_{\rm T}$
of the $\mathchar 28931\relax_{c}^{+}\mu^{-}$ pair is shown in Fig. 8.
As the functional forms of the fragmentation ratios in terms of $p_{\rm T}$
and $\eta$ are not known, we determine the efficiencies for the final states
studied as a function of $p_{\rm T}$ and $\eta$ within the LHCb acceptance.
Figure 9 shows the results.
## 3 Evaluation of the ratios ${f_{s}/(f_{u}+f_{d})}$ and
${f_{\Lambda_{b}}/(f_{u}+f_{d})}$
Perturbative QCD calculations lead us to expect the ratios
$f_{s}/(f_{u}+f_{d})$ and $f_{\Lambda_{b}}/(f_{u}+f_{d})$ to be independent of
$\eta$, while a possible dependence upon the $b$ hadron transverse momentum
$p_{\rm T}$ is not ruled out, especially for ratios involving baryon species
[18]. Thus we determine these fractions in different $p_{\rm T}$ and $\eta$
bins. For simplicity, we use the transverse momentum of the charmed
hadron-$\mu$ pair as the $p_{\rm T}$ variable, and do not try to unfold the
$b$ hadron transverse momentum.
In order to determine the corrected yields entering the ratio
$f_{s}/(f_{u}+f_{d})$, we determine yields in a matrix of three $\eta$ and
five $p_{\rm T}$ bins and divide them by the corresponding efficiencies. We
then use Eq. 5, with the measured lifetime ratio $(\tau_{B^{-}}+\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}})/2\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}_{s}}=1.07\pm 0.02$ [1] to derive
the ratio $f_{s}/(f_{u}+f_{d})$ in two $\eta$ bins. The measured ratio is
constant over the whole $\eta$-$p_{\rm T}$ domain. Figure 10 shows the
$f_{s}/(f_{u}+f_{d})$ fractions in bins of $p_{\rm T}$ in two $\eta$
intervals.
Figure 9: Efficiencies for $D^{0}\mu^{-}\overline{\nu}X$,
$D^{+}\mu^{-}\overline{\nu}X$, $D^{+}_{s}\mu^{-}\overline{\nu}X$, $\mathchar
28931\relax_{c}^{+}\mu^{-}\overline{\nu}X$ as a function of $\eta$ and $p_{\rm
T}$.
$\begin{array}[]{cc}\includegraphics[width=433.62pt]{fig10}\end{array}$
Figure 10: Ratio between $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ and light $B$ meson production fractions as a function of the transverse momentum of the $D^{+}_{s}\mu^{-}$ pair in two bins of $\eta$. The errors shown are statistical only. Table 2: Systematic uncertainties on the relative $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ production fraction. Source | Error (%)
---|---
Bin-dependent errors | 1.0
${\cal B}(D^{0}\rightarrow K^{-}\pi^{+})$ | 1.2
${\cal B}(D^{+}\rightarrow K^{-}\pi^{+}\pi^{+})$ | 1.5
${\cal B}(D^{+}_{s}\rightarrow K^{-}K^{+}\pi^{+})$ | 4.9
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ semileptonic decay modelling | 3.0
Backgrounds | 2.0
Tracking efficiency | 2.0
Lifetime ratio | 1.8
PID efficiency | 1.5
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+}X\mu^{-}\overline{\nu}$ | ${}^{+4.1}_{-1.1}$
${\cal B}((B^{-},\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0})\rightarrow D_{s}^{+}KX\mu^{-}\overline{\nu})$ | 2.0
Total | ${}^{+8.6}_{-7.7}$
Figure 11: $f_{+}/f_{0}$ as a function of $p_{\rm T}$ for $\eta$=(2,3) (a) and
$\eta$=(3,5) (b). The horizontal line shows the average value. The error shown
combines statistical and systematic uncertainties accounting for the detection
efficiency and the particle identification efficiency.
By fitting a single constant to all the data, we obtain
$f_{s}/(f_{u}+f_{d})=0.134\pm 0.004^{+0.011}_{-0.010}$ in the interval
$2<\eta<5$, where the first error is statistical and the second is systematic.
The latter includes several different sources listed in Table 2. The dominant
systematic uncertainty is caused by the experimental uncertainty on
${\cal{B}}(D_{s}^{+}\rightarrow K^{+}K^{-}\pi^{+})$ of 4.9%. Adding in the
contributions of the $D^{0}$ and $D^{+}$ branching fractions we have a
systematic error of 5.5% due to the charmed hadron branching fractions. The
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ semileptonic modelling
error is derived by changing the ratio between different hadron species in the
final state obtained by removing the SU(3) symmetry constrain, and changing
the shapes of the less well known $D^{**}$ states. The tracking efficiency
errors mostly cancel in the ratio since we are dealing only with combinations
of three or four tracks. The lifetime ratio error reflects the present
experimental accuracy [1]. We correct both for the bin-dependent PID
efficiency obtained with the procedure detailed before, accounting for the
statistical error of the calibration sample, and the overall PID efficiency
uncertainty, due to the sensitivity to the event multiplicity. The latter is
derived by taking the kaon identification efficiency obtained with the method
described before, without correcting for the different track multiplicities in
the calibration and signal samples. This is compared with the results of the
same procedure performed correcting for the ratio of multiplicities in the two
samples. The error due to $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{0}K^{+}X\mu^{-}\overline{\nu}$ is obtained by changing the RS/WS background
ratio predicted by the simulation within errors, and evaluating the
corresponding change in $f_{s}/(f_{u}+f_{d})$. Finally, the error due to
$(B^{-},\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0})\rightarrow
D_{s}^{+}KX\mu^{-}\overline{\nu}$ reflects the uncertainty in the measured
branching fraction.
Isospin symmetry implies the equality of $f_{d}$ and $f_{u}$, which allows us
to compare $f_{+}/f_{0}\equiv n_{\rm corr}(D^{+}\mu)/n_{\rm corr}(D^{0}\mu)$
with its expected value. It is not possible to decouple the two ratios for an
independent determination of $f_{u}/f_{d}$. Using all the known semileptonic
branching fractions [1], we estimate the expected relative fraction of the
$D^{+}$ and $D^{0}$ modes from $B^{+/0}$ decays to be $f_{+}/f_{0}=0.375\pm
0.023,$ where the error includes a 6% theoretical uncertainty associated to
the extrapolation of present experimental data needed to account for the
inclusive $b\rightarrow c\mu^{-}\overline{\nu}$ semileptonic rate. Our
corrected yields correspond to $f_{+}/f_{0}=0.373\pm 0.006$ (stat) $\pm$ 0.007
(eff) $\pm$ 0.014, for a total uncertainty of 4.5%. The last error accounts
for uncertainties in $B$ background modelling, in the
$D^{0}K^{+}\mu^{-}\overline{\nu}$ yield, the $D^{0}p\mu^{-}\overline{\nu}$
yield, the $D^{0}$ and $D^{+}$ branching fractions, and tracking efficiency.
The other systematic errors mostly cancel in the ratio. Our measurement of
$f_{+}/f_{0}$ is not seen to be dependent upon $p_{\rm T}$ or $\eta$, as shown
in Fig. 11, and is in agreement with expectation.
Figure 12: Fragmentation ratio $f_{\mathchar 28931\relax_{b}}/(f_{u}+f_{d})$
dependence upon $p_{\rm T}(\mathchar 28931\relax_{c}^{+}\mu^{-})$. The errors
shown are statistical only.
We follow the same procedure to derive the fraction $f_{\mathchar
28931\relax_{b}}/(f_{u}+f_{d})$, using Eq. 7 and the ratio
$(\tau_{B^{-}}+\tau_{\kern
1.25995pt\overline{\kern-1.25995ptB}{}^{0}})/(2\tau_{\mathchar
28931\relax_{b}^{0}})=1.14\pm 0.03$ [1]. In this case, we observe a $p_{\rm
T}$ dependence in the two $\eta$ intervals. Figure 12 shows the data fitted to
a straight line
$\frac{f_{\mathchar 28931\relax_{b}}}{f_{u}+f_{d}}=a[1+b\times p_{\rm T}({\rm
GeV})].$ (8)
Table 3 summarizes the fit results. A corresponding fit to a constant shows
that a $p_{\rm T}$ independent $f_{\mathchar 28931\relax_{b}}/(f_{u}+f_{d})$
is excluded at the level of four standard deviations. The systematic errors
reported in Table 3 include only the bin-dependent terms discussed above.
Table 4 summarizes all the sources of absolute scale systematic uncertainties,
that include several components. Their definitions mirror closely the
corresponding uncertainties for the $f_{s}/(f_{u}+f_{d})$ determination, and
are assessed with the same procedures. The term $\mathchar
28931\relax_{b}\rightarrow D^{0}pX\mu^{-}\overline{\nu}$ accounts for the
uncertainty in the raw $D^{0}pX\mu^{-}\overline{\nu}$ yield, and is evaluated
by changing the RS/WS background ratio (1.4$\pm$0.2) within the quoted
uncertainty. In addition, an uncertainty of 2% is associated with the
derivation of the semileptonic branching fraction ratios from the
corresponding lifetimes, labelled $\Gamma_{\rm sl}$ in Table 4. The
uncertainty is derived assigning conservative errors to the parameters
affecting the chromomagnetic operator that influences the $B$ meson total
decay widths, but not the $\mathchar 28931\relax_{b}^{0}$. By far the largest
term is the poorly known ${\cal B}(\mathchar 28931\relax_{c}^{+}\rightarrow
pK^{-}\pi^{+}$); thus it is quoted separately.
Table 3: Coefficients of the linear fit describing the $p_{\rm T}(\mathchar 28931\relax_{c}^{+}\mu^{-})$ dependence of $f_{\mathchar 28931\relax_{b}}/(f_{u}+f_{d})$. The systematic uncertainties included are only those associated with the bin-dependent MC and particle identification errors. $\eta$ range | $a$ | $b$
---|---|---
2-3 | 0.434$\pm$0.040$\pm$0.025 | -0.036$\pm$0.008$\pm$0.004
3-5 | 0.397$\pm$0.020$\pm$0.009 | -0.028$\pm$0.006$\pm$0.003
2-5 | 0.404$\pm$0.017$\pm$0.009 | -0.031$\pm$0.004$\pm$0.003
Table 4: Systematic uncertainties on the absolute scale of $f_{\Lambda_{b}}/(f_{u}+f_{d})$. Source | Error (%)
---|---
Bin dependent errors | 2.2
${\cal B}(\mathchar 28931\relax_{b}^{0}\rightarrow D^{0}pX\mu^{-}\overline{\nu})$ | 2.0
Monte Carlo modelling | 1.0
Backgrounds | 3.0
Tracking efficiency | 2.0
$\Gamma_{\rm sl}$ | 2.0
Lifetime ratio | 2.6
PID efficiency | 2.5
Subtotal | 6.3
${\cal B}(\mathchar 28931\relax_{c}^{+}\rightarrow pK^{-}\pi^{+})$ | 26.0
Total | $26.8$
In view of the observed dependence upon $p_{\rm T}$, we present our results as
$\left[\frac{f_{\mathchar 28931\relax_{b}}}{f_{u}+f_{d}}\right](p_{\rm
T})=(0.404\pm 0.017\pm 0.027\pm 0.105)\times[1-(0.031\pm 0.004\pm 0.003)\times
p_{\rm T}{\rm(GeV)}],$ (9)
where the scale factor uncertainties are statistical, systematic, and the
error on ${\cal B}(\mathchar 28931\relax_{c}\rightarrow pK^{-}\pi^{+})$
respectively. The correlation coefficient between the scale factor and the
slope parameter in the fit with the full error matrix is $-0.63$. Previous
measurements of this fraction have been made at LEP and the Tevatron [3]. LEP
obtains 0.110$\pm$0.019 [2]. This fraction has been calculated by combining
direct rate measurements with time-integrated mixing probability averaged over
an unbiased sample of semi-leptonic $b$ hadron decays. CDF measures
$f_{\mathchar 28931\relax_{b}}/(f_{u}+f_{d})=0.281\pm
0.012^{+0.011+0.128}_{-0.056-0.086}$, where the last error reflects the
uncertainty in ${\cal B}(\mathchar 28931\relax_{c}^{+}\rightarrow
pK^{-}\pi^{+})$. It has been suggested [3] that the difference between the
Tevatron and LEP results is explained by the different kinematics of the two
experiments. The average $p_{\rm T}$ of the $\mathchar
28931\relax_{c}^{+}\mu^{-}$ system is 10 GeV for CDF, while the $b$-jets, at
LEP, have $p\approx 40$ GeV. LHCb probes an even lower $b$ $p_{\rm T}$ range,
while retaining some sensitivity in the CDF kinematic region. These data are
consistent with CDF in the kinematic region covered by both experiments, and
indicate that the baryon fraction is higher in the lower $p_{\rm T}$ region.
## 4 Combined result for the production fraction $f_{s}/f_{d}$ from LHCb
From the study of $b$ hadron semileptonic decays reported above, and assuming
isospin symmetry, namely $f_{u}=f_{d}$, we obtain
$\left(\frac{f_{s}}{f_{d}}\right)_{\rm sl}=0.268\pm 0.008({\rm
stat})^{+0.022}_{-0.020}({\rm syst}),$
where the first error is statistical and the second is systematic.
Measurements of this quantity have also been made by LHCb by using hadronic
$B$ meson decays [4]. The ratio determined using the relative abundances of
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{+}_{s}\pi^{-}$ to $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow D^{+}K^{-}$ is
$\left(\frac{f_{s}}{f_{d}}\right)_{\rm h1}=0.250\pm 0.024(\rm stat)\pm
0.017({\rm syst})\pm 0.017({\rm theor}),$
while that from the relative abundances of $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{+}_{s}\pi^{-}$ to
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow D^{+}\pi^{-}$ [4]
is
$\left(\frac{f_{s}}{f_{d}}\right)_{\rm h2}=0.256\pm 0.014({\rm stat})\pm
0.019({\rm syst})\pm 0.026({\rm theor}).$
The first uncertainty is statistical, the second systematic and the third
theoretical. The theoretical uncertainties in both cases include non-
factorizable SU(3)-breaking effects and form factor ratio uncertainties. The
second ratio is affected by an additional source, accounting for the
$W$-exchange diagram in the $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}\rightarrow D^{+}\pi^{-}$ decay.
In order to average these results, we consider the correlations between
different sources of systematic uncertainties, as shown in Table 5. We then
utilise a generator of pseudo-experiments, where each independent source of
uncertainty is generated as a random variable with Gaussian distribution,
except for the component $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow
D^{0}K^{+}\mu^{-}\overline{\nu}_{\mu}X$, which is modeled with a bifurcated
Gaussian with standard deviations equal to the positive and negative errors
shown in Table 5. This approach to the averaging procedure is motivated by the
goal of proper treatment of asymmetric errors [21]. We assume that the
theoretical errors have a Gaussian distribution.
Table 5: Summary of the systematic and theoretical uncertainties in the three LHCb measurements of $f_{s}/f_{d}$. Source | Error (%) |
---|---|---
| $(f_{s}/f_{d})_{\rm sl}$ | $(f_{s}/f_{d})_{\rm h1}$ | $(f_{s}/f_{d})_{\rm h2}$ |
Bin dependent error | 1.0 | - | - | Uncorrelated
Semileptonic decay modelling | 3.0 | - | - | Uncorrelated
Backgrounds | 2.0 | - | - | Uncorrelated
Fit model | - | 2.8 | 2.8 | Uncorrelated
Trigger simulation | - | 2.0 | 2.0 | Uncorrelated
Tracking efficiency | 2.0 | - | - | Uncorrelated
${\cal B}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+}X\mu^{-}\overline{\nu})$ | ${}^{+4.1}_{-1.1}$ | - | - | Uncorrelated
${\cal B}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}/B^{-}\rightarrow D^{+}_{s}KX\mu^{-}\overline{\nu})$ | 2.0 | - | - | Uncorrelated
Particle identification calibration | 1.5 | 1.0 | 2.5 | Correlated
$B$ lifetimes | 1.5 | 1.5 | 1.5 | Correlated
${\cal B}(D^{+}_{s}\rightarrow K^{+}K^{-}\pi^{+})$ | 4.9 | 4.9 | 4.9 | Correlated
${\cal B}(D^{+}\rightarrow K^{-}\pi^{+}\pi^{-})$ | 1.5 | 1.5 | 1.5 | Correlated
SU(3) and form factors | - | 6.1 | 6.1 | Correlated
$W$-exchange | - | - | 7.8 | Uncorrelated
We define the average fraction as
$f_{s}/f_{d}=\alpha_{1}(f_{s}/f_{d})_{\rm sl}+\alpha_{2}(f_{s}/f_{d})_{\rm
h1}+\alpha_{3}(f_{s}/f_{d})_{\rm h2},$ (10)
where
$\alpha_{1}+\alpha_{2}+\alpha_{3}=1.$ (11)
The RMS value of $f_{s}/f_{d}$ is then evaluated as a function of $\alpha_{1}$
and $\alpha_{2}$.
We derive the most probable value $f_{s}/f_{d}$ by determining the
coefficients $\alpha_{i}$ at which the RMS is minimum, and the total errors by
computing the boundaries defining the 68% CL, scanning from top to bottom
along the axes $\alpha_{1}$ and $\alpha_{2}$ in the range comprised between 0
and 1. The optimal weights determined with this procedure are
$\alpha_{1}=0.73$, and $\alpha_{2}=0.14$, corresponding to the most probable
value
$f_{s}/f_{d}=0.267^{+0.021}_{-0.020}.$
The most probable value differs slightly from a simple weighted average of the
three measurements because of the asymmetry of the error distribution in the
semileptonic determination. By switching off different components we can
assess the contribution of each source of uncertainty. Table 6 summarizes the
results.
Table 6: Uncertainties in the combined value of $f_{s}/f_{d}$. Source | Error (%)
---|---
Statistical | 2.8
Experimental systematic (symmetric) | 3.3
${\cal B}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}\rightarrow D^{0}K^{+}X\mu^{-}\overline{\nu})$ | ${}^{+3.0}_{-0.8}$
${\cal B}(D^{+}\rightarrow K^{-}\pi^{+}\pi^{-})$ | 2.2
${\cal B}(D^{+}_{s}\rightarrow K^{+}K^{-}\pi^{+})$ | 4.9
$B$ lifetimes | 1.5
${\cal B}(\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}/B^{-}\rightarrow D^{+}_{s}KX\mu^{-}\overline{\nu})$ | 1.5
Theory | 1.9
## 5 Conclusions
We measure the ratio of the $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ production fraction to the sum
of those for $B^{-}$ and $\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}$
mesons $f_{s}/(f_{u}+f_{d})$ = 0.134$\pm$0.004${}^{+0.011}_{-0.010}$, and find
it consistent with being independent of $\eta$ and $p_{\rm T}$. Our results
are more precise than, and in agreement with, previous measurements in
different kinematic regions. We combine the LHCb measurements of the ratio of
$\kern 1.79993pt\overline{\kern-1.79993ptB}{}^{0}_{s}$ to $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ production fractions obtained
using $b$ hadron semileptonic decays, and two different ratios of branching
fraction of exclusive hadronic decays to derive
$f_{s}/f_{d}=0.267^{+0.021}_{-0.020}$. The ratio of the $\Lambda_{b}^{0}$
baryon production fraction to the sum of those for $B^{-}$ and $\kern
1.79993pt\overline{\kern-1.79993ptB}{}^{0}$ mesons varies with the $p_{\rm T}$
of the charmed hadron muon pair. Assuming a linear dependence up to $p_{\rm
T}=14$ GeV, we obtain
$\frac{f_{\mathchar 28931\relax_{b}}}{f_{u}+f_{d}}=(0.404\pm 0.017\pm 0.027\pm
0.105)\times[1-(0.031\pm 0.004\pm 0.003)\times p_{\rm T}{\rm(GeV)}],$ (12)
where the errors on the absolute scale are statistical, systematic and error
on ${\cal B}(\mathchar 28931\relax_{c}^{+}\rightarrow pK^{-}\pi^{+})$
respectively. No $\eta$ dependence is found.
## Acknowledgements
We express our gratitude to our colleagues in the CERN accelerator departments
for the excellent performance of the LHC. We thank the technical and
administrative staff at CERN and at the LHCb institutes, and acknowledge
support from the National Agencies: CAPES, CNPq, FAPERJ and FINEP (Brazil);
CERN; NSFC (China); CNRS/IN2P3 (France); BMBF, DFG, HGF and MPG (Germany); SFI
(Ireland); INFN (Italy); FOM and NWO (the Netherlands); SCSR (Poland); ANCS
(Romania); MinES of Russia and Rosatom (Russia); MICINN, XuntaGal and GENCAT
(Spain); SNSF and SER (Switzerland); NAS Ukraine (Ukraine); STFC (United
Kingdom); NSF (USA). We also acknowledge the support received from the ERC
under FP7 and the Region Auvergne.
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|
arxiv-papers
| 2011-11-09T21:57:41 |
2024-09-04T02:49:24.173970
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "The LHCb Collaboration: R. Aaij, C. Abellan Beteta, B. Adeva, M.\n Adinolfi, C. Adrover, A. Affolder, Z. Ajaltouni, J. Albrecht, F. Alessio, M.\n Alexander, G. Alkhazov, P. Alvarez Cartelle, A. A. Alves Jr, S. Amato, Y.\n Amhis, J. Anderson, R. B. Appleby, O. Aquines Gutierrez, F. Archilli, L.\n Arrabito, A. Artamonov, M. Artuso, E. Aslanides, G. Auriemma, S. Bachmann, J.\n J. Back, D. S. Bailey, V. Balagura, W. Baldini, R. J. Barlow, C. Barschel, S.\n Barsuk, W. Barter, A. Bates, C. Bauer, Th. Bauer, A. Bay, I. Bediaga, S.\n Belogurov, K. Belous, I. Belyaev, E. Ben-Haim, M. Benayoun, G. Bencivenni, S.\n Benson, J. Benton, R. Bernet, M.-O. Bettler, M. van Beuzekom, A. Bien, S.\n Bifani, A. Bizzeti, P.M. Bj{\\o}rnstad, T. Blake, F. Blanc, C. Blanks, J.\n Blouw, S. Blusk, A. Bobrov, V. Bocci, A. Bondar, N. Bondar, W. Bonivento, S.\n Borghi, A. Borgia, T. J. V. Bowcock, C. Bozzi, T. Brambach, J. van den Brand,\n J. Bressieux, D. Brett, S. Brisbane, M. Britsch, T. Britton, N. H. Brook, H.\n Brown, A. B\\\"uchler-Germann, I. Burducea, A. Bursche, J. Buytaert, S.\n Cadeddu, J. M. Caicedo Carvajal, O. Callot, M. Calvi, M. Calvo Gomez, A.\n Camboni, P. Campana, A. Carbone, G. Carboni, R. Cardinale, A. Cardini, L.\n Carson, K. Carvalho Akiba, G. Casse, M. Cattaneo, M. Charles, Ph.\n Charpentier, N. Chiapolini, K. Ciba, X. Cid Vidal, G. Ciezarek, P. E. L.\n Clarke, M. Clemencic, H. V. Cliff, J. Closier, C. Coca, V. Coco, J. Cogan, P.\n Collins, A. Comerma-Montells, F. Constantin, G. Conti, A. Contu, A. Cook, M.\n Coombes, G. Corti, G. A. Cowan, R. Currie, B. D'Almagne, C. D'Ambrosio, P.\n David, I. De Bonis, S. De Capua, M. De Cian, F. De Lorenzi, J. M. De Miranda,\n L. De Paula, P. De Simone, D. Decamp, M. Deckenhoff, H. Degaudenzi, M.\n Deissenroth, L. Del Buono, C. Deplano, D. Derkach, O. Deschamps, F. Dettori,\n J. Dickens, H. Dijkstra, P. Diniz Batista, F. Domingo Bonal, S. Donleavy, F.\n Dordei, A. Dosil Su\\'arez, D. Dossett, A. Dovbnya, F. Dupertuis, R.\n Dzhelyadin, A. Dziurda, S. Easo, U. Egede, V. Egorychev, S. Eidelman, D. van\n Eijk, F. Eisele, S. Eisenhardt, R. Ekelhof, L. Eklund, Ch. Elsasser, D.\n Elsby, D. Esperante Pereira, L. Est\\'eve, A. Falabella, E. Fanchini, C.\n F\\\"arber, G. Fardell, C. Farinelli, S. Farry, V. Fave, V. Fernandez Albor, M.\n Ferro-Luzzi, S. Filippov, C. Fitzpatrick, M. Fontana, F. Fontanelli, R.\n Forty, M. Frank, C. Frei, M. Frosini, S. Furcas, A. Gallas Torreira, D.\n Galli, M. Gandelman, P. Gandini, Y. Gao, J-C. Garnier, J. Garofoli, J. Garra\n Tico, L. Garrido, D. Gascon, C. Gaspar, N. Gauvin, M. Gersabeck, T. Gershon,\n Ph. Ghez, V. Gibson, V. V. Gligorov, C. G\\\"obel, D. Golubkov, A. Golutvin, A.\n Gomes, H. Gordon, M. Grabalosa G\\'andara, R. Graciani Diaz, L. A. Granado\n Cardoso, E. Graug\\'es, G. Graziani, A. Grecu, E. Greening, S. Gregson, B.\n Gui, E. Gushchin, Yu. Guz, T. Gys, G. Haefeli, C. Haen, S. C. Haines, T.\n Hampson, S. Hansmann-Menzemer, R. Harji, N. Harnew, J. Harrison, P. F.\n Harrison, J. He, V. Heijne, K. Hennessy, P. Henrard, J. A. Hernando Morata,\n E. van Herwijnen, E. Hicks, K. Holubyev, P. Hopchev, W. Hulsbergen, P. Hunt,\n T. Huse, R. S. Huston, D. Hutchcroft, D. Hynds, V. Iakovenko, P. Ilten, J.\n Imong, R. Jacobsson, A. Jaeger, M. Jahjah Hussein, E. Jans, F. Jansen, P.\n Jaton, B. Jean-Marie, F. Jing, M. John, D. Johnson, C. R. Jones, B. Jost, M.\n Kaballo, S. Kandybei, M. Karacson, T. M. Karbach, J. Keaveney, I. R. Kenyon,\n U. Kerzel, T. Ketel, A. Keune, B. Khanji, Y. M. Kim, M. Knecht, P.\n Koppenburg, A. Kozlinskiy, L. Kravchuk, K. Kreplin, M. Kreps, G. Krocker, P.\n Krokovny, F. Kruse, K. Kruzelecki, M. Kucharczyk, R. Kumar, T.\n Kvaratskheliya, V. N. La Thi, D. Lacarrere, G. Lafferty, A. Lai, D. Lambert,\n R. W. Lambert, E. Lanciotti, G. Lanfranchi, C. Langenbruch, T. Latham, R. Le\n Gac, J. van Leerdam, J.-P. Lees, R. Lef\\'evre, A. Leflat, J. Lefran\\c{c}ois,\n O. Leroy, T. Lesiak, L. Li, L. Li Gioi, M. Lieng, M. Liles, R. Lindner, C.\n Linn, B. Liu, G. Liu, J. H. Lopes, E. Lopez Asamar, N. Lopez-March, J.\n Luisier, F. Machefert, I. V. Machikhiliyan, F. Maciuc, O. Maev, J. Magnin, S.\n Malde, R. M. D. Mamunur, G. Manca, G. Mancinelli, N. Mangiafave, U. Marconi,\n R. M\\\"arki, J. Marks, G. Martellotti, A. Martens, L. Martin, A. Mart\\'in\n S\\'anchez, D. Martinez Santos, A. Massafferri, Z. Mathe, C. Matteuzzi, M.\n Matveev, E. Maurice, B. Maynard, A. Mazurov, G. McGregor, R. McNulty, C.\n Mclean, M. Meissner, M. Merk, J. Merkel, R. Messi, S. Miglioranzi, D. A.\n Milanes, M.-N. Minard, S. Monteil, D. Moran, P. Morawski, R. Mountain, I.\n Mous, F. Muheim, K. M\\\"uller, R. Muresan, B. Muryn, M. Musy, J.\n Mylroie-Smith, P. Naik, T. Nakada, R. Nandakumar, I. Nasteva, M. Nedos, M.\n Needham, N. Neufeld, C. Nguyen-Mau, M. Nicol, S. Nies, V. Niess, N. Nikitin,\n A. Nomerotski, A. Novoselov, A. Oblakowska-Mucha, V. Obraztsov, S. Oggero, S.\n Ogilvy, O. Okhrimenko, R. Oldeman, M. Orlandea, J. M. Otalora Goicochea, P.\n Owen, K. Pal, J. Palacios, A. Palano, M. Palutan, J. Panman, A. Papanestis,\n M. Pappagallo, C. Parkes, C. J. Parkinson, G. Passaleva, G. D. Patel, M.\n Patel, S. K. Paterson, G. N. Patrick, C. Patrignani, C. Pavel-Nicorescu, A.\n Pazos Alvarez, A. Pellegrino, G. Penso, M. Pepe Altarelli, S. Perazzini, D.\n L. Perego, E. Perez Trigo, A. P\\'erez-Calero Yzquierdo, P. Perret, M.\n Perrin-Terrin, G. Pessina, A. Petrella, A. Petrolini, A. Phan, E. Picatoste\n Olloqui, B. Pie Valls, B. Pietrzyk, T. Pilar, D. Pinci, R. Plackett, S.\n Playfer, M. Plo Casasus, G. Polok, A. Poluektov, E. Polycarpo, D. Popov, B.\n Popovici, C. Potterat, A. Powell, T. du Pree, J. Prisciandaro, V. Pugatch, A.\n Puig Navarro, W. Qian, J. H. Rademacker, B. Rakotomiaramanana, M. S. Rangel,\n I. Raniuk, G. Raven, S. Redford, M. M. Reid, A. C. dos Reis, S. Ricciardi, K.\n Rinnert, D. A. Roa Romero, P. Robbe, E. Rodrigues, F. Rodrigues, P. Rodriguez\n Perez, G. J. Rogers, S. Roiser, V. Romanovsky, M. Rosello, J. Rouvinet, T.\n Ruf, H. Ruiz, G. Sabatino, J. J. Saborido Silva, N. Sagidova, P. Sail, B.\n Saitta, C. Salzmann, M. Sannino, R. Santacesaria, C. Santamarina Rios, R.\n Santinelli, E. Santovetti, M. Sapunov, A. Sarti, C. Satriano, A. Satta, M.\n Savrie, D. Savrina, P. Schaack, M. Schiller, S. Schleich, M. Schmelling, B.\n Schmidt, O. Schneider, A. Schopper, M.-H. Schune, R. Schwemmer, B. Sciascia,\n A. Sciubba, M. Seco, A. Semennikov, K. Senderowska, I. Sepp, N. Serra, J.\n Serrano, P. Seyfert, B. Shao, M. Shapkin, I. Shapoval, P. Shatalov, Y.\n Shcheglov, T. Shears, L. Shekhtman, O. Shevchenko, V. Shevchenko, A. Shires,\n R. Silva Coutinho, H. P. Skottowe, T. Skwarnicki, A. C. Smith, N. A. Smith,\n E. Smith, K. Sobczak, F. J. P. Soler, A. Solomin, F. Soomro, B. Souza De\n Paula, B. Spaan, A. Sparkes, P. Spradlin, F. Stagni, S. Stahl, O. Steinkamp,\n S. Stoica, S. Stone, B. Storaci, M. Straticiuc, U. Straumann, N. Styles, V.\n K. Subbiah, S. Swientek, M. Szczekowski, P. Szczypka, T. Szumlak, S.\n T'Jampens, E. Teodorescu, F. Teubert, C. Thomas, E. Thomas, J. van Tilburg,\n V. Tisserand, M. Tobin, S. Topp-Joergensen, N. Torr, E. Tournefier, M. T.\n Tran, A. Tsaregorodtsev, N. Tuning, M. Ubeda Garcia, A. Ukleja, P. Urquijo,\n U. Uwer, V. Vagnoni, G. Valenti, R. Vazquez Gomez, P. Vazquez Regueiro, S.\n Vecchi, J. J. Velthuis, M. Veltri, K. Vervink, B. Viaud, I. Videau, X.\n Vilasis-Cardona, J. Visniakov, A. Vollhardt, D. Voong, A. Vorobyev, H. Voss,\n K. Wacker, S. Wandernoth, J. Wang, D. R. Ward, N. K. Watson, A. D. Webber, D.\n Websdale, M. Whitehead, D. Wiedner, L. Wiggers, G. Wilkinson, M. P. Williams,\n M. Williams, F. F. Wilson, J. Wishahi, M. Witek, W. Witzeling, S. A. Wotton,\n K. Wyllie, Y. Xie, F. Xing, Z. Xing, Z. Yang, R. Young, O. Yushchenko, M.\n Zavertyaev, F. Zhang, L. Zhang, W. C. Zhang, Y. Zhang, A. Zhelezov, L. Zhong,\n E. Zverev, A. Zvyagin",
"submitter": "Marina Artuso",
"url": "https://arxiv.org/abs/1111.2357"
}
|
1111.2380
|
# Age and structure parameters of a remote M31 globular cluster B514 based on
HST, 2MASS, GALEX and BATC observations
Jun Ma11affiliation: National Astronomical Observatories, Chinese Academy of
Sciences, Beijing 100012, P.R. China; majun@nao.cas.cn 22affiliation: Key
Laboratory of Optical Astronomy, National Astronomical Observatories, Chinese
Academy of Sciences, Beijing, 100012, China , Song Wang11affiliation: National
Astronomical Observatories, Chinese Academy of Sciences, Beijing 100012, P.R.
China; majun@nao.cas.cn 33affiliation: Graduate University of Chinese Academy
of Sciences, A19 Yuquan Road, Shijingshan District, Beijing 100049, China ,
Zhenyu Wu11affiliation: National Astronomical Observatories, Chinese Academy
of Sciences, Beijing 100012, P.R. China; majun@nao.cas.cn , Zhou
Fan11affiliation: National Astronomical Observatories, Chinese Academy of
Sciences, Beijing 100012, P.R. China; majun@nao.cas.cn , Tianmen
Zhang11affiliation: National Astronomical Observatories, Chinese Academy of
Sciences, Beijing 100012, P.R. China; majun@nao.cas.cn , Jianghua
Wu11affiliation: National Astronomical Observatories, Chinese Academy of
Sciences, Beijing 100012, P.R. China; majun@nao.cas.cn , Xu Zhou11affiliation:
National Astronomical Observatories, Chinese Academy of Sciences, Beijing
100012, P.R. China; majun@nao.cas.cn , Zhaoji Jiang11affiliation: National
Astronomical Observatories, Chinese Academy of Sciences, Beijing 100012, P.R.
China; majun@nao.cas.cn and Jiansheng Chen11affiliation: National Astronomical
Observatories, Chinese Academy of Sciences, Beijing 100012, P.R. China;
majun@nao.cas.cn
###### Abstract
B514 is a remote M31 globular cluster which locating at a projected distance
of $R_{p}\simeq 55$ kpc. Deep observations with the Advanced Camera for
Surveys (ACS) on the Hubble Space Telescope (HST) are used to provide the
accurate integrated light and star counts of B514. By coupling analysis of the
distribution of the integrated light with star counts, we are able to reliably
follow the profile of the cluster out to $\sim 40\arcsec$. Based on the
combined profile, we study in detail its surface brightness distribution in
F606W and F814W filters, and determine its structural parameters by fitting a
single-mass isotropic King model. The results showed that, the surface
brightness distribution departs from the best-fit King model for
$r>10^{\prime\prime}$. B514 is quite flatted in the inner region, and has a
larger half-light radius than majority of normal globular clusters of the same
luminosity. It is interesting that, in the $M_{V}$ versus $\log R_{h}$ plane,
B514 lies nearly on the threshold for ordinary globular clusters as defined by
Mackey & van den Bergh. In addition, B514 was observed as part of the Beijing-
Arizona-Taiwan-Connecticut (BATC) Multicolor Sky Survey, using 13
intermediate-band filters covering a wavelength range of 3000–8500 Å. Based on
aperture photometry, we obtain its SEDs as defined by the 13 BATC filters. We
determine the cluster’s age and mass by comparing its SEDs (from 2267 to
20000Å, comprising photometric data in the near-ultraviolet of GALEX, 5 SDSS
bands, 13 BATC intermediate-band, and 2MASS near-infrared $JHK_{\rm s}$
filters) with theoretical stellar population synthesis models, resulting in
age of $11.5\pm 3.5$ Gyr. This age confirms the previous suggestion that B514
is an old GC in M31. B514 has a mass of $0.96-1.08\times 10^{6}\rm M_{\odot}$,
and is a medium-mass globular cluster in M31.
###### Subject headings:
galaxies: evolution – galaxies: individual (M31) – star clusters: B514
††slugcomment: AJ, in press
## 1\. Introduction
In hierarchical cosmological models, galaxies are built up through the
continual accretion and merger of smaller ones. The signature of these system
assembly processes is expected to be seen in the outskirts of a galactic halo.
Globular clusters (GCs), as luminous compact objects that are found out to
distant radii in the haloes of massive galaxies, can serve as excellent
tracers of substructures in the outer region of their parent galaxy. So,
detailed studies on GCs in the outer haloes of the local galaxies are very
important.
M31, with a distance modulus of 24.47 (Holland, 1998; Stanek & Garnavich,
1998; McConnachie et al., 2005), is an ideal nearby galaxy for studying GCs,
since it is so near, and contains more GCs than all other Local Group galaxies
combined (Battistini et al., 1987; Racine, 1991; Harris, 1991; Fusi Pecci et
al., 1993). The study of GCs in M31 was initiated by Hubble (1932), who
discovered 140 GC candidates with $m_{pg}\leq 18$ mag. And then, a number of
catalogs of GC candidates were published. For example, the Bologna Group
(Battistini et al., 1980, 1987, 1993) did independent searches of GC
candidates and compiled them with their own Bologna number. Bologna catalog
contains a total of 827 objects, and all the objects were classified into five
classes by authors’ degree of confidence. 353 of these candidates were
considered as class A or class B by high level of confidence, and the others
fell into class C, D, or E. V magnitude and $B-V$ color for most candidates
were given in the Bologna catalog. There are recent works dealing with the
searched and the catalogs of M31 GCs (e.g., Mochejska et al., 1998; Barmby &
Huchra, 2001; Galleti et al., 2004, 2006, 2007; Huxor et al., 2005, 2008,
2011; Kim et al., 2007; Mackey et al., 2006, 2007, 2010; Martin et al., 2006;
Caldwell et al., 2009, 2011; Peacock et al., 2010). The continued importance
of the study of GCs in this galaxy has been reviewed by Barmby et al. (2000).
From the spatial structure and internal stellar kinematics of GCs, we can get
information on both their formation conditions and dynamical evolution within
the tidal fields of their host galaxies. For example, the structural
parameters of GCs indicate on what timescales the cluster is bound to
dissolve. However, the integrated properties of GCs, such as age and
metallicity, are believed to reflect conditions in the early stages of galaxy
formation (Brodie & Strader 2006).
The most direct method to determine a cluster’s age is by employing main-
sequence photometry, since the absolute magnitude of the main-sequence turnoff
is predominantly affected by age (see Puzia et al., 2002, and references
therein). However, until recently (cf. Perina et al., 2009), this method was
only applied to the star clusters in the Milky Way and its satellites (e.g.,
Rich et al., 2001), although Brown et al. (2004) estimated the age of an M31
GC using extremely deep images observed with the HST’s Advanced Camera for
Surveys (ACS). Generally, the ages of extragalactic star clusters are
determined by comparing their observed spectral-energy distributions (SEDs)
and/or spectroscopy with the predictions of SSP models (Williams & Hodge,
2001a, b; de Grijs et al., 2003b, c, a; Bik et al., 2003; Jiang et al., 2003;
Beasley et al., 2004; Puzia et al., 2005; Fan et al., 2006; Ma et al., 2006a,
2007a, 2009a, 2009b, 2011; Caldwell et al., 2009, 2011; Wang et al., 2010;
Perina et al., 2011).
M31 GC B514 (B for ‘Bologna’, see Battistini 1987), which was detected by
Galleti et al. (2005) based on the XSC sources of the All Sky Data Release of
the Two Micron All Sky Survey (2MASS) within a $\sim 9^{\circ}\times
9^{\circ}$ area centered on M31, is the outermost cluster known in M31 at that
time, which locating at a projected distance of $R_{p}\simeq 55$ kpc. Now,
many new members of M31 halo GC system, which extending to very large radii,
have been discovered (e.g., Huxor et al., 2005, 2008, 2011; Mackey et al.,
2006, 2007, 2010; Martin et al., 2006).
Galleti et al. (2006) presented deep color-magnitude diagram for B514 in F606W
and F814W photometry obtained with the ACS/HST, which reveals a steep red
giant branch (RGB) and a horizontal branch (HB) extending blueward of the
instability strip showing that B514 is a classical old metal-poor GC. Federici
et al. (2007) studied the density profile of B514 based on the HST/ACS
observations as Galleti et al. (2006) had, and they found that the light and
the star-count profiles show a departure from the best-fitted empirical models
of King (1962) for $r\geq 8^{\prime\prime}$ – as a surface brightness excess
at large radii, and the star-count profile shows a clear break in the
correspondence of the estimated tidal radius; they also found that B514 has a
significantly larger half-light radius than ordinary GCs of the same
luminosity. Clementini et al. (2009) identified a rich harvest of RR Lyrae
stars in B514, based on HST Wide Field Planetary Camera 2 (WFPC2) and HST/ACS
time-series observations.
Since B514 is located in the halo of M31, i.e., far away from the galaxy’s
disk, it is (for all practical purposes) only affected by the Galactic
foreground extinction. About the foreground Galactic reddening in the
direction of M31, it was discussed by many authors (e.g., van den Bergh, 1969;
McClure & Racine, 1969; Frogel et al., 1980; Fusi Pecci et al., 2005), and
nearly similar values were determined such as $E(B-V)=0.08$ by van den Bergh
(1969), 0.11 by McClure & Racine (1969) and Hodge (1992), 0.08 by Frogel et
al. (1980). We argue that the reddening value of B514 should not smaller than
the foreground Galactic reddening in the direction of M31. In this paper, we
adopt the reddening value of $E(B-V)=0.10$ from Galleti et al. (2006). The
reddening law from Cardelli et al. (1989) is employed in this paper. In
addition, throughout this paper we adopt a distance to M31 of $783\pm 25$ kpc
($1{\arcsec}$ subtends 3.8 pc), corresponding to a distance modulus of
$(m-M)_{0}=24.47\pm 0.07$ mag (McConnachie et al., 2005).
In this paper, we describe the details of the observations and our approach to
the data reduction with the HST/ACS and the BATC system in §2 and §3. We will
study the surface distribution of B514 with the King models of King (1966) in
details, which are developed by Michie (1963) and King (1966) based on the
assumption that the GCs are formed by single-mass, isotropic, lowered
isothermal spheres (hereafter ‘King models’). We determine the age and mass of
B514 by comparing observational SEDs with population synthesis models in §4.
We provide a summary in §5.
## 2\. Observation and photometric data with HST/ACS
The images of B514 used in this paper were observed with the ACS/Wide Field
Camera (WFC) in the F606W and F814W filters on 2005 July 19 (program ID GO
10394, PI: N. Tanvir), covering the period 2005, July 19–20 in F606W (total
$t_{\rm exp}=1776$ s) and F814W (total $t_{\rm exp}=2505$ s). Upon retrieval
from the STScI archive, all images were processed by the standard ACS
calibration pipeline, in which bias and dark subtractions, flatfield division,
and the masking of known bad pixels are included. Subsequently, photometric
header keywords are populated. In the final stage of the pipeline, the
MultiDrizzle software is used to correct the geometric distortion presented in
the images. Finally, any cosmic rays are rejected while individual images in
each band are combined into a final single image. We checked the images, and
did not find saturated cluster stars. Figure 1 show the images observed with
the ACS/WFC in the F606W and F814W. The ACS/WFC spatial resolution is
$0.05\arcsec$ pixel-1.
### 2.1. Ellipticity, position angle, and surface brightness profile
Surface photometry of the cluster is obtained from the drizzled images, using
the iraf task ellipse. Its center position was fixed at a value derived by
object locator of ellipse task, however an initial center position was
determined by centroiding. Elliptical isophotes were fitted to the data, with
no sigma clipping. We ran two passes of ellipse task, the first pass was run
in the usual way, with ellipticity and position angle allowed to vary with the
isophote semimajor axis. In the second pass, surface brightness profiles on
fixed, zero-ellipticity isophotes were measured, since we choose to fit
circular models for the intrinsic cluster structure and the point spread
function (PSF) as Barmby et al. (2007) did (see §2.3 for details). The
background value was derived as the mean of a region of $100\times 100$ pixels
in “empty” areas far away from the cluster. We performed the photometric
calibration using the results of Sirianni et al. (2005): 26.398 in F606W and
25.501 in F814W zero-points. Magnitudes are derived in the ACS/WFC vegamag
system.
Tables 1 and 2 list the ellipticity, $\epsilon=1-b/a$, and the position angle
(P.A.) as a function of the semi-major axis length, $a$, from the center of
annulus in the F606W and F814W filters, respectively. These observables have
also been plotted in Figure 2; the errors were generated by the iraf task
ellipse, in which the ellipticity errors are obtained from the internal errors
in the harmonic fit, after removal of the first and second fitted harmonics.
From Table 1 and Figure 2, we can see that, the values of ellipticity and
position angle cannot be obtained beyond $0.9744\arcsec$ in the F606W filter
because of low signal-to-noise ratio. In addition, Figure 2 shows that the
ellipticity varies significantly with position along the semimajor axis radius
smaller than $\sim 0.1752\arcsec$. Beyond $\sim 0.1752\arcsec$, the
ellipticity does not vary significantly as a function of the cluster s
semimajor axis. The P.A. does not vary as a function of the cluster s
semimajor axis within $\sim 0.1752\arcsec$ because of high signal-to-noise
ratio; however, beyond this position, it varies significantly with great
errors because of low signal-to-noise ratio.
Tables 3 and 4 list the surface brightness profile, $\mu$, of B514, and its
integrated magnitude, $m$, as a function of radius in the F606W and F814W
filters, respectively. The errors in the surface brightness were generated by
the iraf task ellipse, in which they are obtained directly from the root mean
square scatter of the intensity data along the zero-ellipticity isophotes. In
addition, the surface photometries at radii where the ellipticity and position
angle cannot be measured, are obtained based on the last ellipticity and
position angle as the iraf task ellipse is designed.
In order to derive the surface brightness profile of B514 in its outer region,
we use the profile from star counts. We used the
DOLPHOT111http://americano.dolphinsim.com/dolphot/ photometry software
(Dolphin 2000a), specifically the ACS module, to photometer our images.
DOLPHOT performs PSF fitting using PSFs especially tailored to the ACS camera.
Photometry was done simultaneously on all the flat-fielded images from the
STScI archive (both filters), relative to a deep reference image–we used the
drizzled combination of the F606W image. DOLPHOT accounts for the hot-pixel
and cosmic-ray masking information attached to each flat-fielded image, fits
the sky locally around each detected source, and automatically applies the
correction for the charge-transfer efficiency (CTE, Dolphin 2000b). It then
transforms instrumental magnitude to the vegamag system (Dolphin 2000b). A
variety of quality information is listed with each detected object, including
the object type (stellar, extended, etc.), $\chi^{2}$ of the PSF fit,
sharpness and roundness of the object, and a “crowding” parameter which
describes how much brighter an object would have been had neighboring objects
not been fitted simultaneously. We used the quality information provided by
DOLPHOT to clean the resulting detection lists, selecting only stellar
detections, with valid photometry on all input images, global sharpness
parameter between $-0.3$ and 0.3 in each filter, and crowding parameter less
than 0.25 in each filter.
We joined the two profiles into one based on the method of Federici et al.
(2007). This involves matching the intensity scales of the two profiles by
fitting both profiles to smooth curves in the region $r=9-16^{\prime\prime}$.
The star count profile is listed in Table 5. The errors for the star counts
take account of Poisson statistical uncertainties. The joined profile covered
the full $0<r\leq 40^{\prime\prime}$ range, as shown in Figure 3.
Figure 1.— The images of GC B514 observed in the F606W and F814W filters of
ACS/HST. The image size is $20\arcsec\times 20\arcsec$ for each panel. Figure
2.— Ellipticity and P.A. as a function of the semimajor axis in the F606W and
F814W filters of ACS/HST.
### 2.2. Point spread function
At a distance of 783 kpc, the ACS/WFC has a scale of
$\rm{0.05^{\prime\prime}=0.19~{}pc~{}pixel^{-1}}$, and thus M31 clusters are
clearly resolved with it. Their observed core structures, however, are still
affected by the PSF. We chose not to deconvolve the data, instead fitting
structural models after convolving them with a simple analytic description of
the PSF as Barmby et al. (2007) and McLaughlin et al. (2008) did (see Barmby
et al., 2007; McLaughlin et al., 2008; Ma et al., 2011, for details). In
addition, since this PSF formula is radially symmetric and the models of King
(1966) we fit are intrinsically spherical, the convolved models to be fitted
to the data are also circularly symmetric.
### 2.3. Models and fits
#### 2.3.1 Structural models
After elliptical galaxies, GCs are the best understood and most thoroughly
modelled class of stellar systems. For example, a large majority of the $\sim
150$ Galactic GCs have been fitted by the simple models of single-mass,
isotropic, lowered isothermal spheres developed by Michie (1963) and King
(1966) (i.e. King models), yielding comprehensive catalogs of cluster
structural parameters and physical properties (see McLaughlin & van der Marel,
2005, and references therein). For extragalactic GCs, HST imaging data have
been used to fit King models to a large number of GCs in M31 (e.g., Barmby et
al., 2002, 2007, 2009, and references therein), in M33 (Larsen et al., 2002),
and in NGC 5128 (e.g., Harris et al., 2002; McLaughlin et al., 2008, and
references therein). In addition, there are other models used to fit the
surface profile of GCs, including Wilson (1975), Elson et al. (1987), and
Sérsic (1968). In this paper, we fit King models to the density profile of
B514 observed with ACS/WFC.
#### 2.3.2 Fits
Our fitting procedure involves computing in full large numbers of King
structural models, spanning a wide range of fixed values of the appropriate
shape parameter $W_{0}$ (see McLaughlin & van der Marel, 2005, in detail). And
then the models are convolved with the ACS/WFC PSF for the F606W and F814W
filters (see Barmby et al., 2007, for details):
${\widetilde{I}_{\rm
mod}^{*}(R|r_{0})=\int\\!\\!\\!\int_{-\infty}^{\infty}\widetilde{I}_{\rm
mod}(R^{\prime}/r_{0})\times\widetilde{I}_{\rm
PSF}\left[(x-x^{\prime}),(y-y^{\prime})\right]}$
$\ dx^{\prime}\,dy^{\prime}\ ,$ (1)
where $\widetilde{I}_{\rm mod}\equiv I_{\rm mod}/I_{0}$ (see McLaughlin et
al., 2008, in detail). We changed the luminosity density to surface brightness
$\widetilde{\mu}_{\rm mod}^{*}=-2.5\,\log\,[\widetilde{I}_{\rm mod}^{*}]$
before fitting them to the observed surface brightness profile of B514,
$\mu=\mu_{0}-2.5\,\log\,[I(R/r_{0})/I_{0}]$, finding the radial scale $r_{0}$
and central surface brightness $\mu_{0}$ which minimizing $\chi^{2}$ for every
given value of $W_{0}$. The $(W_{0},r_{0},\mu_{0})$ combination that yields
the global minimum $\chi_{\rm min}^{2}$ over the grid used defines the best-
fit model of that type:
$\chi^{2}=\sum_{i}{\frac{\left[\mu_{\rm obs}(R_{i})-\widetilde{\mu}_{\rm
mod}^{*}(R_{i}|r_{0})\right]^{2}}{\sigma_{i}^{2}}},$ (2)
in which $\sigma_{i}$ is the error in the surface brightness. Estimates of the
one-sigma uncertainties on these basic fit parameters follow from their
extreme values over the subgrid of fits with $\chi^{2}/\nu\leq\chi_{\rm
min}^{2}/\nu+1$, here $\nu$ is the number of free parameters. Figure 3 shows
our best King fits to B514. In Figure 3, open squares are ellipse data points
included in the least-squares model fitting, and the asterisks are points not
used to constrain the fit; black circles are star-counts points included in
the $\chi^{2}$ model fitting, and red circles are star-counts points not used
to constrain the fits. These observed data points shown by asterisks are
included in the radius of $R<2~{}\rm{pixels}=0\farcs 1$, and the isophotal
intensity is dependent on its neighbors. As Barmby et al. (2007) pointed out
that, the ellipse output contains brightnesses for 15 radii inside 2 pixel,
but they are all measured from the same 13 central pixels and are not
statistically independent. So, to avoid excessive weighting of the central
regions of B514 in the fits, we only used intensities at radii $R_{\rm min}$,
$R_{\rm min}+(0.5,1.0,2.0)~{}{\rm pixels}$, or $R>2.5$ as Barmby et al. (2007)
used. Table 6 summarizes the results obtained in this paper.
From Figure 3, we can note that the surface brightness distribution departs
from the best-fit King model for $r>10^{\prime\prime}$, which can be
interpreted as the presence of a population of extratidal stars around the
cluster. In fact, Federici et al. (2007) have reported this population of
extratidal stars (see their Fig. 5 and their discussions).
Figure 3.— Surface brightness profile of B514 measured in the F606W and F814
filters. Dashed curves (blue) trace the PSF intensity profiles and solid (red)
curves are the PSF-convolved best-fit models. Open squares are ellips data
points and black circles are star-counts profiles included in the $\chi^{2}$
model fitting, and the asterisks are ellips data points and red circles are
star-counts profiles not used to constrain the fits (see the text in detail).
Figure 4.— $M_{V}$ vs. $R_{h}$ for interesting stellar systems. The plotted
line is the threshold for ordinary cluster in this plane as defined by Mackey
& van den Bergh (2005), $\log R_{h}({\rm pc})=0.25M_{V}({\rm mag})+2.95$. Red
circles are Galactic globular clusters from the on-line data base of Harris
(1996) (2010 update); green crosses are M31 globular clusters from Barmby et
al. (2007); green filled triangles are M31 young massive clusters from Barmby
et al. (2009); green filled circles are M31 extended clusters from Mackey et
al. (2006); green open triangles are outer M31 GCs from Mackey et al. (2007);
black filled circles are M33 outer halo clusters from Cockcroft et al. (2011);
the black square is B514 derived here.
### 2.4. Distribution of B514 in the $M_{V}$ vs. $\log R_{h}$ plane
The distribution of stellar systems in the $M_{V}$ vs. $\log R_{h}$ plane can
provide interesting information on the evolutionary history of these objects
(e.g., van den Bergh & Mackey, 2004; Mackey & van den Bergh, 2005). In this
plane, the half-light radius is an important parameter, which can be used to
trace the initial size of a cluster, since it changes little in evolution
process (see Spitzer & Thuan, 1972; Henon, 1973; Lightman & Shapiro, 1978;
Murphy et al., 1990, for details).
Recently, van den Bergh & Mackey (2004) and Mackey & van den Bergh (2005)
showed that in a plot of luminosity versus half-light radius, the overwhelming
majority of normal Galactic GCs lie below (or to the right) of the line:
$\log R_{h}({\rm pc})=0.25M_{V}({\rm mag})+2.95.$ (3)
Exceptions to this rule are massive clusters, such as M54 and $\omega$
Centauri in the Milky Way, and G1 in M31, which are widely believed to be the
remnant cores of now defunct dwarf galaxies (Zinnecker et al., 1988; Freeman,
1993; Meylan et al., 2001). Because the well-known giant GC NGC 2419 (van den
Bergh & Mackey, 2004) in the Galaxy and 037-B327 (Ma et al., 2006b) in M31
also lie above this line, it has been speculated that these two objects might
also be the remnant cores of dwarf galaxies (but see de Grijs et al., 2005,
for doubts regarding NGC 2419).
With the value of $R_{h}$ (i.e. $r_{h}$) in the F606W filter obtained in this
paper, we plot the relationship of $M_{V}$ versus $\log R_{h}$ in Figure 4, in
which $M_{V}=-9.02$ which being derived based on $m_{V}=15.76$ from Huxor et
al. (2008). It is interesting that, on this plot B514 is seen to lie nearly on
the line defined by equation (3). Considering the uncertainties of $R_{h}$ and
$M_{V}$, a certain conclusion may not be presented here. However, we argued
that, B514 is a medium-mass GC in M31 (see §4.4 for details), and is not as
massive as G1 and 037–B327 (see Ma et al., 2006a, b, 2007b, 2009b, for
details). Furthermore, and for completeness, in Figure 4 we have also included
GCs in the Milky Way, M31 and M33. Galactic GCs are from the on-line data base
of Harris (1996) (2010 update). This new revision of the McMaster catalog of
Galactic GCs is the first update since 2003 and the biggest single revision
since the original version of the catalog published in 1996. The starting
points for the present list of structural parameters are the major
compilations of McLaughlin & van der Marel (2005) and Trager et al. (1995).
McLaughlin & van der Marel (2005) used the same raw data as Trager et al.
(1995), and derived structural parameter values from King (1966) dynamical
profile models. M31 GCs are from recent compilations of data by Barmby et al.
(2007, 2009), Mackey et al. (2006, 2007). M33 GCs are from Cockcroft et al.
(2011). Barmby et al. (2007) derived structural parameters for 34 GCs in M31
based on ACS/HST observations, and the derived structural parameters are
combined with corrected versions of those measured in an earlier survey in
order to construct a comprehensive catalog of structural and dynamical
parameters for 93 M31 GCs. Barmby et al. (2009) measured structural parameters
for 23 bright young clusters in M31 based on the HST/WFPC2 observations, and
suggested that on average they are larger and more concentrated than typical
old clusters. Mackey et al. (2006) determined structural parameters for 4
extended, luminous globular clusters in the outskirts of M31 based on ACS/HST
observations. These objects were discovered by Huxor et al. (2005) and Martin
et al. (2006). Mackey et al. (2007) derived structural parameters for 10
classical GCs in the far outer regions of M31 based on ACS/WFC observations.
Cockcroft et al. (2011) searched for outer halo star clusters in M33 as part
of the Pan-Andromeda Archaeological Survey using the images taken with the
Canada-France-Hawaii Telescope (CFHT)/MegaCam, and found one new unambiguous
star cluster in addition to the five previously known in the M33 outer halo,
and determined structural parameters for these 6 outer halo clusters.
From Figure 4, we can see that, for the data of Galactic GCs which updated in
2010, in addition to the clusters already noted by Mackey & van den Bergh
(2005), i.e. M 54 (NGC 6715), ${\omega{\rm~{}Cen}}$ (NGC 5139), NGC 2419,
there are three other bright GCs (NGC 104, NGC 5272, and NGC 5024) lying above
the “ordinary globular clusters” threshold. For M31 star clusters, in addition
to G1, there are 26 other bright clusters lying above the line defined by
equation (3). All of these objects are classified as GCs (Barmby et al., 2007;
Mackey et al., 2007) or bright young clusters (Barmby et al., 2009). For 4
extended, luminous GCs in the outskirts of M31, they all lie above the line
defined by equation (3).
Based on F606W and F814W images of B514 obtained with the ACS/HST (program ID
GO 10565, PI: S. Galleti), Federici et al. (2007) also studied in detail its
surface brightness distribution in F606W and F814W filters, and determine its
structural parameters by fitting a King (1962) model to a surface brightness
profile. Comparing the results of Federici et al. (2007) with Table 6 of this
paper, we find that our model fits produce smaller tidal radii, which
resulting in smaller half-light, or effective, radii of a model. In addition,
Federici et al. (2007) adopted $M_{V}=-9.1$ being brighter than $M_{V}=-9.02$
adopted here. So, in Federici et al. (2007), B514 lied above and brightward of
the line defined by equation (3).
## 3\. Archival images of the BATC Multicolor Sky Survey, 2MASS and GALEX and
photometric data of SDSS
In this section, we will determine the magnitudes of B514 based on the
archival images of the BATC Multicolor Sky Survey, 2MASS and GALEX using a
standard aperture photometry approach, i.e., the phot routine in daophot
(Stetson, 1987). In addition, we will introduce the photometric data of B514
from the Sloan Digital Sky Survey (SDSS) obtained by Peacock et al. (2010)
### 3.1. Intermediate-band photometry of B514
Observations of B514 were also obtained with the BATC 60/90cm Schmidt
telescope located at the Xinglong station of the National Astronomical
Observatory of China (NAOC). This telescope is equipped with 15 intermediate-
band filters covering the optical wavelength range from 3000 to 10000 Å(see
Fan et al., 2009, for details). Figure 5 shows a finding chart of B514 in the
BATC $b$ band (centered at 5795 Å).
Figure 5.— Image of B514 in the BATC $b$ band, obtained with the NAOC 60/90cm
Schmidt telescope. B514 is circled using an aperture with a radius of
$13^{\prime\prime}$. The field of view of the image is $4.3^{\prime}\times
4.3^{\prime}$.
The BATC survey team obtained 47 images of B514 in 13 BATC filters between
2005 March 1 and 2006 December 9. Table 7 contains the observation log,
including the BATC filter names, the central wavelength and bandwidth of each
filter, the number of images observed through each filter, and the total
observing time per filter. Multiple images through the same filter were
combined to improve image quality (i.e., increase the signal-to-noise ratio
and remove spurious signal).
Calibration of the magnitude zero level in the BATC photometric system is
similar to that of the spectrophotometric AB magnitude system. For flux
calibration, the Oke-Gunn (Oke & Gunn, 1983) primary flux standard stars HD
19445, HD 84937, BD +26∘2606, and BD +17∘4708 were observed during photometric
nights (Yan et al., 2000). Column (6) of Table 7 gives the zero-point errors
in magnitude for the standard stars through each filter. The formal errors
obtained for these stars in the 13 BATC filters used are $\lesssim 0.02$ mag,
which implies that we can define photometrically the BATC system to an
accuracy of better than 0.02 mag.
We determined the intermediate-band magnitudes of B514 on the combined images.
The (radial) photometric asymptotic growth curves, in all BATC bands, flatten
out at a radius of $\sim 13^{\prime\prime}$. Inspection ensured that this
aperture is adequate for photometry, i.e., B514 does not show any obvious
signal beyond this radius. Therefore, we use an aperture with $r\approx
13^{\prime\prime}$ for integrated photometry. Since B514 is located in the M31
halo, contamination from background fluctuations can be neglected. We adopted
annuli for background subtraction spanning between 14 to $20^{\prime\prime}$.
The calibrated photometry of B514 in 13 filters is summarized in column (7) of
Table 7, in conjunction with the $1\sigma$ magnitude uncertainties, which
include uncertainties from the calibration errors in magnitude from daophot.
### 3.2. Near-infrared 2MASS photometry of B514
B514 was detected by Galleti et al. (2005) based on the XSC sources of the All
Sky Data Release of 2MASS within a $\sim 9^{\circ}\times 9^{\circ}$ area
centered on M31. In order to obtain accurate photometry for B514 in $JHK_{s}$,
we download the images in $JHK_{s}$ filters including B514. The image in each
filter is combined using 6 frames of 1.3 seconds, so the total exposure time
of image in each filter is 7.8 seconds. The mosaic pixel scale of the final
atlas image is resampled to $1^{\prime\prime}$ (see Skrutskie et al., 2006,
for details). The relevant zero-points for photometry are 20.9210, 20.7089,
and 20.0783 in $J$, $H$, and $K_{s}$ magnitudes, respectively, which are
presented in photometric header keywords. We use an aperture with
$r=13^{\prime\prime}$ for integrated photometry, and annuli for background
subtraction spanning between $14^{\prime\prime}$ to $19^{\prime\prime}$. The
calibrated photometry of B514 in $J$, $H$, and $K_{s}$ filters is summarized
in Table 8, in conjunction with the $1\sigma$ magnitude uncertainties obtained
from daophot.
### 3.3. GALEX Ultraviolet photometry of B514
While the principle science goal of the Galaxy Evolution Explorer (GALEX,
Martin et al. (2005); Morrissey et al. (2007)) has been the study of star
formation in the local and intermediate-redshift universe, nearby galaxies
such as M31 have also been surveyed, taking advantage of the wide
($1.2^{\circ}$) field of view of GALEX. The B514 images were obtained as part
of the guest program carried out by GALEX in two UV bands: far-ultraviolet
(FUV) ($\mbox{$\lambda_{\rm eff}$}=1539$ Å, FWHM $\approx 270$ Å), and near-
ultraviolet (NUV) ($\mbox{$\lambda_{\rm eff}$}=2316$ Å, FWHM $\approx 615$ Å)
with resolution $4.2^{\prime\prime}$ (FUV) and $5.3^{\prime\prime}$ (NUV)
(Morrissey et al., 2007). The exposure times are 1616 seconds in FUV and 1704
seconds in NUV. The images are sampled with $1.5^{\prime\prime}$ pixels. The
data was downloaded from the MAST archive. The relevant zero-points for
photometry are 20.08 and 18.82 in NUV and FUV magnitudes, respectively
(Morrissey et al., 2007). We use an aperture with $r=12^{\prime\prime}$ for
integrated photometry, and annuli for background subtraction spanning between
$13.5^{\prime\prime}$ to $19.5^{\prime\prime}$. The calibrated photometry of
B514 in NUV and FUV filters is summarized in Table 8. From Table 8, we can see
that the $1\sigma$ magnitude uncertainties are great, especially the magnitude
uncertainty in FUV is very great (2.3 magnitude), i.e. the signal-to-noise
ratios of these images are low, especially the signal-to-noise ratio of the
image in FUV is very low. Since the magnitude uncertainty in FUV is very
great, we will not use it when fitting to derive the age of B514 in §4.
### 3.4. Photometric data of B514 from SDSS
Peacock et al. (2010) presented an updated catalog of M31 GCs based on images
from the Wide Field Camera (WFCAM) on the United Kingdom Infrared Telescope
and from the SDSS, in which $ugriz$ and $K$-band photometry are determined. In
this catalog, B514 is named H6 from Huxor et al. (2008), and $ugriz$
photometry is presented.
## 4\. Stellar population of B514
### 4.1. Metallicity of B514
Cluster SEDs are determined by the combination of their ages and
metallicities, which is often referred to as the age-metallicity degeneracy.
Therefore, the age of a cluster can only be constrained accurately if the
metallicity is known with confidence, from independent determinations. There
exist four metallicity determinations for B514: namely, $\rm{[Fe/H]}=-1.8\pm
0.3$ (spectroscopic from Galleti et al. 2005), $-1.8\pm 0.15$ (from the CMD;
Galleti et al. 2006), $-2.14\pm 0.15$ (from the CMD; Mackey et al. 2007), and
$-2.06\pm 0.16$ (spectroscopic from Galleti et al. 2009), which are
consistent. In order to adopt a reasonable value of metallicity for B514, the
mean value of these four independent determinations, i.e. $\rm{[Fe/H]=-1.95}$,
is adopted in this paper.
### 4.2. Stellar populations and synthetic photometry
To determine the age and mass of B514, we compared its SEDs with theoretical
stellar population synthesis models. The SEDs consist of photometric data in
NUV of GALEX, 13 BATC intermediate-band and 2MASS near-infrared $JHK_{s}$
filters obtained in this paper, and of the photometric data in 5 SDSS filters
obtained by Peacock et al. (2010). We will not include the photometric datum
in the FUV band when constraining the age of B514 because of its large
photometric error (2.3 magnitude), i.e. the photometric datum is not accurate.
B514 is a very metal poor GC (see discussions above). So, we use the SSP
models of Bruzual & Charlot (2003) (hereafter BC03), which have been upgraded
from the earlier Bruzual & Charlot (1993, 1996) versions, and now provide the
evolution of the spectra and photometric properties for a wide range of
stellar metallicities. For example, BC03 SSP models based on the Padova-1994
evolutionary tracks include six initial metallicities,
$Z=0.0001,0.0004,0.004,0.008,0.02\,(Z_{\odot})$, and 0.05, corresponding to
${\rm[Fe/H]}=-2.25$, $-1.65$, $-0.64$, $-0.33$, $+0.09$, and $+0.56$. BC03
provides 26 SSP models (both of high and low spectral resolution) using the
Padova-1994 evolutionary tracks, half of which were computed based on the
Salpeter (1955) IMF with lower and upper-mass cut-offs of $m_{\rm
L}=0.1~{}M_{\odot}$ and $m_{\rm U}=100~{}M_{\odot}$, respectively. The other
thirteen were computed using the Chabrier (2003) IMF with the same mass cut-
offs. In addition, BC03 provide 26 SSP models using the Padova-2000
evolutionary tracks which including six partially different initial
metallicities, $Z=0.0004$, 0.001, 0.004, 0.008, 0.019 $(Z_{\odot})$, and 0.03,
i.e., ${\rm[Fe/H]}=-1.65,-1.25,-0.64,-0.33,+0.07$, and $+0.29$. In this paper,
we adopt the high-resolution SSP models using the Padova-1994 evolutionary
tracks to determine the most appropriate age for B514 since its metallicity is
$\rm[Fe/H]=-1.95$, and a Salpeter (1955) IMF is used. These SSP models contain
221 spectra describing the spectral evolution of SSPs from $1.0\times 10^{5}$
yr to 20 Gyr. The evolving spectra include the contribution of the stellar
component at wavelengths from 91Å to $160~{}\mu$m.
Since our observational data are integrated luminosities through a given set
of filters, we convolved the theoretical SSP SEDs of BC03 with the GALEX NUV,
SDSS $ugriz$, BATC $a-n$ and 2MASS $JHK_{\rm s}$ filter response curves to
obtain synthetic optical and NIR photometry for comparison (see Ma et al.,
2009a, b; Wang et al., 2010; Ma et al., 2011, for details).
### 4.3. Fit results
We use a $\chi^{2}$ minimization approach to examine which SSP models are most
compatible with the observed SEDs, following
$\chi^{2}=\sum_{i=1}^{22}{\frac{[m_{\lambda_{i}}^{\rm
intr}-m_{\lambda_{i}}^{\rm mod}(t)]^{2}}{\sigma_{i}^{2}}},$ (4)
where $m_{\lambda_{i}}^{\rm mod}(t)$ is the integrated magnitude in the $i{\rm
th}$ filter of a theoretical SSP at age $t$, $m_{\lambda_{i}}^{\rm intr}$
represents the intrinsic integrated magnitude in the same filter, and
$\sigma_{i}$ is the magnitude uncertainty, defined as
$\sigma_{i}^{2}=\sigma_{{\rm obs},i}^{2}+\sigma_{{\rm mod},i}^{2}+\sigma_{{\rm
md},i}^{2}.$ (5)
Here, $\sigma_{{\rm obs},i}$ is the observational uncertainty from Tables 7
and 8 of this paper, and Table 1 of Peacock et al. (2010), $\sigma_{{\rm
mod},i}$ is the uncertainty associated with the model itself, and
$\sigma_{{\rm md},i}$ is associated with the uncertainty with the distance
modulus adopted here. Charlot et al. (1996) estimated the uncertainty
associated with the term $\sigma_{{\rm mod},i}$ by comparing the colors
obtained from different stellar evolutionary tracks and spectral libraries.
Following Ma et al. (2009a), Ma et al. (2009b), Wang et al. (2010) and Ma et
al. (2011), we adopt $\sigma_{{\rm mod},i}=0.05$ mag. For $\sigma_{{\rm
md},i}$, we adopt 0.07 from McConnachie et al. (2005).
Before fitting, we obtained the the theoretical SEDs for the metallicity
$\rm[Fe/H]=-1.95$ model by interpolation of between ${\rm[Fe/H]}=-2.25$ and
$-1.65$ models.
Since the observed magnitudes in the 2MASS photometric systems are given in
the Vega system, we transformed them to the AB system for our fits. The
photometric offsets in the 2MASS filters between the Vega and AB systems were
obtained based on equations (7) and (8) in the manual provided by Bruzual &
Charlot (2003) (bc03.ps). The best-reduced $\chi^{2}_{\rm min}/\nu=0.8$ is
achieved with an age of $11.5\pm 3.5$ Gyr (1 $\sigma$ uncertainties), $\nu=21$
is the number of free parameters, i.e., the number of observational data
points minus the number of parameters used in the theoretical model. In Figure
6, we show the intrinsic SEDs of B514, the integrated SEDs of the best-fitting
model, and the spectra of the best-fitting model. From Figure 6, we can see
that the BC03 SSP models cannot fit the photometric data point in $H$ band as
well as the other 21 data points, i.e. the observed magnitude is brighter than
the model one in the $H$ band. However, the photometric data point in the $H$
band from Galleti et al. (2005) can be fitted by BC03 SSP models as well as
the other 21 data points, and the fitting result (the age of B514) is in
agreement with one obtained above ($11.5\pm 3.5$ Gyr) within the uncertainty.
much better fitted by BC03 SSP models than the data value derived in this
paper, and the fitting results are consistent within the uncertainty.
### 4.4. Mass of B514
We next determined the mass of B514. The BC03 SSP models are normalized to a
total mass of $1M_{\odot}$ in stars at age $t=0$. The absolute magnitudes (in
the Vega system) in $V$, SDSS $ugriz$ and 2MASS $JHK_{\rm s}$ filters are
included in the BC03 SSP models. The difference between the intrinsic absolute
magnitudes and those given by the model provides a direct measurement of the
cluster mass. To reduce mass uncertainties resulting from photometric
uncertainties based on only magnitudes in one filter (in general the $V$ band
is used), we estimated the mass of B514 using magnitudes in the $V$, $ugriz$
and $JHK_{\rm s}$ bands. The resulting mass determinations for B514 are listed
in Table 9 with their $1\sigma$ uncertainties. From Table 9, we can see that
the mass of B514 obtained based on the magnitudes in different filters is
consistent except for one in the $H$ band. In fact, the observed magnitude is
brighter than the model one in $H$ band (see discussion in §4.3). So, the mass
of B514 derived based on the magnitude in the $H$ band is more massive than
its true one. The mass of B514 derived based on the magnitude in the $H$ band
from Galleti et al. (2005) is in agreement with ones derived based on the
magnitudes in the other 8 bands. Table 9, we know that the obtained mass of
B514 is between $0.96-1.08\times 10^{6}\rm M_{\odot}$ not including one in the
$H$ band. Comparing with 037-B327 [$\mathcal{M}_{\rm 037-B327}\sim 8.5\times
10^{6}\rm M_{\odot}$ (Barmby et al., 2002) or $\mathcal{M}_{\rm 037-B327}\sim
3.0\pm 0.5\times 10^{7}\rm M_{\odot}$ (Ma et al., 2006a)] and G1
[$\mathcal{M}_{\rm G1}\sim(7-17)\times 10^{6}\rm M_{\odot}$ (Meylan et al.,
2001) or $\mathcal{M}_{\rm G1}\sim(5.8-10.6)\times 10^{6}\rm M_{\odot}$ (Ma et
al., 2009b)] in M31 and $\omega$ Cen
[$\mathcal{M}_{\omega{\rm~{}Cen}}\sim(2.9-5.1)\times 10^{6}$M⊙ (Meylan, 2002)]
in the Milky Way, the most massive clusters in the Local Group, B514 is only a
medium-mass globular cluster.
Figure 6.— Best-fitting, integrated theoretical BC03 SEDs compared to the
intrinsic SED of B514. The photometric measurements are shown as symbols with
error bars. Open circles represent the calculated magnitudes of the model SED
for each filter.
## 5\. Summary
In this paper, we determined the structural parameters of one remote globular
cluster B514 known in M31 based on F606W and F814W images obtained with the
ACS/HST. By performing a fit to the surface brightness distribution of a
single-mass isotropic King model, we derive its parameters: the best-fitting
scale radii
$r_{0}=0.36^{+0.09}_{-0.05}~{}\rm{arcsec}~{}(=1.35^{+0.35}_{-0.19}~{}\rm{pc})$
and $0.36^{+0.09}_{-0.06}~{}\rm{arcsec}~{}(=1.35^{+0.32}_{-0.22}~{}\rm{pc})$,
tidal radii
$r_{t}=16.08^{+2.11}_{-1.35}~{}\rm{arcsec}~{}(=61.11^{+8.00}_{-5.14}~{}\rm{pc})$
and
$16.79^{+1.74}_{-1.48}~{}\rm{arcsec}~{}(=63.78^{+6.62}_{-5.64}~{}\rm{pc})$,
and concentration indexes $c=\log(r_{t}/r_{0})=1.66^{+0.05}_{-0.04}$ and
$1.68^{+0.04}_{-0.04}$ in F606W and F814W, respectively; the central surface
brightnesses are $16.25^{+0.57}_{-0.56}$ mag arcsec-2 and
$15.64^{+0.80}_{-0.64}$ mag arcsec-2 in F606W and F814W, respectively; the
half-light, or effective, radius of a model that contains half the total
luminosity in projection, at
$r_{h}=1.31^{+0.14}_{-0.08}~{}\rm{arcsec}~{}(=5.00^{+0.55}_{-0.32}~{}\rm{pc})$
and $1.36^{+0.13}_{-0.09}~{}\rm{arcsec}~{}(=5.17^{+0.48}_{-0.34}~{}\rm{pc})$
in F606W and F814W, respectively. The results show that, the surface
brightness distribution departs from the best-fit King model for
$r>10^{\prime\prime}$. In addition, B514 was observed as part of the BATC
Multicolor Sky Survey, using 13 intermediate-band filters covering a
wavelength range of 3000–80,000 Å. Based on aperture photometry, we obtain its
SEDs as defined by the 13 BATC filters. We determine the cluster’s age by
comparing its SEDs (from 2267 to 20,000Å, comprising photometric data in the
NUV of GALEX, 13 BATC intermediate-band filters, and 5 SDSS filters, and 2MASS
near-infrared $JHK_{\rm s}$ data) with theoretical stellar population
synthesis models, resulting in an age of $11.5\pm 3.5$ Gyr. This age confirms
previous suggestions that B514 is an old GC in M31. B514 has a mass of
$0.96-1.08\times 10^{6}\rm M_{\odot}$, and is a medium-mass globular cluster
in M31.
We would like to thank the anonymous referee for providing rapid and
thoughtful report that helped improve the original manuscript greatly. This
work is partly based on observations made with the NASA/ESA Hubble Space
Telescope, obtained at the Space Telescope Science Institute, which is
operated by AURA, Inc., under NASA contract NAS 5-26555. These observations
are associated with proposal 10394. This work was supported by the Chinese
National Natural Science Foundation grands No. 10873016, 10633020, 10803007,
11003021, and 11073032, and by National Basic Research Program of China (973
Program), No. 2007CB815403.
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Table 1B514: Ellipticity, $\epsilon$, and position angle (P.A.) as a function of the semimajor axis, $a$, in the F606W filter of HST ACS-WFC $a$ | $\epsilon$ | P.A. | $a$ | $\epsilon$ | P.A.
---|---|---|---|---|---
(arcsec) | | (deg) | (arcsec) | | (deg)
0.0260 | $0.360\pm 0.082$ | $159.5\pm 8.3$ | 0.1752 | $0.598\pm 0.046$ | $147.9\pm 3.4$
0.0287 | $0.380\pm 0.087$ | $158.6\pm 8.3$ | 0.1928 | $0.150\pm 0.150$ | $150.4\pm 32.0$
0.0315 | $0.402\pm 0.092$ | $157.7\pm 8.5$ | 0.2120 | $0.150\pm 0.183$ | $107.3\pm 38.4$
0.0347 | $0.426\pm 0.098$ | $156.9\pm 8.6$ | 0.2333 | $0.171\pm 0.111$ | $84.8\pm 20.6$
0.0381 | $0.449\pm 0.104$ | $155.6\pm 8.8$ | 0.2566 | $0.203\pm 0.066$ | $76.4\pm 10.5$
0.0420 | $0.474\pm 0.112$ | $154.3\pm 9.2$ | 0.2822 | $0.199\pm 0.049$ | $79.3\pm 8.0$
0.0461 | $0.503\pm 0.123$ | $153.1\pm 9.7$ | 0.3105 | $0.205\pm 0.062$ | $94.7\pm 9.8$
0.0508 | $0.535\pm 0.134$ | $152.3\pm 10.2$ | 0.3415 | $0.210\pm 0.059$ | $102.4\pm 9.1$
0.0558 | $0.570\pm 0.133$ | $151.6\pm 9.7$ | 0.3757 | $0.197\pm 0.059$ | $112.5\pm 9.5$
0.0614 | $0.596\pm 0.137$ | $150.6\pm 9.7$ | 0.4132 | $0.171\pm 0.066$ | $129.9\pm 12.1$
0.0676 | $0.622\pm 0.146$ | $149.6\pm 10.1$ | 0.4545 | $0.090\pm 0.094$ | $143.8\pm 31.9$
0.0743 | $0.649\pm 0.159$ | $148.8\pm 10.8$ | 0.5000 | $0.090\pm 0.071$ | $20.8\pm 23.6$
0.0818 | $0.678\pm 0.134$ | $147.8\pm 8.8$ | 0.5500 | $0.123\pm 0.041$ | $95.8\pm 10.1$
0.0899 | $0.694\pm 0.114$ | $147.4\pm 7.4$ | 0.6050 | $0.102\pm 0.024$ | $117.4\pm 7.2$
0.0989 | $0.686\pm 0.092$ | $146.9\pm 6.0$ | 0.6655 | $0.119\pm 0.024$ | $98.6\pm 6.2$
0.1088 | $0.650\pm 0.075$ | $145.4\pm 5.0$ | 0.7321 | $0.158\pm 0.032$ | $93.4\pm 6.2$
0.1197 | $0.636\pm 0.059$ | $146.1\pm 4.0$ | 0.8053 | $0.067\pm 0.049$ | $54.6\pm 21.3$
0.1317 | $0.619\pm 0.050$ | $146.6\pm 3.6$ | 0.8858 | $0.065\pm 0.034$ | $4.8\pm 15.5$
0.1448 | $0.636\pm 0.037$ | $146.6\pm 2.7$ | 0.9744 | $0.067\pm 0.096$ | $132.5\pm 42.6$
0.1593 | $0.610\pm 0.064$ | $147.1\pm 4.5$ | | |
Table 2B514: Ellipticity, $\epsilon$, and position angle (P.A.) as a function of the semimajor axis, $a$, in the F814W filter of HST ACS-WFC $a$ | $\epsilon$ | P.A. | $a$ | $\epsilon$ | P.A.
---|---|---|---|---|---
(arcsec) | | (deg) | (arcsec) | | (deg)
0.0260 | $0.296\pm 0.081$ | $162.5\pm 9.5$ | 0.6655 | $0.139\pm 0.026$ | $103.1\pm 5.6$
0.0287 | $0.308\pm 0.084$ | $161.4\pm 9.5$ | 0.7321 | $0.175\pm 0.028$ | $94.5\pm 5.1$
0.0315 | $0.323\pm 0.087$ | $160.4\pm 9.5$ | 0.8053 | $0.116\pm 0.043$ | $35.2\pm 11.4$
0.0347 | $0.342\pm 0.091$ | $159.5\pm 9.6$ | 0.8858 | $0.099\pm 0.044$ | $11.7\pm 13.4$
0.0381 | $0.363\pm 0.097$ | $158.5\pm 9.7$ | 0.9744 | $0.099\pm 0.075$ | $137.4\pm 22.8$
0.0420 | $0.386\pm 0.104$ | $157.6\pm 9.9$ | 1.0718 | $0.074\pm 0.117$ | $146.5\pm 46.9$
0.0461 | $0.412\pm 0.112$ | $156.6\pm 10.1$ | 1.1790 | $0.246\pm 0.144$ | $111.3\pm 19.3$
0.0508 | $0.437\pm 0.117$ | $155.1\pm 10.2$ | 1.2969 | $0.277\pm 0.066$ | $129.1\pm 8.1$
0.0558 | $0.462\pm 0.128$ | $153.6\pm 10.7$ | 1.4266 | $0.134\pm 0.058$ | $93.4\pm 13.2$
0.0614 | $0.488\pm 0.150$ | $152.5\pm 12.0$ | 1.5692 | $0.134\pm 0.062$ | $41.3\pm 14.2$
0.0676 | $0.529\pm 0.146$ | $152.0\pm 11.1$ | 1.7261 | $0.162\pm 0.082$ | $70.6\pm 16.1$
0.0743 | $0.571\pm 0.138$ | $151.5\pm 10.0$ | 1.8987 | $0.086\pm 0.124$ | $144.4\pm 43.2$
0.0818 | $0.606\pm 0.127$ | $151.0\pm 8.9$ | 2.0886 | $0.155\pm 0.073$ | $0.1\pm 14.8$
0.0899 | $0.633\pm 0.129$ | $150.9\pm 8.8$ | 2.2975 | $0.384\pm 0.069$ | $121.7\pm 6.8$
0.0989 | $0.663\pm 0.114$ | $151.1\pm 7.6$ | 2.5272 | $0.131\pm 0.103$ | $146.7\pm 24.2$
0.1088 | $0.697\pm 0.087$ | $151.0\pm 5.6$ | 2.7800 | $0.131\pm 0.074$ | $161.3\pm 17.5$
0.1197 | $0.703\pm 0.077$ | $148.9\pm 5.0$ | 3.0580 | $0.196\pm 0.114$ | $108.0\pm 18.6$
0.1317 | $0.704\pm 0.157$ | $148.3\pm 10.1$ | 3.3638 | $0.130\pm 0.072$ | $117.6\pm 16.9$
0.1448 | $0.702\pm 0.064$ | $149.1\pm 4.4$ | 3.7001 | $0.209\pm 0.075$ | $124.5\pm 11.4$
0.1593 | $0.759\pm 0.055$ | $149.1\pm 3.6$ | 4.0701 | $0.199\pm 0.041$ | $97.7\pm 6.5$
0.1752 | $0.772\pm 0.033$ | $147.9\pm 2.2$ | 4.4772 | $0.067\pm 0.051$ | $110.7\pm 22.6$
0.1928 | $0.726\pm 0.106$ | $147.9\pm 6.9$ | 4.9249 | $0.078\pm 0.048$ | $94.6\pm 18.3$
0.2120 | $0.380\pm 0.149$ | $150.9\pm 14.4$ | 5.4174 | $0.128\pm 0.061$ | $105.3\pm 14.5$
0.2333 | $0.100\pm 0.159$ | $121.0\pm 48.4$ | 5.9591 | $0.011\pm 0.043$ | $115.7\pm 118.0$
0.2566 | $0.125\pm 0.125$ | $113.3\pm 30.5$ | 6.5550 | $0.007\pm 0.059$ | $65.8\pm 244.4$
0.2822 | $0.232\pm 0.125$ | $116.3\pm 17.7$ | 7.2105 | $0.078\pm 0.046$ | $57.5\pm 17.6$
0.3105 | $0.299\pm 0.094$ | $118.4\pm 10.9$ | 7.9316 | $0.096\pm 0.041$ | $89.4\pm 12.8$
0.3415 | $0.265\pm 0.056$ | $118.4\pm 7.2$ | 8.7247 | $0.218\pm 0.042$ | $89.4\pm 6.1$
0.3757 | $0.229\pm 0.064$ | $119.4\pm 9.1$ | 9.5972 | $0.160\pm 0.039$ | $94.7\pm 7.6$
0.4132 | $0.226\pm 0.072$ | $120.4\pm 10.5$ | 10.5569 | $0.333\pm 0.029$ | $87.5\pm 3.0$
0.4545 | $0.146\pm 0.096$ | $139.6\pm 20.5$ | 11.6126 | $0.307\pm 0.023$ | $87.5\pm 2.6$
0.5000 | $0.184\pm 0.246$ | $110.0\pm 42.6$ | 12.7738 | $0.171\pm 0.072$ | $90.6\pm 13.3$
0.5500 | $0.325\pm 0.071$ | $102.3\pm 7.5$ | 14.0512 | $0.234\pm 0.058$ | $82.6\pm 8.1$
0.6050 | $0.127\pm 0.030$ | $108.9\pm 7.3$ | 15.4564 | $0.234\pm 0.064$ | $82.6\pm 7.4$
Table 3B514: Surface brightness, $\mu$, and integrated magnitude, $m$, as a function of the radius in the F606W filter of HST ACS-WFC $R$ | $\mu$ | $m$ | $R$ | $\mu$ | $m$
---|---|---|---|---|---
(arcsec) | ($\rm{mag/arcsec^{2}}$) | (mag) | (arcsec) | ($\rm{mag/arcsec^{2}}$) | (mag)
0.0260 | $16.438\pm 0.034$ | 22.868 | 0.6655 | $18.091\pm 0.055$ | 17.008
0.0287 | $16.445\pm 0.038$ | 22.868 | 0.7321 | $18.212\pm 0.071$ | 16.900
0.0315 | $16.452\pm 0.043$ | 22.868 | 0.8053 | $18.374\pm 0.080$ | 16.804
0.0347 | $16.461\pm 0.048$ | 22.868 | 0.8858 | $18.588\pm 0.079$ | 16.705
0.0381 | $16.470\pm 0.054$ | 22.868 | 0.9744 | $18.799\pm 0.086$ | 16.615
0.0420 | $16.480\pm 0.062$ | 22.868 | 1.0718 | $18.808\pm 0.134$ | 16.531
0.0461 | $16.491\pm 0.070$ | 22.868 | 1.1790 | $18.953\pm 0.289$ | 16.434
0.0508 | $16.503\pm 0.080$ | 21.221 | 1.2969 | $19.283\pm 0.175$ | 16.362
0.0558 | $16.517\pm 0.090$ | 21.221 | 1.4266 | $19.545\pm 0.170$ | 16.283
0.0614 | $16.534\pm 0.102$ | 21.221 | 1.5692 | $19.839\pm 0.141$ | 16.212
0.0676 | $16.553\pm 0.114$ | 21.221 | 1.7261 | $19.890\pm 0.243$ | 16.149
0.0743 | $16.574\pm 0.126$ | 20.628 | 1.8987 | $20.258\pm 0.141$ | 16.087
0.0818 | $16.598\pm 0.135$ | 20.628 | 2.0886 | $20.308\pm 0.232$ | 16.024
0.0899 | $16.623\pm 0.146$ | 20.628 | 2.2975 | $20.770\pm 0.147$ | 15.968
0.0989 | $16.650\pm 0.159$ | 20.628 | 2.5272 | $20.905\pm 0.220$ | 15.917
0.1088 | $16.680\pm 0.156$ | 20.281 | 2.7800 | $21.196\pm 0.162$ | 15.866
0.1197 | $16.708\pm 0.148$ | 19.809 | 3.0580 | $21.565\pm 0.252$ | 15.819
0.1317 | $16.713\pm 0.143$ | 19.809 | 3.3638 | $21.967\pm 0.094$ | 15.785
0.1448 | $16.725\pm 0.132$ | 19.638 | 3.7001 | $22.090\pm 0.274$ | 15.751
0.1593 | $16.764\pm 0.117$ | 19.254 | 4.0701 | $22.568\pm 0.117$ | 15.720
0.1752 | $16.788\pm 0.102$ | 19.254 | 4.4772 | $22.904\pm 0.102$ | 15.695
0.1928 | $16.810\pm 0.090$ | 19.060 | 4.9249 | $23.113\pm 0.115$ | 15.673
0.2120 | $16.858\pm 0.075$ | 18.837 | 5.4174 | $23.359\pm 0.147$ | 15.649
0.2333 | $16.900\pm 0.054$ | 18.652 | 5.9591 | $23.743\pm 0.096$ | 15.622
0.2566 | $16.959\pm 0.049$ | 18.418 | 6.5550 | $24.078\pm 0.101$ | 15.604
0.2822 | $17.028\pm 0.057$ | 18.341 | 7.2105 | $24.331\pm 0.099$ | 15.580
0.3105 | $17.109\pm 0.074$ | 18.150 | 7.9316 | $24.692\pm 0.082$ | 15.560
0.3415 | $17.199\pm 0.076$ | 17.991 | 8.7247 | $25.101\pm 0.116$ | 15.545
0.3757 | $17.285\pm 0.066$ | 17.835 | 9.5972 | $25.328\pm 0.118$ | 15.534
0.4132 | $17.366\pm 0.060$ | 17.662 | 10.5569 | $25.526\pm 0.127$ | 15.518
0.4545 | $17.439\pm 0.064$ | 17.532 | 11.6126 | $25.840\pm 0.164$ | 15.496
0.5000 | $17.490\pm 0.074$ | 17.382 | 12.7738 | $26.212\pm 0.170$ | 15.483
0.5500 | $17.630\pm 0.072$ | 17.249 | 14.0512 | $26.527\pm 0.176$ | 15.469
0.6050 | $17.881\pm 0.058$ | 17.108 | 15.4564 | $26.951\pm 0.233$ | 15.462
Table 4B514: Surface brightness, $\mu$, and integrated magnitude, $m$, as a function of the radius in the F814W filter of HST ACS-WFC $R$ | $\mu$ | $m$ | $R$ | $\mu$ | $m$
---|---|---|---|---|---
(arcsec) | ($\rm{mag/arcsec^{2}}$) | (mag) | (arcsec) | ($\rm{mag/arcsec^{2}}$) | (mag)
0.0260 | $15.810\pm 0.050$ | 22.224 | 0.6655 | $17.332\pm 0.068$ | 16.278
0.0287 | $15.819\pm 0.055$ | 22.224 | 0.7321 | $17.451\pm 0.077$ | 16.166
0.0315 | $15.828\pm 0.061$ | 22.224 | 0.8053 | $17.601\pm 0.105$ | 16.068
0.0347 | $15.839\pm 0.068$ | 22.224 | 0.8858 | $17.837\pm 0.096$ | 15.966
0.0381 | $15.850\pm 0.076$ | 22.224 | 0.9744 | $18.048\pm 0.098$ | 15.875
0.0420 | $15.862\pm 0.085$ | 22.224 | 1.0718 | $18.007\pm 0.179$ | 15.791
0.0461 | $15.874\pm 0.095$ | 22.224 | 1.1790 | $18.194\pm 0.245$ | 15.686
0.0508 | $15.888\pm 0.107$ | 20.601 | 1.2969 | $18.447\pm 0.336$ | 15.613
0.0558 | $15.903\pm 0.118$ | 20.601 | 1.4266 | $18.755\pm 0.190$ | 15.533
0.0614 | $15.917\pm 0.126$ | 20.601 | 1.5692 | $19.082\pm 0.151$ | 15.460
0.0676 | $15.931\pm 0.135$ | 20.601 | 1.7261 | $19.078\pm 0.335$ | 15.394
0.0743 | $15.947\pm 0.144$ | 20.013 | 1.8987 | $19.444\pm 0.186$ | 15.328
0.0818 | $15.965\pm 0.152$ | 20.013 | 2.0886 | $19.540\pm 0.232$ | 15.261
0.0899 | $15.982\pm 0.160$ | 20.013 | 2.2975 | $19.965\pm 0.199$ | 15.203
0.0989 | $15.999\pm 0.169$ | 20.013 | 2.5272 | $20.101\pm 0.259$ | 15.150
0.1088 | $16.015\pm 0.169$ | 19.656 | 2.7800 | $20.402\pm 0.164$ | 15.096
0.1197 | $16.034\pm 0.159$ | 19.170 | 3.0580 | $20.718\pm 0.310$ | 15.045
0.1317 | $16.024\pm 0.157$ | 19.170 | 3.3638 | $21.153\pm 0.126$ | 15.012
0.1448 | $16.022\pm 0.151$ | 18.989 | 3.7001 | $21.253\pm 0.322$ | 14.975
0.1593 | $16.055\pm 0.137$ | 18.589 | 4.0701 | $21.778\pm 0.116$ | 14.944
0.1752 | $16.072\pm 0.128$ | 18.589 | 4.4772 | $22.117\pm 0.088$ | 14.920
0.1928 | $16.078\pm 0.118$ | 18.384 | 4.9249 | $22.297\pm 0.119$ | 14.898
0.2120 | $16.120\pm 0.098$ | 18.153 | 5.4174 | $22.484\pm 0.174$ | 14.874
0.2333 | $16.171\pm 0.070$ | 17.959 | 5.9591 | $22.916\pm 0.113$ | 14.845
0.2566 | $16.225\pm 0.059$ | 17.715 | 6.5550 | $23.213\pm 0.123$ | 14.826
0.2822 | $16.287\pm 0.070$ | 17.639 | 7.2105 | $23.470\pm 0.125$ | 14.800
0.3105 | $16.367\pm 0.092$ | 17.441 | 7.9316 | $23.880\pm 0.088$ | 14.778
0.3415 | $16.454\pm 0.098$ | 17.275 | 8.7247 | $24.263\pm 0.121$ | 14.763
0.3757 | $16.558\pm 0.089$ | 17.116 | 9.5972 | $24.457\pm 0.132$ | 14.752
0.4132 | $16.646\pm 0.067$ | 16.944 | 10.5569 | $24.607\pm 0.158$ | 14.735
0.4545 | $16.708\pm 0.070$ | 16.814 | 11.6126 | $24.935\pm 0.155$ | 14.710
0.5000 | $16.740\pm 0.082$ | 16.659 | 12.7738 | $25.294\pm 0.197$ | 14.697
0.5500 | $16.878\pm 0.098$ | 16.524 | 14.0512 | $25.649\pm 0.231$ | 14.681
0.6050 | $17.137\pm 0.079$ | 16.380 | 15.4564 | $26.034\pm 0.379$ | 14.674
Table 5B514: Surface brightness profiles $\mu$ from star counts $R$ | $\mu_{\rm{F606W}}$ | $R$ | $\mu_{\rm{F814W}}$
---|---|---|---
(arcsec) | ($\rm{mag/arcsec^{2}}$) | (arcsec) | ($\rm{mag/arcsec^{2}}$)
| | 7.9625 | $23.831\pm 0.066$
9.1875 | $24.961\pm 0.072$ | 9.1875 | $24.181\pm 0.072$
10.4125 | $25.297\pm 0.079$ | 10.4125 | $24.517\pm 0.079$
11.6375 | $25.622\pm 0.087$ | 11.6375 | $24.842\pm 0.087$
12.8625 | $25.905\pm 0.095$ | 12.8625 | $25.125\pm 0.095$
14.0875 | $26.361\pm 0.111$ | 14.0875 | $25.581\pm 0.111$
15.3125 | $26.694\pm 0.125$ | 15.3125 | $25.914\pm 0.125$
16.5375 | $26.998\pm 0.138$ | 16.5375 | $26.218\pm 0.138$
17.7625 | $27.310\pm 0.154$ | 17.7625 | $26.530\pm 0.154$
18.9875 | $27.360\pm 0.152$ | 18.9875 | $26.580\pm 0.152$
20.2125 | $27.517\pm 0.158$ | 20.2125 | $26.737\pm 0.158$
21.4375 | $27.703\pm 0.168$ | 21.4375 | $26.923\pm 0.168$
22.6625 | $27.816\pm 0.172$ | 22.6625 | $27.036\pm 0.172$
23.8875 | $28.050\pm 0.186$ | 23.8875 | $27.270\pm 0.186$
25.1125 | $28.240\pm 0.198$ | 25.1125 | $27.460\pm 0.198$
26.3375 | $28.292\pm 0.198$ | 26.3375 | $27.512\pm 0.198$
27.5625 | $28.630\pm 0.226$ | 27.5625 | $27.850\pm 0.226$
28.7875 | $28.725\pm 0.231$ | 28.7875 | $27.945\pm 0.231$
30.0125 | $28.722\pm 0.226$ | 30.0125 | $27.942\pm 0.226$
31.2375 | $28.864\pm 0.237$ | 31.2375 | $28.084\pm 0.237$
32.4625 | $29.201\pm 0.271$ | 32.4625 | $28.421\pm 0.271$
33.6875 | $29.242\pm 0.271$ | 33.6875 | $28.462\pm 0.271$
34.9125 | $29.350\pm 0.280$ | 34.9125 | $28.570\pm 0.280$
36.1375 | $29.543\pm 0.301$ | 36.1375 | $28.763\pm 0.301$
37.3625 | $29.761\pm 0.327$ | 37.3625 | $28.981\pm 0.327$
38.5875 | $30.014\pm 0.362$ | 38.5875 | $29.234\pm 0.362$
39.8125 | $30.321\pm 0.410$ | 39.8125 | $29.541\pm 0.410$
41.0375 | $30.208\pm 0.384$ | 41.0375 | $29.428\pm 0.384$
Table 6 Structural parameters of B514 Parameters | F606W | F814W
---|---|---
$r_{0}$ | $0.36^{+0.09}_{-0.05}~{}\rm{arcsec}~{}(=1.35^{+0.35}_{-0.19}~{}\rm{pc})$ | $0.36^{+0.09}_{-0.06}~{}\rm{arcsec}~{}(=1.35^{+0.32}_{-0.22}~{}\rm{pc})$
$r_{t}$ | $16.08^{+2.11}_{-1.35}~{}\rm{arcsec}~{}(=61.11^{+8.00}_{-5.14}~{}\rm{pc})$ | $16.79^{+1.74}_{-1.48}~{}\rm{arcsec}~{}(=63.78^{+6.62}_{-5.64}~{}\rm{pc})$
$c=\log(r_{t}/r_{0})$ | $1.66^{+0.05}_{-0.04}$ | $1.68^{+0.04}_{-0.04}$
$r_{h}$ | $1.31^{+0.14}_{-0.08}~{}\rm{arcsec}~{}(=5.00^{+0.55}_{-0.32}~{}\rm{pc})$ | $1.36^{+0.13}_{-0.09}~{}\rm{arcsec}~{}(=5.17^{+0.48}_{-0.34}~{}\rm{pc})$
$\mu_{0}$ (${\rm mag~{}arcsec^{-2}}$) | $16.25^{+0.57}_{-0.56}$ | $15.64^{+0.80}_{-0.64}$
Table 7BATC photometry of B514 Filter | Central wavelength | Bandwidth | Number of images | Exposure time | rms | Magnitude
---|---|---|---|---|---|---
| (Å) | (Å) | | (hours) | (mag) |
$a$ | 3360 | 222 | 6 | 2:00 | 0.010 | $17.59\pm 0.05$
$b$ | 3890 | 187 | 6 | 2:00 | 0.010 | $16.82\pm 0.02$
$c$ | 4210 | 185 | 4 | 1:00 | 0.002 | $16.52\pm 0.01$
$d$ | 4550 | 222 | 4 | 1:20 | 0.015 | $16.25\pm 0.02$
$e$ | 4920 | 225 | 3 | 1:00 | 0.007 | $16.05\pm 0.01$
$f$ | 5270 | 211 | 3 | 1:00 | 0.014 | $15.85\pm 0.02$
$g$ | 5795 | 176 | 3 | 1:00 | 0.010 | $15.64\pm 0.01$
$h$ | 6075 | 190 | 3 | 0:50 | 0.005 | $15.56\pm 0.01$
$i$ | 6660 | 312 | 3 | 0:50 | 0.004 | $15.44\pm 0.01$
$j$ | 7050 | 121 | 3 | 1:00 | 0.006 | $15.33\pm 0.01$
$k$ | 7490 | 125 | 3 | 1:00 | 0.011 | $15.25\pm 0.01$
$m$ | 8020 | 179 | 3 | 1:00 | 0.003 | $15.19\pm 0.01$
$n$ | 8480 | 152 | 3 | 1:00 | 0.005 | $15.12\pm 0.01$
Table 82MASS and $GALEX$ photometry of B514 Filter | Magnitude
---|---
$J$ | $14.23\pm 0.07$
$H$ | $13.32\pm 0.06$
$K_{s}$ | $13.63\pm 0.10$
NUV | $19.57\pm 0.86$
FUV | $20.44\pm 2.30$
Table 9Mass estimates (and uncertainties) of B514 based on the BC03 models $V$ | $u$ | $g$ | $r$ | $i$ | $z$ | $J$ | $H$ | $K_{\rm s}$
---|---|---|---|---|---|---|---|---
| Mass $(10^{6}~{}M_{\odot})$ |
$1.08\pm 0.03$ | $0.99\pm 0.03$ | $1.06\pm 0.03$ | $1.02\pm 0.03$ | $1.0\pm 0.03$ | $0.98\pm 0.03$ | $0.98\pm 0.06$ | $1.41\pm 0.08$ | $0.96\pm 0.09$
|
arxiv-papers
| 2011-11-10T01:33:41 |
2024-09-04T02:49:24.186116
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jun Ma (1,2), Song Wang (1,3), Zhenyu Wu (1), Zhou Fan (1), Tianmeng\n Zhang (1) and Jianghua Wu (1) ((1) National Astronomical Observatories,\n Chinese Academy of Sciences, (2) Key Laboratory of Optical Astronomy,\n National Astronomical Observatories, Chinese Academy of Sciences, Beijing,\n China, (3) Graduate University of Chinese Academy of Sciences, Shijingshan\n District, Beijing, China)",
"submitter": "Jun Ma",
"url": "https://arxiv.org/abs/1111.2380"
}
|
1111.2620
|
# Search for the rare decays $B^{0}_{(s)}\to\mu^{+}\mu^{-}$ at LHCb
CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
E-mail On behalf of the LHCb collaboration
###### Abstract:
A search for the $B_{s}\to\mu^{+}\mu^{-}$ and $B_{d}\to\mu^{+}\mu^{-}$ decays
is presented using $\sim 300$ pb-1 of $pp$ collisions at $\sqrt{s}$ = 7 TeV
collected by the LHCb experiment at the Large Hadron Collider at CERN. The
measured upper limit for the branching ratio of the $B_{s}\to\mu^{+}\mu^{-}$
decay is $\cal B$($B_{s}\to\mu^{+}\mu^{-}$)$<1.3\ (1.6)\times 10^{-8}$ at 90 %
(95 %) confidence level (CL), while in the case of the
$B_{d}\to\mu^{+}\mu^{-}$ decay the measured upper limit is $\cal
B$($B_{d}\to\mu^{+}\mu^{-}$)$<4.2\ (5.1)\times 10^{-9}$ at 90 % (95 %) CL. A
combination with the 2010 dataset results in $\cal
B$($B_{s}\to\mu^{+}\mu^{-}$)$<1.2\ (1.5)\times 10^{-8}$ at 90 % (95 %) CL.
## 1 Introduction
Measurements at low energies may provide interesting indirect constraints on
the masses of particles that are too heavy to be produced directly. This is
particularly true for Flavour Changing Neutral Currents (FCNC) processes which
are highly suppressed in the Standard Model (SM) and can only occur through
higher order diagrams. The SM prediction for the branching ratios ($\cal B$)
of the FCNC decays $B_{s}\to\mu^{+}\mu^{-}$ and $B_{d}\to\mu^{+}\mu^{-}$ have
been computed to be $(3.2\pm 0.2)\times 10^{-9}$ and $(0.10\pm 0.01)\times
10^{-9}$ respectively [1]. However New Physics (NP) contributions can
significantly enhance these values.
The best published limits from the Tevatron at $95\%$ CL are obtained using
6.1 fb-1 by the D0 collaboration [2], and using 2 fb-1 by the CDF
collaboration [3]. The CDF collaboration has also presented a preliminary
result [4] with 6.9 fb-1 in which an excess of
$B_{s}\to\mu^{+}\mu^{-}$candidates is reported, compatible with a $\cal
B$($B_{s}\to\mu^{+}\mu^{-}$)= $(1.8^{+1.1}_{-0.9})\times 10^{-8}$.
The LHCb collaboration has previously obtained the limits $\cal
B$($B_{s}\to\mu^{+}\mu^{-}$)$<5.4\times 10^{-8}$ and $\cal
B$($B_{d}\to\mu^{+}\mu^{-}$)$<1.5\times 10^{-8}$ at 95$\%$ CL based on 37 pb-1
of luminosity collected in the 2010 run [5]. We present here a measurement
based on 300 pb-1 of integrated luminosity collected between March and June
2011.
## 2 Analysis strategy
The general structure of the analysis is similar to the one described in Ref.
[5] and is detailed in Ref. [6].
The selection procedure treats signal and control/normalization channels in
the same way in order to minimize the systematic uncertainties. Assuming the
SM branching ratio and the $b\bar{b}$ cross-section, measured within the LHCb
acceptance, of $\sigma_{b\overline{b}}=75\pm 14\,\mu$b [7], approximately
$3.4~{}(0.32)$ $B^{0}_{s}\to\mu^{+}\mu^{-}$ ($B^{0}\to\mu^{+}\mu^{-}$) events
are expected to be reconstructed and selected in the analysed sample.
After the selection, each event is given a probability to be signal or
background in a two-dimensional space defined by two independent likelihoods:
the invariant mass and the output of a Boosted Decision Tree (BDT) from the
TMVA package [8]. The combination of variables entering the BDT is optimized
using Monte Carlo (MC) simulation. The following variables have been used: the
$B$ lifetime, impact parameter and transverse momentum of the $B$, the minimum
impact parameter significance of the muons, the distance of closest approach
between the two muons, the degree of isolation of the two muons with respect
to any other track in the event, the cosine of the polarization angle, the $B$
isolation and the minimum $p_{T}$ of the muons. The BDT distribution is then
transformed in order to be flat for the signal and peaked at 0 for the
background.
The calibration of the invariant mass and the BDT likelihoods are obtained
from data using control samples. The signal BDT shape is obtained from
$B^{0}_{(s)}\to h^{+}h^{-}$ events free from trigger biases while the
background shape is obtained using sideband $B^{0}_{(s)}\to\mu^{+}\mu^{-}$
candidates. The resulting distributions are shown in Fig. 1.
Figure 1: BDT calibration for signal and background.
The parameters describing the invariant mass line shape of the signal are
extracted from data using control samples. The average mass values are
obtained from $B^{0}\to K^{+}\pi^{-}$ and $B^{0}_{s}\to K^{+}K^{-}$ exclusive
samples. The $B^{0}_{s}$ and $B^{0}$ mass resolutions are estimated by
interpolating the ones obtained with the dimuon resonances ($J/\psi,\psi(2S)$
and $\Upsilon(1S,2S,3S)$) and cross-checked via a fit to the invariant mass
distribution of the $B^{0}_{(s)}\to h^{+}h^{-}$ inclusive decays and of the
$B^{0}\to K^{+}\pi^{-}$ exclusive decay. The interpolation yields
$\sigma(B)=24.6\pm 0.2_{\rm stat}\pm 1.0_{\rm syst}$.
The number of expected signal events is obtained by normalizing to channels of
known branching ratios, $B^{+}\\!\to J/\psi K^{+}$, $B^{0}_{s}\\!\to
J/\psi\phi$, and $B^{0}\\!\to K^{+}\pi^{-}$, that are selected in a way as
similar as possible to the signal.
The probability for a background event to have a given BDT and invariant mass
value is obtained by a fit of the mass distribution of events in the mass
sidebands, in bins of BDT. Different fit functions and mass ranges are used to
compute the systematics uncertainties. The two-dimensional space formed by the
invariant mass and BDT is binned, and for each bin we compute how many events
are observed in data, how many signal events are expected for a given $\cal B$
hypothesis and luminosity, and how many background events are expected for a
given luminosity. The compatibility of the observed distribution of events in
all bins with the expected one for a given $\cal B$ hypothesis is computed
using the CLs method [9], which allows to exclude a given hypothesis at a
given confidence level.
In order to avoid unconscious biases, the data in the mass region defined by
$M_{B^{0}}-60\,{\rm MeV}/c^{2}$ and $M_{B^{0}_{s}}+60\,{\rm MeV}/c^{2}$ have
been blinded until the completion of the analysis.
## 3 Results
Figure 2: Distribution of selected dimuon events in the invariant mass vs BDT
plane. The orange short-dashed (green long-dashed) lines indicate the $\pm
60\,{\rm MeV}/c^{2}$ search window around the $B^{0}_{s}$ ($B^{0}$).
The distribution of events in the invariant mass versus BDT plane is reported
Fig. 2. The expected limit at 90 (95) % CL for the $B_{s}\to\mu^{+}\mu^{-}$ is
$0.8~{}(1.0)\times 10^{-8}$ in the case of background only hypothesis. When
adding signal events according to the SM branching fraction, these limits
become $1.2~{}(1.5)\times 10^{-8}$. The observed values for the
$B_{s}\to\mu^{+}\mu^{-}$channel is $1.3~{}(1.6)\times 10^{-8}$ with a CLb
value of 0.80. The observed events are in good agreement with the background
expectations and the presence of $B_{s}\to\mu^{+}\mu^{-}$events according to
SM predictions.
For the $B_{d}\to\mu^{+}\mu^{-}$, the expected limit at 90 (95) % CL is
$2.4~{}(3.1)\times 10^{-9}$ in the case of background only hypothesis. The
observed values is $4.2~{}(5.2)\times 10^{-9}$ with a CLb value of 0.79. The
comparison of the observed distribution of events with the expected background
distribution results in a p-value (1-$\textrm{CL}_{\textrm{b}}$) of 20 % (21
%) for the $B_{s}\to\mu^{+}\mu^{-}$ ($B_{d}\to\mu^{+}\mu^{-}$) decays. In the
case of $B_{d}\to\mu^{+}\mu^{-}$, the slightly low p-value is due to an excess
of the observed events in the most sensitive BDT bin with respect to the
background expectations. A larger data sample will allow to clarify the
situation. In the case of $B_{s}\to\mu^{+}\mu^{-}$, when a signal is included
at the level expected in the Standard Model, the p-value increases to 50 %.
Finally, the $B_{s}\to\mu^{+}\mu^{-}$ limit is combined with the one published
from the 2010 data to obtain $\cal B$($B_{s}\to\mu^{+}\mu^{-}$) ¡
$1.2~{}(1.5)\times 10^{-8}$ at 90 % (95 %) CL. This 90 % CL upper limit is
still 3.8 times above the standard model prediction.
## References
* [1] A.J. Buras, G. Isidori and P. Paradisi, “EDMs vs CPV in $B_{s,d}$ mixing in two Higgs doublet models with MFV”, arXiv:1007.5291 [hep-ph] (2010).
* [2] V. Abazov et al. [D0 Collaboration], “Search for the rare decay $B^{0}_{s}\to\mu^{+}\mu^{-}$”, Phys. Lett. B 693 (2010) 539.
* [3] T. Aaltonen et al. [CDF Collaboration], “Search for $B_{s}\to\mu^{+}\mu^{-}$and $B_{d}\to\mu^{+}\mu^{-}$Decays with 2$\mbox{\,fb}^{-1}$ of $p\bar{p}$ Collisions” Phys. Rev. Lett. 100 (2008) 101802.
* [4] T. Aaltonen et al. [CDF Collaboration], “Search for $B_{s}\to\mu^{+}\mu^{-}$and $B_{d}\to\mu^{+}\mu^{-}$Decays with CDF-II” , arXiv:1107.2304v1 [hep-ex] (2011).
* [5] R. Aaji et al. [LHCb Collaboration], “Search for the rare decays $B^{0}_{s}\to\mu^{+}\mu^{-}$ and $B^{0}\to\mu^{+}\mu^{-}$”, Phys. Lett. B 699 (2011) 330.
* [6] R. Aaji et al. [LHCb Collaboration], LHCb-CONF-2011-037.
* [7] R. Aaij et al. [LHCb Collaboration], “Measurement of $\sigma$(pp $\rightarrow b\bar{b}$ X) at $\sqrt{s}$=7 TeV in the forward region”, Phys. Lett. B 694 (2010) 209.
* [8] TMVA, Toolkit for Multi Variate analysis with ROOT, http://tmva.sourceforge.net.
* [9] A. Read, “Presentation of Search Results: The $\textrm{CL}_{\textrm{s}}$ Technique”, J. Phys. G 28 (2002) 2693.
|
arxiv-papers
| 2011-11-10T21:38:38 |
2024-09-04T02:49:24.201240
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Justine Serrano",
"submitter": "Justine Serrano",
"url": "https://arxiv.org/abs/1111.2620"
}
|
1111.2724
|
Temperature anisotropy and differential streaming of solar wind ions –
Correlations with transverse fluctuations
Sofiane Bourouaine1, Eckart Marsch1 and Fritz M. Neubauer2
1 Max-Planck-Institut für Sonnensystemforschung, 37191 Katlenburg-Lindau,
Germany 2 Institut für Geophysik und Meteorologie, Universität zu Köln,
Albertus-Magnus-Platz, Köln, 50923, Germany Abstract. We study correlations of
the temperature ratio (which is an indicator for perpendicular ion heating)
and the differential flow of the alpha particles with the power of transverse
fluctuations that have wave numbers between $0.01$ and $0.1$ (normalized to
$k_{p}=1/l_{p}$, where $l_{p}$ is the proton inertial length). We found that
both the normalized differential ion speed, $V_{\alpha p}/V_{\mathrm{A}}$
(where $V_{\mathrm{A}}$ is the Alfvén speed) and the proton temperature
anisotropy, $T_{\perp p}/T_{\parallel p}$, increase when the relative wave
power is growing. Furthermore, if the normalized differential ion speed stays
below 0.5, the alpha-particle temperature anisotropy,
$T_{\perp\alpha}/T_{\parallel\alpha}$, correlates positively with the relative
power of the transverse fluctuations. However, if $V_{\alpha
p}/V_{\mathrm{A}}$ is higher than 0.6, then the alpha-particle temperature
anisotropy tends to become lower and attain even values below unity despite
the presence of transverse fluctuations of relatively high amplitudes. Our
findings appear to be consistent with the expectations from kinetic theory for
the resonant interaction of the ions with Alfvén/ion-cyclotron waves and the
resulting wave dissipation.
## 1 Introduction
The in-situ measurements of protons and heavier ions in fast solar wind
revealed distinct non-thermal kinetic features, such as the core temperature
anisotropy and beam of the protons, or the preferential heating and
acceleration of alpha particles and other minor species (see, e.g., the review
of Marsch (2006) on this subject). Similarly, solar remote-sensing
observations of coronal holes, which are known as the sources of the fast
solar wind provided evidence via significant broadenings of ultraviolet
emission lines, for strong perpendicular heating of oxygen and heavy ions in
the solar corona (see, e.g., the review of Cranmer (2009)).
These observed kinetic features have usually been interpreted as signatures of
heating by ion-cyclotron waves (Isenberg et al. 2001; Marsch and Tu 2001;
Hollweg and Isenberg 2002; Matteini et al. 2007). Their origin remains
unclear, though, and is still a matter of debate. Some authors suggested that
a nonlinear parallel cascade (via parametric decay) of low-frequency Alfvén
waves may ultimately generate ion-cyclotron waves, since their numerical
simulations showed that protons and heavy ions can be heated perpendicularly
by wave absorption through cyclotron resonance (Araneda et al. 2008, 2009).
Moreover, it has been argued theoretically that in low-beta plasma condition,
the ion temperature anisotropy and the preferential (stochastic) heating of
heavy species could even be caused by Alfvénic fluctuations with frequencies
well below the local ion-cyclotron frequency (Wu and Yoon 2007; Bourouaine et
al. 2008; Chandran 2010).
On the other hand, it has also been argued that the energy of Alfvénic
turbulence does not cascade effectively to high-frequency parallel cyclotron
waves, but rather is transferred into low-frequency and highly oblique kinetic
Alfvén waves. Then dissipation would still take place at the proton gyroscale,
but via Landau damping acting mostly on electrons (Howes et al. 2008).
In this letter, while benefiting from the high-resolution magnetic field data
and the detailed proton and alpha-particle velocity distribution functions
(VDFs) obtained (Marsch et al. 1982) simultaneously by the Helios 2 wave and
plasma experiments, we analysed typical solar wind data from a heliocentric
distance of about 0.7 AU. While studying them in detail we found that the
prominent kinetic features of the alpha particles and protons are closely
correlated with the power of the transverse high-frequency waves. The results
of our study place new empirical constraints on solar wind models and further
limit the theoretical assumptions involved. We will first present a
statistical analysis of the relevant plasma and field parameters, and then
discuss them in the light of kinetic plasma theory.
## 2 Observations
Here we analyse plasma and magnetic field data provided by the Helios 2
spacecraft in 1976. For the purpose of a solid statistical study, we selected
a week of continuous measurements, i.e., the days with DOY (day of year)
numbers 67 to 73, on which the spacecraft was at a solar distance of about 0.7
AU.
As is obvious from Figure 1, the selected set covers data from measurements
made in slow and fast solar wind. The slightly inclined dashed line in that
figure separates the data obtained in fast solar wind from those in slow wind.
As we would expect, the domain with a relatively high number of collisions
mainly corresponds to slow solar wind (with $A_{c}>0.2$), and that with a low
number of collisions corresponds to fast solar wind which is weakly
collisional (with $A_{c}<0.1$). Here collisionality is quantified by the so-
called collision age (see e.g., Salem et al. 2003),
$A_{c}=l/(V_{\mathrm{sw}}\tau_{\alpha p})$, where, $l$ is the distance from
the sun and $\tau_{\alpha p}$ is the collision-exchange time between helium
ions and protons.Moreover, for this data set the proton parallel plasma beta,
$\beta_{\parallel p}$, is comparatively small and mainly varies between 0.1
and 0.7.
Because the solar wind speed is clearly higher than the Alfvén speed, we can
safely assume the so-called Taylor hypothesis to be valid. Therefore, the
spacecraft frequency of the magnetic field fluctuations (obtained from fast
Fourier transform) is simply given by, $2\pi f_{\mathrm{sc}}\simeq
kV_{\mathrm{sw}}$, where $V_{\mathrm{sw}}$ is the solar wind speed, and $k$ is
the wave number of the magnetic field fluctuations. Here, we deal with the
average wave power density, $\delta B^{2}$, which is obtained in the
spacecraft frame by integration over a frequency range that corresponds to the
wave-number range (0.01–0.1)$k_{p}$ (where $k_{\rho}=1/l_{p}$, $l_{p}$ is the
proton inertial length). Therefore, we are dealing with turbulent fluctuations
that still belong to the inertial range but are near the dissipation range of
the turbulence. We have chosen this frequency domain to avoid the inclusion of
higher-frequency fluctuations that might stem from local wave excitation
(e.g., owing to a large temperature anisotropy or high normalized ion
differential speed) below but near the proton gyro-kinetic scale. Consistently
with our choice of this narrow frequency range, we define the mean field
(required as the reference value for the superposed fluctuations) by the
average of the full magnetic-field vector over the short time period of about
one hundred times $l_{p}/V_{\mathrm{sw}}$.
Figure 1: Distributions of the number of data (presented in colour with the
coding indicated by the right-hand bar), which are plotted as a function of
the collision age, $A_{c}$, and the solar wind speed, $V_{sw}$. The inclined
dashed line separates the data from fast and slow solar wind. Figure 2:
Distributions (presented in colour with coding indicated by the right-hand
bars) of the number of our data (a) and the compressibility, $\delta
B_{\parallel}^{2}/(\delta B_{\perp}^{2}+\delta B_{\parallel}^{2})$ (b),
plotted in the parameter plane of the normalized wave power versus the proton
temperature anisotropy, $T_{\perp p}/T_{\parallel p}$. (c) Weighted mean value
of the proton temperature anisotropy for $A_{c}<0.1$ (black), and $A_{c}>0.2$
(green) displayed versus normalized wave power. The short vertical bars
indicate the small uncertainties of the weighted mean values.
In Figure 2 a the distribution of the normalized wave power, $\delta
B^{2}/B^{2}_{0}$, is plotted versus the proton temperature anisotropy,
$T_{\perp p}/T_{\parallel p}$. It turns out that the latter parameter ranges
between 0.4 and 2 and correlates positively with the normalized wave power.
According to Figure 2 b, the relative compressive component (coded in colour)
of the fluctuations is fairly small and not higher than $5\%$ of the overall
wave amplitude, which means that the fluctuations are mainly transverse and
essentially incompressible. It also appears that there is no striking
correlation between magnetic compressibility and normalized wave power.
Unfortunately, we do not have the simultaneous small-scale velocity-
fluctuations owing to the lack of an adequate time resolution of the plasma
experiment, and therefore we cannot scrutinize the nature of the magnetic
fluctuations or confirm by polarization study whether they are Alfvénic or
not. But presumably they are Alfvén waves, because these waves with periods
longer than 40 s are observed widely in fast solar wind (Tu and Marsch 1995)
streams at all distances between 0.3 and 1 AU.
In Figure 2c we plot the weighted mean value (expectation value) of the proton
temperature anisotropy as a function of the normalized wave power for
different solar wind regimes. In fast solar wind (i.e., where $A_{c}<0.1$,
black curve), the proton temperature anisotropy is below unity whenever the
value of the normalized wave power is below 0.01. However, above 0.01 the
temperature anisotropy steeply increases until the normalized wave power
reaches a value of about 0.1, where the temperature anisotropy tends to
saturate at a value of about 1.5. In contrast, when collisions are relatively
strong (i.e., where $A_{c}>0.2$, green curve), the proton temperature
anisotropy also increases with the wave amplitude but reaches lower values
than those when collisions are weak. The temperature anisotropy ranges between
0.6 and about 1.1, while the normalized wave power varies between 0.01 and
0.1.
Our study shows that the proton temperature anisotropy correlates positively
with the normalized wave power.Moreover, a positive correlation between the
wave power and the proton temperature anisotropy was found earlier by
Bourouaine et al. (2010), but for another set of Helios data restricted to
fast solar wind and at various solar distances. However, the Helios
measurements selected for the present study were made at a fixed heliocentric
distance which is about 0.7 AU. Therefore, the new correlation established
here between the proton temperature anisotropy and normalized wave power is
certainly of local nature and not affected by radial evolution.
According to Figure 2c, it seems that collisions play a major role in reducing
the effects waves have on the local heating of the protons. Therefore, protons
can become heated (if wave energy is available) more easily in regions with
relatively weak collisions than in regions where collisions are relatively
strong.
Figure 3 a shows the data distribution of the normalized wave power versus the
normalized ion differential speed. It hardly exceeds unity, a result which is
consistent with the prediction of kinetic theory for a linear plasma
instability. The theory says that, whenever the ion differential speed exceeds
the Alfvén speed, the plasma should become unstable and excite magnetosonic
waves (see, e.g., Li and Habbal (2000)). As Bourouaine et al. (2011) found
previously, the present study also indicates a positive correlation between
the normalized ion differential speed and the normalized wave power (see the
white curve, which represents the weighted mean values of the normalized wave
power shown in Figure 3 a. )
Figure 3: Distributions (presented in colour with coding indicated by the
right-hand bars) of the number of our data (a), the collision age, $A_{c}$ (b)
and the ratio between the temperature anisotropies of alpha particles and
protons, $(T_{\perp\alpha}T_{\parallel p})/(T_{\parallel\alpha}T_{\perp p})$,
(c). All parameter are plotted in the plane of the normalized wave-power,
$\delta B^{2}/B_{0}^{2}$, versus normalized ion differential speed, $V_{\alpha
p}/V_{\mathrm{A}}$; (d) Weighted mean values of the collisional age, $A_{c}$,
and the ratio between the temperature anisotropies of alpha particles and
protons, $(T_{\perp\alpha}T_{\parallel p})/(T_{\parallel\alpha}T_{\perp p})$,
displayed versus normalized ion differential speed, $V_{\alpha
p}/V_{\mathrm{A}}$. The white dots in panel (a) represent the weighted mean
values. The vertical dotted line in panel (d) separates the regimes of fast
and slow solar wind.
It is commonly believed (see, e.g., the review of Marsch (2006)) that low
collisional friction would permit a relatively high differential speed to
occur between the two main ion species in the solar wind. This notion is
confirmed by the results of our Figure 3 b, which shows that higher values of
the normalized ion differential speed correspond to lower collision ages. In
slow solar wind, when the value of the collision age is higher than 0.2, the
corresponding normalized differential speed is low (i.e., $V_{\alpha
p}/V_{A}<0.3$), but it is higher than 0.3 for comparatively low values of the
collision age (as is indicated later in Figure 3 d).
In Figure 3c the coloured pixels represent the ratio of the alpha-particle-to-
proton temperature anisotropy, $(T_{\perp\alpha}T_{\parallel p})/(T_{\perp
p}T_{\parallel\alpha})$, plotted as a function of the relative ion
differential speed and the normalized wave power. Interestingly, this figure
clearly shows that when $V_{\alpha p}/V_{A}\leq 0.4$, the temperature
anisotropy of the alpha particles, $T_{\perp\alpha}/T_{\parallel\alpha}$, is
higher than the anisotropy of the protons, $T_{\perp p}/T_{\parallel p}$.
However, the ratio of the ion temperature anisotropies tends to decrease to
lower values of about 0.6 when $V_{\alpha p}/V_{\mathrm{A}}>0.6$, as is
quantitatively shown in Figure 3 d.
The curve in Figure 3d, which represents the weighted mean value of the ratio
of the ion temperature anisotropies (black symbols), clearly indicates that
alpha particles are preferentially heated (perpendicularly to the mean
magnetic field) with respect to the protons whenever $V_{\alpha
p}/V_{\mathrm{A}}\leq 0.4$, and this is true even for a relatively low wave
energy (indicated by the white curve in Figure 3a) and at high collision rates
(green symbols).
One would expect that the plasma tends to thermal equilibrium, in coincidence
with the lowest values of the normalized ion differential flow speed, if a
relatively high collision rate. Then the temperature ratios of the ion species
should also be near unity. However, observationally it seems that preferential
perpendicular heating of the alpha particles with respect to the protons can
persist even in regions where collision rates are high (with $V_{\alpha
p}/V_{A}\leq 0.4$). This is possible because a wave-related local ion heating
mechanism may be acting on time scales much lower than the long cumulative
collision time. Such a fast wave-heating mechanism can drive the plasma far
away from thermal equilibrium, and therefore may cause a significant ion
temperature anisotropy, because on the other side collisions are not fast
enough to enforce thermal equilibrium.
We found in our previous paper (Bourouaine et al. (2011)) that the helium ion
abundance for this selected data set varies mainly between 0.02 and 0.04. The
helium abundance does not show a clear dependence on the normalized
differential ion speed. However, we showed that there is an anti-correlation
between the alpha-to-proton temperature ratio and the helium abundance at a
fixed $V_{\alpha p}/V_{\mathrm{A}}$.
Figure 4: Left ordinate: The mean values of the temperature anisotropy of the
protons, $T_{\perp p}/T_{\parallel p}$ (red dots), and of the alpha particles,
$T_{\perp\alpha}/T_{\parallel\alpha}$ (black diamonds). Right ordinate: Mean
normalized wave power (blue squares). The bars indicate the uncertainties of
the mean values. All quantities are plotted in bins versus the relative ion
differential speed, $V_{\alpha p}/V_{\mathrm{A}}$. The vertical dotted line
separates fast from slow solar wind regimes.
Figure 4 is a plot of the mean values of the alpha-particle temperature
anisotropy, $T_{\perp\alpha}/T_{\parallel\alpha}$, the proton temperature
anisotropy, $T_{\perp p}/T_{\parallel p}$, and the mean relative wave power
versus $V_{\alpha p}/V_{\mathrm{A}}$. The mean values are obtained by
averaging the data within bins of a width $\Delta(V_{\alpha
p}/V_{\mathrm{A}})=0.1$, and the vertical error bars indicate the related
uncertainties of the mean values.
As mentioned above, it turns out that the proton temperature ratio is strictly
correlated with the normalized wave power, which indicates perpendicular
heating of the protons whenever the power of the transverse fluctuations is
enhanced.
An interesting behaviour of the temperature ratio of the alpha particles can
be inferred from Figure 4. This ratio increases with increasing normalized
wave power as long as the normalized ion differential speed stays below about
0.5. Beyond this value of $V_{\alpha p}/V_{\mathrm{A}}$, the alpha-particle
temperature ratio becomes roughly constant, until $V_{\alpha
p}/V_{\mathrm{A}}$ exceeds a value of about 0.7, but then it decreases towards
a value below unity when $V_{\alpha p}/V_{\mathrm{A}}$ reaches one. In the
slow solar wind region, where the collisions are expected to be relatively
high, the proton temperature anisotropy is ranging between 0.8 and 1, however,
the alpha temperature anisotropy is not in the same range but higher, and
varies between 0.9 and 1.2.
Most likely strong collisionality at low $V_{\alpha p}/V_{\mathrm{A}}$ tends
to isotropize the alphas, whereas the weakening of the resonance at higher
alpha/proton speeds leads to a decreasing of the alpha anisotropy, as in Gary
et al. (2005).
The monotonic increase of the proton temperature anisotropy with the
normalized differential ion speed appears to be consistent with the findings
of Kasper et al. (2008) and with the ACE observations of Gary et al. (2005)).
However, there is clear difference between the trend of the alpha temperature
anisotropy in Figure 4 and the results found by Kasper et al. (2008). In their
paper, the alpha temperature anisotropy reaches a minimum value at $V_{\alpha
p}\sim 0.5V_{\mathrm{A}}$, whereas our Figure 3 shows that the alpha
temperature anisotropy reaches its maximum value when $V_{\alpha p}\sim
0.5V_{\mathrm{A}}$. Moreover, the results of Gary et al. (2005) show that the
average alpha temperature anisotropy is monotonically decreasing with
increasing alpha/proton relative speed.
Unlike what Kasper et al. (2008) and Gary et al. (2005) found previously, our
Figure 4 shows that the perpendicular heating is reduced when $V_{\alpha
p}/V_{A}$ is near zero. This reduced ion heating corresponds to an observed
concurrent decrease in the wave power of the transverse waves. This results
would be expected when the heating ultimately rests with the energy in those
waves. If $V_{\alpha p}/V_{A}$ is near zero, we expect that the alpha-particle
temperature ratio increases resulting in a strong perpendicular ion heating
for sufficient wave power. However, although the wave power is empirically
found to be weak when the normalized ion differential speed is low, the alpha
particles are still heated perpendicularly much more than the protons. For
both ion species the interaction with transverse waves is expected to work
against the radial trend caused by the solar wind expansion in a magnetic
mirror, which tends to build up a much larger parallel than perpendicular
temperature, and accordingly a strong fire-hose-type anisotropy.
## 3 Discussion and conclusions
If we assume that the transverse fluctuations provide the energy input for the
observed preferential perpendicular heating of the alpha particles, then their
temperature-ratio profile as given in Figure 4 could be a signature of
cyclotron-wave heating of the alpha particles.
It has been claimed that the alpha particles can only be heated through ion-
cyclotron wave dissipation if the differential speed between alpha-particles
and protons is approximately less than $0.5V_{A}$ (see e.g., Gary et al.
2005). Therefore, alpha-particles can be heated in the perpendicular direction
as long as they stay in resonance with the ion-cyclotron waves. Moreover, the
wave energy is an important input parameter that controls the heating of
alpha-particles. We expect that high wave energy can cause strong
perpendicular heating of those alpha particles that are in resonance with ion-
cyclotron waves.
Figure 4 shows that the perpendicular heating is reduced when $V_{\alpha
p}/V_{\mathrm{A}}$ is near zero. This reduced ion heating corresponds to a
decrease in the observed wave power of the transverse fluctuations, a result
that is expected because the potential for heating ultimately rests with the
energy contained in those waves. If $V_{\alpha p}/V_{\mathrm{A}}$ is near
zero, we expect that the alpha-particle temperature ratio increases, resulting
in strong perpendicular ion heating for sufficient wave power. However,
although the wave power is found empirically to be weak when the normalized
ion differential speed is small, the alpha particles can still be
perpendicularly heated much more than the protons. Yet, for both ion species
the interaction with transverse fluctuations works against the trend caused by
the solar wind expansion in a magnetic-field mirror configuration, which tends
to build up a high parallel temperature anisotropy. Furthermore, Coulomb
collisions tend to thermalize the solar wind plasma and to reduce the
differential speed between alpha particles and protons. But collisions are
effective in removing the non-thermal ion features merely in the comparatively
cold and dense slow solar wind, in which the collision age is found to be high
and the average wave power observed to be weak.
Another possible scenario that may occur as well is that the long-wavelength
and high-amplitude fluctuations of the inertial-range turbulence may
stochastically heat (Chandran 2010) the ions in the perpendicular direction
with respect to the background magnetic field (or non-resonantly drive a
slowly varying $T_{\perp}/T_{\parallel}$ on both ion species). The presence of
this anisotropy can give rise to Alfvén-cyclotron instabilities, which lead to
the growth of relatively high-frequency modes (with frequencies
$\omega\sim\Omega_{p}$, where $\Omega_{p}$ is the proton cyclotron frequency).
In other words, the ion temperature anisotropies could first be caused by low-
frequency fluctuations in the inertial range, and then this thermal energy may
be exchanged between ions and waves at the proton kinetic scale.
## References
* Araneda et al. (2009) Araneda, J. A., Y. Maneva, and E. Marsch (2009), Preferential Heating and Acceleration of $\alpha$ Particles by Alfvén-Cyclotron Waves, Phys. Rev. Lett. 102, 175001
* Araneda et al. (2008) Araneda, J. A., E. Marsch, and A. F.-Viñas (2008), Proton Core Heating and Beam Formation via Parametrically Unstable Alfvén-Cyclotron Waves, Phys. Rev. Lett. 100, 125003
* Bourouaine et al. (2011) Bourouaine, S., E. Marsch, and F. M. Neubauer (2011), On the relative speed and temperature ratio of solar wind alpha particles and protons: Collisions versus wave effects, ApJ. 728, L3
* Bourouaine et al. (2008) Bourouaine, S., E. Marsch, and C. Vocks (2008), On the Efficiency of Nonresonant Ion Heating by Coronal Alfvén Waves, ApJ. 684, L119
* Bourouaine et al. (2010) Bourouaine, S., E. Marsch, and F. M. Neubauer (2010), Correlations between the proton temperature anisotropy and transverse high-frequency waves in the solar wind, Geophys. Res. Lett., 37, L14104
* Chandran (2010) Chandran, Benjamin D. G. (2010), Alfvén Wave Turbulence and Perpendicular Ion Temperatures in Coronal Holes, ApJ, 720, 548
* Cranmer (2009) Cranmer, S. R., (2009) Coronal Holes, Living Rev. Solar Phys, 6, 3
* Gary et al. (2005) Gary, P. S, C. W. Smith, and R. M. Skoug (2005), Signatures of Alfvén-cyclotron wave-ion scattering: Advanced Composition Explorer (ACE) solar wind observations, J. Geophys. Res. 110, A07 108
* Howes et al. (2008) Howes, G. G. W. Dorland, S. C. Cowley, G. W. Hammett, E. Quataert, A. A. Schekochihin, and T. Tatsuno (2008), Kinetic Simulations of Magnetized Turbulence in Astrophysical Plasmas, Phys. Rev. Lett., 100, 065004
* Matteini et al. (2007) Matteini, L., S. Landi, P. Hellinger, F. Pantellini, M. Maksimovic, M. Velli, B. E. Goldstein, and E. Marsch (2007), Evolution of the solar wind proton temperature anisotropy from 0.3 to 2.5 AU, Geophys. Res. Lett., 34, L20105
* Hollweg and Isenberg (2002) Hollweg, Joseph V. and Isenberg, Philip A (2002), Generation of the fast solar wind: A review with emphasis on the resonant cyclotron interaction, J. Geophys. Res. 107 (A7), 1147
* Isenberg et al. (2001) Isenberg, Philip A., Lee, Martin A. and Hollweg, Joseph V (2001), The kinetic shell model of coronal heating and acceleration by ion cyclotron waves 1. Outward propagating waves, J. Geophys. Res. 106, 5649
* Kasper et al. (2008) Kasper, J. C., A. J. Lazarus, and S. P. Gary (2008), Hot Solar-Wind Helium: Direct Evidence for Local Heating by Alfvén-Cyclotron Dissipation, Phys. Rev. Lett. 101, 261103
* Li and Habbal (2000) Li. X and S. R. Habbal (2000), Proton/alpha magnetosonic instability in the fast solar wind, Geophys. Res. 105, 7483
* Marsch et al. (1982) Marsch, E., R. Schwenn, H. Rosenbauer, K. H. Mühlhäuser, W. Pilipp, and F. M. Neubauer (1982b), Solar wind protons - Three-dimensional velocity distributions and derived plasma parameters measured between 0.3 and 1 AU, J. Geophys. Res. 87, 52
* Marsch and Tu (2001) Marsch, E., and C.-Y. Tu (2001), Heating and acceleration of coronal ions interacting with plasma waves through cyclotron and Landau resonance, J. Geophys. Res. 106, 227
* Marsch (2006) Marsch, E. (2006), Kinetic Physics of the Solar Corona and Solar Wind, Living Rev. Solar Phys. 3, 1
* Salem et al. (2003) Salem, C., D. Hubert, C. Lacombe, S. D. Bale, A. Mangeney, D. E. Larson , and R. P. Lin (2003), Electron Properties and Coulomb Collisions in the Solar Wind at 1 AU: Wind Observations, ApJ. 585, 1147
* Tu and Marsch (1995) Tu, C.-Y. and E. Marsch (1995), MHD structures, waves and turbulence in the solar wind: Observations and theories, Space Science Rev. 73, 1
* Wu and Yoon (2007) Wu, C. S., and P. H. Yoon (2007), Proton Heating via Nonresonant Scattering Off Intrinsic Alfvénic Turbulence, Phys. Rev. Lett. 99, 075001
|
arxiv-papers
| 2011-11-11T12:40:14 |
2024-09-04T02:49:24.208178
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "S. Bourouaine, E. Marsch and F. M. Neubauer",
"submitter": "Sofiane Bourouaine",
"url": "https://arxiv.org/abs/1111.2724"
}
|
1111.2881
|
# Semileptonic decays of the Higgs boson at the Tevatron
MisterX
Joseph D. Lykken , Adam O. Martin
Theoretical Physics Department
Fermi National Accelerator Laboratory, Batavia, IL 60510, USA
Jan Winter
PH-TH, CERN, CH-1211 Geneva 23, Switzerland
lykken@fnal.govaomartin@fnal.govjwinter@cern.ch
###### Abstract
We examine the prospects for extending the Tevatron reach for a Standard Model
Higgs boson by including the semileptonic Higgs boson decays $h\to
WW\to\ell\nu_{\ell}\,jj$ for $M_{h}\gtrsim 2\,M_{W}$, and
$h\to~{}Wjj\to\ell\nu_{\ell}\,jj$ for $M_{h}\lesssim 2\,M_{W}$, where $j$ is a
hadronic jet. We employ a realistic simulation of the signal and backgrounds
using the SHERPA Monte Carlo event generator. We find kinematic selections
that enhance the signal over the dominant $W$+jets background. The resulting
sensitivity could be an important addition to ongoing searches, especially in
the mass range $120\lesssim M_{h}\lesssim 150\ \mathrm{GeV}$. The techniques
described can be extended to Higgs boson searches at the Large Hadron
Collider.
FERMILAB-PUB-11-499-T
CERN-PH-TH-2011-237
###### Contents
1. 1 Introduction
2. 2 Strategy and key observables
3. 3 Inclusive cross sections and event generation
1. 3.1 Standard Model Higgs boson production and decay
2. 3.2 Relevant background processes
1. 3.2.1 $W$ boson plus jets background
3. 3.3 Monte Carlo simulation of signal and backgrounds using SHERPA
1. 3.3.1 Generation of signal events
2. 3.3.2 Generation of background events for $W$ boson plus jets production
3. 3.3.3 Generation of background events for electroweak and top-pair production
4. 4 Signal versus background studies based on Monte Carlo simulations using SHERPA
1. 4.1 Baseline selection
2. 4.2 Higgs boson reconstruction based on invariant masses
1. 4.2.1 Reconstruction below the on-shell diboson mass threshold
2. 4.2.2 Effect of the subdominant backgrounds
3. 4.3 More realistic Higgs boson reconstruction methods
1. 4.3.1 More realistic reconstruction below the on-shell diboson mass threshold
4. 4.4 Optimized selection – analyses refinements and (further) significance improvements
5. 5 Conclusions, caveats, and prospects
6. A Appendix: Monte Carlo event generation
1. A.1 Leading-order cross sections
2. A.2 NLO calculations versus CKKW ME+PS merging
7. B Appendix: Analysis side studies and additional material
1. B.1 Ideal Higgs boson reconstruction analyses
2. B.2 More realistic Higgs boson reconstruction analyses
3. B.3 Directions for additional improvements
## 1 Introduction
The Standard Model (SM) predicts a neutral Higgs boson particle whose
couplings to other particles are proportional to the particle masses, and that
couples to photons and gluons via one-loop-generated effective couplings.
While the Higgs boson mass is not predicted, the relation between the Higgs
boson mass and its width is fixed from the predicted couplings. Virtual Higgs
boson contributions to electroweak precision observables have been computed,
and precision data favor $M_{h}<169\ \mathrm{GeV}$ at 95% confidence level
[1]. Searches at the CERN Large Hadron Collider have produced 95% confidence
level exclusion of a SM Higgs boson for a broad mass range above 145 GeV [2,
3].
Because of small signal cross sections and large backgrounds, the search for
the Higgs boson in experiments at the Fermilab Tevatron is very challenging,
even with the large final datasets approaching $10\ \mathrm{fb}^{-1}$ per
experiment. Nevertheless both the CDF and DØ experiments have achieved steady
improvements in their sensitivities in multiple channels to a SM Higgs boson,
and their individual results already exclude a SM Higgs boson in the mass
range $156.7$–$173.8\ \mathrm{GeV}$ and $162$–$170\ \mathrm{GeV}$,
respectively, at 95% confidence level [4, 5, 6]. This exclusion relies mainly
on sensitivity to the dilepton final-state decay chain analyzed in Refs. [7,
8, 9] : $h\to W^{+}W^{-}\\!\to\ell^{+}\nu_{\ell}\;\ell^{-}\bar{\nu}_{\ell}$
where $\ell^{\pm}=e^{\pm}$ or $\mu^{\pm}$.
Here we examine the prospects for extending the Tevatron reach by including a
search for the semileptonic Higgs boson decays $h\to WW\to\ell\nu_{\ell}\;jj$
for $M_{h}\gtrsim 2\,M_{W}$, and $h\to Wjj\to\ell\nu_{\ell}\;jj$ for
$M_{h}\lesssim 2\,M_{W}$, where $j$ is a hadronic jet. This process was first
considered as a potential Higgs boson discovery channel for the SSC [10, 11,
12, 13], emphasizing the case of a very heavy Higgs boson, where the $h\to
ZZ\to 4\,\ell$ “golden mode” becomes limited by its small branching fraction
and the broad Higgs boson width. Similar to the golden mode, the semileptonic
$h\to WW$ modes are (almost) fully reconstructible: assuming that the leptonic
$W$ is close to on-shell, the mass constraint gives an estimate of the
unmeasured longitudinal momentum of the neutrino, up to a two-fold ambiguity
[13]. For $M_{h}\gtrsim 140\ \mathrm{GeV}$ the overall decay rate is 6 times
larger than any other SM Higgs boson decay mode with a triggerable lepton.
Including these semileptonic channels thus offers the distinct possibility of
significantly extending the Tevatron reach over a rather broad mass range.
This channel suffers from large backgrounds from SM processes with a
leptonically decaying $W$ boson. These include diboson production, top quark
production, and direct inclusive $W$+2-jet production. There is also a purely
QCD background that is difficult to estimate absent a dedicated analysis with
data. The dominant background is inclusive $W$+2-jets; from this background
alone we have estimated a signal to background ratio ($S/B$) of $3\times
10^{-4}$, after nominal preselections. Though worrisome, this is not smaller
than the analogous $S/B\simeq 4\times 10^{-5}$ for the $e^{+}e^{-}$ and
$\mu^{+}\mu^{-}$ modes after preselection in the successful Tevatron analyses
of $h\to W^{+}W^{-}\\!\to\ell^{+}\nu_{\ell}\;\ell^{-}\bar{\nu}_{\ell}$ [14,
15, 16, 17, 18, 19].
A drastic reduction in both the $W$+2-jet and diboson backgrounds to
semileptonic Higgs boson decay can be achieved by forward jet tagging, i.e. by
restricting to Higgs boson production from vector boson fusion (VBF) [10, 20];
it is estimated that the additional requirement of forward jet tagging then
gives a factor of $\sim$ 100 reduction in these backgrounds. However the
reduction in the Higgs signal, versus inclusive Higgs boson production, is
also severe: a factor of $\sim$ 10 at the Tevatron [21]. Looking at the
similar trade-off for the dilepton $h\to
W^{+}W^{-}\\!\to\ell^{+}\nu_{\ell}\;\ell^{-}\bar{\nu}_{\ell}$ channel, a
Tevatron study [22] concluded that the overall sensitivity does not improve by
restricting to VBF Higgs boson versus inclusive Higgs boson production. We do
not know of any comparable analysis for the semileptonic channel.
For inclusive Higgs boson production at the Tevatron, the semileptonic
channels were first studied by Han and Zhang [7, 8]. In a parton-level study
with some jet smearing they found that, after basic acceptance cuts together
with a veto on extra energetic jets designed to suppress the $t\bar{t}$
background, the remaining background is completely dominated by $W$+2-jets.
Han and Zhang then made additional kinematic selections that enhance the
signal to background ratio $S/B$. For $M_{h}=140\ (160)\ \mathrm{GeV}$ they
thus obtained a significance estimate of $S/\sqrt{B}=1.0\ (3.3)$ for $30\
\mathrm{fb}^{-1}$ of integrated Tevatron luminosity. The fully differential
Higgs boson decay width for this process was exhibited by Dobrescu and Lykken
[23], who analyzed the basic kinematics and angular distributions that
characterize the Higgs signal.
We improve on these studies by including realistic parton showering (since
parton-level jet smearing is inadequate), an NLO-rate improved treatment of
the Higgs boson decays (including off-shell effects), and a resummed NNLO
estimate of the $gg\to h$ production cross section. The first two improvements
are incorporated by the use of SHERPA [24, 25], a general purpose showering
Monte Carlo program, for simulation of both the signal and the inclusive
$W$+2-jets background. The NNLO signal cross section is modeled by a
$K$-factor.
Our purpose is to study these semileptonic Higgs boson decay channels in a
systematic way, but not to mimic a fully-optimized experimental analysis. The
DØ experiment has already reported on a semileptonic Higgs boson search using
$5.4\ \mathrm{fb}^{-1}$ of Tevatron data [26, 27]; this analysis uses
multivariate decision trees to enhance the significance of the result. Here we
will limit ourselves to simple cuts, in order to make the features of the
analysis and the underlying physics more explicit. We study the Higgs signal
in the mass range $110$–$220\ \mathrm{GeV}$ to reasonably cover the below,
near and above threshold regions for Higgs boson decay to two on-shell $W$
bosons.
In Section 2 we outline the strategy and define several useful observables. In
Section 3 we discuss inclusive Higgs boson production from the dominant
gluon–gluon fusion mechanism, and implement a $K$-factor correction to the
SHERPA result. In Section 4 we introduce basic preselections and develop cuts
implemented in SHERPA to enhance $S/\sqrt{B}$; this section also contains our
main results. We conclude in Section 5 with caveats about the limitations of
our analysis and suggestions for further improvements. Cross-checks and
additional material are presented in the appendices.
## 2 Strategy and key observables
We are interested in the Higgs boson decay
$h\to W^{*}W^{*}\to e\nu_{e}\;jj\ ,$ (1)
and the similar decay with a muon in the final state. In general we take both
$W$ bosons off shell. We will write $e\nu_{e}\,jj$ as $e\nu_{e}\,jj^{\prime}$
where $j$ is the jet with higher transverse momentum ($p_{T}$), and noting
that physical observables will be symmetric under $j\leftrightarrow
j^{\prime}$. We use a baseline selection adopted from the DØ analysis to
define reconstructed jets and leptons and impose realistic acceptance cuts. We
will assume that events with more than one reconstructed lepton are vetoed,
but we want to allow the possibility of extra jets in order to increase signal
efficiency. When more than two jets are present there is a combinatorial
problem; we will define a Higgs boson candidate selection algorithm that
assigns which two jets to use in the Higgs boson reconstruction; these jets
may or may not correspond to the two leading jets in the event. It is
important to note that this algorithm is chosen so as to optimize the signal
sensitivity after the full selection, which is not equivalent to maximizing
the number of correctly reconstructed signal events.
When the leptonically decaying $W$ boson is (close to) on-shell, these decays
are fully reconstructible up to a two-fold ambiguity in the neutrino momentum
without making any assumption about the Higgs boson mass. Here we are assuming
that the transverse momentum of the neutrino is well estimated given a
measurement of the missing transverse energy (MET), as has been demonstrated
by both Tevatron experiments in the determination of the $W$ boson mass.
The semileptonic channel’s advantage of being, in principle, completely
reconstructible offers a great way to separate signal from backgrounds.
However, when the leptonically decaying $W$ boson is far off shell, a
straightforward full reconstruction is not possible. There are then three
generic possibilities for how to proceed:
> * •
>
> Use only transverse observables.
>
> * •
>
> Perform an approximate event-by-event reconstruction using an estimate of
> the off-shell $W$ boson mass.
>
> * •
>
> Perform an approximate event-by-event reconstruction using a (hypothesized)
> Higgs boson mass constraint.
>
>
Since it is not clear a priori which of these approaches maximizes the Higgs
boson sensitivity, we will pursue all three and compare the results.
Given an event-by-event approximate combinatorial full reconstruction of the
putative decaying Higgs boson, one can approximately reproduce the kinematics
in the Higgs boson rest frame. The true Higgs boson rest frame is given by a
longitudinal boost from the lab frame together with a transverse boost defined
by the transverse momentum $p_{T,h}$ of the Higgs boson. An explicit
representation for the four-momenta in the Higgs boson rest frame is given by:
$\displaystyle p_{e}$ $\displaystyle=$
$\displaystyle\frac{1}{2}\;m_{e\nu_{e}}\biggl{(}\gamma_{e\nu_{e}}(1+\beta_{e\nu_{e}}{\rm
cos}\,\theta_{\ell}),\;{\rm sin}\,\theta_{\ell}\,{\rm
cos}\,\varphi_{\ell},\;{\rm sin}\,\theta_{\ell}\,{\rm
sin}\,\varphi_{\ell},\;\gamma_{e\nu_{e}}(\beta_{e\nu_{e}}+{\rm
cos}\,\theta_{\ell})\biggr{)}\;,$ $\displaystyle p_{\nu_{e}}$ $\displaystyle=$
$\displaystyle\frac{1}{2}\;m_{e\nu_{e}}\biggl{(}\gamma_{e\nu_{e}}(1-\beta_{e\nu_{e}}{\rm
cos}\,\theta_{\ell}),\;-{\rm sin}\,\theta_{\ell}\,{\rm
cos}\,\varphi_{\ell},\;-{\rm sin}\,\theta_{\ell}\,{\rm
sin}\,\varphi_{\ell},\;-\gamma_{e\nu_{e}}(\beta_{e\nu_{e}}-{\rm
cos}\,\theta_{\ell})\biggr{)}\;,$ $\displaystyle p_{j}$ $\displaystyle=$
$\displaystyle\frac{1}{2}\;m_{jj^{\prime}}\biggl{(}\gamma_{jj^{\prime}}(1+\beta_{jj^{\prime}}{\rm
cos}\,\theta_{j}),\;{\rm
sin}\,\theta_{j},\;0,\;-\gamma_{jj^{\prime}}(\beta_{jj^{\prime}}+{\rm
cos}\,\theta_{j})\biggr{)}\;,$ $\displaystyle p_{j^{\prime}}$ $\displaystyle=$
$\displaystyle\frac{1}{2}\;m_{jj^{\prime}}\biggl{(}\gamma_{jj^{\prime}}(1-\beta_{jj^{\prime}}{\rm
cos}\,\theta_{j}),\;-{\rm
sin}\,\theta_{j},\;0,\;-\gamma_{jj^{\prime}}(\beta_{jj^{\prime}}-{\rm
cos}\,\theta_{j})\biggr{)}\;,$ (2)
where we have chosen the dijet plane to coincide to the $x$–$z$ plane, and
have chosen the positive $z$-axis to be the direction of the leptonically
decaying $W$ boson. The boost factors of the two $W$ bosons relative to the
Higgs boson rest frame are given by
$\displaystyle\gamma_{jj^{\prime}}$ $\displaystyle=$
$\displaystyle\frac{M_{h}}{2\,m_{jj^{\prime}}}\left(1+\frac{m_{jj^{\prime}}^{2}-m_{e\nu_{e}}^{2}}{M_{h}^{2}}\right)\;,$
$\displaystyle\gamma_{e\nu_{e}}$ $\displaystyle=$
$\displaystyle\frac{M_{h}}{2\,m_{e\nu_{e}}}\left(1-\frac{m_{jj^{\prime}}^{2}-m_{e\nu_{e}}^{2}}{M_{h}^{2}}\right)\;,$
(3)
and we note the identities
$\displaystyle M_{h}$ $\displaystyle=$ $\displaystyle
m_{jj^{\prime}}\,\gamma_{jj^{\prime}}+m_{e\nu_{e}}\,\gamma_{e\nu_{e}}\;,$
$\displaystyle m_{jj^{\prime}}\,\beta_{jj^{\prime}}\,\gamma_{jj^{\prime}}$
$\displaystyle=$ $\displaystyle
m_{e\nu_{e}}\,\beta_{e\nu_{e}}\,\gamma_{e\nu_{e}}\;.$ (4)
Note that $\theta_{j}$ is the angle between jet $j$ and the direction of the
hadronic $W$ boson, as seen in the $W$ rest frame, while $\theta_{e}$ is the
angle between the charged lepton and the direction of the leptonic $W$ boson
as seen in the $W$ rest frame. The azimuthal angle $\varphi_{e}$ is the angle
between the dilepton and dijet planes. Defining
$\displaystyle r_{jj^{\prime}}$ $\displaystyle=$
$\displaystyle\beta_{jj^{\prime}}^{2}\,\gamma_{jj^{\prime}}^{2}\,{\rm
sin}^{2}\theta_{j}\ ,$ (5)
we can calculate the angle $\theta_{jj^{\prime}}$ between the two jets as seen
in the Higgs boson rest frame:
$\displaystyle{\rm cos}\,\theta_{jj^{\prime}}$ $\displaystyle=$
$\displaystyle\frac{r_{jj^{\prime}}-1}{r_{jj^{\prime}}+1}\ .$ (6)
Signal events have a minimum opening angle between the jets as seen in the
Higgs boson rest frame:
$\displaystyle-1\;\leq\;{\rm
cos}\,\theta_{jj^{\prime}}\;\leq\;2\,\beta_{jj^{\prime}}^{2}-1\ .$ (7)
Figure 1: Examples of observables discriminating between Higgs signal and
(the dominant $W$+jets) backgrounds. Distributions are shown from the upper
left to the lower right, for the reconstructed masses of the
$\\{e,\nu_{e},j,j^{\prime}\\}$ final states, their transverse momenta, the
boost of the dijet subsystem with respect to the parent system and the scalar
sum of the $p_{T}$ of the two jets $j$ and $j^{\prime}$. All results are
obtained after the combinatorial Higgs boson candidate selection based on an
ideal mass reconstruction of the potential resonance. The selection is
facilitated by symmetric mass window constraints for both the 4-object and
dijet mass: $|m_{e\nu_{e}jj^{\prime}}-M_{h}|<\Delta$ and
$|m_{jj^{\prime}}-M_{W}|<\delta$, respectively. More detailed explanations
will follow in Section 4. Note that the background differential cross sections
are scaled down by the respective factors indicated in the legends.
In the approximate reconstructions that we will employ in our analysis, the
Higgs boson mass $M_{h}$ is approximated by a 4-object invariant mass
$m_{e\nu_{e}jj^{\prime}}$. The transverse momentum of the Higgs boson is
approximated by the 4-object transverse momentum $p_{T,e\nu_{e}jj^{\prime}}$.
The dijet boost $\gamma_{jj^{\prime}}$ defined in the Higgs boson rest frame
is approximated by $\gamma_{jj^{\prime}|e\nu_{e}}$, which is the dijet boost
defined in the 4-object rest frame, in which we can compute this boost factor
via $\gamma_{jj^{\prime}|e\nu_{e}}=E_{jj^{\prime}}/m_{jj^{\prime}}$. This is
equivalent to using Eqs. (2) where one inserts for each invariant mass its
reconstructed counterpart.
All three of these observables discriminate between Higgs signal and
backgrounds. As seen in Figure 1 (top left), $m_{e\nu_{e}jj^{\prime}}$ for
signal events peaks strongly near the true Higgs boson mass, with a width
determined primarily by parton shower effects. Thus a simple mass window
selection significantly enhances the signal, and since we are interested in
Higgs boson exclusion there is no “look-elsewhere” effect associated with
imposing mass windows [28]. Note that the backgrounds are not necessarily flat
in the mass windows: as seen in Figure 1 the dominant $W$+jets background is
rather flat in the high mass window, but is steeply rising in the lower mass
window because of the underlying kinematics. In Figure 1 (top right) one sees
that the 4-object transverse momentum $p_{T,e\nu_{e}jj^{\prime}}$ has a harder
spectrum for Higgs signal events than for the $W$+jets backgrounds,
independent of the Higgs boson mass.
The reconstructed dijet boost $\gamma_{jj^{\prime}|e\nu_{e}}$ has
qualitatively different behaviour depending on the underlying Higgs boson
mass. When the Higgs boson mass is close to $2\,M_{W}$, the distribution of
the dijet boost for signal events is strongly peaked near one compared to the
distribution for $W$+jets, as seen in Figure 1 (lower left). For larger Higgs
boson masses the signal distribution of the dijet boost is instead rather
strongly peaked around the value $M_{h}/(2\,M_{W})$, as expected from Eq. (2).
Other physical observables of interest for signal versus background
discrimination are defined directly in the lab frame. This includes the
4-object pseudo-rapidity $\eta_{e\nu_{e}jj^{\prime}}$, the pseudo-rapidity
difference of the two jets, $\Delta\eta_{j,j^{\prime}}$, and the scalar sum of
the two selected jet transverse momenta $H_{T,jj^{\prime}}$. As seen in Figure
1 (lower right), the distribution of $H_{T,jj^{\prime}}$ for signal events is
harder than for the $W$+jets background, independent of the Higgs boson mass.
From the distributions shown in Figure 12 one sees that the dijet pseudo-
rapidity difference has a maximum at zero for signal events (which tend to be
confined to the central region), but peaks at a larger value for the $W$+jets
background. Similarly from Figure 13 one notes that the 4-object pseudo-
rapidity distribution is more central for signal than for the $W$+jets
background.
We will also employ various transverse masses, both as signal discriminators
and as inputs to the algorithms for the approximate reconstruction of the
Higgs boson candidates. One class of $m_{T}$ observables is solely constructed
out of transverse degrees of freedom, $\vec{p}_{T,i}=(p_{x,i},p_{y,i})$; we
define these $m_{T}$ observables as
$m_{T,ij..}\;=\;\sqrt{(\lvert\vec{p}_{T,i}\rvert+\lvert\vec{p}_{T,j}\rvert+\ldots)^{2}-(\vec{p}_{T,i}+\vec{p}_{T,j}+\ldots)^{2}}\;\leq\;m_{ij..}\
,$ (8)
where, for the purposes of this study, the labels $ij..$ refer to two, three,
or four final-state physics objects (charged lepton, MET, and the two selected
jets). We also investigate how our results change when we adopt a slightly
different definition that includes the full information from the invariant
mass of the visible subset:
$\left(m^{(k..)}_{T,ik..l..}\right)^{2}\;=\;m^{2}_{il..}\,+\;2\,\bigl{(}E_{T,il..}\,p_{T,k..}-\vec{p}_{T,il..}\cdot\vec{p}_{T,k..}\bigr{)}\;\leq\;m_{ik..l..}\
,$ (9)
where $E^{2}_{T,il..}=p^{2}_{T,il..}+m^{2}_{il..}$. Here, we have separated
the event into a “visible” ($il..$) and “invisible” ($k..$) part. The
transverse masses are all approximately bounded from above by a kinematic
edge; this gives us another handle when fully reconstructing the event.
Schematically, we have
$m_{T,i[k..]l..}\;\leq\;m^{(k..)}_{T,ik..l..}\;\leq\;m_{ik..l..}$ (10)
where $m_{T,i[k..]l..}$ just indicates that in the presence of multiple
invisible objects the $[k..]$ subsystem enters as a whole when computing Eq.
(8). In the 2-particle case, the two transverse-mass definitions coincide
provided the single objects are massless.
Given the above arsenal of kinematic discriminators and approximate
reconstruction techniques, our basic strategy will be to find the most
promising combinations of selections as a function of the Higgs boson mass.
Since we are only performing a cut and count analysis, and are lacking a
realistic detector description, there is no point in attempting a complete
optimization. Instead we will concentrate on providing a comprehensive look at
the physics that distinguishes signal from background.
## 3 Inclusive cross sections and event generation
We use the multi-purpose Monte Carlo event generator SHERPA [24, 25] to pursue
our analysis of semileptonic Higgs boson decays in Higgs boson production via
gluon fusion. This way we can easily include all (hard and soft) initial-state
radiation and final-state radiation (ISR and FSR) effects and arrive at a
fairly realistic description of final states as used for detector simulations.
Furthermore, sophisticated cuts can be implemented in a straightforward manner
owing to the convenient analysis features that come with the SHERPA package.
We also want to make use of SHERPA’s capabilities in providing an enhanced
modeling of multi-jet final states with respect to a treatment by parton
showers only. Apart from a handful of key processes, the SHERPA Monte Carlo
program evaluates cross sections at leading order/tree level utilizing its
integrated automated matrix-element generators AMEGIC++ [29] and/or COMIX
[30]. However, in a number of studies, SHERPA has been shown to generate
predictions that are in sufficient, often good agreement with the shapes of
kinematic distributions obtained from measurements as well as higher-order
calculations; we will give more details in the respective subsections that
follow up. Hence, we wish to identify appropriate constant $K$-factors between
the most accurate theory results and the leading-order predictions. We in turn
want to apply these $K$-factors to correct SHERPA’s predictions for the
inclusion of exact higher-order rate effects. Therefore, we study signal and
background fixed-order and resummed cross sections at Tevatron Run II energies
for the processes
$P\bar{P}\to\ell\nu_{\ell}\;pp$ (11)
leading to final states consisting of an isolated lepton, missing transverse
energy and at least two jets. The label $\ell$ denotes electrons, $e$, and
muons, $\mu$; the parton label $p$ contains light/massless-quark flavours and
gluons. Note that final-state gluons may only occur in background hard
processes or through the inclusion of I/FSR effects. For the respective
$K$-factors, it is not clear a priori at which level of cuts they are defined
most accurately. The most convenient definition should be given in terms of
the total inclusive cross sections, while one may define more exclusive
$K$-factors, if the higher-order tools allow for the specification of the
desired cuts. We do not expect a strong dependence on the exact $K$-factor
definition provided the shapes are comparable. These issues are examined in
more detail below with the goal of determining reasonable signal and
background $K$-factors that can be used to rescale the respective leading-
order cross sections $\sigma^{(0)}$ of the SHERPA predictions, which we take
to pursue our signal versus background studies.
### 3.1 Standard Model Higgs boson production and decay
The signal processes for the final states of interest are summarized by
$P\bar{P}\to h\to W^{(\ast)}W^{(\ast)}\to\ell\nu_{\ell}\;pp$ (12)
where the Higgs particles are produced through gluon–gluon fusion and decay
into $W$ boson pairs that split further into the desired semileptonic final
states. There are other Higgs boson production mechanisms that can contribute.
In particular, the associated production $P\bar{P}\to Vh$ with the additional
vector boson $V$ decaying hadronically and the production via vector-boson
fusion (VBF) ought to be mentioned in this context. For all production
channels, up-to-date theory predictions for the total inclusive cross sections
of these events are needed to arrive at reliable acceptance estimates for
various Higgs boson masses. Refs. [4, 31, 32] give the most recent overview of
the theory calculations and results that are used as input for the ongoing
Tevatron (and LHC) SM Higgs boson searches. For the Higgs boson masses we are
interested in, it is appropriate to separate the Higgs boson production from
its subsequent decays and multiply the production rates by the respective
branching fractions, which we obtain from Hdecay [33, 34, 35, 36] for $h\to
W^{\ast}W^{\ast}$ and the Particle Data Group (PDG) listings [37] for the
subsequent decays of the $W$ bosons.
For our main production channel, the Standard Model Higgs boson production via
gluon fusion, we want to use the most precise theoretical predictions that
have become available over the last few years, for a great review, we refer to
Ref. [32]. Using an effective theory approach, this production channel is
known at NNLO including electroweak and mixed QCD–electroweak contributions
[38, 39, 40, 41]. For a wide range of Higgs boson masses, these NNLO cross
sections have been shown to reproduce the latest results obtained from soft-
gluon resummation up to NNLL accuracy, cf. Refs. [42, 43, 44]. To mimic the
resummation effects, the optimal scale choice at NNLO is found to be
$\mbox{$\mu_{\mathrm{F}}$}=\mbox{$\mu_{\mathrm{R}}$}=M_{h}/2$, while for the
NNLL calculation one employs common scales of $\mu=M_{h}$. Both higher-order
calculations take the most recent parametrization of PDFs at next-to-next-to-
leading order into account where the corresponding PDF sets have been provided
by the MSTW group in 2008 [45]. The Tevatron Higgs boson searches use these
higher-order $gg\to h$ cross section predictions to report their combined CDF
and DØ upper limits on Standard Model Higgs boson production in the
$W^{+}W^{-}$ decay mode [14, 15, 16, 4, 46, 17, 18, 19]. Hence, it makes sense
to input the same theory cross sections in our studies to guarantee a
reasonable level of compatibility between our work and the experimental
searches. However, one should keep in mind that different viewpoints exist
concerning the determination of the best $gg\to h$ cross section numbers. For
example, in Ref. [47] the authors argue that the 10–15% enhancement seen in
the inclusive rates is unlikely to survive the cuts applied in the Tevatron
analyses and, therefore, should not be included in the calculation of the
limits [48, 49]. On the contrary, the renormalization-group improved resummed
NNLO cross sections discussed by Ahrens, Becher, Neubert and Yang in Refs.
[50, 51, 52] would yield a further 5–6% increase of the NNLL $gg\to h$ rates.
This is because in their approach Ahrens et al. do not only resum threshold
logarithms from soft-gluon emission but also $\pi^{2}$-enhanced terms, which
arise in the analytic continuation of the gluon form factor to timelike
momentum transfer.
$M_{h}$ | $\Gamma_{h}$ | $\sigma^{\mathrm{NNLL}}_{ggh}$ | $B^{\mbox{\tiny Ref.~{}\cite[cite]{[\@@bibref{}{:2011cb}{}{}]}}}_{W^{\ast}W^{\ast}}$ | $\ \sigma^{\mathrm{NNLO}}_{ggh}$ | $B_{W^{\ast}W^{\ast}}$ | $\ \;\sigma^{(0)}_{S,\mathrm{all}}$ | $\sigma^{(0)}_{S,\mathrm{NNLO}}$ | $\ \sigma^{(0)}_{S}$ | $K_{S}$ | $\ \sigma^{(0)}_{S,\mathrm{66}}$ | $K^{\mathrm{66}}_{S}$
---|---|---|---|---|---|---|---|---|---|---|---
$\mathrm{[GeV]}$ | $\mathrm{[GeV]}$ | $\mathrm{[fb]}$ | | $\mathrm{[fb]}$ | | $\mathrm{[fb]}$ | $\mathrm{[fb]}$ | $\mathrm{[fb]}$ | | $\mathrm{[fb]}$ |
$110$ | $0.002939$ | $1385.0$ | $0.0482$ | $1428$ | $0.04633$ | $11.81$ | $9.660$ | $3.350$ | 2.88 | $2.550$ | 3.79
$120$ | $0.003595$ | $1072.3$ | $0.143$ | $1102$ | $0.1380$ | $27.11$ | $22.21$ | $7.717$ | 2.88 | $5.990$ | 3.71
$130$ | $0.004986$ | $842.9$ | $0.305$ | $863$ | $0.2976$ | $45.83$ | $37.50$ | $13.06$ | 2.87 | $10.15$ | 3.69
$140$ | $0.008222$ | $670.6$ | $0.504$ | $685$ | $0.4959$ | $60.53$ | $49.60$ | $17.29$ | 2.87 | $13.46$ | 3.68
$150$ | $0.01726$ | $539.1$ | $0.699$ | $550$ | $0.6927$ | $67.88$ | $55.63$ | $19.33$ | 2.88 | $15.08$ | 3.69
$165$ | $0.2429$ | $383.7$ | $0.960$ | $389$ | $0.9595$ | $66.72$ | $54.50$ | $19.35$ | 2.82 | $15.13$ | 3.60
$170$ | $0.3759$ | $344.0$ | $0.965$ | $347$ | $0.9642$ | $60.02$ | $48.85$ | $17.77$ | 2.75 | $13.91$ | 3.51
$180$ | $0.6290$ | $279.2$ | $0.932$ | $283$ | $0.9327$ | $47.33$ | $38.54$ | $14.19$ | 2.72 | $11.15$ | 3.46
$190$ | $1.036$ | $228.0$ | $0.786$ | $229$ | $0.7871$ | $32.50$ | $26.32$ | $9.862$ | 2.67 | $7.775$ | 3.39
$200$ | $1.426$ | $189.1$ | $0.741$ | $190$ | $0.7426$ | $25.40$ | $20.60$ | $7.827$ | 2.63 | $6.191$ | 3.33
$210$ | $1.841$ | | | $159$ | $0.7250$ | | $16.83$ | $6.473$ | 2.60 | $5.131$ | 3.28
$220$ | $2.301$ | | | $134$ | $0.7160$ | | $14.01$ | $5.420$ | 2.58 | $4.317$ | 3.25
Table 1: Signal cross sections $\sigma^{(0)}_{S}$ at NNLO and LO plus the
resulting signal $K$-factors for several different SM Higgs boson masses. The
Higgs boson widths are found from Hdecay calculations. Columns 5 and 6 show
our input NNLO cross sections taken from recent Fehip calculations [53, 54,
55, 38] and input branching fractions for $h\to W^{\ast}W^{\ast}$ obtained
from Hdecay [34, 35], respectively. They are comparable to the values given in
the 3rd and 4th columns that are used by the Tevatron experimentalists in
their ongoing Higgs boson searches [4]. The signal cross sections labelled
“all” account for contributions stemming not only from gluon–gluon fusion but
also from Higgs strahlung and VBF Higgs boson production processes. All
calculations use MSTW2008 parton distributions except those to extract the LO
values denoted by “$\mathrm{66}$”, which result from using CTEQ6.6 PDFs.
For various Higgs boson masses, Table 1 summarizes the signal cross sections
$\sigma^{(0)}_{S}$ that are of relevance for our studies. Separated from the
other entries, the left and right parts of the table show parameters, which we
use as input to our analysis. We have used the NNLO $gg\to h$ inclusive cross
sections, $\sigma^{\mathrm{NNLO}}_{ggh}$, given in the 5th column and
multiplied them with the branching ratios listed in the column to the right of
it, which we have computed with Hdecay version 3.51.111To obtain the values in
Table 1, the NNLO MSTW fit result for the strong coupling,
$\alpha_{\mathrm{s}}(M_{Z})=0.11707$, has been employed, the running of quark
masses at NNLO has been enabled and the top and bottom quark masses have been
set to $m_{t}=173.1$ and $m_{b}=4.8\ \mathrm{GeV}$, respectively. The Higgs
boson widths shown in the 2nd column are also taken from the Hdecay
calculation; they slightly differ from the older values given in [21]. To
arrive at the higher-order prediction of our signal cross sections depicted in
the 8th column, $\sigma^{(0)}_{S,\mathrm{NNLO}}$, we furthermore have
accounted for the $W$ decays described by the branching ratios for one light
lepton species, $B(W\to\ell\nu_{\ell})=0.108$, and for jets, $B(W\to
pp)=0.676$, and a combinatorial factor of 2 reflecting that either of the $W$
bosons may decay leptonically. The NNLO cross sections used here [55] are
updated values with respect to the ones published in [38]. The difference can
be traced back to the addition of the electroweak real-radiation corrections
as encoded in Ref. [39] and the change to top masses of $m_{t}=173.1\
\mathrm{GeV}$. In the 3rd and 4th columns we respectively also give the $gg\to
h$ cross sections and $h\to W^{\ast}W^{\ast}$ branching fractions as used by
the Tevatron experimentalists in their ongoing searches [4]. These values are
in good agreement with the respective numbers used in our study. The other
signal cross sections listed in the rightmost part of Table 1 are the LO rates
obtained from SHERPA, where the one labelled $\sigma^{(0)}_{S,\mathrm{66}}$
refers to the use of CTEQ6.6 PDFs [56]. We will discuss these LO results and
the corresponding $K$-factors in Section 3.3.1. In the 7th column we show an
upper estimate for the signal cross sections $\sigma^{(0)}_{S,\mathrm{all}}$
if one were to include the contributions from the $Wh$, $Zh$ and VBF
production channels. Over the considered Higgs boson mass range the extra
processes would enhance the signal rates resulting from gluon–gluon fusion by
about 22–23%. We have determined these estimates by adding to the $gg\to h$
rates the theory predictions for $\sigma_{Wh}$, $\sigma_{Zh}$ and
$\sigma_{\mathrm{VBF}}$ as presented in the July 2010 CDF and DØ Higgs boson
searches combination paper [46] (for updated values, cf. [4]) including all
necessary branching fractions to arrive at the $pp\;\ell\nu_{\ell}\,pp$ final
states.222For all production channels, we consider the Higgs boson decay as
specified in (12), since our selection and reconstruction procedures are
tailored to this decay mode, see Section 4. In our analysis we deal with
$\ell$, MET and multiple jets (from the decays as well as I/FSR), thus, the
presence of additional jets stemming from other hard decays or VBF does not
alter our analysis procedures considerably, in other words the Higgs boson
reconstruction and selection procedures are designed in a robust way with
respect to additional jet activity. The dominant Higgs strahlung contributions
to the semileptonic final states arise from the hadronic decays of the
associated vector bosons, where we have used the PDG values $B(W\to pp)=0.676$
and $B(Z\to pp)=0.6991$. All other combinations are suppressed by about an
order of magnitude. Also, we do not consider more than one lepton, i.e. we
implicitly assume an exclusive one-lepton cut. We also neglect the $Vh$ cases
where the $Z$ boson decays invisibly and the associated $W$ boson splits up
leptonically while $h\to pppp$. These modes will fail our $h$ boson
reconstruction. The additional 20% increase resulting from these
considerations should be born in mind when acceptances and significances are
evaluated in a signal versus background study. For the purpose of the analysis
we are pursuing here, we however want to be on the conservative side and
solely concentrate on the gluon–gluon fusion events.
### 3.2 Relevant background processes
Processes that give rise to major background contributions are $V$+jets and
multi-jet production, where the latter comes into play because of jets faking
isolated leptons and/or MET. The muon channel suffers less from jet fakes
reducing the multi-jet background by a factor of 5 with respect to the
electron channel. The $Z$+jets background can also contribute in cases where
one lepton goes missing or a jet mimics a lepton while the $Z$ decays
invisibly. Minor contributions stem from $VV$ and $t\bar{t}$ production. A
first DØ search in the semileptonic Higgs boson decay channel shows how these
backgrounds compare to each other after basic selection cuts, see Ref. [26].
The largest fraction of 83% occurs from $V$+jets, followed by multi-jets,
$t\bar{t}$ and $VV$ contributing with 12%, 3% and 2% to the overall
background. The $V$+jets contribution is totally dominated by $W$+jets
production; the $Z(\to\ell^{+}\ell^{-})$+jets background, where one of the
leptons is missed, is small and makes up less than 1% of the total background.
To gain a better understanding of the backgrounds, we will have a closer look
at the major contributor $W$+jets. The multi-jet background cannot be
simulated straightforwardly, since it requires detailed knowledge of the
experiments and measured fake rates etc. Regarding the minor background
contributors, we will study the $t\bar{t}$ as well as the $WW$ – or, more
exactly, electroweak – background. Even though they enter at a rather low
level after the basic selection compared to $W$+jets, it is necessary to
cross-check what number of events remain after more selective cuts have been
applied, as will be discussed in Section 4.
Another background contribution that has been discussed is gluon-initiated
vector boson pair production [57, 58, 59]. This (quark-loop-induced) process
occurs at ${\mathcal{O}}(\alpha^{4}_{\mathrm{ew}}\alpha^{2}_{\mathrm{s}})$,
the same order as the signal. This background formally arises at NNLO, but
under realistic experimental cuts this production channel has been shown to
significantly increase e.g. the $WW\to 2\,\ell\;2\,\nu_{\ell}$ background at
the LHC. At the Tevatron the gluon densities are small, so the impact of
$gg\to WW$ is expected to be negligible. This expectation was confirmed in
Ref. [60] (a 4‰ effect with respect to the NLO cross section for this decay
channel). A more important effect also recently pointed out by Campbell et al.
in Ref. [60] is the interference between $gg\to WW$ and $gg\to h\to WW$, which
can result in ${\mathcal{O}}(0.1)$ corrections to the Higgs boson signal cross
section. However, interference effects are considerably reduced by requiring
the transverse mass of the leptons plus MET system to be smaller than $M_{h}$.
This type of transverse cut is frequently used in our analyses, so we can
safely neglect interference effects in our study.
#### 3.2.1 $W$ boson plus jets background
For our first study of $W$+jets production, we explore the dependence of
inclusive $W$+jet cross sections on the number of jets and the variation of
the common scale $\mu$ used to specify the factorization and renormalization
scales, $\mu_{\mathrm{F}}$ and $\mu_{\mathrm{R}}$, respectively. This
information will help us identify an optimal definition of the $W$+jets
$K$-factor, which we take to improve the rates of the SHERPA predictions. We
calculate inclusive $W$+$\,n\leq 2$-jet cross sections with MCFM version 5.8
to obtain results that are accurate at NLO in the strong-coupling constant
[61, 62, 63]. We also run MCFM at LO to determine explicit NLO-to-LO
theoretical $K$-factors.333We only consider $W^{+}$ bosons decaying into
$e^{+}\nu_{e}$ pairs; the charge conjugated process will just double the cross
section owing to the $P\bar{P}$ initial states at the Tevatron. We employ the
LO and NLO MSTW2008 PDFs [45] with $\alpha_{\mathrm{s}}(M_{Z})=0.13939$ and
$\alpha_{\mathrm{s}}(M_{Z})=0.12018$, respectively, and impose cuts according
to the parameters given in Section 4.1. Note that we do not account for the
so-called triangle cut relating the transverse mass of the $W$ boson and the
missing energy. Other parameters, such as the electroweak input values of the
Standard Model, have been taken according to the MCFM default settings.
| Inclusive $W^{+}$+$\,n$-jet cross sections in pb.
---|---
$n$ | | $\mu=M_{\perp,W}/2$ | $=M_{\perp,W}$ | $=2\,M_{\perp,W}$ | % | | $\mu=\hat{H}_{T}/2$ | $=\hat{H}_{T}$ | $=2\,\hat{H}_{T}$ | %
$0$ | LO | $457$ | $465$ | $469$ | ${}^{-1.7}_{+0.9}$ | | $453$ | $463$ | $468$ | ${}^{-2.2}_{+1.1}$
| NLO | $625$ | $619$ | $616$ | ${}^{+1.0}_{-0.5}$ | | $606$ | $602$ | $602$ | ${}^{+0.7}_{-0.0}$
| $K$ | 1.37 | 1.33 | 1.31 | | | 1.34 | 1.30 | 1.29 |
$1$ | LO | $66.2$ | $55.6$ | $47.3$ | ${}^{+19.1}_{-14.9}$ | | $62.5$ | $52.7$ | $45.1$ | ${}^{+18.6}_{-14.4}$
| NLO | $79.8$ | $74.6$ | $69.3$ | ${}^{+7.0}_{-7.1}$ | | $74.2$ | $70.2$ | $65.7$ | ${}^{+5.7}_{-6.4}$
| $K$ | 1.21 | 1.34 | 1.47 | | | 1.19 | 1.33 | 1.46 |
| $R_{\mathrm{LO}}^{(1,0)}$ | $0.145$ | $0.120$ | $0.101$ | | | $0.138$ | $0.114$ | $0.096$ |
| $R_{\mathrm{NLO}}^{(1,0)}$ | $0.128$ | $0.121$ | $0.113$ | | | $0.122$ | $0.117$ | $0.109$ |
$2$ | LO | $14.4$ | $10.1$ | $7.40$ | ${}^{+42.6}_{-26.7}$ | | $10.9$ | $7.89$ | $5.89$ | ${}^{+38.1}_{-25.3}$
| NLO | $12.8$ | $11.7$ | $10.4$ | ${}^{+9.4}_{-11.1}$ | | $12.0$ | $10.1$ | $8.95$ | ${}^{+18.8}_{-11.4}$
| $K$ | 0.89 | 1.16 | 1.41 | | | 1.10 | 1.28 | 1.52 |
| $R_{\mathrm{LO}}^{(2,1)}$ | $0.218$ | $0.182$ | $0.156$ | | | $0.174$ | $0.150$ | $0.131$ |
| $R_{\mathrm{NLO}}^{(2,1)}$ | $0.160$ | $0.157$ | $0.150$ | | | $0.162$ | $0.144$ | $0.136$ |
Table 2: Inclusive $W^{+}$+$\,n$-jet cross sections $\sigma_{n}$ in pb at LO
and NLO in QCD for different scale choices and jet multiplicities using
MSTW2008 PDFs. The variations with respect to the nominal choices,
$M_{\perp,W}$ with $M^{2}_{\perp,W}=M^{2}_{W}+p^{2}_{T,W}$ and $\hat{H}_{T}$,
are given in the columns labelled by “%”. Numerical integration uncertainties
are not displayed, since they are at least one order of magnitude below the
accuracy indicated here. NLO-to-LO $K$-factors and $n$-to-($n-1$)-jet cross
section ratios are also shown for all possible instances.
We display our MCFM results in Table 2 for different inclusive jet bins $n$
and scale choices $\mu$. As expected, for each $n$-jet multiplicity, the NLO
cross sections are more stable under scale variations with the largest
deviations occurring for the more complex $W^{+}$+2-jet processes. This is
also reflected by the various NLO-to-LO $K$-factors, which vary from about 0.9
to 1.5 for $n=2$ while they are rather constant for $n=0$ ranging from about
1.3 to 1.4 only. For illustrative purposes, we also list the LO and NLO
inclusive jet-rate ratios $R^{(n,n-1)}=\sigma_{n}/\sigma_{n-1}$ starting with
$n=1$. The $W^{+}$+$\,n\leq 1$-jet cross sections do not deviate substantially
for the two nominal scales chosen, $\mu\sim M_{\perp,W}$ where
$M^{2}_{\perp,W}=M^{2}_{W}+p^{2}_{T,W}$ and $\mu\sim\hat{H}_{T}$, which are
determined dynamically for each event. Note that $\hat{H}_{T}$ is the scalar
sum of the transverse momenta of all particles (partons) in the event, i.e. no
jet clustering has taken place. The $\mu\sim M_{\perp,W}$ scales lead to
slightly larger rates when compared to those obtained for
$\mu\sim\hat{H}_{T}$. This can be traced back to the occurrence of
$\mu$-values that are on average larger in the latter case, since
$\langle\hat{H}_{T}\rangle\gtrsim\langle M_{\perp,W}\rangle$ for $n\geq 1$.
For the same reason, the cross section differences become more manifest for
$n=2$. The presence of the second jet gives an extra $p_{T}$ contribution to
$\hat{H}_{T}$ per event whereas $M_{\perp,W}$ is less affected. This further
enhances the deviation of the $\hat{H}_{T}$ and $M_{\perp,W}$ averages.
Given the numbers of Table 2 we can conclude that our knowledge of the
$W$+2-jet background is accurate on the level of $\lesssim$ 20%. A $K$-factor
of about $1.5$ should be viewed as the upper limit for correcting LO results;
in Section 3.3.2 we will however compare the SHERPA background rates more
closely with the results of Table 2 and determine a $K$-factor accordingly.
### 3.3 Monte Carlo simulation of signal and backgrounds using SHERPA
For reasons outlined at the beginning of Section 3, we use SHERPA version
1.1.3 [24, 25] to generate the $\ell\nu_{\ell}$+jets signal and background
events that are needed to understand the potential of a Standard Model Higgs
boson analysis in the lepton + MET + jets channel.444Version 1.1.3 was the
last of the previous SHERPA generation; for all our purposes, it models the
necessary physics equally well compared to the upgraded versions of the
current (1.3.x) generation. Cross-comparisons have confirmed this result. We
will employ the results of the previous two subsections to settle the
inclusive $K$-factors needed to re-scale SHERPA’s LO predictions and include
higher-order rate effects.
The signal and background simulations share a number of common parameters and
options that have been set as follows: we simulate all events at the parton-
shower level, i.e. we include initial- and final-state QCD radiation, but do
not account for hadronization effects and corrections owing to the underlying
event, since their impact is considerably smaller with respect to additional
QCD radiation arising from the hard processes. The intrinsic transverse motion
of quarks and gluons inside the colliding hadrons is however modeled by an
intrinsic Gaussian $k_{T}$-smearing of $\mu(k_{T})=0.2$ and
$\sigma(k_{T})=0.8\ \mathrm{GeV}$. The electroweak parameters are explicitly
given: $M_{W}=80.419$, $\Gamma_{W}=2.06$, $M_{Z}=91.188$, $\Gamma_{Z}=2.49\
\mathrm{GeV}$; the Higgs boson masses and widths are mutable, taken according
to Table 1; the couplings are specified by
$\alpha_{\mathrm{ew}}(0)=1/137.036$, $\sin^{2}_{\mathrm{W}}=0.2222$ and the
Higgs field vacuum expectation value and its quartic coupling are given as
$246\ \mathrm{GeV}$ and $0.47591$, respectively. The CKM matrix is simply
parametrized by the identity matrix. The bottom and top quark masses are set
to $m_{b}=4.8$ and $m_{t}=173.1\ \mathrm{GeV}$, respectively, and all other
quark masses are zero. To avoid any bias owing to the utilization of different
PDFs and in order to develop a consistent picture, signal and background
events are generated using the same parton distributions. Our first choice of
PDFs is the LO MSTW set MSTW2008lo90cl [45], because its NNLO version has been
the preferred PDF set used for the recent calculations of the gluon–gluon
fusion Higgs boson production cross sections. The strong coupling is
determined by one-loop running with $\alpha_{\mathrm{s}}(M_{Z})=0.13939$,
which is the advertised fit value of the LO MSTW2008 set.
To gain some understanding of PDF effects, we compare our MSTW2008 results
against predictions generated with a different PDF set. To fully establish the
comparison on the same level as for the MSTW2008 PDFs, signal and background
rates have to be predicted from theory using the alternative PDF libraries. We
cannot follow this approach here, instead we start out from the same
normalization that has been used for the SHERPA predictions calculated with
MSTW2008 PDFs. After the application of our cuts we then focus on the
differences induced by the alternative PDF set. As our second choice we employ
the CTEQ6.6 PDF libraries [56] where the strong coupling is set by
$\alpha_{\mathrm{s}}(M_{Z})=0.118$ and the running of the coupling is again
computed at one loop. Notice that SHERPA invokes a 6-flavour running for all
strong-coupling evaluations.
#### 3.3.1 Generation of signal events
We simulate signal events with electrons or positrons in the final state
according to
$P\bar{P}\to h\to e\nu_{e}\;pp\to e\nu_{e}+\mathrm{jets}\ .$ (13)
The hard process composed as $gg\to h\to e\nu_{e}\,pp$ is calculated at LO.
The incoming gluons and the quarks arising from the decay undergo further
parton showering, which automatically is taken care of by the SHERPA
simulation. One ends up with the $e\nu_{e}+\mathrm{jets}$ final states
generated at shower level. The hard-process tree-level matrix elements and
subsequent parton showers needed for the simulation are provided by the SHERPA
modules AMEGIC++ and APACIC++, respectively. For our purpose, it is sufficient
to treat the muon final states in exactly the same manner as the electron
final states, i.e. the muon decay channel is included by multiplication with
the lepton factor $f_{\ell}=2$ at the appropriate places.
The Higgs boson production occurs through gluon–gluon fusion via intermediate
heavy-quark loops. In SHERPA this is modeled at LO by an effective $gg\to h$
coupling where the top quarks have been integrated out. The EHC (Effective
Higgs Couplings) implementation of SHERPA includes all interactions up to
5-point vertices that result from the effective-theory Lagrangian. These
effective vertices can simply be added to the Standard Model. We do not work
in the infinite top-mass limit, because we also want to consider Higgs bosons
heavier than the top quark, the approximation however is well applicable only
as long as $m_{t}>M_{h}$. The Higgs boson decays are described by $1\to 4$
processes, i.e. we directly consider $h\to e\nu_{e}\,pp$. We thereby make use
of SHERPA’s feature to decompose processes on the amplitude level into the
production and decays of unstable intermediate particles while the colour and
spin correlations are fully preserved between the production and decay
amplitudes [25]. This way one can focus on certain resonant contributions
instead of calculating the full set of diagrams contributing to a given final
state, which in our case would lead to the inclusion of contributions from the
backgrounds. The intermediate propagators are allowed to be off-shell, such
that finite-width effects are naturally incorporated into the simulation. This
comes in handy especially for Higgs boson masses below the $WW$ mass threshold
as the $1\to 4$ decays moreover guarantee the inclusion of off-shell $W$-boson
effects. A consistent LO treatment would require the use of total Higgs boson
widths as computed at LO. We instead put in the values from the Hdecay
calculations [34, 35] as listed in Table 1. This modifies the Higgs boson
propagators and one arrives at a more accurate description of the finite-width
effects of the Higgs boson decays. The effect on the total rate,
$\sigma^{(0)}_{S}\;=\;\frac{\Gamma(h\to
e\nu_{e}\,pp)}{\Gamma_{h}}\;\sigma^{\mathrm{LO}}_{ggh}\ ,$ (14)
is nullified, since we eventually correct for the NNLO rates
$\sigma^{(0)}_{S,\mathrm{NNLO}}$ worked out in Section 3.1.
In Ref. [35] a comparative study has been presented for Higgs boson production
via gluon fusion at the LHC. Amongst a variety of predictions including those
given by HNNLO [64, 65], the SHERPA versions 1.1.3 and 1.2.1 have been
validated to produce very reasonable results for the shapes of distributions
like the rapidity and transverse momentum of the Higgs boson, pseudo-
rapidities and transverse momenta of associated jets and jet–jet $\Delta R$
separations. We hence rely on a well validated approach that works not only
for pure parton showering in addition to the Higgs boson production and
decays, but also beyond in the context of merging higher-order tree-level
matrix elements with parton showers. Nevertheless, we have carried out a
number of cross-checks to convince ourselves of the correctness of the SHERPA
calculations; for the details, we refer the reader to Appendix A.1.
Finally we turn to the discussion of the $K$-factors. Recalling our findings
of Section 3.1, we want to re-scale SHERPA’s leading-order signal cross
sections $\sigma^{(0)}_{S}$ to the fixed-order NNLO predictions given by Fehip
for Higgs boson production in $gg\to h$ fusion via intermediate heavy-quark
loops [53, 54, 38, 55]. To be consistent, the renormalization and
factorization scales of the LO hard-process evaluations are chosen as for the
higher-order calculations, which employ
$\mu=\mbox{$\mu_{\mathrm{R}}$}=\mbox{$\mu_{\mathrm{F}}$}=M_{h}/2$. The
resulting cross sections ultimately define our signal $K$-factors:
$K_{S}\;=\;\frac{\sigma^{(0)}_{S,\mathrm{NNLO}}}{\sigma^{(0)}_{S}}\ .$ (15)
We have determined two sets of $K$-factors for our two choices of PDFs where
the $K$-factors and LO cross sections labelled by “$\mathrm{66}$” refer to the
case of utilizing the CTEQ6.6 libraries when calculating the LO cross
sections. Our results have already been summarized in Section 3.1, they are
presented in the right part of Table 1. The $K$-factors are remarkably stable
varying slowly from 2.8 to 2.6 over the entire Higgs boson mass range when
relying on MSTW2008 PDFs. In the CTEQ6.6 case, where we have employed
$\mu=\mbox{$\mu_{\mathrm{R}}$}=\mbox{$\mu_{\mathrm{F}}$}=\sqrt{\hat{s}}/2\approx
M_{h}/2$, they are larger due to the smaller LO rates but their magnitude
still remains $\lesssim$ 3.6.555The LO rates calculated with the CTEQ PDF
libraries are diminished for two reasons mainly, the value of
$\alpha_{\mathrm{s}}$ at $M_{Z}$ is considerably lower and the altered scale
choice entails a further reduction of the cross sections.
In addition to the default scale choice of $\mu=M_{h}/2$ that we used for the
MSTW runs, we have explored other options by essentially varying this default
setting for $\mu$ by factors of 2. We obtained results for $\mu=M_{h}/4$,
$\mu=\sqrt{\hat{s}}/2\approx M_{h}/2$ and $\mu=\sqrt{\hat{s}}\approx M_{h}$
with the effect that the LO rates were varied by +20% to -15% but – as
expected – no shape changes were induced.
#### 3.3.2 Generation of background events for $W$ boson plus jets production
We restrict ourselves to the Monte Carlo simulation of the $e^{\pm}$ channels.
Their final states are generated through
$P\bar{P}\to e\nu_{e}+0,1,2\,p\to e\nu_{e}+\mathrm{jets}$ (16)
using an inclusive $W$+2-jets sample obtained from the
Catani–Krauss–Kuhn–Webber (CKKW) merging of the corresponding tree-level
matrix elements with the parton showers (ME+PS) [66, 67]. In these $W$+2-jets
calculations the electroweak order is tied to $\alpha^{2}_{\mathrm{ew}}$.
Unlike the NLO calculation we do include matrix elements where the extra
partons may occur as $b$ quarks; effectively, they are however treated as
massless quarks in the evaluation of the matrix elements and generation of the
radiation pattern. The events are corrected for the $b$-quark mass after the
parton showering. This approach generates slightly harder $p_{T}$ spectra but
as part of being more conservative in estimating this background it is totally
reasonable. Similarly, we simply assume no effect of a $b$-jet veto in
removing $W$+jets events.
The parameters of the matrix-element parton-shower merging are the jet
separation scale $Q_{\mathrm{jet}}$ and the $D$-parameter, which is used to
fix the minimal separation of the parton jets. These parameters are
respectively set to $Q_{\mathrm{jet}}=20\ \mathrm{GeV}$ and $D=0.4$ in
correspondence to the jet $p_{T}$ threshold and cone definitions of our
analysis, see Section 4.1. $Q_{\mathrm{jet}}$ denotes the scale at which –
according to the internal $k_{T}$-jet measure incorporating the $D$-parameter
– the multi-jet phase space is divided into the two domains of
$Q>Q_{\mathrm{jet}}$ where the jets are produced through exact tree-level
matrix elements and $Q_{\mathrm{jet}}>Q>Q_{\textrm{cut-off}}\sim 1\
\mathrm{GeV}$ where the parton-shower intra-jet evolution takes place. We
generate predictions from samples that merge matrix elements with up to
$n^{\mathrm{max}}_{p}=2$ partons. Although we could increase this maximum
number, at this point we do not want to include matrix elements with more than
two partons in order to be consistent with our signal event generation where
the jets beyond those arising from the $W$-boson decays are produced by parton
showers only. If one wishes to further improve on the description of
additional hard jets, both background and signal simulations should be
extended on the same footing.
$n$ | | $\mu=M_{\perp,W}/2$ | $=M_{\perp,W}$ | $=2\,M_{\perp,W}$ | | $\mu=\hat{H}_{T}/2$ | $=\hat{H}_{T}$ | $=2\,\hat{H}_{T}$ | | $\sigma_{\mathrm{{CKKW}}}/\mathrm{pb}$
---|---|---|---|---|---|---|---|---|---|---
LO | $0$ | 0.92 | 0.94 | 0.95 | | 0.91 | 0.93 | 0.94 | | $496$
| $2$ | 1.45 | 1.02 | 0.75 | | 1.10 | 0.80 | 0.59 | | $9.90$
NLO | $0$ | 1.26 | 1.25 | 1.24 | | 1.22 | 1.21 | 1.21 | |
| $2$ | 1.29 | 1.18 | 1.05 | | 1.21 | 1.02 | 0.90 | |
Table 3: Ratios at LO and QCD NLO taken between rates of MCFM and SHERPA CKKW
(rightmost column) for inclusive $W^{+}$+$\,n$-jet production at different
choices of scales in MCFM using MSTW2008 PDFs in all cases. The MCFM cross
sections are listed in Table 2.
The $V$+jets predictions of SHERPA have been extensively studied and
validated over the last few years. Studies exist for comparisons against other
Monte Carlo tools [68, 69, 70, 71, 72], NLO calculations [68, 69, 73] and
Tevatron Run I and II data [68, 74, 75, 25, 76, 77, 78, 79]. They have helped
improve SHERPA gradually and provided evidence that SHERPA gives a good
description of the shapes of the $V$+jet final-state distributions missing a
global scaling factor only, which can be extracted from the data [25] or
higher-order calculations [73].In Appendix A.2 we briefly highlight to what
extent the CKKW ME+PS merging includes important features of NLO computations.
We use the results of Table 2 to identify a reasonable $K$-factor for our
simulated $W$+jets backgrounds. Relying on MSTW2008 PDFs, the SHERPA numbers
for the inclusive $W^{+}$ and $W^{+}$+2-jet cross sections are 496 pb and 9.90
pb, respectively. The 0-jet SHERPA rate thereby is about 7% larger than the
corresponding LO rates given by MCFM. The differences occur because on the one
hand MCFM by default invokes a non-diagonal CKM matrix and a somewhat larger
$W$-boson width 666Switching to an unity CKM matrix and using SHERPA’s input
parameters, one finds 486 pb at $\mu=M_{W}$., on the other hand SHERPA’s
merged-sample generation relies on a very different scale-setting procedure
compared to the leading fixed-order calculations. These differences have no
effect on the kinematic distributions – and are fully absorbed by the
$K$-factor, i.e. CKM effects may eventually enter through the correction of
SHERPA’s rate. Table 3 summarizes the ratios between the MCFM predictions of
Table 2 and SHERPA’s CKKW cross sections mentioned above. This overview neatly
points to the two options that give the most stable ratios; they are found at
NLO for $\mu=M_{\perp,W}/2$ and $\mu=\hat{H}_{T}/2$ where the latter scale
choice has been reported to be well suitable for even higher jet
multiplicities [80, 73, 81]. Based on these observations, we can hence
conclude that it is fair to apply a $K$-factor of
$K_{B}\;=\;1.25$ (17)
to the $W$+jets backgrounds employed in our study. The number found here
compares well to global $K$-factors as reported throughout the literature.
As outlined at the beginning of Section 3.3, we want to normalize the
backgrounds obtained with CTEQ6.6 to those computed with MSTW2008 PDFs. In the
CTEQ case the SHERPA CKKW cross sections amount to 544 pb and 8.13 pb for the
inclusive $W^{+}$ and $W^{+}$+2-jet final states, respectively. Since the
latter selection of $W$+2-jet events is more exclusive, we re-scale the CTEQ
backgrounds according to $K^{\mathrm{66}}_{B}\times 8.13\
\mathrm{pb}=K_{B}\times 9.90\ \mathrm{pb}$ and arrive at
$K^{\mathrm{66}}_{B}\;=\;1.52\ .$ (18)
#### 3.3.3 Generation of background events for electroweak and top-pair
production
The $WW$ background enters at ${\mathcal{O}}(\alpha^{4}_{\mathrm{ew}})$ of the
electroweak coupling constant $\alpha_{\mathrm{ew}}$, i.e. it is suppressed by
more than two orders of magnitude with respect to the $W$+2-jets contribution
occurring at ${\mathcal{O}}(\alpha^{2}_{\mathrm{ew}}\alpha^{2}_{\mathrm{s}})$.
Still, without running the simulation we cannot say for sure whether the
continuum $WW$ production remains an 1% effect after application of the
analysis cuts and – if necessary – what handles exist to distinguish it from
the signal. Because of the large resemblance between the topologies of the
Higgs boson decay and the dominant $WW$ production channels, we anticipate
some of the cuts to be equally efficient for both signal and minor background.
This makes it hard to estimate a priori the extent to which the Higgs boson
signal will be diluted by the electroweak production type of processes. For
the same reasons, the $t\bar{t}$ production final states can be expected to
enhance the signal dilution on a similar level. Certainly, whether we end up
with an 1% or 10% effect, this time it is sufficient to apply $K$-factors
taken from the literature.
For the simulation of the diboson production background, we take the complete
set of electroweak diagrams occurring at
${\mathcal{O}}(\alpha^{4}_{\mathrm{ew}})$ into account including interference
effects. This way we comprise physics effects beyond the plain $WW$ production
with subsequent decays of the gauge bosons.777Relying on the full set of
electroweak processes is more conservative: the rate increases by about 20%;
the effect on the shapes is rather small in general, although we observe
slightly harder tails in $p_{T}$ distributions. As before we only generate the
processes regarding the first lepton family:
$P\bar{P}\to e\nu_{e}\;pp\to e\nu_{e}+\mathrm{jets}$ (19)
where additional jets are produced by the parton shower. Similar setups have
been validated for SHERPA in [82] and more recently in [83, 84]. Here, we
employ a dynamic choice, $\mu=\sqrt{\hat{s}}\sim 2\,M_{V}$, to calculate the
scales of the LO processes. Parton-level jets are generated as in Section
3.3.2 using the same jet-finder algorithm and the same parameters
($Q_{\mathrm{jet}}=20\ \mathrm{GeV}$ and $D=0.4$). Processes with bottom
quarks are included; just as in the $W$+2-jets case, they are treated as
massless.
The $t\bar{t}$ background events are generated according to
$P\bar{P}\to t\bar{t}\to b\bar{b}\;e\nu_{e}\;pp\to e\nu_{e}+\mathrm{jets}$
(20)
again utilizing the parton shower to describe any additional jet activity
beyond that generated by the top quark decays. We only consider the
semileptonic channel. The fully hadronic channel has to be considered together
with the QCD background, and the fully leptonic channel will suffer from
smaller branching fractions, the single isolated-lepton requirement and any
dijet mass window that we impose around the $W$ mass. The LO processes are
calculated at the scale $\mu=m_{t}$, the mass of the $b$ quarks is fully taken
into account and the partonic phase-space generation is subject to the same
jet-finding constraints as used for the compilation of the electroweak
background. In addition we place mild generation cuts on the $b$ quarks:
$p_{T,b}>10\ \mathrm{GeV}$ and $\Delta R_{b,p}>0.3$.
We also examined the impact of $Z$+jets production on our analyses, and found
that this contribution makes up less than 1% of the total background. Since
$Z$+jets has kinematics similar to $W$+jets, we will not study it further.
In SHERPA the minor backgrounds are computed at LO. As in all other cases, we
correct the total inclusive cross section for NLO effects by multiplying with
global $K$-factors, which for both electroweak and $t\bar{t}$ production are
larger than 1. Tevatron diboson searches like [85, 86] measure cross sections
in good agreement with the prediction given by Campbell and Ellis ($16.1\pm
0.9\ \mathrm{pb}$ for $WW$+$WZ$). From their work [87] (Table III) we infer
an NLO-to-LO $K$-factor ranging from $1.30$ to $1.35$. For our analysis, we
will then use the conservative estimate 888NLO corrections to $VV$ production
can become large, for a recent example, see [88] where $K$-factors as large as
$1.77$ have been reported; taken this value, we would certainly overestimate
the electroweak contribution, since the CDF-type cuts employed in [88] are
more exclusive. As for the shapes, we found them reliably described in a
cross-check against an electroweak $VV$+1-jet merged sample, including
matrix-element contributions at
$\mathcal{O}(\alpha^{4}_{\mathrm{ew}}\alpha_{\mathrm{s}})$.
$K_{B,\mathrm{ew}}\;=\;1.35\ .$ (21)
For the inclusive $t\bar{t}$ production, we can safely estimate a conservative
$K$-factor of $1.30$ by comparing the cross section results given for the
Tevatron in Ref. [89]. Adopting a $b$-tagging efficiency of the order of 50%
would give us a 75% chance of vetoing $t\bar{t}$ events with at least one
$b$-quark jet, i.e. we were able to remove about 3/4 of the $t\bar{t}$
background; again, we will be more conservative here and assume that about 40%
of the $t\bar{t}$ events will pass; hence, for our purpose, we finally assign
$K^{b\textrm{-}\mathrm{veto}}_{B,t\bar{t}}\;=\;0.52\ .$ (22)
## 4 Signal versus background studies based on Monte Carlo simulations using
SHERPA
We report the successive improvements of the $S/\sqrt{B}$ significances when
applying a series of cuts that preserve most of the signal and reduce the
inclusive $W$+2-jets background significantly.
### 4.1 Baseline selection
We follow the event-selection procedure as used by the DØ collaboration [26]:
hadronic jets $j$ are identified by a seeded midpoint cone algorithm using the
$E$-scheme for recombining the momenta [90]. The cone size is taken as $R=0.5$
and selection cuts of $p^{\mathrm{jet}}_{T}>20\ \mathrm{GeV}$ and
$\lvert\eta^{\mathrm{jet}}\rvert<2.5$ are imposed. Additionally, we require a
lepton–jet isolation of $\Delta R^{\mathrm{lep}\textrm{--}\mathrm{jet}}>0.4$.
For the leptonic sector, we apply transverse-momentum and pseudo-rapidity cuts
of $p^{\mathrm{lep}}_{T}>15\ \mathrm{GeV}$ and
$\lvert\eta^{\mathrm{lep}}\rvert<1.1$, respectively, supplemented by a
missing-energy cut via $\not{p}_{T}>15\ \mathrm{GeV}$. In addition, we also
account for $M_{T,W}+\not{E}_{T}/2>40\ \mathrm{GeV}$, which is known as
triangle cut.999The cut is applied to the leptonic $W$ boson where
$M_{T,W}=\sqrt{(\lvert\vec{p}_{T,\ell}\rvert+\lvert\not{\vec{p}}_{T}\rvert)^{2}-(\vec{p}_{T,\ell}+\not{\vec{p}}_{T})^{2}}\equiv
m_{T,\ell\nu_{\ell}}$, cf. Eq. (8). For Higgs boson masses above the $WW$
threshold, the rate reduction and shape changes induced by this cut are
marginal.
### 4.2 Higgs boson reconstruction based on invariant masses
After the application of the basic cuts, we identify the best-fit
$\\{e,\nu_{e},j,j^{\prime}\\}$ set from all possible candidates allowed by
combinatorics. The algorithm we use to identify the best-fit object is
referred to as the Higgs boson candidate selection. Several different
selection algorithms are possible, however for now, we will use an invariant
mass (or invm) selection: the four particles (reconstructed in a more or less
ideal way) whose combined mass $m_{e\nu_{e}jj^{\prime}}$ is closest to a
“test” Higgs boson mass $M_{h}$ are chosen. Of course, in the context of the
analysis, the Higgs boson mass enters as a hypothesis and, thus, is treated as
a parameter. Regardless of the selection algorithm, we refer to $j$ and
$j^{\prime}$ as the two selected jets, which are not necessarily the hardest
jets in the event.
After selection, we impose a requirement on the absolute difference between
$m_{e\nu_{e}jj^{\prime}}$ and the hypothesized Higgs boson mass; events are
kept only if they reconstruct a mass that lies within the window
$M_{h}-\Delta<m_{e\nu_{e}jj^{\prime}}<M_{h}+\Delta$. This completes our
combinatorial Higgs boson reconstruction, which we label as “comb. $h$-reco”
in our tables. On top of this selection, we may include an additional dijet
mass constraint of $M_{W}-\delta<m_{jj^{\prime}}<M_{W}+\delta$ (marked by
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ in the
tables). The selection procedure will certainly shape – to some extent – the
remaining background to look like the signal, however the primary effect we
are interested in concerns the reduction of the background rate while we want
to preserve as many signal events as possible.
One may ask whether the reconstruction of the Higgs particle candidate can be
achieved more easily by selecting the set containing the respective hardest
particles, in particular, by choosing the two hardest jets, $j_{1}$ and
$j_{2}$, to reconstruct the hadronically decaying $W$ boson. We will refer to
this approach as the naive Higgs boson reconstruction, denoted as “naive
$h$-reco” later on (as before we use
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ in the tables
to indicate that a dijet mass constraint has been imposed in addition). There
are no combinatorial issues in the naive scheme. However, as we show in our
tables, it yields poorer significances than the selection based on
combinatorics.
We calculate the number of $S$ signal and $B$ background events for different
Higgs boson masses assuming a total integrated luminosity of
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$. This seems to be a good ${\mathcal{L}}$
estimate for what each of the two Tevatron experiments, CDF and DØ, were able
to collect before the eventual Run II shutdown in September 2011. We compute
the numbers according to
$\displaystyle S$ $\displaystyle=$ $\displaystyle
K_{S}\,\varepsilon_{S}\,\sigma^{(0)}_{S}\;\times\;2\,f_{\ell}\,{\mathcal{L}}\;=\;K_{S}\,\sigma_{S}\;\times\;2\,f_{\ell}\,{\mathcal{L}}\
,$ $\displaystyle B$ $\displaystyle=$ $\displaystyle
K_{B}\,\varepsilon_{B}\,\sigma^{(0)}_{B}\;\times\;2\,f_{\ell}\,{\mathcal{L}}\;=\;K_{B}\,\sigma_{B}\;\times\;2\,f_{\ell}\,{\mathcal{L}}$
(23)
where $\varepsilon$ and $K$ respectively denote the total cut efficiencies and
the $K$-factors, which we have worked out in Section 3.1, cf. Table 1, and
Section 3.3, cf. Eqs. (17), (18), (21) and (22). The total efficiencies are a
product of single-step efficiencies, i.e.
$\varepsilon=\prod_{i}\varepsilon_{i}$. The factor $f_{\ell}=2$ accounts for
including the decay channels that involve muons and their associate neutrinos.
Notice that the Higgs boson mass enters in our simulation in two, potentially
different ways. In practice, the Higgs boson mass that we used to generate the
signal need not be the same as the Higgs boson mass we use to formulate the
analysis. We refer to the former as the injected mass $M^{\mathrm{inj}}_{h}$
in the text, while we have already introduced the terminology of the latter as
the “test” or “hypothesis” Higgs boson mass $M_{h}$. However, for simplicity
we take the generation level Higgs boson mass and the analysis level Higgs
boson mass to be equal, $M^{\mathrm{inj}}_{h}=M_{h}$. A discussion on how
different generation versus analysis masses would change our results can be
found in the Appendix B.1.
cuts & | $2\,\Delta/$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$
---|---|---|---|---|---|---|---|---|---|---
selections | $\mathrm{GeV}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $165\quad[20]$ | $170\quad[20]$ | $180\quad[20]$
$\sigma^{(0)}$ | | $19.35$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 20$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $17.77$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $14.19$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 14$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$
| | 1.0 | 1.0 | 0.21 | 1.0 | 1.0 | 0.19 | 1.0 | 1.0 | 0.15
lepton & | | $10.66$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 24$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $9.869$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 22$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $7.946$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$
MET cuts | | 0.551 | 0.45 | 0.17 | 0.555 | 0.45 | 0.15 | 0.560 | 0.45 | 0.12
as above & | | $8.572$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 0.0010 | $7.967$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 92$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $6.471$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 74$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
$\geq 2$ jets | | 0.443 | 0.0087 | 0.99 | 0.448 | 0.0087 | 0.90 | 0.456 | 0.0087 | 0.72
as above & | | $5.195$ | $6997$ | 0.0017 | $4.735$ | $6997$ | 0.0015 | $3.691$ | $6997$ | 0.0011
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.269 | 0.0032 | 0.99 | 0.266 | 0.0032 | 0.88 | 0.260 | 0.0032 | 0.68
naive $h$-reco | $50$ | $5.422$ | $6492$ | 0.0019 | $4.983$ | $6492$ | 0.0017 | $3.911$ | $6749$ | 0.0013
| | 0.280 | 0.0030 | 1.07 | 0.280 | 0.0030 | 0.96 | 0.276 | 0.0031 | 0.73
naive $h$-reco | $30$ | $3.948$ | $4108$ | 0.0022 | $3.897$ | $4108$ | 0.0021 | $3.039$ | $4199$ | 0.0016
| | 0.204 | 0.0019 | 0.98 | 0.219 | 0.0019 | 0.95 | 0.214 | 0.0019 | 0.72
naive $h$-reco | $48$ | $4.657$ | $2965$ | 0.0035 | $4.214$ | $3210$ | 0.0029 | $3.232$ | $3539$ | 0.0020
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.241 | 0.0013 | 1.36 | 0.237 | 0.0015 | 1.16 | 0.228 | 0.0016 | 0.84
naive $h$-reco | $20$ | $3.080$ | $1374$ | 0.0051 | $2.876$ | $1512$ | 0.0042 | $2.219$ | $1676$ | 0.0029
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.159 | 62$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.33 | 0.162 | 69$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.15 | 0.156 | 76$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.83
comb. $h$-reco | $50$ | $7.105$ | $6816$ | 0.0024 | $6.557$ | $7117$ | 0.0020 | $5.241$ | $7396$ | 0.0015
| | 0.367 | 0.0031 | 1.37 | 0.369 | 0.0032 | 1.21 | 0.369 | 0.0034 | 0.94
comb. $h$-reco | $20$ | $4.827$ | $3094$ | 0.0035 | $4.577$ | $3191$ | 0.0032 | $3.657$ | $3255$ | 0.0024
| | 0.249 | 0.0014 | 1.38 | 0.258 | 0.0015 | 1.26 | 0.258 | 0.0015 | 0.99
comb. $h$-reco | $50$ | $6.346$ | $3336$ | 0.0043 | $5.884$ | $3697$ | 0.0035 | $4.679$ | $4098$ | 0.0025
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.328 | 0.0015 | 1.75 | 0.331 | 0.0017 | 1.51 | 0.330 | 0.0019 | 1.12
comb. $h$-reco | $30$ | $5.586$ | $2217$ | 0.0057 | $5.159$ | $2488$ | 0.0046 | $4.083$ | $2756$ | 0.0032
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.289 | 0.0010 | 1.89 | 0.290 | 0.0011 | 1.61 | 0.288 | 0.0013 | 1.20
comb. $h$-reco | $20$ | $4.616$ | $1525$ | 0.0068 | $4.280$ | $1731$ | 0.0054 | $3.404$ | $1933$ | 0.0038
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.239 | 69$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.89 | 0.241 | 79$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.60 | 0.240 | 88$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.19
comb. $h$-reco | $16$ | $4.075$ | $1235$ | 0.0074 | $3.784$ | $1396$ | 0.0060 | $3.017$ | $1575$ | 0.0042
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.211 | 56$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.85 | 0.213 | 63$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.58 | 0.213 | 72$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.17
comb. $h$-reco | $10$ | $3.025$ | $787.9$ | 0.0087 | $2.624$ | $905.0$ | 0.0064 | $2.103$ | $1006$ | 0.0046
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.156 | 36$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.72 | 0.148 | 41$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.36 | 0.148 | 46$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.02
Table 4: Impact of the different levels of cuts on the $e\nu_{e}$+jets final
states for the $gg\to h\to WW$ production and decay signal and the $W$+jets
background as obtained from SHERPA. Cross sections $\sigma_{S}$, $\sigma_{B}$,
acceptances $\varepsilon_{S}$, $\varepsilon_{B}$ and $S/B$, $S/\sqrt{B}$
ratios are shown for Higgs boson masses of $M_{h}=165\ \mathrm{GeV}$,
$M_{h}=170\ \mathrm{GeV}$ and $M_{h}=180\ \mathrm{GeV}$. Note that
$\tilde{m}_{ij}=m_{ij}/\mathrm{GeV}$ and $\tilde{\delta}=\delta/\mathrm{GeV}$.
Significances were calculated using Eqs. (4.2) assuming ${\mathcal{L}}=10\
\mathrm{fb}^{-1}$ of integrated luminosity, counting both electrons and muons
and combining Tevatron experiments.
We can now go ahead and calculate the $S/B$ ratios and $S/\sqrt{B}$
significances. For various Higgs boson mass hypotheses, Table 4 and Tables
7–10 of Appendix B.1 list signal and $W$+jet-background cross sections,
acceptances, $S/B$ ratios and significances at different levels of cuts for
the selection procedures discussed in this subsection. The SHERPA simulation
runs obtained with the MSTW2008 LO PDFs have been used to extract the results
of all tables except those of Table 9 presented in Appendix B.1 which are
based on a set of runs taken with the CTEQ6.6 PDFs. In Appendix B.1 we will
then also briefly discuss the differences that can be seen between the
predictions for the two PDF sets.
We now turn to the discussion of the tables. Their setup is as follows: the
rows represent different stages in the cut-flow, Higgs boson reconstruction
strategies, and mass window cuts, while the third through fifth columns
contain the outcomes for different Higgs boson masses. The second column
indicates the mass window cut (in $\mathrm{GeV}$, referred to as $\Delta$ in
the text), which has been applied to all reconstructed Higgs boson candidates.
Similarly, the number in square brackets next to each Higgs boson mass is the
dijet mass window cut (referred to as $\delta$, also in $\mathrm{GeV}$). At
every analysis level, six numbers are displayed for each Higgs boson mass. The
top row displays the LO signal cross section (in $\mathrm{fb}$), the LO
$W$+jets cross section (in $\mathrm{fb}$) and $S/B$ at $10\ \mathrm{fb}^{-1}$
of integrated luminosity, calculated including $K$-factors and factors of 2
following Eqs. (4.2). The bottom three numbers in each table entry are the
signal and background efficiencies and $S/\sqrt{B}$. Of these entries,
$S/\sqrt{B}$ is displayed in bold. For the first set of tables, Table 4 and 7
(see Appendix B.1), we concentrate on Higgs boson masses greater than $\approx
162\ \mathrm{GeV}$ – above the $WW$ threshold. Higgs boson masses below the
$WW$ threshold have additional challenges, which we explore in a later
subsection.
The rows are divided into three groups. In the first group, rows 1–4, the
baseline selection cuts, as described in Section 4.1, are applied.101010Note
that the “lepton & MET cuts” level also includes a lepton–jet separation of
$\Delta R^{\mathrm{lep}\textrm{--}\mathrm{jet}}>0.4$ in the presence of jets.
In the second group, rows 5–8, events are selected using the “naive” criteria,
then retained if their reconstructed sum falls within various Higgs boson and
dijet mass windows. Finally, in the last set of rows, 9–15, we select events
with the “comb. $h$-reco” algorithm, then apply several different mass
windows. The effect of the mass window cuts, with either the “naive” or “comb.
$h$-reco" selection scheme, are fairly intuitive; mass windows always help
because they emphasize the peaks in the signal in comparison to a featureless
$W$+jets background. Tighter mass windows are usually, but not always,
better. Clearly, among the three groups the combinatorial selections give the
best significances, followed by the naive ones, which already improve over the
baseline selection cuts.
Comparing rows with identical cuts but different selections (“naive” versus
“comb. $h$-reco”), such as rows 5 and 9, or 7 and 11, the combinatorial Higgs
boson reconstruction is better across all Higgs boson masses by roughly 30%.
The difference can be traced to events where one of the hardest jets comes
from I/FSR rather than from one of the jets of the $W$ decay. Had we truncated
our treatment of the background at the matrix-element level (or even at
matrix-element level plus some Gaussian smearing, as in Ref. [7], additional
jet activity arising from I/FSR would be absent and the “comb. $h$-reco”
scheme would give the same result as the “naive $h$-reco” scheme.
Incorporating these relevant I/FSR jets using a complete, matrix-element plus
parton-shower treatment of the background, we notice that the “naive” scheme
is no longer the best option. The ME+PS merging thereby allows us a fully
inclusive description of $W$+2-jet events on almost equal footing with the
related NLO calculation, however with the advantage of accounting for multiple
parton emissions at leading-logarithmic accuracy. These effects are pivotal to
obtain reliable results for the combinatorial selections.
Showering effects are not just limited to the background. In particular, the
width of the Higgs boson candidates reconstructed from showered events is much
broader than the reconstructed width derived from parton level. In fact, after
showering, the reconstructed Higgs boson peak is typically so broad that the
tightest mass windows used in the tables ($\Delta=5\ \mbox{and}\ 8\
\mathrm{GeV}$) cut out some of the signal and yield worse significances than
broader windows. For example, the combinatorial selections supplemented by a
dijet mass window yield FWHM of about $10\ \mathrm{GeV}$ at the shower level,
while the FWHM at the parton level are reduced down to $2\ \mathrm{GeV}$ –
that basically is the width of one bin. If we relied on the matrix-element
level results, we would obtain far too promising $S/\sqrt{B}$. Focusing on the
$M_{h}=180\ \mathrm{GeV}$ test point and the “comb. $h$-reco” with a dijet
mass window, we would find the significances increasing from $1.5$ for
$\Delta=25\ \mathrm{GeV}$, $2.2$ for $\Delta=10\ \mathrm{GeV}$ to $3.0$ for
$\Delta=5\ \mathrm{GeV}$. These numbers should be contrasted with those in
Table 4, namely $1.12$, $1.19$ and $1.02$, respectively.
Figure 2: $S/\sqrt{B}$ significances for Higgs boson masses varying from
$M_{h}=110$ to $220\ \mathrm{GeV}$ after different levels of cuts. The numbers
are taken from Tables 4 and 7–10, which reflect in more detail the outcome of
the analysis based on the invm selection procedure for $e\nu_{e}$+jets final
states originating from the $gg\to h\to WW$ signal and the $W$+jets
background. Results are shown for Higgs boson masses below and above the $WW$
mass threshold; the threshold region has been left out though. All
significances were calculated according to Eqs. (4.2) under the assumption of
an integrated luminosity of ${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ and including
electron and muon channels, i.e. $f_{\ell}=2$.
To conclude this discussion, it is illustrative to show a plot of the
significances versus Higgs boson masses for various selections as presented in
the tables (Table 4 and Tables 7–10 in Appendix B.1), all of which is
summarized in Figure 2. The significance, at least after this level of
analysis, reaches a maximum of $\sim 2.0$. The highest significance occurs, as
expected, near the $WW$ threshold. For heavier Higgs bosons, as the $WW$ decay
mode becomes subdominant to the $ZZ$ mode, the significance drops slowly,
reaching $\sim 1.0$ at $M_{h}\sim 185\ \mathrm{GeV}$. By gradually enhancing
the selections the gain in significance remains approximately equal over the
whole region of large $M_{h}$; this is indicated by the parallel shifts of the
respective significance curves. Hence, the differences seen in the
significances per $M_{h}$ test point are mainly driven by the behaviour of the
total inclusive cross section for the signal. Looking back at Tables 4 and 7,
we in fact realize that the acceptances $\varepsilon_{S}$ and
$\varepsilon_{B}$ are rather similar at any selection step (for each row),
only mildly varying across the different Higgs boson test mass points. For
Higgs boson masses below threshold (Table 8 in Appendix B.1), as we will
discuss in the next sections, the drop-off is more severe. Not only does the
branching fraction to $WW^{*}$ fall rapidly, but the signal becomes more
background-like once the two $W$ bosons from the Higgs boson decay cannot both
be produced on-shell. The significances shown in Figure 2 reflect our best
estimates, however, we have also performed several checks on the stability of
these significances under slight variations in the analysis. These checks not
only include – as mentioned earlier – varying parton distribution functions,
but also varying jet definitions, etc. and are summarized in Appendix B.1.
#### 4.2.1 Reconstruction below the on-shell diboson mass threshold
In the above-threshold case there is good hope that the idealized approach of
considering the neutrino as a fully measurable particle will not lead to
results, which are sizably different from those obtained by a realistic
treatment of neutrinos. This is based on the fact that in most cases the
leptonic $W$ will be on its mass shell. The approximation $m_{e\nu_{e}}\approx
M_{W}$ can in principle be used to determine the neutrino’s longitudinal
momentum – up to a twofold ambiguity – by employing the lepton and missing
transverse energy measurements. Below the $WW$ mass threshold one of the $W$
bosons will be off-shell, so that the simple ansatz in calculating
$p_{\parallel,\nu_{e}}$ will be rather inaccurate. Hence, it a priori is not
clear whether an event selection based on invariant-mass windows will give an
overall picture that can be maintained in more realistic scenarios.
Nevertheless here we briefly establish what kind of significances may be
achievable assuming we had knowledge about the off-shellness (the actual mass)
of the leptonic $W$ boson. This will give us a benchmark, which we may use to
assess more realistic reconstruction approaches.
When we apply the same analysis as above the $WW$ threshold, we find
significances as summarized in Table 8 of Appendix B.1. They are visualized in
Figure 2. The numbers demonstrate that we quickly lose sensitivity below the
$WW$ threshold, in particular for test points $M_{h}\lesssim 130\
\mathrm{GeV}$. This happens for three reasons (which apply to the signal
only): one factor is the decline of the total inclusive signal cross section
$\sigma^{(0)}_{S,\mathrm{NNLO}}$ towards lower $M_{h}$, which actually is
comparable to that seen for large $M_{h}$. As shown in Table 1, this effect is
not as drastic as one would assume from the drop in the $h\to W^{*}W^{*}$
branching ratios; it is partly compensated by the rising gluon–gluon fusion
rate for low $M_{h}$. In contrast to the above-threshold case, there are yet
two more factors coming into the equation. Firstly, the basic selection cuts
affect the signal more severely;111111For $M_{h}=110\ \mathrm{GeV}$, only
about 7% of the events survive, while 45–49% of the signal is kept above
threshold (cf. the respective 1st rows in Table 8 and 3rd rows in Tables 4 and
7). secondly, the low $M_{h}$ signals that pass the baseline selection are
often penalized because of substantial off-shell effects. In particular, the
Higgs boson propagator can be pushed far off-shell and the Higgs boson
reconstruction will fall outside the mass window, such that the event will be
discarded. The tendency for lighter Higgs bosons to go off-shell increases,
since the basic cuts make it extremely unlikely for the leptonic and hadronic
$W$ masses to drop below $\sim 30\ \mathrm{GeV}$.
Figure 3: Mass spectra $m_{e\nu_{e}jj^{\prime}}$ after the combinatorial
reconstruction of Higgs boson candidates for very wide Higgs boson mass
windows. Results are shown for $M_{h}=130\ \mathrm{GeV}$ and $M_{h}=180\
\mathrm{GeV}$ and $e\nu_{e}$+jets final states originating from the $gg\to
h\to WW$ signal (peaked distributions) and the $W$+jets background (flat
distributions).
Figure 3 shows the $m_{e\nu_{e}jj^{\prime}}$ spectra including shower effects
for signals and backgrounds at $M_{h}=130$ and $180\ \mathrm{GeV}$ after the
combinatorial Higgs boson reconstruction has been applied using wide Higgs
boson mass windows ($\Delta\equiv M_{h}$). The parton showering washes out the
peaks, therefore reduces and broadens them. Both signal distributions develop
a softer tail above $M_{h}$ as a result of the jet combinatorics. For
$M_{h}=130\ \mathrm{GeV}$, the tail plateaus due to the off-shell effects
mentioned earlier. Figure 3 also illustrates why the value of the significance
jumps up significantly (as shown in Figure 2) when we tighten the Higgs boson
mass window from $\Delta=25$ to $10\ \mathrm{GeV}$ for $M_{h}=130\
\mathrm{GeV}$. This effect arises because we place our window cuts in a
steeply rising $W$+jets background.
When we studied which choice of mass window gives us the best results in terms
of separating signal from background, it came as somewhat of a surprise that
we did not have to alter the additional dijet mass constraint of
$M_{W}-\delta<m_{jj^{\prime}}<M_{W}+\delta$, $\delta=20\ \mathrm{GeV}$. Our
studies indicate that it is helpful to have the hadronically decaying $W$
boson to be close to its on-shell mass $M_{W}$. The $W$ boson decaying
leptonically is then forced to go off-shell ($m_{e\nu_{e}}<M_{W}$), a
kinematic configuration at odds with most $W$+jets events.Cutting on
$m_{jj^{\prime}}$ therefore helps suppress the dominant background and,
moreover, should also be convenient to demote the production of multi-jets
efficiently.
For the tighter Higgs boson mass windows, our results show that a simple one-
sided lower cut on $m_{jj^{\prime}}$, i.e. $m_{jj^{\prime}}>M_{W}-\delta$ is
slightly more efficient than using any type of dijet mass window. The one-
sided cut improves the significances as given in Table 8 by 1–2%. The removal
of the upper bound on $m_{jj^{\prime}}$ has however negligible effects on
selections using broad Higgs boson mass windows. As a consequence of keeping
an $m_{jj^{\prime}}$ constraint the leptonically decaying $W$ will almost
always be off-shell, such that the reconstruction of the longitudinal
component of the neutrino’s four-momentum cannot succeed without a good guess
of the mass of the $e\nu_{e}$ pair. We will address this issue in Section 4.3.
#### 4.2.2 Effect of the subdominant backgrounds
In this section, we examine to what extent the significances of the ideal
Higgs boson reconstruction will be diluted by contributions from the
electroweak and top-pair production of the $e\nu_{e}$+jets final states. To
this end we apply the analysis as established so far, without any
modification.
Figure 4: Cut efficiencies $\varepsilon_{B}$ and $\varepsilon_{S}$ for
$e\nu_{e}$+jets final states and different invm naive and combinatorial Higgs
boson reconstructions taking the mass points $M_{h}=130$, $170$ and $210\
\mathrm{GeV}$. The selections are labelled by the row numbers as assigned in
Tables 4 and 7–10; row 3 marks the baseline selection (used as benchmark). The
left pane exhibits the efficiencies found for the minor backgrounds –
electroweak and top–antitop pair production (dashed and solid lines,
respectively) – whereas the right pane displays the $\varepsilon_{B}$ for the
$W$+jets background (dashed lines) as well as the $\varepsilon_{S}$ of the
$gg\to h\to WW$ signal (solid lines). The two plots to the right compared with
each other nicely visualize why the combinatorial outperforms the naive
selection: the signal cut efficiencies get increased, while, for $W$+jets,
the cuts remain about as effective as for the naive approach. Also notice the
drop of the $M_{h}=130\ \mathrm{GeV}$ signal curves – they show the penalty in
employing the same basic cuts as above the $WW$ threshold. Figure 5: Single-
background and total significances as a function of $M_{h}$ for three
different Higgs boson candidate invm selections as denoted on top of each
panel. The $e\nu_{e}$+jets final states are generated from the $gg\to h\to WW$
signal, $W$+jets, electroweak and $t\bar{t}$ production backgrounds. All
$S/\sqrt{B_{i}}$ were calculated according to Eqs. (4.2) assuming an
integrated luminosity of ${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ and including
electron and muon channels, i.e. $f_{\ell}=2$. The total-background
significances were obtained with Eq. (24). The lower plots hold the ratios of
$S/\sqrt{B_{\mathrm{tot}}}$ over $S/\sqrt{B}$ for $W$+jets only. Note that
the large $t\bar{t}$ significances obtained by the naive selection (left
panel) result from the failure of the relatively harder leading-jet pairs to
satisfy the mass window constraints.
The first thing to notice is the total inclusive LO cross sections for these
minor backgrounds are ${\mathcal{O}}(1)\ \mathrm{pb}$ – substantially smaller
than the $W$ production contribution. After the application of the basic cuts,
the inclusive $e\nu_{e}$+2-jet cross sections drop to about $0.5\
\mathrm{pb}$, a factor of 40 below the major background. Including all the
various $K$-factors, see Table 1 and Eqs. (17), (21), (22), we find that the
total significance,
$\frac{S}{\sqrt{B_{\mathrm{tot}}}}\;=\;\frac{1}{\sqrt{\,\sum_{i}\left(\frac{S}{\sqrt{B_{i}}}\right)^{-2}}\,}\;=\;\left(\left(\frac{S}{\sqrt{B}}\right)^{-2}+\left(\frac{S}{\sqrt{B_{\mathrm{ew}}}}\right)^{-2}+\left(\frac{S}{\sqrt{B^{b\textrm{-}\mathrm{veto}}_{t\bar{t}}}}\right)^{-2}\right)^{-\frac{1}{2}}\
,$ (24)
at the basic selection level is only 2% smaller compared to the significance
$S/\sqrt{B}$ using only $W$+jets.121212The cross sections stated are LO-like
cross sections as obtained with SHERPA: before (after) the basic cuts, we find
$1.21$ and $0.886\ \mathrm{pb}$ ($540$ and $508\ \mathrm{fb}$) for the
electroweak and $t\bar{t}$ backgrounds, respectively; the resulting single-
background significances turn out to be almost 6 and 10 times larger than the
$W$+jets $S/\sqrt{B}$.
Switching from the baseline selection to a combinatorial selection and finite
dijet mass window, the single minor-background significances improve by up to
50% above (100% below) the $WW$ threshold. So, for beyond-baseline $h$
reconstructions, the significance corrections owing to the inclusion of the
minor backgrounds will be of the same order as before. This is documented in
both Figures 4 and 5. In the former we compare the cut efficiencies between
all backgrounds and the signal (shown together with the $W$+jets background
in the plots to the right in Figure 4) for various naive and combinatorial
Higgs boson candidate selections. Firstly, no background cut efficiency ever
exceeds any of the $\varepsilon_{S}$. In all cases the background curves
decrease more strongly when tightening the selection. Secondly, the pattern we
observe for the minor-background cut efficiency curves resembles by and large
those of the major background.131313In particular, the minor-background cut
efficiencies show more pronounced drops, if one enhances the baseline to a
Higgs boson candidate selection, introduces the dijet mass window $\delta$ or
tightens the $\Delta$ Higgs boson mass range. Notice that the naive Higgs
boson reconstruction very efficiently beats down the $t\bar{t}$ background.
This is because the two leading jets turn out harder compared to all other
cases. The presence of a sufficient number of subleading jets however makes
the selection based on jet combinatorics pick a pair of soft jets, and, on the
contrary almost too effective for Higgs boson masses above the $WW$ threshold.
For these reasons, the single-background significances plotted versus $M_{h}$
follow the trend found for the $W$+jets contribution, but remain well above
the $W$+jets significances.141414Both the electroweak and $t\bar{t}$
production significances show the same strong enhancement around $M_{h}=130\
\mathrm{GeV}$ as a result of the effect discussed around Figure 3 which is due
to the use of a tight Higgs boson mass window. As for $W$+jets, the minor
backgrounds fall rapidly for decreasing $M_{h}$. All of which is exemplified
in Figure 5 using three Higgs boson candidate selections, which impose a dijet
mass window (corresponding to the rows 8, 11 and 13 in Tables 4, 7 and 10),
namely the naive method with $\Delta=10\ \mathrm{GeV}$ (left panel) and the
combinatorial method with the same and broader window of $\Delta=25\
\mathrm{GeV}$ (middle panel). The total significances
$S/\sqrt{B_{\mathrm{tot}}}$ resulting from combining the three single
backgrounds are also shown. In fact, they only decrease by 1–5% as
demonstrated by the ratio plots in the lower part of Figure 5. The high
$M_{h}$ region is found to receive the larger, ${\mathcal{O}}(\mbox{5\%})$
corrections once the electroweak and $t\bar{t}$ contributions are included in
the overall background. As noted early in Section 3.2, experimenters have
estimated this fraction of events with 5% implying a 2.5% drop in
significance. It is reassuring to be able to confirm this expectation with our
results. Slightly contrary to the expectation, we identify the electroweak as
the leading minor background in all selections.
Based on these results it is easy to conclude; at this stage of our analysis
we do not have to worry about contributions from minor backgrounds. Although
additional handles exist to further reduce these backgrounds or supplement the
(here conservatively chosen) $b$-jet veto, it is of far more importance to
find ways to diminish the $W$+jets background. We postpone this discussion
until Section 4.4.
### 4.3 More realistic Higgs boson reconstruction methods
Up to this point we have ignored one big problem, namely the neutrino problem.
In our selection based on the reconstruction of invariant masses – which we
dubbed invm approach – we currently treat neutrinos as if we were able to
measure them like leptons. This is, of course, unrealistic and before we can
talk about further significance improvements, we have to investigate in which
way our analysis may fall short when switching to more experimentally
motivated Higgs boson candidate selections. Under experimental conditions,
missing energy is taken from the $\vec{p}_{T}$ imbalance in the event.
However, in our analysis we then make a small simplification and identify the
missing energy with the neutrino’s transverse momentum as given by the Monte
Carlo simulation.
There are multiple choices for how to proceed.Recall that whatever method we
pick acts as a selection criterion; we decide which two jets to keep in the
event based on these variables, therefore we want to design variables, which
are best at correctly picking out the jets from a Higgs boson decay. One way
to proceed is to give up complete reconstruction and to work solely with
transverse quantities; this is clean and unambiguous, but we throw out
information. The second approach is to attempt to guess the longitudinal
neutrino momentum by requiring that some or all of the final-state objects
reconstruct an object we expect, such as a $W$ or Higgs boson. Full
reconstruction then gives access to a larger set of observables, therefore
keeps more handles and information, but it is also more ambiguous.
$M_{h}/M_{W}$ | | pzmw | | pzmh | | mt | | mtp
---|---|---|---|---|---|---|---|---
$<2$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$
$>2$ | | $m_{T,e\nu_{e}jj^{\prime}}$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$ | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$ | | $m_{T,e\nu_{e}jj^{\prime}}$
Table 5: The preferred choice of definition for the 4-particle transverse mass
shown for each of the more realistic Higgs boson candidate selections. The
$m_{T}$ definitions are given in Eqs. (8) and (9).
To remove some of the combinatorial headache, we use $e\nu_{e}jj^{\prime}$ and
$jj^{\prime}$ (transverse) mass windows as before; moreover, we can impose
criteria on subsets of the event. For example, if the (3-particle) mass of the
visible system $m_{ejj^{\prime}}$ is greater than the test Higgs boson mass,
that particular choice of jets is unphysical and we can move on to the next
choice. A second constraint we often impose is that the 4-particle transverse
mass does not exceed the upper bound on the Higgs boson mass window:
$m_{T,e\nu_{e}jj^{\prime}}\leq M_{h}+\Delta$. As to the definition of $m_{T}$,
we generally use the definitions stated in Section 2, see Eqs. (8) and (9).
For our selections, we found that the distinction of the two $m_{T}$
definitions in fact only matters when we calculate the 4-particle transverse
masses. Accordingly, each selection comes in two versions either using
$m_{T,e\nu_{e}jj^{\prime}}$ or $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$. In
Table 5 we summarize for each type which version is more appropriate to use
and in what context. Whenever we refer to a specific selection in due course,
we understand it according to the findings listed in Table 5.
With these criteria in hand, the different selection methods are specified as
follows:
* •
mt: we want to test a method where the final state of the Higgs boson decay
will be identified purely with the help of transverse masses rather than
invariant masses. To this end we calculate $m_{T,e\nu_{e}jj^{\prime}}$
according to Eqs. (8) or (9) and prefer the final state giving us the
4-particle transverse mass closest to $M_{h}$. Owing to
$m_{T,e\nu_{e}jj^{\prime}}\leq m_{e\nu_{e}jj^{\prime}}$ the test mass window
is placed on the lower side only,
$M_{h}-2\,\Delta<m_{T,e\nu_{e}jj^{\prime}}<M_{h}$, with double the size as
compared to the other selections. The advantage, but also disadvantage of the
method is there is no reconstruction. Avoiding reconstruction eliminates
uncertainties owing to constraining masses plus resolving ambiguities, but
means we have no access to longitudinal and invariant-mass observables
involving the neutrino.
For the next two selections, we aim to approximately determine the
longitudinal momentum of the neutrino, $p_{\parallel,\nu_{e}}$, using
knowledge about which value for $m_{e\nu_{e}}$ should likely be reconstructed
by the combined system,$p_{e}+p_{\nu_{e}}$.151515Provided the MET cut was
passed, we assume here that all MET in the event has been produced by a single
neutrino. When we write
$\displaystyle m^{2}_{\ast,i\nu_{e}k..}\;\approx\;m^{2}_{i\nu_{e}k..}\;=\;$
(25) $\displaystyle
m^{2}_{ik..}\,+\,2\left(\sqrt{m^{2}_{ik..}+\,\vec{p}^{2}_{ik..}}\sqrt{\vec{p}^{2}_{T,\nu_{e}}+p^{2}_{\parallel,\nu_{e}}}-\vec{p}_{T,ik..}\cdot\vec{p}_{T,\nu_{e}}-p_{\parallel,ik..}\,p_{\parallel,\nu_{e}}\right)\
,$
using the “(in)visible” subsystem notation, we note that such problems can be
solved up to a twofold ambiguity. The difference among the two selections lies
in how particles are grouped in Eq. (25) and how the twofold ambiguity is
resolved.161616If the solutions are complex-valued, we only assign the real
part to describe $p_{\parallel,\nu_{e}}$ with no ambiguity left to resolve.
* •
pzmw: in this selection we use the $W$ mass constraint to solve for the
neutrino momentum: $m^{2}_{*,e\nu_{e}}=M^{2}_{W}$. The ambiguity is then
resolved by picking the neutrino $p_{z}$ solution, which brings the
reconstructed mass $m_{e\nu_{e}jj^{\prime}}$ more closely to the Higgs boson
test mass $M_{h}$. For true signal events, it then is more likely to find the
reconstructed $m_{e\nu_{e}jj^{\prime}}$ matching the Higgs boson mass.
The tricky part is to pick the best choice for the $m^{2}_{*,e\nu_{e}}$
constraint – meaning is $M^{2}_{W}$ always optimal given that the $W$ boson
may be off-shell? For pzmw, we do the following: first, we inspect the
transverse mass, $m_{T,e\nu_{e}}$, of the $e\nu_{e}$ subsystem in each event.
If $m_{T,e\nu_{e}}\geq M_{W}$ we choose $m_{*,e\nu_{e}}=m_{T,e\nu_{e}}$,
otherwise we pick $m_{*,e\nu_{e}}=M_{W}$ as long as $M_{h}>2\,M_{W}$ or
$0.9<m_{T,e\nu_{e}}/M_{W}<1.0$. That is, above and around the $WW$ threshold,
we take $m_{e\nu_{e}}$ towards $M_{W}$. If below threshold and
$m_{T,e\nu_{e}}/M_{W}<0.9$, the mass estimate is chosen taking various
subsystem invariant and transverse masses into account but enforcing
$m_{*,e\nu_{e}}$ to lie between $m_{T,e\nu_{e}}$ and $M_{W}$. For example, if
$m_{jj^{\prime}}>2\,m_{T,jj^{\prime}}$ we set
$m_{*,e\nu_{e}}=m_{T,e\nu_{e}jj^{\prime}}-m_{jj^{\prime}}$ while otherwise
$m_{*,e\nu_{e}}=m_{T,e\nu_{e}jj^{\prime}}-m_{T,jj^{\prime}}$ unless
$m_{ejj^{\prime}}>m_{T,e\nu_{e}jj^{\prime}}$ where we say
$m_{*,e\nu_{e}}=m_{T,e\nu_{e}}$.
* •
pzmh: we again infer the neutrino’s longitudinal momentum from mass
constraints. Although technically similar to pzmw – with the “visible”
subsystem entering Eq. (25) now being $\\{e,j,j^{\prime}\\}$ – we here turn
the idea around and already require $m_{e\nu_{e}jj^{\prime}}\approx M_{h}$ in
order to solve for $p_{\parallel,\nu_{e}}$. That is to say we enforce the
combined system, $p_{ejj^{\prime}}+p_{\nu_{e}}$, to mimic a Higgs boson signal
mass while leaving us with reasonable leptonic $W$ masses $m_{e\nu_{e}}$ at
the same time. When reconstructing the signal these observables are likely
correlated, while for the background they are uncorrelated apart from
kinematic constraints.
The details of the method are: we specify the target mass via
$m_{*,e\nu_{e}jj^{\prime}}=M_{h}$ unless we find
$m_{T,e\nu_{e}jj^{\prime}}/M_{h}\geq 0.94$, i.e. the 4-particle transverse
mass turns out too large already so that
$m^{2}_{*,e\nu_{e}jj^{\prime}}=m^{2}_{T,e\nu_{e}jj^{\prime}}/0.95$ is the more
appropriate choice. We approximate the leptonic $W$ boson mass by
$m_{*,e\nu_{e}}=M_{h}-m_{jj^{\prime}}$ freezing it at $m_{*,e\nu_{e}}=M_{W}$
if this difference exceeds $M_{W}$. We however require
$\min(m_{*,e\nu_{e}})=m_{T,e\nu_{e}}$. Taking this estimate, we can form the
absolute difference $\delta m_{e\nu_{e}}=|m_{e\nu_{e}}-m_{*,e\nu_{e}}|$ using
the reconstructed mass $m_{e\nu_{e}}$ for each possible neutrino solution. In
the presence of two solutions we define, as a measure of the longitudinal
activity, $b_{ij}=m_{\perp,ij}\exp|y_{ij}|=\max\\{E_{ij}\pm
p_{\parallel,ij}\\}$ with $m^{2}_{\perp,ij}=m^{2}_{ij}+p^{2}_{T,ij}$ and pick
the solution that generates the smaller $b_{e\nu_{e}}$, i.e. the $e\nu_{e}$
subsystem less likely going forward. We do so unless the other solution’s
$\delta m_{e\nu^{\prime}_{e}}$ drops below $\delta m_{e\nu_{e}}$ and
$(b_{jj^{\prime}}+b_{e\nu^{\prime}_{e}})/(m_{jj^{\prime}}+m_{e\nu^{\prime}_{e}})<(b_{jj^{\prime}}+b_{e\nu_{e}})/(m_{jj^{\prime}}+m_{e\nu_{e}})+\delta
x$ is satisfied; this is when we pick conversely ($\delta x=0.5$ if
$M_{h}<2\,M_{W}$ otherwise $\delta x=1.0$). Finally, we ensure that the
$\\{e,\nu_{e},j,j\\}$ set minimizing $\delta m_{e\nu_{e}}$ will be preferred
by the overall selection among all sets reconstructing the same 4-particle
mass. Note that we do not reject the selected ensemble if the $\delta
m_{e\nu_{e}}$ deviation becomes too large; we leave this potential to be
exploited by supplemental cuts, which we discuss in Section 4.4.
Figure 6: Single $W$+jets background significances as a function of $M_{h}$
for 4 different realistic Higgs boson candidate selections using jet
combinatorics. The selection types are denoted on top of each panel. Results
are shown for 4 different mass window parameter settings each, overlaying the
sort of optimal case defined by the invm combinatorial $h$ reconstruction for
$\Delta=10$, $\delta=20\ \mathrm{GeV}$, which also serves as the main
reference. The $e\nu_{e}$+jets final states are generated from the signal,
$gg\to h\to WW$, and the $W$+jets background. All $S/\sqrt{B}$ were
calculated according to Eqs. (4.2) assuming an integrated luminosity of
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ and including electron and muon channels,
i.e. $f_{\ell}=2$. The lower plots hold the ratios of realistic over ideal
$S/\sqrt{B}$. The thick lines are drawn with respect to the main reference,
while the thin lines are taken from comparing to the invm combinatorial
selection relying on the same window parameters. Figure 7: $S/B$ ratios as a
function of $M_{h}$ for 4 types of realistic combinatorial Higgs boson
candidate selections, each with their window parameters chosen as to reach
maximal significance. The $e\nu_{e}$+jets final states are generated from the
$gg\to h\to WW$ signal and $W$+jets background. The $S/B$ were calculated
from the $\sigma_{S}$ and $\sigma_{B}$ as obtained after the selection, the
signal $K$-factors of Table 1 and $K_{B}$ as given in Eq. (17).
We explored in detail how each of these selection criteria compare to the
ideal case. Figure 6 shows the significances per $M_{h}$ test point that we
can achieve running the different combinatorial selections for various,
reasonable window parameters. In the upper panels we display them directly on
top of the (quasi optimal) ideal case, i.e. the invm combinatorial $h$
candidate selection. The plots on the right and in the center respectively
exhibit the results of the mt and pzmh methods for the whole $M_{h}$ test
range. The pzmw method yields similar, yet slightly worse results below the
$WW$ mass threshold compared to pzmh. Therefore, we split the leftmost pane
into two subplots: on the right, one finds the pzmw results for masses above
threshold; on the left we then already reveal the outcome of the mtp
selection, whose discussion we postpone until the next subsection.
The bottom-row plots of Figure 6 depict the respective survival fractions: the
significance ratios for each selection and different window parameters always
taken with respect to the quasi optimal case (cf. the thick lines with
symbols). To indicate the net effect of the more realistic approaches, we also
show in each case the significance loss relative to the invm selection using
exactly the same mass windows as the realistic one (cf. the thin lines with no
symbols). We observe serious losses, larger than 50%, if one uses selections
with tighter window parameters and/or runs for $M_{h}$ points further away
from $2\,M_{W}$. The reconstruction methods (pzwm, pzmh) do not work better
than 90% of the quasi optimal case. This only is improved by the mt approach
making use of broad Higgs boson mass windows where one can reach up to about
100%. However, lowering $\Delta$ here quickly results in losing sensitivity.
Related to the quasi optimal invm selection, the various results point us to
work with medium-sized Higgs boson mass windows always imposing the dijet mass
cuts. Tighter $\Delta$ constraints may help improve the outcome of the
reconstruction types, but are of disadvantage in the measurement.
In all selections the below-threshold region is especially problematic. One
might settle for 65–80% of the ideal significances, but if we want to get a
better handle on the low $h$ boson masses, in particular include the
$M_{h}=130\ \mathrm{GeV}$ mass point, we have to push further – which we do so
in Section 4.3.1.
Focusing on the above-threshold region, we see that pzmw slightly outperforms
pzmh over the whole range; only for the near-threshold region up to
$M_{h}=180\ \mathrm{GeV}$ this is topped by the mt selection for medium-sized
$\Delta$ windows. This is somewhat surprising, but seems plausible, if one
considers that signal events are central in rapidity and manifest themselves
in larger transverse activity on average. This has to be opposed to the
$W$+2-jets background whose events tend to populate phase space more along
the beam direction, which generates $y_{e\nu_{e}jj^{\prime}}$ distributions
peaking about half an unit away from zero rapidity. Nevertheless the
differences between the 3 methods are not conclusive per se; to some extent
the selections will shape distributions differently and it is easy to imagine
the picture changing if additional cuts are imposed. But, one has to bear in
mind, there is a second, very important criterion, the ratio $S/B$, which one
wants to maximize. For their optimal window parameters, we show in Figure 7
the $S/B$ curves of the more realistic and ideal selections as functions of
$M_{h}$. Even more surprisingly than before, the mt outperforms the neutrino
reconstruction methods and, moreover, the ideal $S/B$ are also beaten unless
$M_{h}<2\,M_{W}$. Based on this observation one may prefer the methods where
the selection utilizes transverse masses, with the only drawback of having no
$p_{\parallel,\nu_{e}}$ estimate available.
#### 4.3.1 More realistic reconstruction below the on-shell diboson mass
threshold
We have just seen that the significances achievable in more realistic
scenarios drop off considerably below the $WW$ mass threshold, amplifying the
loss already present in the ideal case. Therefore, it is of great importance
to learn how the reconstruction methods described above can be applied more
efficiently in the below-threshold region.
We noticed that the mt selection picks up background events, which often fall
outside (mostly above) the Higgs boson mass window. As a result, a somewhat
different class of background events survives the mt selection procedure
compared to utilizing the invm, ideal, approach. This is no surprise since we
have already argued that the Higgs boson decays yield an enhanced transverse
production with regard to the $W$+2-jet background. Using $M_{h}=130\
\mathrm{GeV}$, Figure 8 exemplifies this by means of the
$m_{e\nu_{e}jj^{\prime}}$ and
$m_{T,e\nu_{e}jj^{\prime}}/m_{e\nu_{e}jj^{\prime}}$ ratio distributions. We
clearly see the large impact on the $W$+jets results being a consequence of
enforcing a transverse- rather than invariant-mass window. Imposing the
constraint $80\leq m_{T,e\nu_{e}jj^{\prime}}/\mathrm{GeV}\leq 130$ on the
background is fairly equivalent to choosing events with larger longitudinal
components. This drives the associated invariant masses to higher values
whereas the $m_{T}/m$ ratios are shifted to lower ones.
Figure 8: Differential 4-particle mass spectra $m_{e\nu_{e}jj^{\prime}}$ (see
left panel) and
$m_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}/m_{e^{-}\bar{\nu}_{e}jj^{\prime}}$ ratio
distributions after combinatorial selection of Higgs boson candidates using a
$50\ \mathrm{GeV}$ symmetric mass window centered at a mass of $M_{h}=130\
\mathrm{GeV}$. The mt selection characteristics is compared to the ideal case
of choosing candidates according to the invm criteria. Eq. (8) was used to
compute the mt criteria. All results are shown for $e\nu_{e}$+jets final
states originating from the $gg\to h\to WW$ signal (solid lines) and the
$W$+jets background (dashed lines).
To give a more quantitative example, we consider the case $M_{h}=130\
\mathrm{GeV}$ for broad (tight) Higgs boson mass windows and $\delta=20\
\mathrm{GeV}$. If we select events using mt, then discard all those events
with an invariant mass $m_{e\nu_{e}jj^{\prime}}$ outside $M_{h}\pm 25\ (8)\
\mathrm{GeV}$, we find a quite impressive gain of 82% (360%). Similarly, if we
use the invm selection and apply further cuts removing events with
$m_{T,e\nu_{e}jj^{\prime}}$ values greater than $M_{h}$ or less than
$M_{h}-2\,\Delta=M_{h}-50\ (16)\ \mathrm{GeV}$, we observe that $S/\sqrt{B}$
improves by 12% (drops by 20%). As the invm selection has a better starting
$S/\sqrt{B}$ than mt, the final significances are similar in all cases,
however the essential point is we can improve the significance by combining
the transverse- and invariant-mass selections.
To exploit this potential in a more realistic scenario, we use (in a first
phase) the mt selection to pick out the 4-particle system, then (in a second
phase) reconstruct the neutrino momentum following the pzmw procedure.
Compared to the ideal case, this reconstruction works rather inefficiently in
identifying background events that yield invariant masses
$m_{e\nu_{e}jj^{\prime}}>M_{h}+\Delta$. As it is optimized to the features of
the Higgs boson decay signal, pzmw generates $m_{*,e\nu_{e}}$ estimates by
assuming $m_{T,e\nu_{e}}/m_{e\nu_{e}}$ ratios close to 1. For $W$+jets, these
choices often turn out to be sufficiently smaller than the actual masses of
the leptonically decaying $W$ boson. The $W$+jets background usually contains
an on-shell $W$, a lower $m_{*,e\nu_{e}}$ is not ideal and tends to bring down
the deconstructed, associated 4-particle masses $m_{e\nu_{e}jj^{\prime}}$. As
a result a good fraction of background events possessing true invariant masses
exceeding the upper bound on $m_{e\nu_{e}jj^{\prime}}$ (see left plot of
Figure 8) are shuffled back into the Higgs boson mass window.
Hence, we make adjustments to the pzmw reconstruction used in this symbiosis
of selections so that it performs better below threshold. The basic idea is to
maximally exploit the differences of signal and background in the leptonic and
hadronic $W$ mass distributions. If, after initial mt selection, we constrain
the transverse masses of the $e\nu_{e}$ and dijet subsystems by placing cuts
that favor $m_{T,e\nu_{e}}\approx M_{W}/2$ and $m_{T,jj^{\prime}}\approx
M_{W}$, we enforce off-shell $W\to e\nu_{e}$ decays while keeping those of
$W\to jj^{\prime}$ on-shell. This is beneficial to the signal and suppresses,
at the same time, the $W$+jets background. Here, we only add requirements on
the hadronic subsystem by imposing a minimum value for $m_{T,jj}$. We leave
the $m_{T,e\nu_{e}}$ potential to be exploited by the $M_{h}$-dependent cut
optimization (which we discuss in Section 4.4), since demanding an upper bound
on $m_{T,e\nu_{e}}$ would not only enhance mt but all realistic selections.
With this major adjustment, we can now design a promising mt+pzmw selection,
which we call mtp:
* •
mtp: the very first part of computing $m_{T,e\nu_{e}jj^{\prime}}$ is identical
to the mt selection. Subsequently, we require $m_{T,jj^{\prime}}$ to be
greater than
$m^{\mathrm{min}}_{T,jj^{\prime}}=3.51\,M_{h}-0.015\,M^{2}_{h}-135.4\
\mathrm{GeV}$, which we have parametrized in terms of $M_{h}$ for convenience.
This concludes the phase of testing the transverse criteria. Accepted
$\\{e,\nu_{e},j,j^{\prime}\\}$ candidates are subject to the neutrino
reconstruction, whose implementation deviates from pzmw to some extent: we
choose $m_{*,e\nu_{e}}=0.55\,m_{T,e\nu_{e}}+0.45\,Q$ freezing it below
$m_{T,e\nu_{e}}$; here we employ $Q=M_{h}-m_{jj^{\prime}}$ but keep it
constant above $M_{W}+\delta x$ ($\delta x=4\ \mathrm{GeV}$). As in all other
reconstruction methods we then solve for the longitudinal momentum of the
neutrino, cf. Eq. (25), and reject the particular
$\\{e,\nu_{e},j,j^{\prime}\\}$ choice if the solutions are degenerate or
generate a mean $m_{e\nu_{e}jj^{\prime}}$ mass deviating from $M_{h}$ by more
than $\Delta^{\prime}=\max\\{\Delta,20\ \mathrm{GeV}\\}$.171717Here, we do not
adopt tight $\Delta$ choices owing to the uncertainties intrinsic to the
reconstruction of the $p_{\parallel,\nu_{e}}$ component. In all other cases,
we pick the solution giving the smaller $|m_{e\nu_{e}jj^{\prime}}-M_{h}|$ and
accept it if the reconstructed $h$ mass falls inside the $\Delta^{\prime}$
window around $M_{h}$.
The mtp results we can reach in terms of significance and sensitivity are
known already from Figures 6 and 7. The increase in $S/\sqrt{B}$ (leftmost
plot in Figure 6) and $S/B$ (Figure 7) is highly visible for all of the $h$
test masses below threshold. For the mtp analyses with broadest $M_{h}$
windows, the significance effect is huge compared to the respective ideal
selections using the same $\Delta$ parameter. Likewise the signal over
background ratio turns out similarly or even better than in the quasi optimal
case.
Not only does mtp profit from the better selection performance, but we end up
with a fully reconstructed neutrino vector, giving us access to a larger set
of observables. However, because of the hard cut on $m_{T,jj^{\prime}}$, we
notice that the mtp-selected backgrounds are strongly sculpted, washing out a
number of shape differences. The effect of the $m_{T,jj^{\prime}}$ cut
moreover deteriorates as soon as $M_{h}>2\,M_{W}$. We obtain significances
similar to the other methods and since they sculpt the background less, there
is no real advantage to applying mtp above threshold, hence, we restrict its
use to the low-mass $h$ region. On that note, it remains to be studied whether
the $S/\sqrt{B}$, $S/B$ improvements came at the expense of further handles
for the $M_{h}$-dependent optimization. As a possible consequence the mtp
selection actually might be superseded by an – at this level – inferior
selection once enhanced by appropriately designed cuts.
### 4.4 Optimized selection – analyses refinements and (further) significance
improvements
Having established the more realistic overall picture, we here discuss steps
to achieve better signal over background discrimination. After baseline and
combinatorial selections, we are interested in cuts that help further increase
the significances obtained so far, i.e.
$\varepsilon_{B,i}<\varepsilon^{2}_{S,i}$. Our aim is to identify observables
for each $M_{h}$ test point that are sufficiently uncorrelated such that
simultaneous selections yield a total significance gain in the range:
$\max\left\\{\frac{\varepsilon_{S,1}}{\sqrt{\varepsilon_{B,1}}}\,,\,\frac{\varepsilon_{S,2}}{\sqrt{\varepsilon_{B,2}}}\right\\}\;<\;\frac{\varepsilon_{S}}{\sqrt{\varepsilon_{B}}}\;\leq\;\frac{\varepsilon_{S,1}\,\varepsilon_{S,2}}{\sqrt{\varepsilon_{B,1}\,\varepsilon_{B,2}}}\
.$ (26)
Above relation is written out for the example of two extra handles, but easily
extensible to multi-cut scenarios. As we have shown, the subdominant
backgrounds have negligible effects at this level; we therefore concentrate on
reducing the $W$+jets background, although the other backgrounds are still
included in computing the total significance. In order to be conservative
about mass resolutions for hadronic final states at the Tevatron, we will fix
the mass window parameters as $\Delta=\delta=20\ \mathrm{GeV}$.
As we have done in earlier sections, we divide the optimized selection into
three broad Higgs boson mass ranges: below threshold, near threshold, and
above threshold. As we vary the Higgs boson mass, we probe different kinematic
configurations for the background. For the lowest Higgs boson mass, many
background distributions ($H_{T}$, $p_{T}$, etc.) are steeply rising in the
region of interest, cut off from below by baseline kinematics. For
intermediate Higgs boson masses, the backgrounds tend to be flatter or peaked,
while the higher Higgs boson mass region overlaps with backgrounds that are
sharply falling. This basic shape behind (many) background distributions
drives which cuts are optimal for a given Higgs boson mass. We present a
number of distributions in Appendix B.2 to back up the many findings presented
in this subsection.
The optimized analysis employs only observables constructed out of the momenta
of the selected 4-particle system $\\{e,\nu_{e},j,j^{\prime}\\}$. More
inclusive observables may be useful: for example, the scalar sum of the two
selected-jet $p_{T}$ versus the two hardest-jet $p_{T}$,
$H_{T,jj^{\prime}}\leftrightarrow H_{T,12}$. However such observables may also
be subject to larger uncertainties from, e.g. the modeling of hard initial
state radiation. To reduce these uncertainties one could extend the ME+PS
program here, e.g. to inclusive $W$+3-jets, however this is beyond the scope
of this study.
In the optimized analysis the longitudinal observables,
$\Delta\eta_{j,j^{\prime}}$ or $\eta_{e\nu_{e}jj^{\prime}}$, offer only
moderate gains in significance (typically, we obtain gains on the order of
3–10% with larger gains occurring for heavy $M_{h}$), but they are insensitive
to the exact value of the Higgs boson mass and are therefore more broadly
applicable. The total 4-particle pseudo-rapidity $\eta_{e\nu_{e}jj^{\prime}}$
can only be used in the pzmw, pzmh or mtp selections, since reconstruction of
the neutrino is necessary. By the same logic, because the pseudo-rapidity
difference between the selected jets is independent of the neutrino, cuts on
$\Delta\eta_{j,j^{\prime}}$ can be used with all selections (though they are
not efficient for mtp). We find that the requirements
$|\Delta\eta_{j,j^{\prime}}|\lesssim 1.5$ and
$|\eta_{e\nu_{e}jj^{\prime}}|\lesssim 3.0$ work well for the entire range of
Higgs boson masses we are interested in, so we include these cuts into our
optimized (pzmh/w and mt) selection. In Appendix B.2 we present examples of
$|\Delta\eta_{j,j^{\prime}}|$ distributions (Figure 12) and
$\eta_{e\nu_{e}jj^{\prime}}$ spectra (Figure 13) after un-optimized
combinatorial selections, documenting the usefulness of these cuts.
The other useful observables we have found are all transverse, or in some
cases based on invariant masses. Unlike the longitudinal variables, the
optimal transverse variables and cut values depend strongly on the Higgs boson
mass. We also find that transverse and longitudinal observables are largely
uncorrelated in this study, so any gains in significance from selections in
the transverse observables will add to the gains from
$|\Delta\eta_{j,j^{\prime}}|$ and $|\eta_{e\nu_{e}jj^{\prime}}|$. We refer to
Figures 14 and 15 of Appendix B.2 conveniently illustrating this decorrelation
for the case $|\Delta\eta_{j,j^{\prime}}|$ versus $m_{e\nu_{e}jj^{\prime}}$.
$M_{h}$ | comb. $h$-reco | leading = major cut | gain | subleading cut | gain | | minor cut | gain
---|---|---|---|---|---|---|---|---
$\mathrm{[GeV]}$ | selection | $\mathrm{[range\ in\ GeV]}$ | $\mathrm{[\%]}$ | $\mathrm{[range\ in\ GeV]}$ | $\mathrm{[\%]}$ | | $\mathrm{[range\ in\ GeV]}$ | $\mathrm{[\%]}$
| mtp | $m_{jj^{\prime}}\ \ [75,\infty]$ | 17 | $H_{T,jj^{\prime}}\ \ [76,\infty]$ | 9 | | $p_{T,j}\ \ [38,\infty]$ | 4
$120$ | mt | $m_{T,jj^{\prime}}\ \ [72,\infty]$ | 70 | $H_{T,jj^{\prime}}\ \ [72,\infty]$ | 52 | | $H_{T,e\nu_{e}jj^{\prime}}\ \ [108,\infty]$ | 6
| pzmw | $m_{T,e\nu_{e}}\ \ [0,40]$ | 73 | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,120]$ | 38 | | $H_{T,jj^{\prime}}\ \ [64,\infty]$ | 10
| pzmh | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,120]$ | 63 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,124]$ | 59 | | $m_{\perp,jj^{\prime}}\ \ [76,\infty]$ | 48
| mtp | $p_{T,e\nu_{e}jj^{\prime}}\ \ [9,\infty]$ | 6 | $m_{T,e\nu_{e}}\ \ [0,48]$ | 4 | | $m_{\perp,jj^{\prime}}\ \ [76,\infty]$ | 4
$130$ | mt | $H_{T,jj^{\prime}}\ \ [72,\infty]$ | 38 | $m_{T,jj^{\prime}}\ \ [68,\infty]$ | 31 | | $m_{jj^{\prime}}\ \ [73,\infty]$ | 4
| pzmw | $m_{T,e\nu_{e}}\ \ [0,48]$ | 53 | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,130]$ | 42 | | $m_{\perp,jj^{\prime}}\ \ [72,\infty]$ | 11
| pzmh | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,130]$ | 48 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,134]$ | 46 | | $H_{T,jj^{\prime}}\ \ [68,\infty]$ | 25
| mtp | $p_{T,e\nu_{e}jj^{\prime}}\ \ [12,\infty]$ | 8 | $H_{T,jj^{\prime}}\ \ [68,\infty]$ | 3 | | $m_{T,e\nu_{e}}\ \ [0,60]$ | 3
$140$ | mt | $H_{T,jj^{\prime}}\ \ [68,\infty]$ | 24 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [16,\infty]$ | 15 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [15,\infty]$ | 6
| pzmw | $m_{T,e\nu_{e}}\ \ [0,56]$ | 30 | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,140]$ | 30 | | $m_{\perp,jj^{\prime}}\ \ [70,\infty]$ | 6
| pzmh | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,140]$ | 29 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,144]$ | 28 | | $H_{T,jj^{\prime}}\ \ [68,\infty]$ | 14
| mtp | $p_{T,e\nu_{e}jj^{\prime}}\ \ [17,\infty]$ | 14 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [116,\infty]$ | 3 | | $H_{T,e\nu_{e}jj^{\prime}}\ \ [116,\infty]$ | *
$150$ | mt | $p_{T,e\nu_{e}jj^{\prime}}\ \ [20,\infty]$ | 18 | $H_{T,jj^{\prime}}\ \ [60,\infty]$ | 10 | | $H_{T,jj^{\prime}}\ \ [56,\infty]$ | *
| pzmw | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,150]$ | 20 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 18 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 9
| pzmh | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [0,150]$ | 19 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,154]$ | 19 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 9
| mt | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 18 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [136,\infty]$ | 12 | | $\Delta\phi_{e,\nu_{e}}\geq 1.9$ | 3
$165$ | pzmw | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 18 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,170]$ | 17 | | $\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.06$ | 20
| pzmh | $p_{T,e\nu_{e}jj^{\prime}}\ \ [18,\infty]$ | 18 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,170]$ | 13 | | $\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.09$ | 15
| mt | $p_{T,e\nu_{e}jj^{\prime}}\ \ [21,\infty]$ | 20 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [140,\infty]$ | 16 | | $\Delta\phi_{e,\nu_{e}}\geq 1.7$ | *
$170$ | pzmw | $p_{T,e\nu_{e}jj^{\prime}}\ \ [19,\infty]$ | 20 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,176]$ | 13 | | $\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.12$ | 11
| pzmh | $p_{T,e\nu_{e}jj^{\prime}}\ \ [20,\infty]$ | 20 | $m_{e\nu_{e}jj^{\prime}}\ \ [0,176]$ | 9 | | $\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.16$ | 9
| mt | $p_{T,e\nu_{e}jj^{\prime}}\ \ [22,\infty]$ | 24 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [148,\infty]$ | 22 | | $H_{T,e\nu_{e}jj^{\prime}}\ \ [140,\infty]$ | *
$180$ | pzmw | $p_{T,e\nu_{e}jj^{\prime}}\ \ [21,\infty]$ | 23 | $H_{T,jj^{\prime}}\ \ [64,\infty]$ | 11 | | $1.06\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.22$ | 5
| pzmh | $p_{T,e\nu_{e}jj^{\prime}}\ \ [22,\infty]$ | 24 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [140,\infty]$ | 10 | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [136,182]$ | 5
| mt | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 28 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [156,\infty]$ | 27 | | $H_{T,e\nu_{e}jj^{\prime}}\ \ [148,\infty]$ | *
$190$ | pzmw | $p_{T,e\nu_{e}jj^{\prime}}\ \ [23,\infty]$ | 24 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [148,\infty]$ | 15 | | $1.12\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.30$ | 5
| pzmh | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 29 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [144,\infty]$ | 17 | | $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}\ \ [142,194]$ | 3
| mt | $H_{T,e\nu_{e}jj^{\prime}}\ \ [164,\infty]$ | 31 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 28 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [15,\infty]$ | 9
$200$ | pzmw | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 25 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [156,\infty]$ | 20 | | $1.18\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.40$ | 6
| pzmh | $p_{T,e\nu_{e}jj^{\prime}}\ \ [27,\infty]$ | 32 | $H_{T,e\nu_{e}jj^{\prime}}\ \ [156,\infty]$ | 25 | | $H_{T,e\nu_{e}jj^{\prime}}\ \ [144,\infty]$ | 4
| mt | $H_{T,e\nu_{e}jj^{\prime}}\ \ [172,\infty]$ | 36 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [25,\infty]$ | 27 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [15,\infty]$ | 8
$210$ | pzmw | $H_{T,e\nu_{e}jj^{\prime}}\ \ [160,\infty]$ | 24 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 23 | | $1.25\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.45$ | 14
| pzmh | $H_{T,e\nu_{e}jj^{\prime}}\ \ [162,\infty]$ | 36 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [30,\infty]$ | 36 | | $1.25\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.54$ | 7
| mt | $H_{T,e\nu_{e}jj^{\prime}}\ \ [180,\infty]$ | 39 | $m_{T,e\nu_{e}jj^{\prime}}\ \ [174,\infty]$ | 26 | | $p_{T,e\nu_{e}jj^{\prime}}\ \ [12,\infty]$ | 8
$220$ | pzmw | $H_{T,e\nu_{e}jj^{\prime}}\ \ [168,\infty]$ | 29 | $p_{T,e\nu_{e}jj^{\prime}}\ \ [24,\infty]$ | 22 | | $1.31\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.53$ | 15
| pzmh | $H_{T,e\nu_{e}jj^{\prime}}\ \ [172,\infty]$ | 49 | $p_{T,e\nu_{e}}\ \ [56,\infty]$ | 43 | | $1.30\leq\gamma_{jj^{\prime}|e\nu_{e}}\leq 1.58$ | 8
Table 6: Leading (optimal/major) and subleading cuts for each Higgs boson
mass and selection criteria (all selections are combinatorial with window
parameters $\Delta=\delta=20\ \mathrm{GeV}$). Gain in $S/\sqrt{B}$, in
percent, is shown after each cut. Having used the major discriminators
including pseudo-rapidity cuts (see text), next in the cut hierarchy are the
minor cuts shown in the two rightmost columns. The significance gains
associated with them are understood in addition to the major cut improvements.
Cuts marked with an asterisk have less than 2% improvement.
##### Below-threshold region:
For pzmw and pzmh, the reconstruction selections, the transverse mass of the
4-particle system, $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$, and (for pzmw at
low $M_{h}$) the transverse mass of the leptonic system, $m_{T,e\nu_{e}}$, are
the best observables. This result is really just a reiteration of the fact
that simple reconstruction selections work poorly for below-threshold Higgs
bosons. In our effort to make pzmh/w more flexible and apply them to below-
threshold scenarios, we have allowed the possibility for background
configurations that are inconsistent with a single parent resonance – such as
4-particle systems with $m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}>M_{h}$.
Removing this inconsistent region results in gains of
${\mathcal{O}}(\mbox{40\%})$.
For mt, the transverse mass of the dijet system or the scalar $p_{T}$ sum
formed with the selected jets, $H_{T,jj^{\prime}}$, are the most optimal, with
improvements of ${\mathcal{O}}(\mbox{50\%})$. The former cut takes advantage
of the fact that the signal jets originate in a (real or virtual) $W$ boson,
while the background jets come primarily from ISR – information that is not
exploited in the initial mt selection. Somewhat smaller gains come from
cutting on the total selected system’s transverse momentum,
$p_{T,e\nu_{e}jj^{\prime}}$.
For mtp, there is little optimization to be done. Much of the physics that is
behind the optimal cuts in the pzmh/w or mt cases has already been
incorporated into the selection process. Some improvement is possible by
cutting out the region with very low 4-particle transverse momentum,
$p_{T,e\nu_{e}jj^{\prime}}$.
We illustrate our findings by showing and commenting on a small collection of
spectra resulting from the baseline combinatorial selections; for more
details, we refer the reader to the discussion around Figure 16 of Appendix
B.2.
##### Near-threshold region:
For Higgs boson masses close to the $WW$ threshold, the 4-particle $p_{T}$ is
the most powerful additional handle. In Higgs boson production, as in other
colour-singlet resonance production, the $p_{T,e\nu_{e}jj^{\prime}}$
distribution is cut off at low values by soft-gluon resummation, and falls off
at high values because of parton kinematics. The result is a peaked
distribution. The hard scale of the process, dictated by the Higgs boson mass,
sets the initial ISR scale, thereby influencing where
$p_{T,e\nu_{e}jj^{\prime}}$ peaks. Since the Higgs boson is heavier than the
$W$ boson, $p_{T,e\nu_{e}jj^{\prime}}$ always peaks at higher values for the
signal compared to $W$ production. Even though the dominant background for our
study is $W$+2-jets, rather than $W$+0-jets, the argument still holds. The
peak in $p_{T,e\nu_{e}jj^{\prime}}$ for $W$+2-jets is still governed by
$M_{W}$ and continues to peak at lower values than the Higgs boson signal.
Selection criteria change the tails of the background
$p_{T,e\nu_{e}jj^{\prime}}$ distribution, but do not affect the location of
the peak. Cutting out the low-$p_{T,e\nu_{e}jj^{\prime}}$ region, improvements
on the order of 20% are possible. Distributions such as $H_{T,jj^{\prime}}$,
the scalar $p_{T}$ sum of the two selected jets, or
$H_{T,e\nu_{e}jj^{\prime}}$, the $H_{T}$ of the 4-particle system, also have
potential discriminating power. Signal versus background distributions in the
relevant variables after baseline combinatorial selections are shown in Figure
17, see Appendix B.2.
##### Above-threshold region:
For higher Higgs boson masses, the total amount of (transverse) energy in the
$W$+$jj^{\prime}$ system becomes the most powerful discriminator between the
signal and the background. Specifically, once $M_{h}\gtrsim 200\
\mathrm{GeV}$, the $H_{T}$ of the selected 4-particle system peaks at
significantly higher values than the background, regardless of the selection
technique. By cutting away the low-$H_{T}$ region, we find gains of order 25%
are possible. The 4-object $p_{T}$ remains a very useful observable, as does
the 4-particle transverse mass. Examples of signal versus background
distributions are shown in Figure 18, see Appendix B.2.
##### Results:
The optimal or “leading” or “major” cuts for the different Higgs boson mass
categories and selection methods are summarized in Table 6. To give some idea
how useful the single best discriminator is compared to other observables, we
also show the percent increase in significance for the second best or
“subleading” single discriminator. Separately we have also determined which
combinations of the leading discriminator (supplemented by the respective
pseudo-rapidity cuts discussed above) with a second observable give the
largest (additional) increase in significance. The second “minor” cuts of
these optimal two-variable selections are summarized in the last two columns
of Table 6. Note that in most cases these minor cuts do not involve the same
observables as the subleading cut; this is because the subleading
discriminator is typically strongly correlated with the leading discriminator,
and thus does not add much to the combined significance. Some ideas beyond the
application of minor cuts exist; we comment in Appendix B.3 on a few possible
routes one can take to enhance the optimized analyses presented here.
When looking for minor cuts, we found in addition to variables we have already
discussed, such as $\gamma_{jj^{\prime}|e\nu_{e}}$ (in Section 2),
$H_{T,(e\nu_{e})jj^{\prime}}$ and the 4-object $p_{T}$, a few other
observables, namely $p_{T,j}$, $\Delta\phi_{e,\nu_{e}}$ and
$m_{\perp,jj^{\prime}}$ to be beneficial. The first two, $p_{T,j}$ and
$\Delta\phi_{e,\nu_{e}}$ are common, so we do not repeat their definitions
here. The last minor cut observable, $m_{\perp,jj^{\prime}}$ is defined
through $m^{2}_{\perp,jj^{\prime}}=m^{2}_{jj^{\prime}}+p^{2}_{T,jj^{\prime}}$
exhibiting yet another way of defining a transverse mass. The additional gains
from the minor cuts are typically small, except close to the $WW$ threshold
and for the largest Higgs boson masses considered here. In particular for
$M_{h}>2\,M_{W}$, the boost of the $jj^{\prime}$ system in the reconstructed
4-object rest frame stands out as a helpful discriminator of secondary order;
for more details, we refer again to Appendix B.2 and the discussion around
Figure 19. Also, the dijet-system based handles, $m_{\perp,jj^{\prime}}$ and
$H_{T,jj^{\prime}}$ yield fairly substantial extra gains, but only at low
$M_{h}$ if we rely on the pzmh method. Once jets are picked stemming from the
backgrounds, the strict $M_{h}$ mass reconstruction of the selected 4-object,
as encoded in pzmh and amplified by the major cut given through
$m^{(\nu_{e})}_{T,e\nu_{e}jj^{\prime}}$, gives rise to the selection of less
energetic $j$, $j^{\prime}$ jets with preferred pair masses of
$m_{jj^{\prime}}\sim M_{W}-\delta$. The observables $m_{\perp,jj^{\prime}}$
and $H_{T,jj^{\prime}}$ exploit this fact conveniently, hence facilitate such
secondary improvements, as shown in Figure 19 for the example of $M_{h}=130\
\mathrm{GeV}$.
Figure 9: Total significances as a function of $M_{h}$ after including the
subdominant backgrounds and all major plus minor cuts (as specified in the
text and Table 6). The results (lines with squares) are shown for the four
types of more realistic Higgs boson candidate selections, which have been
advertised in this work. They are denoted on top of each panel. In all cases
the combinatorial approach using window parameters $\Delta=\delta=20\
\mathrm{GeV}$ has been applied for selecting the candidate set of particles.
As before in Figure 6, the ideal case (i.e. the invm combinatorial
reconstruction of the Higgs boson candidates using $2\,\Delta=\delta=20\
\mathrm{GeV}$) is taken as the main reference to compare the different
results. For each combinatorial selection, the outcome (lines with circles)
with no cuts applied (but using a slightly smaller mass window, $\Delta=15\
\mathrm{GeV}$) is also displayed to emphasize the effect of the cut
optimization. Note that the effect of the minor backgrounds has been neglected
in computing each of these reference curves. The $e\nu_{e}$+jets final states
are generated from the signal, $gg\to h\to WW$, the $W$+jets, electroweak and
$t\bar{t}$ backgrounds. All $S/\sqrt{B_{i}}$ were calculated using Eqs. (4.2)
and combined according to Eq. (24) assuming an integrated luminosity of
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ and including electron and muon channels,
i.e. $f_{\ell}=2$. The lower plots depict for each selection the ratios of
optimized over un-optimized $S/\sqrt{B}$ (lines with squares) and optimized
over ideal-case reference $S/\sqrt{B}$ (lines with circles). The thin blue
lines visualize in how far the final $S/\sqrt{B_{\mathrm{(tot)}}}$ results
suffer from the presence of the minor backgrounds. Figure 10: The $S/B$
ratios associated with the total significances presented in Figure 9. The
ratios are shown as a function of $M_{h}$ for the final selections: the
combinatorial pzmw (triangles), pzmh (circles), mt (squares) and mtp
(diamonds) selections using the window parameters $\Delta=\delta=20\
\mathrm{GeV}$ each supplemented by their specific optimization cuts, see Table
6. The reference curve (dashed line with diamonds) is the same as used in
Figure 7 representing the invm combinatorial reconstruction with
$2\,\Delta=\delta=20\ \mathrm{GeV}$, no extra cuts applied and neglecting the
impact of the minor backgrounds. The $e\nu_{e}$+jets final states were
generated from $gg\to h\to WW$ signal events, the $W$+jets, electroweak and
$t\bar{t}$ backgrounds. The $S/B_{i}$ were calculated using the $\sigma_{S}$
and $\sigma_{B}$ as obtained after the final selections, the signal
$K$-factors of Table 1 and the $K_{B}$ as given in Eqs. (17), (21) and (22).
Single ratios were combined according to
$S/B_{\mathrm{(tot)}}=1/\sum_{i}(S/B_{i})^{-1}$.
Finally, the baseline combinatorial and optimized combinatorial significances
are displayed as a function of the Higgs boson mass in Figure 9. The ratio
plots associated with each selection in the lower part of the figure visualize
the significance increase achieved by the optimization. Independent of the
selection, they also indicate a $\mathcal{O}(\mbox{10\%})$ drop of
significances caused by the subdominant backgrounds. Focusing on the
$S/\sqrt{B}$ ratios taken with respect to the ideal case (orange lines with
circles), these ratio plots emphasize that the optimized mt(p) (transverse)
and pzmw selections work best below and above the $WW$ mass threshold,
respectively. The related $S/B$ ratios presented in Figure 10 confirm these
findings. They turn out to be rather small, as a consequence of maximizing the
significance and trying to preserve most of the signal; both of which does not
allow for imposing too restrictive cuts. Advantageously, the actual number of
signal events, $S$, present in this $h\to WW$ channel is not small. Except for
the Higgs signal at $M_{h}=120\ \mathrm{GeV}$ (with $\mathcal{O}(\mbox{4})$
expected events), the optimized analyses usually leave us with hundreds of
signal events (50–300), if we assume an integrated luminosity of
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$. Even the $\sim$ 25 signal events for
$M_{h}=130\ \mathrm{GeV}$ are sufficient, particularly as $S/B$ increases to
$\sim 0.04$.
Clearly, as seen from Figure 9, the optimized significances for the four
different Higgs boson reconstruction methods are very similar. The best
significance is for Higgs boson masses close to the $WW$ threshold, suggesting
that it would be possible to achieve 95% confidence level exclusion in a
stand-alone analysis. For a Higgs boson mass in the range $130\lesssim
M_{h}\lesssim 150\ \mathrm{GeV}$, the optimized significance is between $0.7$
and $1.4$. Given the additional improvements expected from a full multivariate
analysis, this indicates that the semileptonic Higgs boson decay channel can
make a significant contribution to Higgs boson exclusion in this mass range.
## 5 Conclusions, caveats, and prospects
We have presented a systematic study of the prospects for extending the
Tevatron exclusion reach for a Standard Model Higgs boson by including the
final states arising from semileptonic Higgs boson decays. We have used a
realistic simulation of the Higgs signal and the relevant Standard Model
background processes to exhibit the kinematic differences between the signal
and background. We have used three qualitatively different approaches to
extracting the event kinematics, one based on transverse observables and the
others based on approximate even-by-event full reconstruction. We have shown
that all three approaches give similar results when one optimizes selections
based on several discriminating observables. The details of the optimization
depend on the Higgs boson mass, and in particular on whether it is below,
near, or above the threshold for decay to two on-shell $W$ bosons.
The optimized significances that we have achieved are not sufficient for
stand-alone Higgs boson exclusion except in the most favorable case where the
Higgs boson mass is close to threshold. However the sensitivities shown here
are certainly promising as ingredients to a combined multi-channel analysis.
One important caveat is that the signal to background ratios for this type of
analysis are fairly small, on the order of a percent, as illustrated in Figure
10. This means that Higgs boson exclusion is sensitive to relatively small
systematic errors in the modeling of the backgrounds, notably the dominant
background from $W$+jets. However we have shown here that an experimental
analysis has multiple over-constrained handles on the kinematic features of
the data, providing extra cross-checks. In addition, the Higgs boson candidate
selection employed here to identify the two jets from the Higgs boson decay is
by design rather stable against effects from extra hard radiation; this
reduces the uncertainty in the background modeling.
The techniques described here are applicable to Higgs boson searches at the
CERN Large Hadron Collider. At 7 TeV center-of-mass collision energy, the
Higgs boson production cross section increases by a factor of $\sim$ 30, while
the $W$+jets background increases by a factor of $\sim$ 20\. With no hard
upper limit of the amount of data available, one can use more restrictive
selections and improve both the signal to background ratio and the overall
significance. A recent analysis by the ATLAS collaboration [91] employed a
4-object invariant mass reconstruction and a dijet mass window similar to the
baseline naive analysis used here, but applied to the heavy Higgs boson mass
region $M_{h}>240\ \mathrm{GeV}$. A CMS study [92] looked at the enhanced
signal to background ratio provided by focusing on VBF production (and thus
requiring two extra forward jets). Because of the large branching fraction,
the $h\to WW$ semileptonic mode might prove to be observable at the LHC over a
larger mass range than the $h\to WW$ dilepton mode; thus it could play a
critical role in establishing the nature of electroweak symmetry breaking, in
the event that the Higgs boson mass is $\sim$ 120 GeV.
## Acknowledgments
We would like to thank Bogdan Dobrescu for many lively and stimulating
discussions. We would also like to thank Thomas Becher, Li Lin Yang, Tanju
Gleisberg, Frank Petriello, John Campbell and Ciaran Williams for their help
in accomplishing necessary cross-checks. We are grateful to Frank Krauss and
Frank Siegert, as well as all other members of the SHERPA collaboration for
their continuous support. We gratefully acknowledge useful discussions with
Bob Hirosky and Lidija Zivkovic, as well as Lance Dixon, Gavin Salam, Walter
Giele, Rakhi Mahbubani, Patrick Fox and Agostino Patella.
Fermilab is operated by Fermi Research Alliance, LLC, under contract DE-
AC02-07CH11359 with the United States Department of Energy.
## Appendix A Appendix: Monte Carlo event generation
### A.1 Leading-order cross sections
Here we briefly report on the tests we have done to convince ourselves of the
correctness of the SHERPA leading-order cross section calculations. We found
satisfactory or better agreement in all our cross-checks, which we briefly
summarize here:
* •
For the case of $h\to e^{-}\bar{\nu}_{e}\,pp$ decays, we have compared
SHERPA’s branching ratios with those obtained by multiplying the Hdecay
results for $B_{W^{\ast}W^{\ast}}$ given in Table 1 times the PDG literature
numbers for $B(W\to e\nu_{e})\times B(W\to pp)=0.1075\times 0.676$. Using
AMEGIC++’s mode of calculating partial widths of $1\to N$ processes, we
determined $B_{e^{-}\bar{\nu}_{e}pp}=\Gamma(h\to
e^{-}\bar{\nu}_{e}\,pp)/\Gamma_{h}$ for the various Higgs boson masses and
respective widths of Table 1. The differences seen are at most on the few-
percent level.
* •
With the explicit knowledge of the SHERPA branching fractions we were able to
extract Higgs boson production rates at LO from the SHERPA signal cross
section calculations according to
$\sigma^{\mathrm{LO}}_{ggh}\;=\;\frac{\Gamma_{h}}{\Gamma(h\to
e^{-}\bar{\nu}_{e}\,pp)}\;\;\sigma^{(0)}_{e^{-}\bar{\nu}_{e}pp}\;=\;\frac{\sigma^{(0)}_{e^{-}\bar{\nu}_{e}pp}}{B_{e^{-}\bar{\nu}_{e}pp}}\
.$ (27)
The numbers that we obtained from this procedure compare well to numbers of
other LO calculations, for example the LO rates as evaluated in MCFM or
provided by Becher and Yang for verification purposes [93].
* •
SHERPA LO rates were computed for both finite top masses and in the infinite
top-mass limit. The ratio of the former over the latter cross section given as
$\frac{\sigma^{(0)}_{S}}{\sigma^{(0)}_{S,m_{t}\to\infty}}\;=\;\left|\,I\left(\frac{m^{2}_{t}}{M^{2}_{h}}\right)\right|^{2}$
(28)
singles out the dependence on the top mass versus Higgs boson mass ratio,
which is encoded by the function
$\displaystyle I(x)\;=\;6x+3x\,(4x-1)$ $\displaystyle\Biggl{\\{}$
$\displaystyle\frac{\Theta(1-4x)}{2}\left[\ln\left(\frac{1+\sqrt{1-4x}}{1-\sqrt{1-4x}}\right)-i\pi\right]^{2}-$
(29) $\displaystyle
2\,\Theta(4x-1)\arcsin^{2}\left(\frac{1}{2\sqrt{x}}\right)\;\,\Biggl{\\}}\ .$
Note that $|I(x)|^{2}$ attains unity as $x\to\infty$, while it vanishes for
$x\to 0$. Comparing the numerical cross section ratios with the analytical
values for $|I(x)|^{2}$, we found excellent agreement over the entire Higgs
boson mass range considered in this study.
### A.2 NLO calculations versus CKKW ME+PS merging
When compared to NLO calculations, it is evident that SHERPA’s CKKW merging
approach does not account for the virtual corrections to $V$+jets in their
entirety.181818For a brief summary of the basics of the CKKW merging, see Ref.
[70]. For the current generation of SHERPA Monte Carlo programs, the ME+PS
facilities have been extended to allow for truncated showering, which is a
major refinement over the CKKW approach, see Refs. [76, 94, 95]. The only
contributions enter through Sudakov form-factor terms at leading-logarithmic
accuracy used in the parton shower and to reweight the tree-level matrix
elements. The real-emission corrections however are included on a fairly
comparable level with respect to full NLO calculations.191919More recent
versions of SHERPA, from version 1.2.3 on, have been enhanced by the means to
generate, for a number of important processes, events at the hadron level with
a rate correct at next-to-leading order in $\alpha_{\mathrm{s}}$ [83, 84, 96].
To make this work, SHERPA relies on interfacing external one-loop amplitude
generators like MCFM [61, 97], BlackHat [98, 99] or, more generally, via the
Binoth Les Houches Accord [100]. Because of the complexity of the procedure,
such improvements are not yet available for arbitrary processes, in particular
the multi-jet final states we are interested in. Unless one decides for a
fixed-scale choice at NLO, both approaches determine the strong-coupling
scales dynamically, i.e. on an event-by-event basis, taking the kinematic
configuration of the event into account. For all these reasons, it then occurs
that the CKKW shapes of distributions emerge in many cases quite similarly to
those evaluated at NLO, making an application of global $K$-factors
feasible.202020For example, in [25] a global $K$-factor of magnitude $1.33$
with respect to the total inclusive cross section as measured by CDF [101] was
applied to achieve a good agreement between the data and the SHERPA
predictions for inclusive jet multiplicity and transverse momentum
distributions. The treatment to fix the strong couplings is different when
multiple scales are present. While at NLO the scales are set uniformly such
that all $\alpha_{\mathrm{s}}$ factors obtain the same value, in the CKKW
method they are set locally by the procedure itself, cf. [68, 102, 70] for
example. This is known as $\alpha_{\mathrm{s}}$ reweighting, constituting the
second component of the matrix-element reweighting of the CKKW method. The
assignment of the scales proceeds hierarchically based on the splitting
history, which is identified by the $k_{T}$-jet cluster algorithm when applied
to the initial matrix-element configuration considering physical parton
combinations only. The nodal $k_{T}$ values found by the clustering can be
interpreted as the relative transverse momenta of the identified splittings.
They are then used as the scales for the strong-coupling constants replacing
the predefined choice of the initial matrix-element generation. It would be
interesting to see if a hierarchical scale setting can further stabilize NLO
results, but no such $\alpha_{\mathrm{s}}$ reweighting has been completely
worked out yet for NLO calculations.
## Appendix B Appendix: Analysis side studies and additional material
cuts & | $2\,\Delta/$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$
---|---|---|---|---|---|---|---|---|---|---
selections | $\mathrm{GeV}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $190\quad[20]$ | $200\quad[20]$ | $210\quad[20]$
$\sigma^{(0)}$ | | $9.862$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 96$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$ | $7.827$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 75$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$ | $6.473$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 61$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$
| | 1.0 | 1.0 | 0.10 | 1.0 | 1.0 | 0.08 | 1.0 | 1.0 | 0.06
lepton & | | $5.561$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 12$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $4.433$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 95$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$ | $3.689$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 78$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$
MET cuts | | 0.564 | 0.45 | 0.08 | 0.566 | 0.45 | 0.07 | 0.570 | 0.45 | 0.05
as above & | | $4.586$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $3.709$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 41$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $3.128$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 34$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
$\geq 2$ jets | | 0.465 | 0.0087 | 0.50 | 0.474 | 0.0087 | 0.40 | 0.483 | 0.0087 | 0.33
as above & | | $2.533$ | $6997$ | 77$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $2.007$ | $6997$ | 60$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.671$ | $6997$ | 50$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.257 | 0.0032 | 0.46 | 0.256 | 0.0032 | 0.36 | 0.258 | 0.0032 | 0.29
naive $h$-reco | $50$ | $2.699$ | $6644$ | 87$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $2.146$ | $6303$ | 72$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.789$ | $5837$ | 64$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.274 | 0.0030 | 0.50 | 0.274 | 0.0029 | 0.40 | 0.276 | 0.0027 | 0.34
naive $h$-reco | $30$ | $2.082$ | $4062$ | 0.0011 | $1.649$ | $3805$ | 91$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.374$ | $3521$ | 81$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.211 | 0.0018 | 0.49 | 0.211 | 0.0017 | 0.40 | 0.212 | 0.0016 | 0.34
naive $h$-reco | $48$ | $2.177$ | $3483$ | 0.0013 | $1.699$ | $3151$ | 0.0011 | $1.397$ | $2649$ | 0.0011
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.221 | 0.0016 | 0.56 | 0.217 | 0.0014 | 0.45 | 0.216 | 0.0012 | 0.40
naive $h$-reco | $20$ | $1.488$ | $1565$ | 0.0020 | $1.159$ | $1312$ | 0.0019 | $0.9518$ | $1080$ | 0.0018
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.151 | 71$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.57 | 0.148 | 60$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.48 | 0.147 | 49$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.43
comb. $h$-reco | $50$ | $3.666$ | $7296$ | 0.0011 | $2.937$ | $6946$ | 89$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $2.464$ | $6465$ | 79$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.372 | 0.0033 | 0.65 | 0.375 | 0.0032 | 0.52 | 0.381 | 0.0029 | 0.45
comb. $h$-reco | $20$ | $2.545$ | $3145$ | 0.0017 | $2.031$ | $2964$ | 0.0014 | $1.699$ | $2756$ | 0.0013
| | 0.258 | 0.0014 | 0.69 | 0.260 | 0.0013 | 0.56 | 0.262 | 0.0013 | 0.48
comb. $h$-reco | $50$ | $3.243$ | $4088$ | 0.0017 | $2.573$ | $3755$ | 0.0014 | $2.138$ | $3211$ | 0.0014
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.329 | 0.0019 | 0.77 | 0.329 | 0.0017 | 0.62 | 0.330 | 0.0015 | 0.55
comb. $h$-reco | $30$ | $2.806$ | $2662$ | 0.0023 | $2.205$ | $2314$ | 0.0020 | $1.823$ | $1925$ | 0.0020
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.284 | 0.0012 | 0.82 | 0.282 | 0.0011 | 0.68 | 0.282 | 88$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.61
comb. $h$-reco | $20$ | $2.333$ | $1830$ | 0.0027 | $1.829$ | $1555$ | 0.0025 | $1.507$ | $1293$ | 0.0024
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.237 | 83$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.82 | 0.234 | 71$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.69 | 0.233 | 59$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.62
comb. $h$-reco | $16$ | $2.066$ | $1484$ | 0.0030 | $1.617$ | $1255$ | 0.0027 | $1.331$ | $1029$ | 0.0027
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.209 | 67$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.81 | 0.207 | 57$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.68 | 0.206 | 47$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.61
comb. $h$-reco | $10$ | $1.436$ | $921.6$ | 0.0033 | $1.121$ | $773.9$ | 0.0030 | $0.9215$ | $632.3$ | 0.0030
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.146 | 42$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.71 | 0.143 | 35$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.60 | 0.142 | 29$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.54
Table 7: Impact of the different levels of cuts on the $e\nu_{e}$+jets final
states for the $gg\to h\to WW$ production and decay signal and the $W$+jets
background as obtained from SHERPA. Cross sections $\sigma_{S}$, $\sigma_{B}$,
acceptances $\varepsilon_{S}$, $\varepsilon_{B}$ and $S/B$, $S/\sqrt{B}$
ratios are shown for Higgs boson masses of $M_{h}=190$, $200$ and $210\
\mathrm{GeV}$. Note that $\tilde{m}_{ij}=m_{ij}/\mathrm{GeV}$ and
$\tilde{\delta}=\delta/\mathrm{GeV}$. Significances were calculated using Eqs.
(4.2) assuming ${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ of integrated luminosity,
counting both electrons and muons and combining Tevatron experiments.
cuts & | $2\,\Delta/$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$
---|---|---|---|---|---|---|---|---|---|---
selections | $\mathrm{GeV}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $110\quad[20]$ | $120\quad[20]$ | $130\quad[20]$
lepton & MET | | $0.2254$ | $19080$ | 27$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $0.8017$ | $19080$ | 97$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $2.260$ | $19080$ | 27$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
cuts & $\geq 2$ jets | | 0.0673 | 0.0087 | 0.027 | 0.104 | 0.0087 | 0.095 | 0.173 | 0.0087 | 0.27
naive $h$-reco | $50$ | $0.02723$ | $321.2$ | 20$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $0.2132$ | $969.2$ | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $0.9580$ | $1963$ | 0.0011
| | 0.00813 | 15$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.025 | 0.0276 | 44$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.11 | 0.0733 | 89$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.35
naive $h$-reco | $48$ | $0.00873$ | $34.23$ | 59$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $0.1160$ | $90.62$ | 0.0029 | $0.5804$ | $363.7$ | 0.0037
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.00261 | 16$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 0.024 | 0.0150 | 41$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 0.20 | 0.0444 | 17$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.49
comb. $h$-reco | $50$ | $0.04236$ | $392.8$ | 25$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $0.3000$ | $1131$ | 61$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.253$ | $2239$ | 0.0013
| | 0.0126 | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.035 | 0.0389 | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.15 | 0.0959 | 0.0010 | 0.43
comb. $h$-reco | $50$ | $0.01252$ | $41.87$ | 69$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $0.1515$ | $124.0$ | 0.0028 | $0.7673$ | $461.3$ | 0.0038
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.00374 | 19$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 0.032 | 0.0196 | 56$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 0.22 | 0.0587 | 21$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.58
comb. $h$-reco | $20$ | $0.00607$ | $8.805$ | 0.0016 | $0.1017$ | $23.97$ | 0.0098 | $0.4995$ | $53.81$ | 0.021
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.00181 | 40$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$ | 0.033 | 0.0132 | 11$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 0.34 | 0.0382 | 24$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | 1.11
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | | $140\quad[20]$ | $150\quad[20]$
lepton & MET | | | | | $4.316$ | $19080$ | 52$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $6.343$ | $19080$ | 77$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
cuts & $\geq 2$ jets | | | | | 0.250 | 0.0087 | 0.51 | 0.328 | 0.0087 | 0.75
naive $h$-reco | $50$ | | | | $2.213$ | $3230$ | 0.0016 | $3.642$ | $4593$ | 0.0018
| | | | | 0.128 | 0.0015 | 0.63 | 0.188 | 0.0021 | 0.88
naive $h$-reco | $48$ | | | | $1.374$ | $924.4$ | 0.0034 | $2.509$ | $1662$ | 0.0035
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | | | | 0.0795 | 42$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.73 | 0.130 | 76$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.00
comb. $h$-reco | $50$ | | | | $2.902$ | $3637$ | 0.0018 | $4.778$ | $5103$ | 0.0022
| | | | | 0.168 | 0.0017 | 0.78 | 0.247 | 0.0023 | 1.09
comb. $h$-reco | $50$ | | | | $1.877$ | $1121$ | 0.0038 | $3.487$ | $1968$ | 0.0041
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | | | | 0.109 | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.91 | 0.180 | 89$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.28
comb. $h$-reco | $20$ | | | | $1.213$ | $234.1$ | 0.012 | $2.462$ | $704.8$ | 0.0080
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | | | | 0.0701 | 11$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.29 | 0.127 | 32$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.51
Table 8: Impact of the different levels of cuts on the $e\nu_{e}$+jets final
states for the $gg\to h\to WW$ production and decay signal and the $W$+jets
background as obtained from SHERPA. Cross sections $\sigma_{S}$, $\sigma_{B}$,
acceptances $\varepsilon_{S}$, $\varepsilon_{B}$ and $S/B$, $S/\sqrt{B}$
ratios are shown for Higgs boson masses below the on-shell diboson mass
threshold from $M_{h}=110\ \mathrm{GeV}$ to $M_{h}=150\ \mathrm{GeV}$. Note
that $\tilde{m}_{ij}=m_{ij}/\mathrm{GeV}$ and
$\tilde{\delta}=\delta/\mathrm{GeV}$. The significances were calculated
according to Eqs. (4.2) assuming ${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ of
integrated luminosity, counting both electrons and muons and combining
Tevatron experiments. The layout of the table is the same as in Tables 4 and
7, however a smaller number of Higgs boson candidate selections is shown.
### B.1 Ideal Higgs boson reconstruction analyses
We first complete the presentation of our Section 4.2 main results by showing
Tables 7 and 8 where we display the numbers associated with the high and low
Higgs boson mass region, respectively. As in Table 4 we list the signal and
$W$+jet background cross sections, selection efficiencies, $S/B$ ratios and
significances at different analysis levels. All cuts, the selection
procedures, the layout of the tables and the interpretation of the results
given in the tables have been discussed in detail in Section 4.2. Note that
the rightmost column of Table 10 carries the outcomes for the test mass point
$M_{h}=220\ \mathrm{GeV}$.
One remark shall be added regarding the magnitude of off-shell effects. The
loss one faces due to off-shell Higgs bosons can be read off Table 8 by
comparing the acceptances after baseline (1st rows) and combinatorial Higgs
boson selection (4th rows). While above the $WW$ mass threshold the loss on
the signal (background) is mild (significant) ranging from $1.2$–$1.3$
($2.6$–$3.2$), it steadily increases for decreasing $M_{h}$, approaching $1.8$
and $5.3$ at $M_{h}=130\ \mathrm{GeV}$ and $M_{h}=110\ \mathrm{GeV}$,
respectively. The background loss factors turn huge (up to 48) because of the
steeply falling $m_{e\nu_{e}jj^{\prime}}$ spectrum (cf. Figure 3), but this
cannot overcome the smallness of $S/\sqrt{B}$ due to the signal reduction.
We now present the results of our side studies, which we decided to put in an
appendix in order to not distract the flow of the main body.
Figure 11: $S/\sqrt{B}$ significances as a function of injected Higgs boson
masses varying from $M^{\mathrm{inj}}_{h}=165\ \mathrm{GeV}$ to $200\
\mathrm{GeV}$ for different mass window parameters $\Delta$ and $\delta$.
Results of two combinatorial analyses based on the invm reconstruction are
shown: one using the default setting $M_{h}=M^{\mathrm{inj}}_{h}$ (dashed
lines), the other where the hypothesized Higgs boson mass is fixed at
$M_{h}=180\ \mathrm{GeV}$ (solid lines). For the different injected Higgs
boson masses, the $e\nu_{e}$+jets final states are generated from the $gg\to
h\to WW$ signal and the $W$+jets production background. All significances
were calculated according to Eqs. (4.2) taking only the dominant background
into account and under the assumption of an integrated luminosity of
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$, including electron and muon channels,
i.e. $f_{\ell}=2$.
cuts & | $2\,\Delta/$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$
---|---|---|---|---|---|---|---|---|---|---
selections | $\mathrm{GeV}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $165\quad[20]$ | $170\quad[20]$ | $180\quad[20]$
naive $h$-reco | $50$ | $4.495$ | $5317$ | 0.0020 | $4.121$ | $5317$ | 0.0018 | $3.232$ | $5527$ | 0.0013
| | 0.297 | 0.0022 | 1.14 | 0.296 | 0.0022 | 1.02 | 0.290 | 0.0023 | 0.77
naive $h$-reco | $48$ | $3.965$ | $2371$ | 0.0040 | $3.585$ | $2571$ | 0.0032 | $2.759$ | $2849$ | 0.0022
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.262 | 0.0010 | 1.50 | 0.258 | 0.0011 | 1.27 | 0.248 | 0.0012 | 0.92
comb. $h$-reco | $50$ | $5.603$ | $5511$ | 0.0024 | $5.149$ | $5758$ | 0.0021 | $4.105$ | $5982$ | 0.0016
| | 0.370 | 0.0023 | 1.39 | 0.370 | 0.0024 | 1.22 | 0.368 | 0.0025 | 0.94
comb. $h$-reco | $50$ | $5.081$ | $2632$ | 0.0046 | $4.690$ | $2923$ | 0.0037 | $3.722$ | $3251$ | 0.0026
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.336 | 0.0011 | 1.83 | 0.337 | 0.0012 | 1.56 | 0.334 | 0.0014 | 1.16
comb. $h$-reco | $20$ | $3.890$ | $1206$ | 0.0076 | $3.592$ | $1357$ | 0.0061 | $2.846$ | $1518$ | 0.0043
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.257 | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 2.07 | 0.258 | 57$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.76 | 0.255 | 64$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.30
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $190\quad[20]$ | $200\quad[20]$ | $210\quad[20]$
naive $h$-reco | $50$ | $2.231$ | $5435$ | 92$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.779$ | $5158$ | 76$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.480$ | $4787$ | 67$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.287 | 0.0023 | 0.53 | 0.287 | 0.0022 | 0.42 | 0.288 | 0.0020 | 0.36
naive $h$-reco | $48$ | $1.867$ | $2821$ | 0.0015 | $1.465$ | $2560$ | 0.0013 | $1.203$ | $2161$ | 0.0012
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.240 | 0.0012 | 0.61 | 0.237 | 0.0011 | 0.49 | 0.234 | 91$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.44
comb. $h$-reco | $50$ | $2.865$ | $5892$ | 0.0011 | $2.302$ | $5607$ | 90$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $1.925$ | $5221$ | 80$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.368 | 0.0025 | 0.65 | 0.372 | 0.0024 | 0.53 | 0.375 | 0.0022 | 0.45
comb. $h$-reco | $50$ | $2.575$ | $3243$ | 0.0018 | $2.052$ | $2988$ | 0.0015 | $1.702$ | $2557$ | 0.0014
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.331 | 0.0014 | 0.79 | 0.331 | 0.0013 | 0.64 | 0.332 | 0.0011 | 0.57
comb. $h$-reco | $20$ | $1.949$ | $1445$ | 0.0030 | $1.536$ | $1242$ | 0.0027 | $1.265$ | $1029$ | 0.0027
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.251 | 61$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.89 | 0.248 | 53$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.74 | 0.247 | 44$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.66
Table 9: Impact of the different levels of cuts on the $e\nu_{e}$+jets final
states for the $gg\to h\to WW$ production and decay signal and the $W$+jets
background as obtained from SHERPA when using the CTEQ6.6 PDF libraries. Cross
sections $\sigma_{S}$, $\sigma_{B}$, acceptances $\varepsilon_{S}$,
$\varepsilon_{B}$ and $S/B$, $S/\sqrt{B}$ ratios are shown for Higgs boson
masses from $M_{h}=165\ \mathrm{GeV}$ to $M_{h}=210\ \mathrm{GeV}$. Note that
$\tilde{m}_{ij}=m_{ij}/\mathrm{GeV}$ and $\tilde{\delta}=\delta/\mathrm{GeV}$.
All significances were calculated according to Eqs. (4.2) assuming
${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ of integrated luminosity, counting both
electron and muon channels and combining Tevatron experiments.
The Higgs boson masses used in the analyses are, of course, hypothetical.
However, we only considered the obvious scenario where the test mass $M_{h}$
has been set equal to the Higgs boson mass $M^{\mathrm{inj}}_{h}$ injected
while generating the signal predictions. It is clear that one scans over a
range of masses when pursuing an analysis, and here we do it in steps of 10
GeV; still the actual Higgs boson mass could deviate as much as 5 GeV from the
assumed mass. Hence, we want to briefly study how strongly the invm Higgs
boson candidate selections and their related significances depend on the match
between the test and injected Higgs boson mass. To this end we generated other
than default signal predictions for Higgs boson masses of
$M^{\mathrm{inj}}_{h}=165,175,185,195\ \mathrm{GeV}$ and input them into the
analyses using $M_{h}=180\ \mathrm{GeV}$. Figure 11 shows the outcome of this
side study where, for both types of analyses, we collected results for several
mass window settings. We learn two things from plotting the significance as a
function of the Higgs boson generation mass. First, the significances that we
attain if we keep the selection parameters ($M_{h}$, $\Delta$, $\delta$)
constant are fairly robust over a broader range of generation masses. Yet the
maximum $S/\sqrt{B}$ occur for $M_{h}=M^{\mathrm{inj}}_{h}$. Secondly we
learn, asymmetric window placements such that $M_{h}<M^{\mathrm{inj}}_{h}$ are
beneficial to achieve significance gains. By focusing on a single generation
mass, e.g. $M^{\mathrm{inj}}_{h}=190\ \mathrm{GeV}$, we see that the response
in significance can easily get as large as 40%. The strong sensitivity can be
understood by comparing the signal and background shapes as visualized in
Figure 3 or the upper left plot of Figure 1.
cuts & | $2\,\Delta/$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$ | $\sigma_{S}/\mathrm{fb}$ | $\sigma_{B}/\mathrm{fb}$ | $S/B$
---|---|---|---|---|---|---|---|---|---|---
selections | $\mathrm{GeV}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$ | $\varepsilon_{S}$ | $\varepsilon_{B}$ | $S/\sqrt{B}$
$M_{h}/\mathrm{GeV}\quad[\delta/\mathrm{GeV}]$ | $180\quad[15]$ | $180\quad[20]\,[p^{\mathrm{jet}}_{T}/\mathrm{GeV}\\!>\\!30]$ | $220\quad[20]$
$\sigma^{(0)}$ | | $14.19$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 14$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $14.19$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 14$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $5.420$ | $220\mbox{$\cdot\;\\!\\!$}10^{4}$ | 51$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$
| | 1.0 | 1.0 | 0.15 | 1.0 | 1.0 | 0.15 | 1.0 | 1.0 | 0.05
lepton & | | $7.946$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $8.127$ | $987\mbox{$\cdot\;\\!\\!$}10^{3}$ | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}6}$ | $3.101$ | $984\mbox{$\cdot\;\\!\\!$}10^{3}$ | 65$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}7}$
MET cuts | | 0.560 | 0.45 | 0.12 | 0.573 | 0.45 | 0.13 | 0.572 | 0.45 | 0.05
as above & | | $6.471$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 74$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | $3.907$ | $6715$ | 0.0013 | $2.661$ | $191\mbox{$\cdot\;\\!\\!$}10^{2}$ | 29$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
$\geq 2$ jets | | 0.456 | 0.0087 | 0.72 | 0.275 | 0.0031 | 0.73 | 0.491 | 0.0087 | 0.28
as above & | | $3.169$ | $5272$ | 0.0013 | $1.932$ | $2194$ | 0.0019 | $1.413$ | $6997$ | 42$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.223 | 0.0024 | 0.67 | 0.136 | 0.0010 | 0.63 | 0.261 | 0.0032 | 0.25
naive $h$-reco | $50$ | $3.911$ | $6749$ | 0.0013 | $1.948$ | $1429$ | 0.0030 | $1.510$ | $5342$ | 58$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.276 | 0.0031 | 0.73 | 0.137 | 65$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.79 | 0.279 | 0.0024 | 0.30
naive $h$-reco | $30$ | $3.039$ | $4199$ | 0.0016 | $1.551$ | $907.0$ | 0.0037 | $1.161$ | $3196$ | 75$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.214 | 0.0019 | 0.72 | 0.109 | 41$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.79 | 0.214 | 0.0015 | 0.30
naive $h$-reco | $48$ | $2.857$ | $2778$ | 0.0022 | $1.631$ | $925.5$ | 0.0038 | $1.169$ | $2167$ | 0.0011
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.201 | 0.0013 | 0.83 | 0.115 | 42$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.82 | 0.216 | 99$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.37
naive $h$-reco | $20$ | $2.118$ | $1351$ | 0.0034 | $1.170$ | $454.3$ | 0.0056 | $0.7991$ | $874.9$ | 0.0019
$\lvert\tilde{m}_{j_{1}j_{2}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.149 | 61$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.89 | 0.0825 | 21$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.84 | 0.147 | 40$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.39
comb. $h$-reco | $50$ | $5.241$ | $7396$ | 0.0015 | $2.484$ | $1534$ | 0.0035 | $2.088$ | $5942$ | 73$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$
| | 0.369 | 0.0034 | 0.94 | 0.175 | 70$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.98 | 0.385 | 0.0027 | 0.40
comb. $h$-reco | $20$ | $3.657$ | $3255$ | 0.0024 | $1.718$ | $671.7$ | 0.0056 | $1.440$ | $2508$ | 0.0012
| | 0.258 | 0.0015 | 0.99 | 0.121 | 31$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.02 | 0.266 | 0.0011 | 0.42
comb. $h$-reco | $50$ | $4.248$ | $3241$ | 0.0029 | $2.185$ | $1024$ | 0.0046 | $1.798$ | $2662$ | 0.0014
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.299 | 0.0015 | 1.15 | 0.154 | 47$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.05 | 0.332 | 0.0012 | 0.51
comb. $h$-reco | $30$ | $3.845$ | $2220$ | 0.0038 | $1.916$ | $699.8$ | 0.0060 | $1.528$ | $1581$ | 0.0020
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.271 | 0.0010 | 1.26 | 0.135 | 32$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.11 | 0.282 | 72$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.56
comb. $h$-reco | $20$ | $3.260$ | $1564$ | 0.0045 | $1.617$ | $492.9$ | 0.0071 | $1.264$ | $1058$ | 0.0025
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.230 | 71$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.27 | 0.114 | 22$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.12 | 0.233 | 48$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.57
comb. $h$-reco | $16$ | $2.895$ | $1273$ | 0.0049 | $1.444$ | $399.0$ | 0.0079 | $1.115$ | $846.6$ | 0.0027
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.204 | 58$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.25 | 0.102 | 18$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.11 | 0.206 | 38$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.56
comb. $h$-reco | $10$ | $2.010$ | $805.0$ | 0.0054 | $1.038$ | $258.5$ | 0.0087 | $0.7722$ | $525.2$ | 0.0030
$\lvert\tilde{m}_{jj^{\prime}}\\!-80\rvert\\!<\\!\tilde{\delta}$ | | 0.142 | 37$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 1.09 | 0.0731 | 12$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.99 | 0.142 | 24$\mbox{$\cdot\;\\!\\!$}10^{\textrm{-}5}$ | 0.49
Table 10: Impact of the different levels of cuts on the $e\nu_{e}$+jets final
states for the $gg\to h\to WW$ production and decay signal and the $W$+jets
background as obtained from SHERPA. Cross sections $\sigma_{S}$, $\sigma_{B}$,
acceptances $\varepsilon_{S}$, $\varepsilon_{B}$ and $S/B$, $S/\sqrt{B}$
ratios are given for Higgs boson masses of $M_{h}=180\ \mathrm{GeV}$ and
$M_{h}=220\ \mathrm{GeV}$. For the former mass point, the central column shows
the values for a jet $p_{T}$ threshold increased by $10\ \mathrm{GeV}$, while
the left column has the values for a smaller dijet mass window of $\delta=15\
\mathrm{GeV}$. A further decrease of the dijet mass window to $\delta=10\
\mathrm{GeV}$ yields $S/B=0.0034$ as well as $S/\sqrt{B}=1.14$ and
$S/B=0.0058$ as well as $S/\sqrt{B}=1.36$ for the combinatorial Higgs boson
reconstruction and mass windows of $\tilde{\Delta}=50$ and
$\tilde{\Delta}=20$, respectively. Note that
$\tilde{m}_{ij}=m_{ij}/\mathrm{GeV}$; all other mass variables denoted by a
tilde are understood in the same way. All significances were calculated
according to Eqs. (4.2) assuming ${\mathcal{L}}=10\ \mathrm{fb}^{-1}$ of
integrated luminosity, counting both electron and muon channels and combining
Tevatron experiments.
In the remainder of this appendix, we outline the impact of PDF variations and
parameter variations other than $M_{h}$, $\Delta$ and $\delta$ on our
analyses.
The SHERPA calculations resulting from using the CTEQ6.6 PDF libraries give a
similar pattern with significances that are about 1–10% larger. This can be
read off Table 9 and seen in Figure 2. By normalizing the Monte Carlo
predictions for both PDF choices, MSTW2008 and CTEQ6.6, to the same respective
theory cross sections, it is altogether reassuring to see that the cut and
event selection procedures only induce deviations of the order of 10% or
below. For the combinatorial Higgs boson candidate selection, one finds almost
the same $S/\sqrt{B}$ ratios, provided a wide Higgs boson mass window is used.
The significances of the CTEQ6.6 calculations outperform those obtained with
MSTW2008, once the mass windows for the Higgs bosons ($\Delta$) and dijets
($\delta$) are tightened. We find the differences being more pronounced for
Higgs boson masses just above the $WW$ threshold.
Speaking of shape differences triggered by the use of different PDFs, we note
that the MSTW2008 PDF set accounts for a larger transverse activity, i.e.
rapidity distributions turn out steeper while the $p_{T}$ spectra develop
10–25% harder tails compared to the CTEQ6.6 predictions. Also, for MSTW2008,
mass peaks are more washed out, again resulting in differences of the order of
25%.
Lastly, apart from the rightmost column, Table 10 gives more details for the
choice $M_{h}=M^{\mathrm{inj}}_{h}=180\ \mathrm{GeV}$ including the cases
where either the dijet mass window has been tightened from $\delta=20\
\mathrm{GeV}$ to $\delta=15\ \mathrm{GeV}$ or the threshold of the jet
transverse momenta has been enhanced to $p^{\mathrm{jet}}_{T}>30\
\mathrm{GeV}$.
### B.2 More realistic Higgs boson reconstruction analyses
We collect here in this appendix additional material to substantiate our
findings presented in Section 4.4. Figures 12–18 display a number of
distributions resulting from the baseline plus combinatorial analyses. In each
plot we show four curves, two predictions each for the signal and the
$W$+jets background as obtained after the ideal (invm) and one of the more
realistic combinatorial selections. While we vary the Higgs boson masses and
the choice of the realistic $h$ reconstruction method, we keep the window
parameters of the combinatorial selections fixed at
$\tilde{\Delta}=\Delta/\mathrm{GeV}=25$ and
$\tilde{\delta}=\delta/\mathrm{GeV}=20$ over the whole set of spectra shown in
these figures.
In Figure 19 we look at distribution taken at the intermediate optimization
level, including the effects of the major cuts but before the application of
the respective minor cuts as listed in Table 6. In these plots we compare the
Monte Carlo predictions for the Higgs boson signals of different $M_{h}$ with
all, dominant and subdominant background predictions. All these predictions
follow from using realistic selection procedures where Higgs boson and dijet
mass windows of $\tilde{\Delta}=\tilde{\delta}=20$ have been employed.
Each one-dimensional distribution is supplemented by an one-minus-ratio
subplot. The first prediction in the respective legend is always taken as the
reference, which we have arranged to be a $W$+jets prediction for all plots
shown here. The ratio subplots nicely visualize why we place the cuts as given
in Table 6. Note that in all plots we only compare the shapes, i.e. all
distributions are normalized to unit area.
Figure 12: Pseudo-rapidity difference between the two selected jets.
Predictions for $gg\to h\to e\nu_{e}$+jets production obtained after invm
(dashed) and more realistic (solid) combinatorial selections (top: pzmh,
center: pzmw, bottom: mt) are compared with each other and to the
corresponding predictions for the $W(\to e\nu_{e})$+jets background. Left
panes show results for $M_{h}=150\ \mathrm{GeV}$; in all other cases
$M_{h}=190\ \mathrm{GeV}$ was used.
Figure 13: Pseudo-rapidity of the selected $\\{e,\nu_{e},j,j^{\prime}\\}$
set, the Higgs boson candidate. The predictions for $gg\to h\to e\nu_{e}$+jets
production (coloured lines) obtained after the combinatorial invm (dashed) and
more realistic $h$ reconstruction (solid) selections (upper: pzmh, lower:
pzmw) are compared with each other and to the corresponding predictions of the
$W(\to e\nu_{e})$+jets background (black lines). Left panes show results for
$M_{h}=150\ \mathrm{GeV}$, while in the right panes the outcomes for
$M_{h}=190\ \mathrm{GeV}$ are shown.
We start by showing the $|\Delta\eta_{j,j^{\prime}}|$ distributions in Figure
12. The differences in the results of the ideal and more realistic selections
are immaterial; there are essentially no differences in the above-threshold
cases. Furthermore, the shapes are very stable under $M_{h}$ variations. We
see that placing a cut around $|\Delta\eta_{j,j^{\prime}}|=1.5$ keeps most of
the signal, while it removes a large fraction of the $W$+jets events. We note
it is only the $W$+jets background featuring a peak location away from zero,
all other backgrounds (not shown here) behave similarly to the signal.
When working with the reconstruction methods, pzmh/w, we have access to
another longitudinal variable: we can include a cut on the $h$ candidate’s
pseudo-rapidity to supplement the constraints from the major and
$|\Delta\eta_{j,j^{\prime}}|$ cuts. Figure 13 displays various
$\eta_{e\nu_{e}jj^{\prime}}$ distributions. Again, all predicted shapes are
rather independent of the choice of the test mass $M_{h}$. The $W$+jets
background (as well as the electroweak background which is not shown here)
tends to preferably populate the forward rapidities while the signal (and the
$t\bar{t}$ contribution also not shown here) shows up more central. This leads
us to require $|\eta_{e\nu_{e}jj^{\prime}}|\lesssim 3.0$ as pointed out in
Section 4.4 to achieve additional significance gains. Deviations between the
ideal and more realistic reconstructions become visible; they now are
$\mathcal{O}(\mbox{25\%})$, this is clearly because one needs information
about the neutrino to form this observable. We observe that the $W$+jets
background receives the larger corrections compared to the signal.
Figure 14: Two-dimensional distributions showing the selected-jet pseudo-
rapidity difference, $|\Delta\eta_{j,j^{\prime}}|$, plotted versus the
reconstructed mass $m_{e\nu_{e}jj^{\prime}}$ of the selected
$\\{e,\nu_{e},j,j^{\prime}\\}$ combinations. The predictions for $gg\to h\to
e\nu_{e}$+jets production obtained after pzmw reconstruction of Higgs boson
candidates are compared with each other and to the corresponding predictions
given by the $W(\to e\nu_{e})$+jets background. The upper (lower) plots
represent the signal (background) predictions, while the left (right) panes
show results for $M_{h}=140\ (180)\ \mathrm{GeV}$.
Figure 15: Two-dimensional distributions showing the selected-jet pseudo-
rapidity difference, $|\Delta\eta_{j,j^{\prime}}|$, plotted versus the
reconstructed mass $m_{e\nu_{e}jj^{\prime}}$ of the selected
$\\{e,\nu_{e},j,j^{\prime}\\}$ combinations. The predictions for $gg\to h\to
e\nu_{e}$+jets production obtained after combinatorial selection according to
the mt procedure are compared with each other and to the corresponding
predictions given by the $W(\to e\nu_{e})$+jets background. The upper (lower)
plots represent the signal (background) predictions, while the left (right)
panes show results for $M_{h}=140\ (180)\ \mathrm{GeV}$. Note that for the
purpose of illustration, the $m_{e\nu_{e}jj^{\prime}}$ quantities are
reconstructed as in the ideal case.
We add one more comment regarding longitudinal quantities. In this study, as
stated in Section 4.4, the pseudo-rapidity variables, which we discussed
above, occur largely uncorrelated with the transverse observables as well as
invariant masses. Schematically, we illustrate this on the basis of two-
dimensional $|\Delta\eta_{j,j^{\prime}}|$ versus $m_{e\nu_{e}jj^{\prime}}$
distributions for both the Higgs boson signal and the $W$+jets background. In
Figure 14 we show these distributions as resulting from the combinatorial pzmw
reconstruction for two different Higgs boson masses, below ($M_{h}=140\
\mathrm{GeV}$) and above ($M_{h}=180\ \mathrm{GeV}$) the diboson mass
threshold. Similarly, Figure 15 exhibits the results obtained with the mt
selection where the $m_{e\nu_{e}jj^{\prime}}$ quantities were reconstructed as
in the ideal case. When confronted with the respective $W$+jets backgrounds,
we notice that the predictions originating from the production of Higgs bosons
cover rather different parameter regions in the $m_{e\nu_{e}jj^{\prime}}$ –
$|\Delta\eta_{j,j^{\prime}}|$ plane. This happens independent of the value of
$M_{h}$ and the chosen combinatorial selection. Based on these observations,
we expect the total $S/\sqrt{B}$ increase to almost completely factorize into
a product of single $S/\sqrt{B}$ improvement factors.
Figure 16: Examples of leading and subleading cut observables below the
$2\,M_{W}$ threshold, for $M_{h}=130\ \mathrm{GeV}$. Predictions are shown for
the $gg\to h$ and $W$ production of $e\nu_{e}$+jets final states using the
invm (dashed) and more realistic (solid) combinatorial selections. Top to
bottom, left panes: $m_{T,e^{-}\bar{\nu}_{e}}$ (pzmw leading),
$m^{(\nu_{e})}_{T,e^{+}\nu_{e}jj^{\prime}}$ and $m_{e^{+}\nu_{e}jj^{\prime}}$
(pzmh leading and subleading); right panes: $H_{T,jj^{\prime}}$,
$m_{T,jj^{\prime}}$ and $m_{jj^{\prime}}$ (mt leading to subsubleading).
Figure 17: Examples of (sub)leading cut observables for Higgs boson masses
$M_{h}\sim 2\,M_{W}$. Predictions are shown for the $gg\to h$ and $W$
production of $e\nu_{e}$+jets final states using the invm (dashed) and more
realistic (solid) combinatorial selections. Top to bottom, left panes:
$p_{T,e^{+}\nu_{e}jj^{\prime}}$ and $m_{e^{-}\bar{\nu}_{e}jj^{\prime}}$ (pzmh
leading and subleading), $H_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ (mt
subleading); right panes: $m^{(\nu_{e})}_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$,
$m_{e^{-}\bar{\nu}_{e}jj^{\prime}}$ and $H_{T,jj^{\prime}}$ (pzmw leading to
subsubleading).
In Figure 16 we depict examples of potential discriminators below the diboson
mass threshold. From the upper left to the lower right we present, for the
choice of $M_{h}=130\ \mathrm{GeV}$, the transverse mass
$m_{T,e^{-}\bar{\nu}_{e}}$ of the leptonically decaying $W^{-}$, the scalar
$p_{T}$ sum $H_{T,jj^{\prime}}$ of the selected jets, the transverse masses of
the selected 4-particle set, $m^{(\nu_{e})}_{T,e^{+}\nu_{e}jj^{\prime}}$, and
of the selected jets, $m_{T,jj^{\prime}}$, as well as their corresponding
invariant masses, $m_{e^{+}\nu_{e}jj^{\prime}}$ and $m_{jj^{\prime}}$. The
first three plots in that order are the leading cut variables for pzmw, mt and
pzmh followed by the subleading cut variables for mt and pzmh, cf. Table 6. In
the lower left plot we display the subsubleading cut variable for the mt
method, $m_{jj^{\prime}}$, which would yield a significance gain of 19% if we
were to demand $m_{jj^{\prime}}\geq 72\ \mathrm{GeV}$. Unlike the transverse
observables depicted in the upper and center panes of Figure 16, the
invariant-mass distributions are affected by the modifications of the ideal
selection owing to a realistic neutrino treatment. The selected 4-object
invariant mass (lower left) exemplifies to what degree shapes can get
distorted by the pzmh approach. The shoulder above $M_{h}=130\ \mathrm{GeV}$
emerges because complex solutions cannot be completely avoided in the
reconstruction of the neutrino momenta; the lower tail arises from the
$m_{T,e\nu_{e}jj^{\prime}}$ constraints on the target mass
$m_{*,e\nu_{e}jj^{\prime}}$, see Section 4.3.
The 2-particle transverse masses shown in the upper left and center right of
Figure 16 together with the $m_{jj^{\prime}}$ spectra document why we exploit
the signal’s preference for on-shell hadronic and off-shell leptonic decays of
the $W$ bosons. All backgrounds considered in this study disfavor this
correlation. The 4-particle transverse mass (center left) exhibits – as
expected – a nice kinematic edge for the signal at
$m^{(\nu_{e})}_{T,e^{+}\nu_{e}jj^{\prime}}=M_{h}$, while all backgrounds are
continuous in this variable reaching their broad maxima above the applied mass
window. The features of the $H_{T,jj^{\prime}}$ handle (upper right) have been
already described in Section 2. For this variable, the electroweak background
turns out signal-like whereas the $t\bar{t}$ background generates (by far) the
hardest tails.
Figure 18: Leading and subleading cut observables for $M_{h}$ choices well
above the $2\,M_{W}$ threshold. Predictions are shown for the $gg\to h$ and
$W$ production of $e\nu_{e}$+jets final states using the invm (dashed) and
more realistic (solid) combinatorial selections. Upper panes:
$H_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ (leading, pzmh (left) and mt); center
panes: $p_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ (subleading, pzmh (left) and
pzmw); lower panes (subleading), left: $p_{T,e^{-}\bar{\nu}_{e}}$ (pzmh),
right: $m_{T,e^{+}\nu_{e}jj^{\prime}}$ (mt).
We now discuss some of the near-threshold discriminators where most of the
examples are given for the test point $M_{h}=170\ \mathrm{GeV}$. Figure 17
shows in its top row the two leading cut variables that we found for this mass
region. In the upper right we have the transverse mass
$m^{(\nu_{e})}_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ of the selected
$\\{e,\nu_{e},j,j^{\prime}\\}$ objects as resulting after the pzmw
reconstruction with $M_{h}=150\ \mathrm{GeV}$. We discover properties very
similar to those discussed for the case of low Higgs boson masses, cf. Figure
16. In the upper left we have depicted the first-rank discriminator for medium
Higgs boson masses – the selected 4-object transverse momentum distribution.
It was introduced early on in Section 2. Here, we present the
$p_{T,e^{+}\nu_{e}jj^{\prime}}$ distribution as obtained after the
combinatorial pzmh selection. We remark that the curves of the other two
realistic approaches, pzmw and mt, (both not shown here) deviate even less
from the respective curves of the ideal selection. Also not shown in Figure 17
but worthwhile to mention, the electroweak background would yield spectra
similar to those of the $W$+jets background whereas the top-pair production
would turn up significantly harder than the signal’s $p_{T}$ spectra.
The other example plots of Figure 17 are chosen from the set of second-leading
cut variables. The center panes display the invariant mass distributions
$m_{e^{-}\bar{\nu}_{e}jj^{\prime}}$ resulting from the pzmh/w selections. As
opposed to the – by construction – sculpted shapes of the pzmh method, we
observe that the pzmw selection reproduces the invariant mass shapes of the
invm ideal reconstruction to a large extent. The lower panes demonstrate the
potential possessed by scalar transverse momentum sums that we exemplify by
means of the selected-set $H_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ observable as
given by the mt selection (lower left) and the selected-jet
$H_{T,jj^{\prime}}$ observable resulting from the pzmw reconstruction (lower
right). Remarkably, based on the 4-object $H_{T}$, we can achieve an even
clearer separation between the signal and the $W$+jets background once we
select $h$ candidates according to the mt method.
Figure 18 summarizes the types of discriminating observables as identified in
Table 6 for the region of large Higgs boson masses; we use $M_{h}=210$ and
$220\ \mathrm{GeV}$ in the example plots. The upper panels represent
$H_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ distributions obtained after pzmh (left)
and mt combinatorial selections. Between these two cases we detect only
marginal differences, and similarly between the predictions of the ideal and
pzmw reconstructions (not shown here). At large Higgs boson masses, the signal
develops the peak at considerably larger $H_{T}$ values compared to the
$W$+jets background. This characterizes the selected-set $H_{T}$ as the
strongest handle we have above the $WW$ threshold. Moreover, the pzmh and mt
realistic selections further enhance the separation between the two peak
regions. For the subdominant backgrounds (not shown here), we noticed a strong
similarity between the $H_{T}$ spectra arising from the electroweak production
and the $W$+jets background. The $t\bar{t}$ background however yields the
hardest spectra both in terms of the peak position as well as the tail of the
$H_{T}$ distributions.
The center panes of Figure 18 show two examples of selected $h$ candidate
$p_{T}$ distributions. These variables do not constitute the best
discriminators anymore, but still quite often rank second best in separating
signal from $W$+jets production in the domain of large $M_{h}$. As the pzmw
reconstruction works extremely well for heavy Higgs boson decays into on-shell
$W$ bosons, it gives $p_{T,e^{-}\bar{\nu}_{e}jj^{\prime}}$ shapes almost
identical with the ideal selection.
For the $M_{h}=220\ \mathrm{GeV}$ point, we present two different subleading
cut variables in the lower pane of Figure 18. To the left, one finds the
$p_{T,e^{-}\bar{\nu}_{e}}$ distribution of the reconstructed $W^{-}$ resulting
after selecting Higgs boson candidates according to pzmh. The purely
transverse mass $m_{T,e^{+}\nu_{e}jj^{\prime}}$ of the mt-selected 4-object
set, cf. Eq. (8), is depicted on the lower right. For these variables, we
recognize similar features as for the $H_{T,e\nu_{e}jj^{\prime}}$ spectra
concerning the minor backgrounds and the comparison with the invm Higgs boson
candidate reconstruction.
Figure 19: Examples of minor cut observables after application of major cuts
and rapidity constraints including spectra resulting from subdominant
backgrounds. Predictions using the more realistic combinatorial $h$ candidate
selections are shown for $e\nu_{e}$+jets final states arising from $gg\to h$
(solid red), $W$+jets (solid black), $t\bar{t}$ (dashed blue) and electroweak
(dashed green) production at Tevatron Run II. See text for more details.
In Figure 19 we give a brief overview of discriminators that lead to a further
increase in significance after implementing all major and rapidity constraints
discussed in Section 4.4. The boost factor $\gamma_{jj^{\prime}|e\nu_{e}}$, as
introduced in Section 2, proves very helpful in separating signal from
backgrounds over a large range of above-threshold Higgs boson masses. Since
$\gamma_{jj^{\prime}|e\nu_{e}}$ develops a peak for the signal only, it is
advantageous to isolate the boost factor peak region in order to exploit that
a not too broad scalar resonance has been produced over a multitude of
continuous backgrounds. This is exemplified in the upper plots of Figure 19
where we present two boost factor $\gamma_{jj^{\prime}|e^{-}\bar{\nu}_{e}}$
example distributions for $M_{h}=165\ \mathrm{GeV}$ (left) and $M_{h}=220\
\mathrm{GeV}$ when selecting via the pzmw method.
For the pzmh method in particular, we identified two transverse variables
that, if constrained from below, are very yielding in the low Higgs boson mass
region, even at this more involved level of the analysis. We exhibit these
variables, $H_{T,jj^{\prime}}$ on the left and $m_{\perp,jj^{\prime}}$ on the
right, in the center panes of Figure 19 choosing $M_{h}=130\ \mathrm{GeV}$. As
in similar cases the significance gain originates from exploiting the
different peak locations associated with the signal, at higher values, and the
dominant background, preferring the low values.
For $M_{h}$ very close above threshold, we can use the azimuthal angle between
the lepton and MET or the two selected jets (as suggested by Han and Zhang in
Ref. [7, 8]) to further suppress the backgrounds. A low-$p_{T}$ Higgs boson
produced at $WW$ threshold gives rise to two longitudinally moving $W$ bosons,
which in turn each decay into two objects oriented almost back-to-back in the
transverse plane. This is demonstrated by the example plot in the lower left
of Figure 19 where signal and background distributions are shown for the
$\Delta\phi_{e^{-},\bar{\nu}_{e}}$ observable when employing the mt selection
for $M_{h}=165\ \mathrm{GeV}$.
Finally, we display in the lower right of Figure 19 an example of a global-
event observable, namely $H_{T}$ as calculated from the entire event, not
vetoed by the selection and major cuts. It illustrates the hierarchy of scales
intrinsic to the heavy Higgs boson signal ($M_{h}=220\ \mathrm{GeV}$) and
different background processes; it also visualizes the leftover potential when
considering the global $H_{T}$ for the implementation of additional cuts, see
Appendix B.3.
### B.3 Directions for additional improvements
The final sets of events surviving our optimized combinatorial selections of
Higgs boson candidates are perfect for use as input to a full multivariate
analysis. This is primarily because the analysis types presented here were
geared towards significance maximization, so that a sufficiently large number
of events can be preserved.
* •
$S/B$ improvements: since we are rather safe from a statistics point of view,
there is in many cases potential to improve $S/B$ by simply requiring more
restrictive constraints accepting (mild) significance losses at the same time.
The simplest way is to cut harder in the tails of the major observables; for
instance for pzmw at $M_{h}=220\ \mathrm{GeV}$, using
$H_{T,e\nu_{e}jj^{\prime}}\geq 188\ \mathrm{GeV}$ (instead of the bound given
in Table 6) results in a 40% gain in $S/B$ while the significance only drops
by 10%. Of course, a change in variable sometimes is more beneficial for
maximizing $S/B$ and keeping a reasonable significance. Consider for example
mt at $M_{h}=140\ \mathrm{GeV}$; hardening the $H_{T,jj^{\prime}}$ constraint
by cutting out the region below $92\ \mathrm{GeV}$ maximizes $S/B$ by doubling
it, but reduces the significance by a factor of $3.4$. In contrast, using
$m_{T,jj^{\prime}}\geq 72\ \mathrm{GeV}$ gives a factor $1.6$ increase in
$S/B$, yet only a factor $1.2$ decrease in significance. Certainly one can opt
for the (other) extreme and totally maximize $S/B$, e.g. for the pzmw case
just mentioned, we find a huge $S/B$ gain of 600% by demanding
$p_{T,j^{\prime}}\geq 64\ \mathrm{GeV}$ but we actually just traded a
reasonable significance associated with half a percent $S/B$ for a $\sim$ 4%
$S/B$ of very low significance diminished by a factor of $6.5$. This behaviour
is typical owing to the limited amount of data taken at the Tevatron.
* •
Overall $H_{T}$ cut: the philosophy of this study is to only constrain
variables involving the candidate set of particles. Allowing cut observables
sensitive to the whole event structure can lead to additional significance
improvements but the related uncertainties are larger since such observables
are more prone to hard radiative corrections that need to be described
appropriately, often beyond parton-shower modeling. At the level of
identifying minor cuts, as given in Table 6, we found that an overall $H_{T}$
cut on the selected events yields in most cases significance gains of the
order of 20–40% near and above threshold. This is a conservative estimate
considering the larger uncertainties on such cuts. An example is shown in the
lower right of Figure 19 where we see why one benefits by constraining $H_{T}$
from below. In addition one could suppress the $t\bar{t}$ background by
introducing an upper $H_{T}$ bound, or equally, exploit the fact that the
leading jets in $t\bar{t}$ production yield a substantially harder
$H_{T,2}=H_{T,j_{1}j_{2}}$ spectrum.
* •
Asymmetric mass windows for reconstruction methods: as touched during the
discussion of Figure 11, the use of asymmetric test mass windows, i.e.
$(M_{h}-\Delta_{\mathrm{low}},M_{h}+\Delta_{\mathrm{up}})$, can improve the
realistic selections that either approximate or set the selected 4-object
mass, $m_{e\nu_{e}jj^{\prime}}$. For all $M_{h}$ values, it is advantageous to
choose $\Delta_{\mathrm{low}}$ larger by a few GeV than $\Delta_{\mathrm{up}}$
where $2\,\Delta=\Delta_{\mathrm{low}}+\Delta_{\mathrm{up}}$. The results
presented in Figure 1 (upper left) and Figure 3 clarify why this is a good
idea: first, the Higgs boson resonance is deformed by radiative losses and
resolution effects (there is no perfect jet finding let alone (sub)event
reconstruction) amounting to a larger portion of cross section below the peak;
second, the backgrounds are either sufficiently flat or – below the $WW$
threshold – steeply falling towards smaller masses amplifying the asymmetry
effect further. We checked the asymmetric window scenario for the pzmw method,
where this led to significance gains up to 10%. For the dijet mass windows,
the effects turned out too weak.
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|
arxiv-papers
| 2011-11-11T23:00:41 |
2024-09-04T02:49:24.221206
|
{
"license": "Public Domain",
"authors": "Joseph D. Lykken, Adam O. Martin, Jan-Christopher Winter",
"submitter": "Jan Winter",
"url": "https://arxiv.org/abs/1111.2881"
}
|
1111.2925
|
# Low Mach number limit for the full compressible magnetohydrodynamic
equations with general initial data
Song Jiang LCP, Institute of Applied Physics and Computational Mathematics,
P.O. Box 8009, Beijing 100088, P.R. China jiang@iapcm.ac.cn , Qiangchang Ju
Institute of Applied Physics and Computational Mathematics, P.O. Box 8009-28,
Beijing 100088, P.R. China qiangchang_ju@yahoo.com , Fucai Li∗ Department of
Mathematics, Nanjing University, Nanjing 210093, P.R. China fli@nju.edu.cn
and Zhouping Xin The Institute of Mathematical Sciences, The Chinese
University of Hong Kong, Shatin, NT, Hong Kong zpxin@ims.cuhk.edu.hk
###### Abstract.
The low Mach number limit for the full compressible magnetohydrodynamic
equations with general initial data is rigorously justified in the whole space
$\mathbb{R}^{3}$. The uniform in Mach number estimates of the solutions in a
Sobolev space are obtained on a time interval independent of the Mach number.
The limit is proved by using the established uniform estimates and a theorem
due to Métiver and Schochet [Arch. Ration. Mech. Anal. 158 (2001), 61-90] for
the Euler equations that gives the local energy decay of the acoustic wave
equations.
###### Key words and phrases:
Full compressible magnetohydrodynamic equations, local smooth solution, low
Mach number limit, general initial data
###### 2000 Mathematics Subject Classification:
76W05, 35B40
∗Corresponding author
## 1\. Introduction
In this paper we study the low Mach number limit of local smooth solutions to
the following full compressible magnetohydrodynamic (MHD) equations with
general initial data in the whole space $\mathbb{R}^{3}$ (see [20, 27, 37,
43]):
$\displaystyle\partial_{t}\rho+{\rm div}(\rho{\mathbf{u}})=0,$ (1.1)
$\displaystyle\partial_{t}(\rho{\mathbf{u}})+{\rm
div}\left(\rho{\mathbf{u}}\otimes{\mathbf{u}}\right)+{\nabla
P}=\frac{1}{4\pi}({\rm curl\,}\mathbf{H})\times\mathbf{H}+{\rm
div}\Psi({\mathbf{u}}),$ (1.2) $\displaystyle\partial_{t}\mathbf{H}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})=-{\rm curl\,}(\nu\,{\rm
curl\,}\mathbf{H}),\quad{\rm div}\mathbf{H}=0,$ (1.3)
$\displaystyle\partial_{t}{\mathcal{E}}+{\rm
div}\left({\mathbf{u}}({\mathcal{E}}^{\prime}+P)\right)=\frac{1}{4\pi}{\rm
div}(({\mathbf{u}}\times\mathbf{H})\times\mathbf{H})$
$\displaystyle\qquad\qquad\qquad\qquad\qquad\,\,+{\rm
div}\Big{(}\frac{\nu}{4\pi}\mathbf{H}\times({\rm
curl\,}\mathbf{H})+{\mathbf{u}}\Psi({\mathbf{u}})+\kappa\nabla\theta\Big{)}.$
(1.4)
Here the unknowns $\rho$,
${\mathbf{u}}=(u_{1},u_{2},u_{3})\in{\mathbb{R}}^{3}$ ,
$\mathbf{H}=(H_{1},H_{2},H_{3})\in{\mathbb{R}}^{3}$, and $\theta$ denote the
density, velocity, magnetic field, and temperature, respectively;
$\Psi({\mathbf{u}})$ is the viscous stress tensor given by
$\Psi({\mathbf{u}})=2\mu\mathbb{D}({\mathbf{u}})+\lambda{\rm
div}{\mathbf{u}}\;\mathbf{I}_{3}$
with
$\mathbb{D}({\mathbf{u}})=(\nabla{\mathbf{u}}+\nabla{\mathbf{u}}^{\top})/2$,
$\mathbf{I}_{3}$ the $3\times 3$ identity matrix, and
$\nabla{\mathbf{u}}^{\top}$ the transpose of the matrix $\nabla{\mathbf{u}}$;
${\mathcal{E}}$ is the total energy given by
${\mathcal{E}}={\mathcal{E}}^{\prime}+|\mathbf{H}|^{2}/({8\pi})$ and
${\mathcal{E}}^{\prime}=\rho\left(e+|{\mathbf{u}}|^{2}/2\right)$ with $e$
being the internal energy, $\rho|{\mathbf{u}}|^{2}/2$ the kinetic energy, and
$|\mathbf{H}|^{2}/({8\pi})$ the magnetic energy. The viscosity coefficients
$\lambda$ and $\mu$ of the flow satisfy $\mu>0$ and $2\mu+3\lambda>0$. The
parameter $\nu>0$ is the magnetic diffusion coefficient of the magnetic field
and $\kappa>0$ the heat conductivity. For simplicity, we assume that
$\mu,\lambda,\nu$ and $\kappa$ are constants. The equations of state
$P=P(\rho,\theta)$ and $e=e(\rho,\theta)$ relate the pressure $P$ and the
internal energy $e$ to the density $\rho$ and the temperature $\theta$ of the
flow.
Multiplying (1.2) by ${\mathbf{u}}$ and (1.3) by $\mathbf{H}/({4\pi})$ and
summing over, one finds that
$\displaystyle\frac{d}{dt}\Big{(}\frac{1}{2}\rho|{\mathbf{u}}|^{2}+\frac{1}{8\pi}|\mathbf{H}|^{2}\Big{)}+\frac{1}{2}{\rm
div}\left(\rho|{\mathbf{u}}|^{2}{\mathbf{u}}\right)+\nabla P\cdot{\mathbf{u}}$
$\displaystyle\quad={\rm div}\Psi\cdot{\mathbf{u}}+\frac{1}{4\pi}({\rm
curl\,}\mathbf{H})\times\mathbf{H}\cdot{\mathbf{u}}+\frac{1}{4\pi}{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})\cdot\mathbf{H}$
$\displaystyle\,\,\qquad\qquad-\frac{\nu}{4\pi}{\rm curl\,}({\rm
curl\,}\mathbf{H})\cdot\mathbf{H}.$ (1.5)
Due to the identities
$\displaystyle{\rm div}(\mathbf{H}\times({\rm curl\,}\mathbf{H}))=|{\rm
curl\,}\mathbf{H}|^{2}-{\rm curl\,}({\rm curl\,}\mathbf{H})\cdot\mathbf{H},$
$\displaystyle{\rm div}(({\mathbf{u}}\times\mathbf{H})\times\mathbf{H})=({\rm
curl\,}\mathbf{H})\times\mathbf{H}\cdot{\mathbf{u}}+{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})\cdot\mathbf{H},$ (1.6)
one can subtract (1) from (1.4) to rewrite the energy equation (1.4) in terms
of the internal energy as
$\partial_{t}(\rho e)+{\rm div}(\rho{\mathbf{u}}e)+({\rm
div}{\mathbf{u}})P=\frac{\nu}{4\pi}|{\rm
curl\,}\mathbf{H}|^{2}+\Psi({\mathbf{u}}):\nabla{\mathbf{u}}+\kappa\Delta\theta,$
(1.7)
where $\Psi({\mathbf{u}}):\nabla{\mathbf{u}}$ denotes the scalar product of
two matrices:
$\Psi({\mathbf{u}}):\nabla{\mathbf{u}}=\sum^{3}_{i,j=1}\frac{\mu}{2}\left(\frac{\partial
u^{i}}{\partial x_{j}}+\frac{\partial u^{j}}{\partial
x_{i}}\right)^{2}+\lambda|{\rm
div}{\mathbf{u}}|^{2}=2\mu|\mathbb{D}({\mathbf{u}})|^{2}+\lambda|\mbox{tr}\mathbb{D}({\mathbf{u}})|^{2}.$
To establish the low Mach number limit for the system (1.1)–(1.3) and (1.7),
in this paper we shall focus on the ionized fluids obeying the following
perfect gas relations
$\displaystyle P=\mathfrak{R}\rho\theta,\quad e=c_{V}\theta,$ (1.8)
where the parameters $\mathfrak{R}>0$ and $c_{V}\\!>\\!0$ are the gas constant
and the heat capacity at constant volume, respectively, which will be assumed
to be one for simplicity of the presentation. We also ignore the coefficient
$1/(4\pi)$ in the magnetic field.
Let $\epsilon$ be the Mach number, which is a dimensionless number. Consider
the system (1.1)–(1.3), (1.7) in the physical regime:
$\displaystyle P\sim P_{0}+O(\epsilon),\quad{\mathbf{u}}\sim
O(\epsilon),\quad\mathbf{H}\sim O(\epsilon),\quad\nabla\theta\sim O(1),$
where $P_{0}>0$ is a certain given constant which is normalized to be
$P_{0}=1$. Thus we consider the case when the pressure $P$ is a small
perturbation of the given state $1$, while the temperature $\theta$ has a
finite variation. As in [2], we introduce the following transformation to
ensure positivity of $P$ and $\theta$
$\displaystyle\ P(x,t)=e^{\epsilon p^{\epsilon}(x,\epsilon
t)},\quad\theta(x,t)=e^{\theta^{\epsilon}(x,\epsilon t)},$ (1.9)
where a longer time scale $t=\tau/\epsilon$ (still denote $\tau$ by $t$ later
for simplicity) is introduced in order to seize the evolution of the
fluctuations. Note that (1.8) and (1.9) imply that $\rho(x,t)=e^{\epsilon
p^{\epsilon}(x,\epsilon t)-\theta^{\epsilon}(x,\epsilon t)}$ since
$\mathfrak{R}\equiv c_{V}\equiv 1$. Set
$\displaystyle{\mathbf{H}}(x,t)=\epsilon\mathbf{H}^{\epsilon}(x,\epsilon
t),\quad{{\mathbf{u}}}(x,t)=\epsilon{\mathbf{u}}^{\epsilon}(x,\epsilon t),$
(1.10)
and
$\displaystyle\mu=\epsilon\mu^{\epsilon},\quad\lambda=\epsilon\lambda^{\epsilon},\quad\nu=\epsilon\nu^{\epsilon},\quad\kappa=\epsilon\kappa^{\epsilon}.$
Under these changes of variables and coefficients, the system, (1.1)–(1.3),
(1.7) with (1.8), takes the following equivalent form:
$\displaystyle\partial_{t}p^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)p^{\epsilon}+\frac{1}{\epsilon}{\rm
div}(2{\mathbf{u}}^{\epsilon}-\kappa^{\epsilon}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla\theta^{\epsilon})$
$\displaystyle\quad\quad\qquad=\epsilon e^{-\epsilon
p^{\epsilon}}[\nu^{\epsilon}|{\rm
curl\,}\mathbf{H}^{\epsilon}|^{2}+\Psi({\mathbf{u}}^{\epsilon}):\nabla{\mathbf{u}}^{\epsilon}]+\kappa^{\epsilon}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla
p^{\epsilon}\cdot\nabla\theta^{\epsilon},$ (1.11) $\displaystyle
e^{-\theta^{\epsilon}}[\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon}]+\frac{\nabla
p^{\epsilon}}{\epsilon}=e^{-\epsilon p^{\epsilon}}[({\rm
curl\,}\mathbf{H}^{\epsilon})\times\mathbf{H}^{\epsilon}+{\rm
div}\Psi^{\epsilon}({\mathbf{u}}^{\epsilon})],$ (1.12)
$\displaystyle\partial_{t}\mathbf{H}^{\epsilon}-{\rm
curl\,}({\mathbf{u}}^{\epsilon}\times\mathbf{H}^{\epsilon})-\nu^{\epsilon}\Delta\mathbf{H}^{\epsilon}=0,\quad{\rm
div}\mathbf{H}^{\epsilon}=0,$ (1.13)
$\displaystyle\partial_{t}\theta^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)\theta^{\epsilon}+{\rm
div}{\mathbf{u}}^{\epsilon}$
$\displaystyle\quad\quad\qquad=\epsilon^{2}e^{-\epsilon
p^{\epsilon}}[\nu^{\epsilon}|{\rm
curl\,}\mathbf{H}^{\epsilon}|^{2}+\Psi^{\epsilon}({\mathbf{u}}^{\epsilon}):\nabla{\mathbf{u}}^{\epsilon}]+\kappa^{\epsilon}e^{-\epsilon
p^{\epsilon}}{\rm div}(e^{\theta^{\epsilon}}\nabla\theta^{\epsilon}),$ (1.14)
where
$\Psi^{\epsilon}({\mathbf{u}}^{\epsilon})=2\mu^{\epsilon}\mathbb{D}({\mathbf{u}}^{\epsilon})+\lambda^{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon}\,\mathbf{I}_{3}$, and the identity ${\rm
curl\,}({\rm curl\,}\mathbf{H}^{\epsilon})=\nabla{\rm
div}\mathbf{H}^{\epsilon}-\Delta\mathbf{H}^{\epsilon}$ and the constraint that
${\rm div}\mathbf{H}^{\epsilon}=0$ have been used.
We shall study the limit as $\epsilon\to 0$ of solutions to the system
(1.11)–(1.14). Formally, as $\epsilon$ goes to zero, if the sequence
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$
converges strongly to a limit $(1,{\mathbf{w}},{\mathbf{B}},\vartheta)$ in
some sense, and
$(\mu^{\epsilon},\lambda^{\epsilon},\nu^{\epsilon},\kappa^{\epsilon})$
converges to a constant vector
$(\bar{\mu},\bar{\lambda},\bar{\nu},\bar{\kappa})$, then taking the limit to
(1.11)–(1.14), we have
$\displaystyle{\rm
div}(2{\mathbf{w}}-\bar{\kappa}\,e^{\vartheta}\nabla\vartheta)=0,$ (1.15)
$\displaystyle
e^{-\vartheta}[\partial_{t}{\mathbf{w}}+({\mathbf{w}}\cdot\nabla){\mathbf{w}}]+\nabla\pi=({\rm
curl\,}{\mathbf{B}})\times{\mathbf{B}}+{\rm div}\Phi({\mathbf{w}}),$ (1.16)
$\displaystyle\partial_{t}{\mathbf{B}}-{\rm
curl\,}({\mathbf{w}}\times{\mathbf{B}})-\bar{\nu}\Delta{\mathbf{B}}=0,\quad{\rm
div}{\mathbf{B}}=0,$ (1.17)
$\displaystyle\partial_{t}\vartheta+({\mathbf{w}}\cdot\nabla)\vartheta+{\rm
div}{\mathbf{w}}=\bar{\kappa}\,{\rm div}(e^{\vartheta}\nabla\vartheta),$
(1.18)
with some function $\pi$, where $\Phi({\mathbf{w}})$ is defined by
$\Phi({\mathbf{w}})=2\bar{\mu}\mathbb{D}({\mathbf{w}})+\bar{\lambda}{\rm
div}{\mathbf{w}}\,\mathbf{I}_{3}.$
The purpose of this paper is to establish the above limit process rigorously.
For this purpose, we supplement the system (1.11)–(1.14) with the following
initial conditions
$\displaystyle(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})|_{t=0}=(p^{\epsilon}_{\rm
in}(x),{\mathbf{u}}^{\epsilon}_{\rm in}(x),\mathbf{H}^{\epsilon}_{\rm
in}(x),\theta^{\epsilon}_{\rm in}(x)),\quad x\in\mathbb{R}^{3}.$ (1.19)
For simplicity of presentation, we shall assume that
$\mu^{\epsilon}\equiv\bar{\mu}>0$, $\nu^{\epsilon}\equiv\bar{\nu}>0$,
$\kappa^{\epsilon}\equiv\bar{\kappa}>0$, and
$\lambda^{\epsilon}\equiv\bar{\lambda}$. The general case
$\mu^{\epsilon}\rightarrow\bar{\mu}>0$,
$\nu^{\epsilon}\rightarrow\bar{\nu}>0$,
$\kappa^{\epsilon}\rightarrow\bar{\kappa}>0$ and
$\lambda^{\epsilon}\rightarrow\bar{\lambda}$ simultaneously as
$\epsilon\rightarrow 0$ can be treated by slightly modifying the arguments
presented here.
As in [2], we will use the notation
$\|v\|_{H_{\eta}^{\sigma}}:=\|v\|_{H^{\sigma-1}}+\eta\|v\|_{H^{\sigma}}$ for
any $\sigma\in\mathbb{R}$ and $\eta\geq 0$. For each $\epsilon>0$, $t\geq 0$
and $s\geq 0$, we will also use the following norm:
$\displaystyle\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})(t)\|_{s,\epsilon}$
$\displaystyle\qquad:=\sup_{\tau\in[0,t]}\big{\\{}\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(\tau)\|_{H^{s}}+\|(\epsilon
p^{\epsilon},\epsilon{\mathbf{u}}^{\epsilon},\epsilon\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})(\tau)\|_{H_{\epsilon}^{s+2}}\big{\\}}$
$\displaystyle\qquad\ \
\quad+\Big{\\{}\int^{t}_{0}[\|\nabla(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\|^{2}_{H^{s}}+\|\nabla(\epsilon{\mathbf{u}}^{\epsilon},\epsilon
H^{\epsilon},\theta^{\epsilon})\|^{2}_{H_{\epsilon}^{s+2}}](\tau)d\tau\Big{\\}}^{1/2}.$
Then, the main result of this paper reads as follows.
###### Theorem 1.1.
Let $s\geq 4$. Assume that the initial data $(p^{\epsilon}_{\rm
in},{\mathbf{u}}^{\epsilon}_{\rm in},\mathbf{H}^{\epsilon}_{\rm
in},\theta^{\epsilon}_{\rm in})$ satisfy
$\displaystyle\|(p^{\epsilon}_{\rm in},{\mathbf{u}}^{\epsilon}_{\rm
in},\mathbf{H}^{\epsilon}_{\rm in})\|_{H^{s}}+\|(\epsilon p^{\epsilon}_{\rm
in},\epsilon{\mathbf{u}}^{\epsilon}_{\rm
in},\epsilon\mathbf{H}^{\epsilon}_{\rm in},\theta^{\epsilon}_{\rm
in}-\bar{\theta})\|_{H_{\epsilon}^{s+2}}\leq L_{0}$ (1.20)
for all $\epsilon\in(0,1]$ and two given positive constants $\bar{\theta}$ and
$L_{0}$. Then there exist positive constants $T_{0}$ and $\epsilon_{0}<1$,
depending only on $L_{0}$ and $\bar{\theta}$, such that the Cauchy problem
(1.11)–(1.14), (1.19) has a unique solution
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$
satisfying
$\displaystyle\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})(t)\|_{s,\epsilon}\leq
L,\qquad\forall\,\,t\in[0,T_{0}],\ \forall\,\epsilon\in(0,\epsilon_{0}],$
(1.21)
where $L$ depends only on $L_{0}$, $\bar{\theta}$ and $T_{0}$. Moreover,
assume further that the initial data satisfy the following conditions
$\displaystyle|\theta^{\epsilon}_{0}(x)-\bar{\theta}|\leq{N}_{0}|x|^{-1-\zeta},\quad|\nabla\theta^{\epsilon}_{0}(x)|\leq
N_{0}|x|^{-2-\zeta},\quad\forall\,\epsilon\in(0,1],$ (1.22)
$\displaystyle\\!\\!\\!\\!\big{(}p^{\epsilon}_{\rm in},{\rm
curl\,}(e^{-\theta^{\epsilon}_{\rm in}}{\mathbf{u}}^{\epsilon}_{\rm
in}),\mathbf{H}^{\epsilon}_{\rm in},\theta^{\epsilon}_{\rm
in}-\bar{\theta}\big{)}\rightarrow(0,{\mathbf{w}}_{0},{\mathbf{B}}_{0},\vartheta_{0}-\bar{\theta})\;\mbox{
in }H^{s}(\mathbb{R}^{3})$ (1.23)
as $\epsilon\rightarrow 0$, where $N_{0}$ and $\zeta$ are fixed positive
constants. Then the solution sequence
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},$
$\theta^{\epsilon}-\bar{\theta})$ converges weakly in
$L^{\infty}(0,T_{0};H^{s}({\mathbb{R}}^{3}))$ and strongly in $L^{2}(0,T_{0};$
$H^{s_{2}}_{\mathrm{loc}}({\mathbb{R}}^{3}))$ for all $0\leq s_{2}<s$ to the
limit $(0,{\mathbf{w}},{\mathbf{B}},\vartheta-\bar{\theta})$, where
$({\mathbf{w}},{\mathbf{B}},\vartheta)$ satisfies the system (1.15)–(1.18)
with initial data
$({\mathbf{w}},{\mathbf{B}},\vartheta)|_{t=0}=({\mathbf{w}}_{0},{\mathbf{B}}_{0},\vartheta_{0})$.
We now give some comments on the proof of Theorem 1.1. The key point in the
proof is to establish the uniform estimates in Sobolev norms for the acoustic
components of solutions, which are propagated by the wave equations whose
coefficients are functions of the temperature. Our main strategy is to bound
the norm of $(\nabla p^{\epsilon},{\rm div}{\mathbf{u}}^{\epsilon})$ in terms
of the norm of
$(\epsilon\partial_{t})(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$
and $(\epsilon
p^{\epsilon},\epsilon{\mathbf{u}}^{\epsilon},\epsilon\mathbf{H}^{\epsilon},\theta^{\epsilon})$
through the density and the momentum equations. This approach is motivated by
the previous works on the compressible Navier-Stokes equations due to Alazard
in [2], and Levermore, Sun and Trivisa [38]. It should be pointed out that the
analysis for (1.11)–(1.14) is complicated and difficult due to the strong
coupling of the hydrodynamic motion and the magnetic fields. Moreover, it is
observed that the terms $({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}$ in the momentum
equations, ${\rm
curl\,}({{\mathbf{u}}^{\epsilon}}\times{\mathbf{H}^{\epsilon}})$ in the
magnetic field equation, and $|\nabla\times{\mathbf{H}^{\epsilon}}|^{2}$ in
the temperature equation change basically the structure of the system. More
efforts should be paid on the estimates involving these terms, in particular,
on the estimate of higher order spatial derivatives. We shall exploit the
special structure of the system to obtain the tamed estimate on higher order
derivatives, so that we can close our estimates on the uniform boundedness of
the solutions. Once the uniform boundedness of the solutions has been
established, one can the convergence result in Theorem 1.1 by applying the
compactness arguments and the dispersive estimates on the acoustic wave
equations in the whole space developed in [40].
###### Remark 1.1.
The positivity of the coefficients $\mu$, $\nu$ and $\kappa$ plays an
fundamental role in the proof of Theorem 1.1. The arguments given in this
paper can not be applied to the case when one of them disappears. Recently,
Jiang, Ju and Li [31] have studied the incompressible limit of the
compressible non-isentropic ideal MHD equations with general initial data in
the whole space $\mathbb{R}^{d}$ ($d=2,3$) when the initial initial data
belong to $H^{s}(\mathbb{R}^{d})$ with $s$ being an even integer. We emphasize
that the restriction on the Sobolev index $s$ to be even plays a crucial role
in the proof since in this case the nonstandard highest order derivative
operators applied to the momentum equations are not intertwined with the
pressure equation, and thus we can apply the same operators to the magnetic
field equations to close the estimates on ${\mathbf{u}}$ and $\mathbf{H}$. On
the other hand, the proof presented in [31] fully exploits the structure of
the ideal MHD equations and can not be directly extended to the full
compressible MHD equations studied in the current paper, where the heat
conductivity is positive.
We point out that the low Mach number limit is an interesting topic in fluid
dynamics and applied mathematics. Now we briefly review some related results
on the Euler, Navier-Stokes and MHD equations. In [48], Schochet obtained the
convergence of the non-isentropic compressible Euler equations to the
incompressible non-isentropic Euler equations in a bounded domain for local
smooth solutions and well-prepared initial data. As mentioned above, in [40]
Métivier and Schochet proved rigorously the incompressible limit of the
compressible non-isentropic Euler equations in the whole space with general
initial data, see also [1, 2, 38] for further extensions. In [41] Métivier and
Schochet showed the incompressible limit of the one-dimensional non-isentropic
Euler equations in a periodic domain with general data. For compressible heat-
conducting flows, Hagstrom and Lorenz established in [18] the low Mach number
limit under the assumption that the variation of the density and temperature
is small. In the case of without heat conductivity, Kim and Lee [33]
investigated the incompressible limit to the non-isentropic Navier-Stokes
equations in a periodic domain with well-prepared data, while Jiang and Ou
[32] investigated the incompressible limit in three-dimensional bounded
domains, also for well-prepared data. The justification of the low Mach number
limit for the non-isentropic Euler or Navier-Stokes equations with general
initial data in bounded domains or multi-dimensional periodic domains is still
open. We refer the interested reader to [6] on formal computations for viscous
polytropic gases, and to [41, 5] for the study on the acoustic waves of the
non-isentropic Euler equations in periodic domains. Compared with the non-
isentropic case, the description of the propagation of oscillations in the
isentropic case is simpler and there are many articles on this topic
(isentropic flows) in the literature, see, for example, Ukai [50], Asano [3],
Desjardins and Grenier [11] in the whole space case; Isozaki [25, 26] in the
case of exterior domains; Iguchi [24] in the half space case; Schochet [47]
and Gallagher [16] in the case of periodic domains; and Lions and Masmoudi
[44], and Desjardins, et al. [12] in the case of bounded domains.
For the compressible isentropic MHD equations, the justification of the low
Mach number limit has been established in several aspects. In [34] Klainerman
and Majda studied the low Mach number limit to the compressible isentropic MHD
equations in the spatially periodic case with well-prepared initial data.
Recently, the low Mach number limit to the compressible isentropic viscous
(including both viscosity and magnetic diffusivity) MHD equations with general
data was studied in [23, 28, 29]. In [23] Hu and Wang obtained the convergence
of weak solutions to the compressible viscous MHD equations in bounded
domains, periodic domains and the whole space. In [28] Jiang, Ju and Li
employed the modulated energy method to verify the limit of weak solutions of
the compressible MHD equations in the torus to the strong solution of the
incompressible viscous or partially viscous MHD equations (zero shear
viscosity but with magnetic diffusion), while in [29] the convergence of weak
solutions of the viscous compressible MHD equations to the strong solution of
the ideal incompressible MHD equations in the whole space was established by
using the dispersion property of the wave equation, as both shear viscosity
and magnetic diffusion coefficients go to zero. For the full compressible MHD
equations, the incompressible limit in the framework of the so-called
variational solutions was established in [35, 36, 42]. Recently, the low Mach
number limit for the full compressible MHD equations with small entropy or
temperature variation was justified rigourously in [30]. Besides the
references mentioned above, the interested reader can refer to the monograph
[14] and the survey papers [9, 45, 49] for more related results on the low
Mach number limit to fluid models.
We also mention that there are a lot of articles in the literatures on the
other topics related to the compressible MHD equations due to theirs physical
importance, complexity, rich phenomena, and mathematical challenges, see, for
example, [4, 7, 8, 20, 39, 10, 13, 15, 19, 43, 21, 22, 52] and the references
cited therein.
This paper is arranged as follows. In Section 2, we describe some notations,
recall basic facts and present commutators estimates. In Section 3 we first
establish a priori estimates on $(\mathbf{H}^{\epsilon},\theta^{\epsilon})$,
$(\epsilon
p^{\epsilon},\epsilon{\mathbf{u}}^{\epsilon},\epsilon\mathbf{H}^{\epsilon},\theta^{\epsilon})$
and on $(p^{\epsilon},{\mathbf{u}}^{\epsilon})$. Then, with the help of these
estimates we establish the uniform boundeness of the solutions and prove the
existence part of Theorem 1.1. Finally, in Section 4 we study the local energy
decay for the acoustic wave equations and prove the convergence part of
Theorem 1.1.
## 2\. Preliminary
In this section, we give some notations and recall basic facts which will be
frequently used throughout the paper. We also present some commutators
estimates introduced in [38] and state the results on local solutions to the
Cauchy problem (1.11)–(1.14), (1.19).
We denote by $\langle\cdot,\cdot\rangle$ the standard inner product in
$L^{2}({\mathbb{R}}^{3})$ with norm $\langle f,f\rangle=\|f\|^{2}_{L^{2}}$ and
by $H^{k}$ the standard Sobolev space $W^{k,2}$ with norm $\|\cdot\|_{H^{k}}$.
The notation $\|(A_{1},\dots,A_{k})\|_{L^{2}}$ means the summation of
$\|A_{i}\|_{L^{2}},i=1,\cdots,k$, and it also applies to other norms. For a
multi-index $\alpha=(\alpha_{1},\alpha_{2},\alpha_{3})$, we denote
$\partial^{\alpha}=\partial^{\alpha_{1}}_{x_{1}}\partial^{\alpha_{2}}_{x_{2}}\partial^{\alpha_{3}}_{x_{3}}$
and $|\alpha|=|\alpha_{1}|+|\alpha_{2}|+|\alpha_{3}|$. We will omit the
spatial domain ${\mathbb{R}}^{3}$ in integrals for convenience. We use
$l_{i}>0$ ($i\in\mathbb{N}$) to denote given constants. We also use the symbol
$K$ or $C_{0}$ to denote generic positive constants, and $C(\cdot)$ to denote
a smooth function which may vary from line to line.
For a scalar function $f$, vector functions $\mathbf{a}$ and $\mathbf{b}$, we
have the following basic vector identities:
$\displaystyle{\rm div}(f\mathbf{a})$ $\displaystyle=f{\rm
div}\mathbf{a}+\nabla f\cdot\mathbf{a},$ (2.1) $\displaystyle{\rm
curl\,}(f\mathbf{a})$ $\displaystyle=f\cdot{\rm curl\,}\mathbf{a}-\nabla
f\times\mathbf{a},$ (2.2) $\displaystyle{\rm div}(\mathbf{a}\times\mathbf{b})$
$\displaystyle=\mathbf{b}\cdot{\rm curl\,}\mathbf{a}-\mathbf{a}\cdot{\rm
curl\,}\mathbf{b},$ (2.3) $\displaystyle{\rm
curl\,}(\mathbf{a}\times\mathbf{b})$
$\displaystyle=(\mathbf{b}\cdot\nabla)\mathbf{a}-(\mathbf{a}\cdot\nabla)\mathbf{b}+\mathbf{a}({\rm
div}\mathbf{b})-\mathbf{b}({\rm div}\mathbf{a}),$ (2.4)
$\displaystyle\nabla(\mathbf{a}\cdot\mathbf{b})$
$\displaystyle=(\mathbf{a}\cdot\nabla)\mathbf{b}+(\mathbf{b}\cdot\nabla)\mathbf{a}+\mathbf{a}\times({\rm
curl\,}\mathbf{b})+\mathbf{b}\times({\rm curl\,}\mathbf{a}).$ (2.5)
Below we recall some results on commutators estimates.
###### Lemma 2.1.
Let $\alpha=(\alpha_{1},\alpha_{2},\alpha_{3})$ be a multi-index such that
$|\alpha|=k$. Then, for any $\sigma\geq 0$, there exists a positive constant
$C_{0}$, such that for all $f,g\in H^{k+\sigma}({\mathbb{R}}^{3})$,
$\displaystyle\|[f,\partial^{\alpha}]g\|_{H^{\sigma}}\leq$ $\displaystyle
C_{0}(\|f\|_{W^{1,\infty}}\|g\|_{H^{\sigma+k-1}}+\|f\|_{H^{\sigma+k}}\|g\|_{L^{\infty}}).$
(2.6)
###### Lemma 2.2.
Let $s>5/2$. Then there exists a positive constant $C_{0}$, such that for all
$\epsilon\in(0,1]$, $T>0$ and multi-index
$\beta=(\beta_{1},\beta_{2},\beta_{3})$ satisfying $0\leq|\beta|\leq s-1$, and
any $f,g\in C^{\infty}([0,T],H^{s}({\mathbb{R}}^{3}))$, it holds that
$\displaystyle\|[f,\partial^{\beta}(\epsilon\partial_{t})]g\|_{L^{2}}\leq$
$\displaystyle\epsilon
C_{0}(\|f\|_{H^{s-1}}\|\partial_{t}g\|_{H^{s-2}}+\|\partial_{t}f\|_{H^{s-1}}\|g\|_{H^{s-1}}).$
(2.7)
Since the system (1.1)–(1.3), (1.7), (1.8) is hyperbolic-parabolic, so the
classical result of Vol’pert and Hudiaev [51] implies that
###### Proposition 2.3.
Let $s\geq 4$. Assume that the initial data
$(\rho_{0},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0})$ satisfy
$\displaystyle\|(\rho_{0}-\underline{\rho},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0}-\underline{\theta})\|_{H^{s}}\leq
C_{0}$
for some positive constants $\underline{\rho}$, $\underline{\theta}$ and
$C_{0}$. Then there exists a $\tilde{T}>0$, such that the system (1.1)–(1.3),
(1.7), and (1.8) with these initial data has a unique classical solution
$(\rho,{\mathbf{u}},\mathbf{H},\theta)$ enjoying $\rho-\underline{\rho}\in
C([0,\tilde{T}],H^{s}({\mathbb{R}}^{3}))$,
$({\mathbf{u}},\mathbf{H},\theta-\bar{\theta})\in C([0,\tilde{T}],\linebreak
H^{s}({\mathbb{R}}^{3}))\cap L^{2}(0,\tilde{T};H^{s+1}({\mathbb{R}}^{3}))$,
and
$\displaystyle\sup_{0\leq
t\leq\tilde{T}}\|(\rho-\underline{\rho},{\mathbf{u}},\mathbf{H},\theta-\bar{\theta})\|^{2}_{H^{s}}$
$\displaystyle+\int^{\tilde{T}}_{0}\Big{\\{}\mu\|\mathbb{D}({\mathbf{u}})\|^{2}_{H^{s}}+\lambda\|{\rm
div}{\mathbf{u}}\|^{2}_{H^{s}}$
$\displaystyle\qquad\quad+\nu\|\nabla\mathbf{H}\|^{2}_{H^{s}}+\kappa\|\nabla\theta\|^{2}_{H^{s}}\Big{\\}}(\tau)d\tau\leq
4C_{0}^{2}.$
It follows from Proposition 2.3, and the transforms (1.9) and (1.10) that
there exists a $T_{\epsilon}>0$, depending on $\epsilon$ and $L_{0}$, such
that for each fixed $\epsilon$ and any initial data (1.19) satisfying (1.20),
the Cauchy problem (1.11)–(1.14), (1.19) has a unique solution
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})$
satisfying
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})\in
C([0,T_{\epsilon}),H^{s}({\mathbb{R}}^{3}))$ and
$({\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})\in
L^{2}(0,T_{\epsilon};H^{s+1}({\mathbb{R}}^{3}))$. Moreover, let
$T_{\epsilon}^{*}$ be the maximal time of existence of such a smooth solution,
then if $T_{\epsilon}^{*}$ is finite, one has
$\displaystyle{\underset{t\rightarrow
T_{\epsilon}^{*}}{\lim\sup}}\,\left\\{\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})(t)\|_{W^{1,\infty}}+\|({\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})(t)\|_{W^{2,\infty}}\right\\}=\infty.$
Therefore, we shall see by the same argument as in [40] that the existence
part of Theorem 1.1 is a consequence of the above assertion and the following
key a priori estimates which will be shown in the next section.
###### Proposition 2.4.
For any given $s\geq 4$ and fixed $\epsilon>0$, let
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$
be the classical solution to the Cauchy problem (1.11)–(1.14) and (1.19).
Denote
$\displaystyle\mathcal{O}(T):=$
$\displaystyle\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}-\bar{\theta})(T)\|_{s,\epsilon},$
$\displaystyle\mathcal{O}_{0}:=$ $\displaystyle\|(p^{\epsilon}_{\rm
in},{\mathbf{u}}^{\epsilon}_{\rm in},\mathbf{H}^{\epsilon}_{\rm
in})\|_{H^{s}}+\|(\epsilon p^{\epsilon}_{\rm
in},\epsilon{\mathbf{u}}^{\epsilon}_{\rm
in},\epsilon\mathbf{H}^{\epsilon}_{\rm in},\theta^{\epsilon}_{\rm
in}-\bar{\theta})\|_{H_{\epsilon}^{s+2}}.$
Then there exist positive constants $\hat{T}_{0}$ and $\epsilon_{0}<1$, and an
increasing positive function $C(\cdot)$, such that for all
$T\in[0,\hat{T}_{0}]$ and $\epsilon\in(0,\epsilon_{0}]$,
$\displaystyle\mathcal{O}(T)\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$
## 3\. Uniform estimates
In this section we shall establish the uniform bounds of the solutions to the
Cauchy problem (1.11)–(1.14) and (1.19) stated in Proposition 2.4 by modifying
the approaches developed in [40, 2, 38] and making careful use of the special
structure of the system (1.11)–(1.14). In the rest of this section, we will
drop the superscripts $\epsilon$ of the variables in the Cauchy problem and
denote
$\Psi({\mathbf{u}})=2\bar{\mu}\mathbb{D}({\mathbf{u}})+\bar{\lambda}{\rm
div}{\mathbf{u}}\;\mathbf{I}_{3}.$
Recall that it has been assumed that $\mu^{\epsilon}\equiv\bar{\mu}>0$,
$\nu^{\epsilon}\equiv\bar{\nu}>0$, $\kappa^{\epsilon}\equiv\bar{\kappa}>0$,
and $\lambda^{\epsilon}\equiv\bar{\lambda}$ independent of $\epsilon$.
### 3.1. $H^{s}$-estimates on $(\mathbf{H},\theta)$ and $(\epsilon
p,\epsilon{\mathbf{u}})$
To prove Proposition 2.4, we first give some estimates derived directly from
the system (1.11)–(1.14). Denoting
$\displaystyle\mathcal{Q}:$
$\displaystyle=\|(p,{\mathbf{u}},\mathbf{H},\theta-\bar{\theta})\|_{H^{s}}+\|(\epsilon
p,\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta-\bar{\theta})\|_{H_{\epsilon}^{s+2}},$
$\displaystyle\mathcal{S}:$ $\displaystyle=\|(\nabla{\mathbf{u}},\nabla
p,\nabla\mathbf{H})\|_{H^{s}}+\|\nabla(\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta)\|_{H_{\epsilon}^{s+2}},$
one has
###### Lemma 3.1.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be a solution to the
problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. There exists an increasing
function $C(\cdot)$ such that, for any $\epsilon\in(0,1]$ and $t\in[0,T]$,
$T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\\{\|(\mathbf{H},\theta)(\tau)\|_{H^{s}}+\|\epsilon\mathbf{H}(\tau)\|_{H_{\epsilon}^{s+1}}\\}$
$\displaystyle\quad+\Big{\\{}\int^{t}_{0}\big{(}\|\nabla(\mathbf{H},\theta)(\tau)\|_{H^{s}}^{2}+\|\nabla(\epsilon\mathbf{H})(\tau)\|^{2}_{H_{\epsilon}^{s+1}}\Big{)}d\tau\Big{\\}}^{1/2}$
$\displaystyle\qquad\leq
C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
###### Proof.
For any multi-index $\alpha=(\alpha_{1},\alpha_{2},\alpha_{3})$ satisfying
$|\alpha|\leq s$, let $\mathbf{H}_{\alpha}=\partial^{\alpha}\mathbf{H}$. Then
$\displaystyle\partial_{t}\mathbf{H}_{\alpha}+({\mathbf{u}}\cdot\nabla)\mathbf{H}_{\alpha}-\bar{\nu}\Delta\mathbf{H}_{\alpha}=-[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\mathbf{H}-\partial^{\alpha}(\mathbf{H}{\rm
div}{\mathbf{u}})+\partial^{\alpha}((\mathbf{H}\cdot\nabla){\mathbf{u}}).$
Taking inner product of the above equations with $\mathbf{H}_{\alpha}$ and
integrating by parts, we have
$\displaystyle\frac{1}{2}\frac{d}{dt}\|\mathbf{H}_{\alpha}\|_{L^{2}}^{2}+\bar{\nu}\|\nabla\mathbf{H}_{\alpha}\|_{L^{2}}^{2}=$
$\displaystyle-\langle({\mathbf{u}}\cdot\nabla)\mathbf{H}_{\alpha},\mathbf{H}_{\alpha}\rangle-\langle[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\mathbf{H},\mathbf{H}_{\alpha}\rangle$
$\displaystyle-\langle\partial^{\alpha}(\mathbf{H}{\rm
div}{\mathbf{u}}),\mathbf{H}_{\alpha}\rangle+\langle\partial^{\alpha}((\mathbf{H}\cdot\nabla){\mathbf{u}}),\mathbf{H}_{\alpha}\rangle.$
By integration by parts we obtain
$\displaystyle-\langle({\mathbf{u}}\cdot\nabla)\mathbf{H}_{\alpha},\mathbf{H}_{\alpha}\rangle=\frac{1}{2}\int{\rm
div}{\mathbf{u}}|\mathbf{H}_{\alpha}|^{2}dx\leq
C(\mathcal{Q})\|\mathbf{H}_{\alpha}\|^{2}_{L^{2}}.$
It follows from the commutator inequality (2.6) that
$\|[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\mathbf{H}\|_{L^{2}}\leq
C_{0}(\|{\mathbf{u}}\|_{W^{1,\infty}}\|\nabla\mathbf{H}\|_{H^{s-1}}+\|{\mathbf{u}}\|_{H^{s}}\|\nabla\mathbf{H}\|_{L^{\infty}})\leq
C(\mathcal{Q}).$
By Sobolev’s inequality, one gets
$-\langle\partial^{\alpha}(\mathbf{H}{\rm
div}{\mathbf{u}}),\mathbf{H}_{\alpha}\rangle+\langle\partial^{\alpha}(\mathbf{H}\cdot\nabla{\mathbf{u}}),\mathbf{H}_{\alpha}\rangle\leq
C_{0}\|\mathbf{H}\|_{H^{s}}^{2}\|{\mathbf{u}}\|_{H^{s+1}}\leq
C(\mathcal{Q})\mathcal{S}.$
Thus, we conclude that
$\displaystyle\ \ \
\sup_{\tau\in[0,t]}\|\mathbf{H}(\tau)\|_{H^{s}}+\bar{\nu}\Big{\\{}\int^{t}_{0}\|\nabla\mathbf{H}\|^{2}_{H^{s}}d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq
C(\mathcal{O}_{0})+C(\mathcal{O}(t)t+C(\mathcal{O}(t))\int_{0}^{t}\mathcal{S}(\tau)d\tau$
$\displaystyle\leq
C(\mathcal{O}_{0})+C(\mathcal{O}(t)t+C(\mathcal{O}(t))\sqrt{t}$
$\displaystyle\leq C(\mathcal{O}_{0})+C(\mathcal{O}(T))\sqrt{T}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
Now denote $\hat{\mathbf{H}}=\epsilon\mathbf{H}$ and
$\hat{\mathbf{H}}_{\alpha}=\partial^{\alpha}(\epsilon\mathbf{H})$ for
$|\alpha|=s+1$. Then, $\hat{\mathbf{H}}_{\alpha}$ satisfies
$\displaystyle\partial_{t}\hat{\mathbf{H}}_{\alpha}+({\mathbf{u}}\cdot\nabla)\hat{\mathbf{H}}_{\alpha}-\bar{\nu}\Delta\hat{\mathbf{H}}_{\alpha}=-\epsilon[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\mathbf{H}-\epsilon\partial^{\alpha}(\mathbf{H}{\rm
div}{\mathbf{u}})+\epsilon\partial^{\alpha}((\mathbf{H}\cdot\nabla){\mathbf{u}}).$
(3.1)
The commutator inequality (2.6) implies that
$\|-\epsilon[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\mathbf{H}\|_{L^{2}}\leq
C_{0}(\|{\mathbf{u}}\|_{W^{1,\infty}}\|\nabla\hat{\mathbf{H}}\|_{H^{s}}+\|\epsilon{\mathbf{u}}\|_{H^{s+1}}\|\nabla\mathbf{H}\|_{L^{\infty}})\leq
C(\mathcal{Q}),$
while an integration by parts and Sobolev’s inequality lead to
$\displaystyle-\langle\epsilon\partial^{\alpha}(\mathbf{H}{\rm
div}{\mathbf{u}}),\hat{\mathbf{H}}_{\alpha}\rangle+\langle\epsilon\partial^{\alpha}((\mathbf{H}\cdot\nabla){\mathbf{u}}),\hat{\mathbf{H}}_{\alpha}\rangle$
$\displaystyle\leq\frac{\bar{\nu}}{2}\|\nabla\hat{\mathbf{H}}_{\alpha}\|_{L^{2}}^{2}+C_{0}\|\mathbf{H}\|_{H^{s}}^{2}\|\epsilon{\mathbf{u}}\|_{H^{s+1}}^{2}$
$\displaystyle\leq\frac{\bar{\nu}}{2}\|\nabla\hat{\mathbf{H}}_{\alpha}\|_{L^{2}}^{2}+C(\mathcal{Q}).$
Hence, we obtain
$\displaystyle\sup_{\tau\in[0,t]}\|\epsilon\mathbf{H}(\tau)\|_{H^{s+1}}+\Big{\\{}\int^{t}_{0}\bar{\nu}\|\nabla(\epsilon\mathbf{H})(\tau)\|^{2}_{H^{s+1}}d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
Similarly, we can obtain
$\displaystyle\sup_{\tau\in[0,t]}\epsilon^{2}\|\mathbf{H}(\tau)\|_{H^{s+2}}+\epsilon^{2}\Big{\\{}\int^{t}_{0}\bar{\nu}\|\nabla\mathbf{H}(\tau)\|^{2}_{H^{s+2}}d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
Next, we estimate $\theta$. Using Sobolev’s inequality, one finds that
$\displaystyle\|\partial^{s}(\epsilon^{2}e^{-\epsilon p}[\bar{\nu}|{\rm
curl\,}\mathbf{H}|^{2}+\Psi({\mathbf{u}}):\nabla{\mathbf{u}}])\|_{L^{2}}$
$\displaystyle\qquad\qquad\qquad\leq C_{0}\|\epsilon
p\|_{H^{s}}(\|\epsilon\nabla\mathbf{H}\|^{2}_{H^{s}}+\|\epsilon\nabla{\mathbf{u}}\|^{2}_{H^{s}})\leq
C(\mathcal{Q}).$
Employing arguments similar to those used for $\mathbf{H}$, we can obtain
$\displaystyle\sup_{\tau\in[0,t]}\|\theta(\tau)\|_{H^{s}}+\Big{\\{}\int^{t}_{0}\bar{\kappa}\|\nabla\theta(\tau)\|^{2}_{H^{s}}d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
Thus, the lemma is proved. ∎
###### Lemma 3.2.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be a solution to the
problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. Then there exists an increasing
function $C(\cdot)$ such that, for any $\epsilon\in(0,1]$ and
$t\in[0,T],T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\|(\epsilon
p,\epsilon{\mathbf{u}})(\tau)\|_{H^{s}}+\Big{\\{}\int^{t}_{0}\bar{\mu}\|\nabla(\epsilon{\mathbf{u}})(\tau)\|^{2}_{H^{s}}d\tau\Big{\\}}^{1/2}\leq
C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
###### Proof.
Let $\check{p}=\epsilon p$, and $\check{p}_{\alpha}=\partial^{\alpha}(\epsilon
p)$ for any multi-index $\alpha=(\alpha_{1},\alpha_{2},\alpha_{3})$ satisfying
$|\alpha|\leq s$. Then
$\displaystyle\partial_{t}\check{p}_{\alpha}+({\mathbf{u}}\cdot\nabla)\check{p}_{\alpha}=$
$\displaystyle-[\partial^{\alpha},{\mathbf{u}}]\cdot(\nabla\check{p})-\partial^{\alpha}[{\rm
div}(2{\mathbf{u}}-\kappa a(\epsilon p)b(\theta)\nabla\theta)]$
$\displaystyle+\partial^{\alpha}\\{a(\epsilon p)[\nu|{\rm
curl\,}(\epsilon\mathbf{H})|^{2}+\Psi(\epsilon{\mathbf{u}}):\nabla(\epsilon{\mathbf{u}})]\\}$
$\displaystyle+\kappa\partial^{\alpha}\\{a(\epsilon p)b(\theta)\nabla(\epsilon
p)\cdot\nabla\theta\\}$ $\displaystyle:=$ $\displaystyle
h_{1}+h_{2}+h_{3}+h_{4},$ (3.2)
where, for simplicity of presentation, we have set
$a(\epsilon p):=e^{-\epsilon p},\quad b(\theta):=e^{\theta}.$
It is easy to see that the energy estimate for (3.2) gives
$\displaystyle\frac{1}{2}\frac{d}{dt}\|\check{p}_{\alpha}\|_{L^{2}}^{2}=-\langle({\mathbf{u}}\cdot\nabla)\check{p}_{\alpha},\check{p}_{\alpha}\rangle+\langle
h_{1}+h_{2}+h_{3}+h_{4},\check{p}_{\alpha}\rangle,$ (3.3)
where we have to estimate each term on the right-hand side of (3.3). First, an
integration by parts yields
$\displaystyle-\langle({\mathbf{u}}\cdot\nabla)\check{p}_{\alpha},\check{p}_{\alpha}\rangle=\frac{1}{2}\int{\rm
div}{\mathbf{u}}|\check{p}_{\alpha}|^{2}dx\leq
C(\mathcal{Q})\|\check{p}_{\alpha}\|^{2}_{L^{2}},$
while the commutator inequality leads to
$\displaystyle\|h_{1}\|\leq
C_{0}(\|{\mathbf{u}}\|_{W^{1,\infty}}\|\nabla\check{p}\|_{H^{s-1}}+\|{\mathbf{u}}\|_{H^{s}}\|\nabla\check{p}\|_{L^{\infty}})\leq
C(\mathcal{Q}).$
Consequently,
$\displaystyle\langle
h_{1},\check{p}_{\alpha}\rangle\leq\|\check{p}_{\alpha}\|_{L^{2}}\|h_{1}\|_{L^{2}}\leq
C(\mathcal{Q}).$
From Sobolev’s inequality one gets
$\|h_{2}\|\leq C_{0}\|{\mathbf{u}}\|_{H^{s+1}}+\|\theta\|_{H^{s}}\|\epsilon
p\|_{H^{s}}\|\theta\|_{H^{s+2}}\leq C(\mathcal{S})(1+C(\mathcal{Q})),$
whence,
$\displaystyle\langle h_{2},\check{p}_{\alpha}\rangle\leq$ $\displaystyle
C(\mathcal{Q})C(\mathcal{S}).$
Similarly, one can prove that
$\displaystyle\langle(h_{3}+h_{4}),\check{p}_{\alpha}\rangle\leq$
$\displaystyle C(\mathcal{Q}).$
Hence, we conclude that
$\displaystyle\sup_{\tau\in[0,t]}\|\epsilon p(\tau)\|_{H^{s}}\leq
C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
In a similar way, we can estimate ${\mathbf{u}}$. Thus the proof of the lemma
is completed. ∎
Next, we control the term $\|({\mathbf{u}},p)\|_{H^{s}}$. The idea is to bound
the norm of $({\rm div}{\mathbf{u}},$ $\nabla p)$ in terms of the suitable
norm of $(\epsilon{\mathbf{u}},\epsilon p,\epsilon\mathbf{H},\theta)$ and
$\epsilon(\partial_{t}{\mathbf{u}},\partial_{t}p)$ by making use of the
structure of the system. To this end, we first estimate
$\|(\epsilon{\mathbf{u}},\epsilon p,\theta)\|_{H^{s+1}}$.
### 3.2. $H^{s+1}$-estimates on $(\epsilon{\mathbf{u}},\epsilon
p,\epsilon\mathbf{H},\theta)$
Following [2], we set
$\displaystyle(\hat{p},\hat{\mathbf{u}},\hat{\mathbf{H}},\hat{\theta}):=(\epsilon
p-\theta,\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta-\bar{\theta}).$
A straightforward calculation results in that
$(\hat{p},\hat{\mathbf{u}},\hat{\mathbf{H}},\hat{\theta})$ solves the
following system:
$\displaystyle\partial_{t}\hat{p}+({\mathbf{u}}\cdot\nabla)\hat{p}+\frac{1}{\epsilon}{\rm
div}\hat{\mathbf{u}}=0,$ (3.4) $\displaystyle
b(-\theta)[\partial_{t}\hat{\mathbf{u}}+({\mathbf{u}}\cdot\nabla)\hat{\mathbf{u}}]+\frac{1}{\epsilon}(\nabla\hat{p}+\nabla\hat{\theta})$
$\displaystyle\qquad\qquad\qquad\qquad\qquad=a(\epsilon p)[({\rm
curl\,}\mathbf{H})\times\hat{\mathbf{H}}+{\rm div}\Psi(\hat{\mathbf{u}})],$
(3.5)
$\displaystyle\partial_{t}\hat{\mathbf{H}}+{\mathbf{u}}\cdot\nabla\hat{\mathbf{H}}+\mathbf{H}{\rm
div}\hat{\mathbf{u}}-\mathbf{H}\cdot\nabla\hat{\mathbf{u}}-\bar{\nu}\Delta\hat{\mathbf{H}}=0,\quad{\rm
div}\hat{\mathbf{H}}=0,$ (3.6)
$\displaystyle\partial_{t}\hat{\theta}+({\mathbf{u}}\cdot\nabla)\hat{\theta}+\frac{1}{\epsilon}{\rm
div}\hat{\mathbf{u}}=\epsilon a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm curl\,}\hat{\mathbf{H}}+\epsilon a(\epsilon
p)\Psi({\mathbf{u}}):\nabla\hat{\mathbf{u}}]$
$\displaystyle\qquad\qquad\qquad\qquad\qquad\ \quad+\bar{\kappa}a(\epsilon
p){\rm div}(b(\theta)\nabla\hat{\theta}).$ (3.7)
We have
###### Lemma 3.3.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be a solution to the
problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. Then there exist a constant
$l_{1}>0$ and an increasing function $C(\cdot)$ such that, for any
$\epsilon\in(0,1]$ and $t\in[0,T]$, $T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\|(\epsilon
q,\epsilon{\mathbf{u}},\theta-\bar{\theta})(\tau)\|_{H^{s+1}}+l_{1}\Big{\\{}\int^{t}_{0}\|\nabla(\epsilon{\mathbf{u}},\theta)\|^{2}_{H^{s+1}}(\tau)d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
(3.8)
###### Proof.
Let $\alpha=(\alpha_{1},\alpha_{2},\alpha_{3})$ be a multi-index such that
$|\alpha|=s+1$. Set
$\displaystyle(\hat{p}_{\alpha},\hat{{\mathbf{u}}}_{\alpha},\hat{H}_{\alpha},\hat{\theta}_{\alpha}):=\left(\partial^{\alpha}(\epsilon
p-\theta),\partial^{\alpha}(\epsilon{\mathbf{u}}),\partial^{\alpha}(\epsilon\mathbf{H}),\partial^{\alpha}(\theta-\bar{\theta})\right).$
Then, $\hat{\mathbf{H}}_{\alpha}$ satisfies (3.1) and
$(\hat{p}_{\alpha},\hat{{\mathbf{u}}}_{\alpha},\hat{\theta}_{\alpha})$ solves
$\displaystyle\partial_{t}\hat{p}_{\alpha}+({\mathbf{u}}\cdot\nabla)\hat{p}_{\alpha}+\frac{1}{\epsilon}{\rm
div}\hat{{\mathbf{u}}}_{\alpha}=g_{1},$ (3.9) $\displaystyle
b(-\theta)[\partial_{t}\hat{{\mathbf{u}}}_{\alpha}+({\mathbf{u}}\cdot\nabla)\hat{{\mathbf{u}}}_{\alpha}]+\frac{1}{\epsilon}(\nabla\hat{p}_{\alpha}+\nabla\hat{\theta}_{\alpha})$
$\displaystyle\quad\quad\qquad\quad=a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\hat{\mathbf{H}}_{\alpha}+a(\epsilon p){\rm
div}\Psi(\hat{{\mathbf{u}}}_{\alpha})+g_{2},$ (3.10)
$\displaystyle\partial_{t}\hat{\theta}_{\alpha}+({\mathbf{u}}\cdot\nabla)\hat{\theta}_{\alpha}+\frac{1}{\epsilon}{\rm
div}\hat{{\mathbf{u}}}_{\alpha}=\epsilon a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm
curl\,}\hat{\mathbf{H}}_{\alpha}+\Psi({\mathbf{u}}):\nabla\hat{\mathbf{u}}_{\alpha}]$
$\displaystyle\quad\quad\qquad\quad+\bar{\kappa}a(\epsilon p){\rm
div}(b(\theta)\nabla\hat{\theta}_{\alpha})+g_{3},$ (3.11)
with initial data
$\displaystyle(\hat{p}_{\alpha},\hat{{\mathbf{u}}}_{\alpha},\hat{\mathbf{H}}_{\alpha},\hat{\theta}_{\alpha})|_{t=0}:=\big{(}$
$\displaystyle\partial^{\alpha}(\epsilon p_{\rm in}(x)-\theta_{\rm
in}(x)),\partial^{\alpha}(\epsilon{\mathbf{u}}_{\rm in}(x)),$
$\displaystyle\partial^{\alpha}\mathbf{H}_{\rm
in}(x)),\partial^{\alpha}(\theta_{\rm in}(x)-\bar{\theta})\big{)},$ (3.12)
where
$\displaystyle g_{1}:=$
$\displaystyle-[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla(\epsilon
p-\theta),$ $\displaystyle g_{2}:=$
$\displaystyle-[\partial^{\alpha},b(-\theta)]\partial_{t}(\epsilon{\mathbf{u}})-[\partial^{\alpha},b(-\theta){\mathbf{u}}]\cdot\nabla(\epsilon{\mathbf{u}})$
$\displaystyle+[\partial^{\alpha},a(\epsilon p){\rm
curl\,}(\epsilon\mathbf{H})]\times\mathbf{H}+[\partial^{\alpha},a(\epsilon
p)]{\rm div}\Psi(\epsilon{\mathbf{u}}),$ $\displaystyle g_{3}:=$
$\displaystyle-[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla\theta+\bar{\nu}\,[\partial^{\alpha},a(\epsilon
p){\rm curl\,}(\epsilon\mathbf{H})]:{\rm curl\,}(\epsilon\mathbf{H})$
$\displaystyle+\epsilon[\partial^{\alpha},a(\epsilon
p)\Psi({\mathbf{u}})]:\nabla(\epsilon{\mathbf{u}})+\bar{\kappa}\partial^{\alpha}\big{(}a(\epsilon
p){\rm div}(b(\theta)\nabla\theta)\big{)}$
$\displaystyle-\bar{\kappa}a(\epsilon p){\rm
div}(b(\theta)\nabla\hat{\theta}_{\alpha}).$
It follows from Proposition 2.3 and the positivity of $a(\cdot)$ and
$b(\cdot)$ that $a(\cdot)$ and $b(\cdot)$ are bounded away from $0$ uniformly
with respect to $\epsilon$, i.e.
$\displaystyle a(\epsilon
p)\geq\underline{a}>0,\;\;\;b(-\theta)\geq\underline{b}>0.$ (3.13)
The standard $L^{2}$-energy estimates for (3.9), (3.10) and (3.11) yield that
$\displaystyle\frac{1}{2}\frac{d}{dt}\big{(}\|\hat{p}_{\alpha}\|_{L^{2}}^{2}+\langle
b(-\theta)\hat{{\mathbf{u}}}_{\alpha},\hat{{\mathbf{u}}}_{\alpha}\rangle+\|\hat{\theta}_{\alpha}\|_{L^{2}}^{2}\big{)}$
$\displaystyle\leq\frac{1}{2}\langle
b_{t}(\theta)\hat{\mathbf{u}}_{\alpha},\hat{\mathbf{u}}_{\alpha}\rangle-\langle({\mathbf{u}}\cdot\nabla)\hat{p}_{\alpha},\hat{p}_{\alpha}\rangle-\langle
b(-\theta)({\mathbf{u}}\cdot\nabla)\hat{{\mathbf{u}}}_{\alpha},\hat{{\mathbf{u}}}_{\alpha}\rangle-\langle({\mathbf{u}}\cdot\nabla)\hat{\theta}_{\alpha},\hat{\theta}_{\alpha}\rangle$
$\displaystyle\quad+\langle a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\hat{\mathbf{H}}_{\alpha},\hat{\mathbf{u}}_{\alpha}\rangle+\langle
a(\epsilon p){\rm
div}\Psi(\hat{{\mathbf{u}}}_{\alpha}),\hat{\mathbf{u}}_{\alpha}\rangle$
$\displaystyle\qquad\langle\epsilon a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm
curl\,}\hat{\mathbf{H}}_{\alpha}+\Psi({\mathbf{u}}):\nabla\hat{\mathbf{u}}_{\alpha}],\theta_{\alpha}\rangle+\langle\bar{\kappa}a(\epsilon
p){\rm div}(b(\theta)\nabla\theta_{\alpha}),\hat{\theta}_{\alpha}\rangle$
$\displaystyle\quad+\langle g_{1},\hat{p}_{\alpha}\rangle+\langle
g_{2},\hat{u}_{\alpha}\rangle+\langle g_{3},\hat{\theta}_{\alpha}\rangle.$
(3.14)
It follows from equation (1.14), and the definition of $\mathcal{Q}$ and
$\mathcal{S}$ that
$\displaystyle\|b_{t}(\theta)\|_{L^{\infty}}\leq\|b(\theta)\|_{H^{s}}\|\theta_{t}\|_{H^{s}}\leq
C(\mathcal{Q})(1+\mathcal{S}).$
Therefore,
$\displaystyle\frac{1}{2}\langle
b_{t}(\theta)\hat{\mathbf{u}}_{\alpha},\hat{\mathbf{u}}_{\alpha}\rangle\leq
C(\mathcal{Q})(1+\mathcal{S}).$
On the other hand, it is easy to see that
$\displaystyle-\langle({\mathbf{u}}\cdot\nabla)\hat{p}_{\alpha},\hat{p}_{\alpha}\rangle-\langle
b(-\theta)({\mathbf{u}}\cdot\nabla)\hat{{\mathbf{u}}}_{\alpha},\hat{{\mathbf{u}}}_{\alpha}\rangle-\langle({\mathbf{u}}\cdot\nabla)\hat{\theta}_{\alpha},\hat{\theta}_{\alpha}\rangle\leq
C(\mathcal{Q})$
and
$\displaystyle\langle a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\hat{\mathbf{H}}_{\alpha},\hat{\mathbf{u}}_{\alpha}\rangle\leq
C(\mathcal{Q}).$
By integration by parts we have
$\displaystyle-\langle a(\epsilon p_{0}){\rm
div}\Psi(\hat{\mathbf{u}}),\hat{\mathbf{u}}\rangle=$ $\displaystyle\int\mu
a(\epsilon p_{0})[|\nabla\hat{\mathbf{u}}_{\alpha}|^{2}+(\mu+\lambda)|{\rm
div}\hat{\mathbf{u}}_{\alpha}|^{2}]dx$ $\displaystyle+\mu\langle(\nabla
a(\epsilon
p)\cdot\nabla)\hat{\mathbf{u}}_{\alpha},\hat{\mathbf{u}}_{\alpha}\rangle$
$\displaystyle+(\mu+\lambda)\langle(\nabla a(\epsilon p){\rm
div}\hat{\mathbf{u}}_{\alpha},\hat{\mathbf{u}}_{\alpha})dx$
$\displaystyle\equiv$ $\displaystyle\,d_{1}+d_{2}+d_{3}.$ (3.15)
Thanks to the assumption that $\bar{\mu}>0$ and $2\bar{\mu}+3\bar{\lambda}>0$,
there exists a positive constant $\xi_{1}$, such that
$\displaystyle d_{1}$
$\displaystyle\geq{\underline{a}}\xi\int|\nabla\hat{\mathbf{u}}_{\alpha}|^{2}dx,$
(3.16)
while Cauchy-Schwarz’s inequality implies
$\displaystyle|d_{2}|+|d_{3}|\leq C(\mathcal{Q})\mathcal{S}.$ (3.17)
Similarly, we can obtain
$\displaystyle-\langle\bar{\kappa}a(\epsilon p){\rm
div}(b(\theta)\nabla\hat{\theta}_{\alpha}),\hat{\theta}_{\alpha}\rangle\geq\bar{\kappa}{\underline{a}}\,{\underline{b}}\|\nabla\hat{\theta}_{\alpha}\|^{2}_{L^{2}}-C(\mathcal{Q})\mathcal{S}.$
(3.18)
Easily, one has
$\displaystyle|\langle\epsilon a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm
curl\,}\hat{\mathbf{H}}_{\alpha}+\Psi({\mathbf{u}}):\nabla\hat{\mathbf{u}}_{\alpha}],\hat{\theta}_{\alpha}\rangle|\leq
C(\mathcal{Q})(1+\mathcal{S}).$
It remains to estimate $\langle g_{1},\hat{p}_{\alpha}\rangle$, $\langle
g_{2},\hat{u}_{\alpha}\rangle$ and $\langle
g_{3},\hat{\theta}_{\alpha}\rangle$ in (3.14). First, an application of
Hölder’s inequality gives
$\displaystyle|\langle\hat{p}_{\alpha},g_{1}\rangle|\leq
C_{0}\|\hat{p}_{\alpha}\|_{L^{2}}\|g_{1}\|_{L^{2}},$
where $\|g_{1}\|_{L^{2}}$ can be bounded, by using (2.6), as follows
$\displaystyle\|g_{1}\|_{L^{2}}=$
$\displaystyle\|[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla(\epsilon
p-\theta)\|_{L^{2}}$ $\displaystyle\leq$ $\displaystyle
C_{0}(\|{\mathbf{u}}\|_{W^{1,\infty}}\|\nabla(\epsilon
p-\theta)\|_{H^{s}}+\|{\mathbf{u}}\|_{H^{s+1}}\|\nabla(\epsilon
p-\theta)\|_{L^{\infty}}).$
It follows from the definition of $\mathcal{Q}$ and Sobolev’s inequalities
that
$\displaystyle\|\nabla(\epsilon
p,\theta)\|_{H^{s}}\leq\mathcal{Q},\;\;\;\;\;\|\nabla(\epsilon
p,\theta)\|_{L^{\infty}}\leq\mathcal{Q}.$
Therefore, we obtain $\|g_{1}\|_{L^{2}}\leq C(\mathcal{Q})(1+\mathcal{S})$,
and
$\displaystyle|\langle p_{\alpha},g_{1}\rangle|\leq
C(\mathcal{Q})(1+\mathcal{S}).$ (3.19)
Next, we turn to the term $|\langle{\mathbf{u}}_{\alpha},g_{2}\rangle|$. Due
to the equation (1.12), one has
$\displaystyle-[\partial^{\alpha},b(-\theta)]\partial_{t}(\epsilon{\mathbf{u}})=$
$\displaystyle[\partial^{\alpha},b(-\theta)]\big{(}({\mathbf{u}}\cdot\nabla)(\epsilon{\mathbf{u}})\big{)}+[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)\nabla
p\big{)}$
$\displaystyle-[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)a(\epsilon
p)({\rm curl\,}\mathbf{H}\times(\epsilon\mathbf{H}))\big{)}$
$\displaystyle-[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)a(\epsilon p){\rm
div}\Psi(\epsilon{\mathbf{u}})\big{)}.$ (3.20)
The inequality (2.6) implies that
$\displaystyle\left|\left\langle{\mathbf{u}}_{\alpha},[\partial^{\alpha},b(-\theta)]\big{(}({\mathbf{u}}\cdot\nabla)(\epsilon{\mathbf{u}})\big{)}\right\rangle\right|$
$\displaystyle\leq
C_{0}\|{\mathbf{u}}_{\alpha}\|_{L^{2}}\|[\partial^{\alpha},b(-\theta)]\big{(}{\mathbf{u}}\cdot\nabla(\epsilon{\mathbf{u}})\big{)}\|_{L^{2}}$
$\displaystyle\leq
C(\mathcal{Q})\big{(}\|b(-\theta)\|_{W^{1,\infty}}\|{\mathbf{u}}\cdot\nabla(\epsilon{\mathbf{u}})\|_{H^{s}}+\|b(-\theta)\|_{H^{s+1}}\|{\mathbf{u}}\cdot\nabla(\epsilon{\mathbf{u}})\|_{L^{\infty}}$
$\displaystyle\leq C(\mathcal{Q}),$
and
$\displaystyle\left|\left\langle{\mathbf{u}}_{\alpha},[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)\nabla
p\big{)}\right\rangle\right|$ $\displaystyle\leq
C_{0}\|{\mathbf{u}}_{\alpha}\|_{L^{2}}\|[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)\nabla
p\big{)}\|_{L^{2}}$ $\displaystyle\leq
C(\mathcal{Q})\big{(}\|b(-\theta)\|_{W^{1,\infty}}\|b(\theta)\nabla
p\|_{H^{s}}+\|b(-\theta)\|_{H^{s+1}}\|b(\theta)\nabla p\|_{L^{\infty}}$
$\displaystyle\leq C(\mathcal{Q})(1+\mathcal{S}).$
The third term on the right-hand side of (3.2) can be treated in a similar
manner, and we obtain
$\displaystyle\left|\left\langle{\mathbf{u}}_{\alpha},[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)a(\epsilon
p)({\rm
curl\,}\mathbf{H}\times(\epsilon\mathbf{H}))\big{)}\right\rangle\right|\leq
C(\mathcal{Q})(1+\mathcal{S}).$
To bound the last term on the right-hand side of (3.2), we use (2.6) to deduce
that
$\displaystyle\left\langle{\mathbf{u}}_{\alpha},[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)a(\epsilon
p){\rm div}\Psi(\epsilon{\mathbf{u}})\big{)}\right\rangle$ $\displaystyle\leq
C_{0}\|{\mathbf{u}}_{\alpha}\|_{L^{2}}\|[\partial^{\alpha},b(-\theta)]\big{(}b(\theta)a(\epsilon
p){\rm div}\Psi(\epsilon{\mathbf{u}})\big{)}\|_{L^{2}}$ $\displaystyle\leq
C(\mathcal{Q})(\|b(-\theta)\|_{W^{1,\infty}}\|b(\theta)a(\epsilon p){\rm
div}\Psi(\epsilon{\mathbf{u}})\|_{H^{s}}$
$\displaystyle\quad+\|b(-\theta)\|_{H^{s+1}}\|b(\theta)a(\epsilon p){\rm
div}\Psi(\epsilon{\mathbf{u}})\|_{L^{\infty}})$ $\displaystyle\leq
C(\mathcal{Q})(1+\mathcal{S}).$
Hence, it holds that
$\displaystyle\left|\left\langle{\mathbf{u}}_{\alpha},g_{2}\right\rangle\right|\leq
C(\mathcal{Q})(1+\mathcal{S}).$ (3.21)
Since $g_{3}$ is similar to $g_{1}$ in structure, we easily get
$\displaystyle\big{|}\big{\langle}\hat{\theta}_{\alpha},g_{3}\big{\rangle}\big{|}\leq
C(\mathcal{Q})(1+\mathcal{S}).$ (3.22)
Therefore, it follows from (3.19), (3.21)–(3.22), the positivity of
$b(-\theta)$, and the definition of $\mathcal{O}$, $\mathcal{O}_{0}$,
$\mathcal{Q}$ and $\mathcal{S}$, that there exists a constant $l_{1}>0$, such
that for $t\in[0,T]$ and $T=\min\\{T_{1},1\\}$,
$\displaystyle\sup_{\tau\in[0,t]}\|(\hat{p}_{\alpha},\hat{\mathbf{u}}_{\alpha},\hat{\theta}_{\alpha})(\tau)\|^{2}_{H^{s+1}}+l_{1}\int^{t}_{0}\|\nabla(\hat{\mathbf{u}}_{\alpha},\hat{\theta}_{\alpha})\|^{2}_{H^{s+1}}(\tau)d\tau$
$\displaystyle\leq
C(\mathcal{O}_{0})+C(\mathcal{O}(t))t+C(\mathcal{O}(t))\int^{t}_{0}\mathcal{S}(\tau)d\tau$
$\displaystyle\leq C(\mathcal{O}_{0})+C(\mathcal{O}(t))\sqrt{t}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T)\\}.$
Summing up the above estimates for all $\alpha$ with $0\leq|\alpha|\leq s+1$,
we obtain the desired inequality (3.3). ∎
In a way similar to the proof of Lemma 3.3, we can show that
###### Lemma 3.4.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be a solution to
(1.11)–(1.14), (1.19) on $[0,T_{1}]$. Then there exist a constant $l_{2}>0$
and an increasing function $C(\cdot)$ such that, for any $\epsilon\in(0,1]$
and $t\in[0,T]$, $T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\|(\epsilon^{2}q,\epsilon^{2}{\mathbf{u}},\epsilon(\theta-\bar{\theta})(\tau)\|_{H^{s+2}}+l_{2}\Big{\\{}\int^{t}_{0}\|\nabla(\epsilon^{2}{\mathbf{u}},\epsilon\theta)\|^{2}_{H^{s+2}}(\tau)d\tau\Big{\\}}^{1/2}$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{(\sqrt{T})C(\mathcal{O}(T))\\}.$
Recalling Lemma 2.2 and the definition of $\mathcal{Q}$ and $\mathcal{S}$, we
find that
$\displaystyle\|\partial_{t}(\epsilon
p,\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta)\|_{H^{s-1}}\leq
C(\mathcal{Q}),$ (3.23) $\displaystyle\|\partial_{t}(\epsilon
p,\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta)\|_{H^{s}}\leq
C(\mathcal{Q})(1+\mathcal{S}),$ (3.24)
$\displaystyle\epsilon\|\partial_{t}(\epsilon
p,\epsilon{\mathbf{u}},\epsilon\mathbf{H},\theta)\|_{H^{s}}\leq
C(\mathcal{Q}).$ (3.25)
Moreover, it follows easily from Lemmas 3.1–3.4 and the equation (1.14) that
for some constant $l_{3}>0$,
$\displaystyle\sup_{\tau\in[0,t]}\|\epsilon\partial_{t}\theta\|_{H^{s}}^{2}+l_{3}\int_{0}^{t}\|\nabla((\epsilon\partial_{t})\theta\|_{H^{s}}^{2}(\tau)d\tau\leq
C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T)\\}.$ (3.26)
### 3.3. $H^{s-1}$-estimates on $({\rm div}{\mathbf{u}},\nabla p)$
To establish the estimates for $p$ and the acoustic part of ${\mathbf{u}}$, we
first control the term $(\epsilon\partial_{t})(p,{\mathbf{u}})$. To this end,
we start with a $L^{2}$-estimate for the linearized system. For a given state
$(p_{0},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0})$, consider the following
linearized system of (1.11)–(1.14):
$\displaystyle\partial_{t}p+({\mathbf{u}}_{0}\cdot\nabla)p+\frac{1}{\epsilon}{\rm
div}(2{\mathbf{u}}-\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\theta)=\epsilon a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm curl\,}\mathbf{H}]$
$\displaystyle\quad\quad\qquad\quad\quad\ \ +\epsilon a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla{\mathbf{u}}+\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla p_{0}\cdot\nabla\theta+f_{1},$ (3.27) $\displaystyle
b(-\theta_{0})[\partial_{t}{\mathbf{u}}+({\mathbf{u}}_{0}\cdot\nabla){\mathbf{u}}]+\frac{\nabla
p}{\epsilon}=a(\epsilon p_{0})[({\rm
curl\,}\mathbf{H}_{0})\times\mathbf{H}+{\rm div}\Psi({\mathbf{u}})]+f_{2},$
(3.28) $\displaystyle\partial_{t}\mathbf{H}-{\rm
curl\,}({\mathbf{u}}_{0}\times\mathbf{H})-\bar{\nu}\Delta\mathbf{H}=f_{3},\quad{\rm
div}\mathbf{H}=0,$ (3.29)
$\displaystyle\partial_{t}\theta+({\mathbf{u}}_{0}\cdot\nabla)\theta+{\rm
div}{\mathbf{u}}=\epsilon^{2}a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm
curl\,}\mathbf{H}+\Psi({\mathbf{u}}_{0}):\nabla{\mathbf{u}}]$
$\displaystyle\quad\quad\qquad\quad\quad\ \ +\bar{\kappa}a(\epsilon p_{0}){\rm
div}(b(\theta_{0})\nabla\theta)+f_{4},$ (3.30)
where we have added the source terms $f_{i}$ ($1\leq i\leq 4$) on the right-
hands sides of (3.27)–(3.30) for latter use, and used the following notations:
$a(\epsilon p_{0}):=e^{-\epsilon p_{0}},\quad b(\theta_{0}):=e^{\theta_{0}}.$
The system (3.27)–(3.30) is supplemented with initial data
$\displaystyle(p,{\mathbf{u}},\mathbf{H},\theta)|_{t=0}=(p_{\rm
in}(x),{\mathbf{u}}_{\rm in}(x),\mathbf{H}_{\rm in}(x),\theta_{\rm
in}(x)),\quad x\in\mathbb{R}^{3}.$ (3.31)
###### Lemma 3.5.
Let $(p,{\mathbf{u}},\mathbf{H},\theta)$ be a solution to the Cauchy problem
(3.27)–(3.31) on $[0,\hat{T}]$. Then there exist a constant $l_{4}>0$ and an
increasing function $C(\cdot)$ such that, for any $\epsilon\in(0,1]$ and
$t\in[0,T]$, $T=\min\\{\hat{T},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\|(p,{\mathbf{u}},\mathbf{H})(\tau)\|^{2}_{L^{2}}+l_{4}\int^{t}_{0}\|\nabla({\mathbf{u}},\mathbf{H})\|_{L^{2}}^{2}(\tau)d\tau$
$\displaystyle\leq
e^{TC(R_{0})}\|(p,{\mathbf{u}},\mathbf{H})(0)\|^{2}_{L^{2}}+C(R_{0})e^{TC(R_{0})}\sup_{\tau\in[0,T]}\|\nabla\theta(\tau)\|^{2}_{L^{2}}$
$\displaystyle\quad+C(R_{0})\int_{0}^{T}\|\nabla(\epsilon{\mathbf{u}},\epsilon\mathbf{H})\|_{L^{2}}^{2}d\tau+C(R_{0})\int^{T}_{0}\|\nabla\theta\|^{2}_{H^{1}}(\tau)d\tau$
$\displaystyle\quad+C(R_{0})\int^{T}_{0}\left\\{\|f_{1}\|^{2}_{L^{2}}+\|f_{2}\|^{2}_{L^{2}}+\|f_{3}\|^{2}_{L^{2}}+\|\nabla
f_{4}\|_{L^{2}}^{2}\right\\}(\tau)d\tau,$ (3.32)
where
$\displaystyle
R_{0}=\sup_{\tau\in[0,T]}\\{\|\partial_{t}\theta_{0}(\tau)\|_{L^{\infty}},\|(p_{0},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0})(\tau)\|_{W^{1,\infty}}\\}.$
(3.33)
###### Proof.
Set
$\displaystyle(\tilde{p},\tilde{\mathbf{u}},\tilde{\mathbf{H}},\tilde{\theta})=(p,2{\mathbf{u}}-\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\theta,\mathbf{H},\theta).$
Then $\tilde{p}$ and $\tilde{\mathbf{H}}$ satisfy
$\displaystyle\partial_{t}\tilde{p}$
$\displaystyle+({\mathbf{u}}_{0}\cdot\nabla)\tilde{p}+\frac{1}{\epsilon}{\rm
div}\tilde{\mathbf{u}}=\epsilon a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm
curl\,}\tilde{\mathbf{H}}]+\frac{\epsilon}{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}}$
$\displaystyle+\frac{\epsilon}{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})+\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla p_{0}\cdot\nabla\tilde{\theta}+f_{1}$ (3.34)
and
$\displaystyle\partial_{t}\tilde{\mathbf{H}}-{\rm
curl\,}({\mathbf{u}}_{0}\times\tilde{\mathbf{H}})-\bar{\nu}\Delta\tilde{\mathbf{H}}=f_{3},\quad{\rm
div}\tilde{\mathbf{H}}=0,$ (3.35)
respectively. One can derive the equation for $\tilde{\mathbf{u}}$ by applying
the operator $\nabla$ to (3.30) to obtain
$\displaystyle\partial_{t}\nabla\tilde{\theta}+({\mathbf{u}}_{0}\cdot\nabla)\nabla\tilde{\theta}+\frac{1}{2}\nabla{\rm
div}\tilde{\mathbf{u}}+\frac{1}{2}\nabla{\rm div}(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})$
$\displaystyle=\nabla\left\\{\epsilon^{2}a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm
curl\,}\tilde{\mathbf{H}}]\right\\}+\frac{1}{2}\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}}\right\\}$
$\displaystyle\quad+\frac{1}{2}\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})\right\\}$
$\displaystyle\quad+\nabla\left\\{\bar{\kappa}a(\epsilon p_{0}){\rm
div}(b(\theta_{0})\nabla\tilde{\theta})\right\\}+[\nabla,{\mathbf{u}}_{0}]\cdot\nabla\tilde{\theta}+\nabla
f_{4}.$ (3.36)
If we multiply (3.3) with $\frac{1}{2}\bar{\kappa}a(\epsilon p_{0})$, we get
$\displaystyle\frac{1}{2}b(-\theta_{0})\left\\{\partial_{t}(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})+({\mathbf{u}}_{0}\cdot\nabla)[\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta}]\right\\}$
$\displaystyle=\frac{\bar{\kappa}}{2}b(-\theta_{0})\partial_{t}\\{a(\epsilon
p_{0})b(\theta_{0})\\}\nabla\tilde{\theta}+\frac{\bar{\kappa}}{2}b(-\theta_{0})\left\\{{\mathbf{u}}_{0}\cdot\nabla[a(\epsilon
p_{0})b(\theta_{0})]\nabla\tilde{\theta}\right\\}$
$\displaystyle\quad-\frac{1}{2}\bar{\kappa}a(\epsilon p_{0})\nabla{\rm
div}\tilde{\mathbf{u}}-\frac{1}{2}\bar{\kappa}a(\epsilon p_{0})\nabla{\rm
div}(\bar{\kappa}a(\epsilon p_{0})b(\theta_{0})\nabla\tilde{\theta})$
$\displaystyle\quad+\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm
curl\,}\tilde{\mathbf{H}}]\right\\}+\frac{1}{4}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}}\right\\}$
$\displaystyle\quad+\frac{1}{4}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})\right\\}$
$\displaystyle\quad+\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\bar{\kappa}a(\epsilon p_{0}){\rm
div}(b(\theta_{0})\nabla\tilde{\theta})\right\\}$
$\displaystyle\quad+\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})[\nabla,{\mathbf{u}}_{0}]\cdot\nabla\tilde{\theta}+\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})\nabla f_{4}.$ (3.37)
Subtracting (3.3) from (3.28) yields
$\displaystyle\frac{1}{2}$ $\displaystyle
b(-\theta_{0})[\partial_{t}\tilde{\mathbf{u}}+{\mathbf{u}}_{0}\cdot\nabla\tilde{\mathbf{u}}]+\frac{\nabla\tilde{p}}{\epsilon}$
$\displaystyle=$
$\displaystyle-\frac{\bar{\kappa}}{2}b(-\theta_{0})\partial_{t}\\{a(\epsilon
p_{0})b(\theta_{0})\\}\nabla\tilde{\theta}-\frac{\bar{\kappa}}{2}b(-\theta_{0})\left\\{{\mathbf{u}}_{0}\cdot\nabla[a(\epsilon
p_{0})b(\theta_{0})]\nabla\tilde{\theta}\right\\}$
$\displaystyle+\frac{1}{4}\bar{\kappa}a(\epsilon p_{0})\nabla{\rm
div}\tilde{\mathbf{u}}+\frac{1}{4}\bar{\kappa}a(\epsilon p_{0})\nabla{\rm
div}(\bar{\kappa}a(\epsilon p_{0})b(\theta_{0})\nabla\tilde{\theta})$
$\displaystyle-\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm curl\,}\tilde{\mathbf{H}}]\right\\}$
$\displaystyle-\frac{1}{4}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}}\right\\}$
$\displaystyle-\frac{1}{4}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\epsilon^{2}a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})\right\\}$
$\displaystyle-\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})\nabla\left\\{\bar{\kappa}a(\epsilon p_{0}){\rm
div}(b(\theta_{0})\nabla\tilde{\theta})\right\\}+a(\epsilon
p_{0})[\bar{\nu}\,{\rm curl\,}\mathbf{H}_{0}:{\rm curl\,}\tilde{\mathbf{H}}]$
$\displaystyle+\frac{1}{2}a(\epsilon
p_{0})[\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}}]+\frac{1}{2}a(\epsilon
p_{0})[\Psi({\mathbf{u}}_{0}):\nabla(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})]$
$\displaystyle-\frac{1}{2}\bar{\kappa}a(\epsilon
p_{0})[\nabla,{\mathbf{u}}_{0}]\cdot\nabla\tilde{\theta}+a(\epsilon
p_{0})[({\rm curl\,}\mathbf{H}_{0})\times\tilde{\mathbf{H}}]$
$\displaystyle+\frac{1}{2}a(\epsilon p_{0}){\rm
div}\Psi(\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla\tilde{\theta})+\frac{1}{2}a(\epsilon p_{0}){\rm
div}\Psi(\tilde{\mathbf{u}})-\frac{1}{2}\bar{\kappa}a(\epsilon p_{0})\nabla
f_{4}+f_{2}$ $\displaystyle:=$
$\displaystyle\sum^{14}_{i=1}h_{i}+\frac{1}{2}a(\epsilon p_{0}){\rm
div}\Psi(\tilde{\mathbf{u}})-\frac{1}{2}\bar{\kappa}a(\epsilon p_{0})\nabla
f_{4}+f_{2}.$ (3.38)
Multiplying (3.34) by $\tilde{p}$, (3.35) by $\tilde{\mathbf{H}}$, and (3.3)
by $\tilde{\mathbf{u}}$ in $L^{2}(\mathbb{R}^{3})$ respectively, and summing
up the resulting equations, we deduce that
$\displaystyle\frac{d}{dt}$
$\displaystyle\Big{\\{}\frac{1}{2}\langle\tilde{p},\tilde{p}\rangle+\frac{1}{4}\langle
b(-\theta_{0})\tilde{\mathbf{u}},\tilde{\mathbf{u}}\rangle+\frac{1}{2}\langle\tilde{\mathbf{H}},\mathbf{H}\rangle\Big{\\}}+\bar{\nu}\|\nabla\tilde{\mathbf{H}}\|^{2}_{L^{2}}$
$\displaystyle=$
$\displaystyle-\langle({\mathbf{u}}_{0}\cdot\nabla)\tilde{p},\tilde{p}\rangle+\frac{1}{4}\langle\partial_{t}b(-\theta_{0})\tilde{\mathbf{u}},\tilde{\mathbf{u}}\rangle-\frac{1}{2}\langle
b(-\theta_{0})({\mathbf{u}}_{0}\cdot\nabla)\tilde{\mathbf{u}},\tilde{\mathbf{u}}\rangle$
$\displaystyle+\langle\epsilon a(\epsilon p_{0})[\bar{\nu}\,{\rm
curl\,}\mathbf{H}_{0}:{\rm
curl\,}\tilde{\mathbf{H}}],\tilde{p}\rangle+\frac{\epsilon}{2}\langle
a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}},\tilde{p}\rangle$
$\displaystyle+\frac{\epsilon}{2}\langle a^{2}(\epsilon
p_{0})b(\theta_{0})\Psi({\mathbf{u}}_{0}):\nabla(\nabla\tilde{\theta}),\tilde{p}\rangle+\langle\bar{\kappa}a(\epsilon
p_{0})b(\theta_{0})\nabla p_{0}\cdot\nabla\tilde{\theta},\tilde{p}\rangle$
$\displaystyle+\sum^{14}_{i=1}\left\langle
h_{i},\tilde{\mathbf{u}}\right\rangle+\frac{1}{2}\langle a(\epsilon p_{0}){\rm
div}\Psi(\tilde{\mathbf{u}}),\tilde{\mathbf{u}}\rangle$
$\displaystyle-\frac{1}{2}\left\langle\bar{\kappa}a(\epsilon p_{0})\nabla
f_{4},\tilde{\mathbf{u}}\right\rangle+\left\langle
f_{2},\tilde{\mathbf{u}}\right\rangle+\langle
f_{3},\tilde{\mathbf{H}}\rangle+\langle f_{1},\tilde{p}\rangle,$ (3.39)
where the singular terms have been canceled out.
Now, the terms on the right-hand side of (3.3) can be estimated as follows.
First, it follows from the regularity of
$(p_{0},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0})$, a partial integration
and Cauchy-Schwarz’s inequality that
$\displaystyle\frac{1}{4}\left|\langle\partial_{t}b(-\theta_{0})\tilde{\mathbf{u}},\tilde{\mathbf{u}}\rangle\right|\leq\frac{1}{4}\|\partial_{t}b(-\theta_{0})\|_{L^{\infty}}\|\tilde{\mathbf{u}}\|^{2}_{L^{2}}\leq
C(R_{0})\|\tilde{\mathbf{u}}\|^{2}_{L^{2}},$
$\displaystyle|\langle({\mathbf{u}}_{0}\cdot\nabla)\tilde{p},\tilde{p}\rangle|=\frac{1}{2}\left|\int({\rm
div}{\mathbf{u}}_{0})|\tilde{p}|^{2}dx\right|\leq
C(R_{0})\|\tilde{p}\|_{L^{2}}^{2},$ $\displaystyle\frac{1}{2}|\langle
b(-\theta_{0})({\mathbf{u}}_{0}\cdot\nabla)\tilde{\mathbf{u}},\tilde{\mathbf{u}}\rangle|\leq
C(R_{0})\|\tilde{\mathbf{u}}\|_{L^{2}}^{2},$ $\displaystyle|\langle\epsilon
a(\epsilon p_{0})[\bar{\nu}\,{\rm curl\,}\mathbf{H}_{0}:{\rm
curl\,}\tilde{\mathbf{H}}],\tilde{p}\rangle|\leq
C(R_{0})(\|\epsilon\nabla\tilde{\mathbf{H}}\|^{2}_{L^{2}}+\|\tilde{p}\|^{2}_{L^{2}}),$
$\displaystyle\frac{\epsilon}{2}|\langle a(\epsilon
p_{0})\Psi({\mathbf{u}}_{0}):\nabla\tilde{\mathbf{u}},\tilde{p}\rangle|\leq
C(R_{0})(\|\epsilon\nabla\tilde{\mathbf{u}}\|^{2}_{L^{2}}+\|\tilde{p}\|^{2}_{L^{2}}),$
$\displaystyle\frac{\epsilon}{2}|\langle a^{2}(\epsilon
p_{0})b(\theta_{0})\Psi({\mathbf{u}}_{0}):\nabla(\nabla\tilde{\theta}),\tilde{p}\rangle|\leq
C(R_{0})\|\tilde{p}\|^{2}_{L^{2}}$
$\displaystyle\qquad\qquad\qquad\qquad\qquad\qquad\qquad+G_{1}(\epsilon
p_{0},\theta_{0})\sum_{|\alpha|=2}\|\partial^{\alpha}(\epsilon\tilde{\theta})\|^{2}_{L^{2}},$
$\displaystyle|\langle\bar{\kappa}a(\epsilon p_{0})b(\theta_{0})\nabla
p_{0}\cdot\nabla\tilde{\theta},\tilde{p}\rangle|\leq
C(R_{0})(\|\nabla\tilde{\theta}\|^{2}_{L^{2}}+\|\tilde{p}\|^{2}_{L^{2}}),$
where $G_{1}(\cdot,\cdot)$ is a smooth function. Similarly, one can bound the
terms involving $h_{i}$ in (3.3) as follows.
$\displaystyle\sum^{14}_{i=1}|\langle h_{i},\tilde{\mathbf{u}}\rangle|\leq$
$\displaystyle\frac{\bar{\nu}}{8}\|\nabla\tilde{\mathbf{H}}\|^{2}_{L_{2}}+\frac{\underline{a}\bar{\mu}}{8}\|\nabla\tilde{\mathbf{u}}\|^{2}_{L_{2}}+\frac{\underline{a}\bar{\nu}}{8}\|{\rm
div}\tilde{\mathbf{u}}\|^{2}_{L_{2}}$
$\displaystyle+C(R_{0})\|\tilde{\mathbf{u}}\|^{2}_{L_{2}}+C(R_{0})\|\nabla(\epsilon\tilde{\mathbf{u}},\tilde{\theta})\|^{2}_{L_{2}}+G_{2}(\epsilon
p_{0},\theta_{0})\|\Delta\tilde{\theta}\|^{2}_{L_{2}},$
where $G_{2}(\cdot,\cdot)$ is a smooth function.
For the dissipative term $\frac{1}{2}\langle a(\epsilon p_{0}){\rm
div}\Psi(\tilde{\mathbf{u}}),\tilde{\mathbf{u}}\rangle$, we can employ
arguments similar to those used in the estimate of the slow motion in
(3.2)–(3.17) to obtain that
$\displaystyle-\frac{1}{2}\langle a(\epsilon p_{0}){\rm
div}\Psi(\hat{\mathbf{u}}),\hat{\mathbf{u}}\rangle\geq\frac{\underline{a}\bar{\mu}}{4}(\|\nabla\hat{\mathbf{u}}\|^{2}_{L^{2}}+\|{\rm
div}\hat{\mathbf{u}}\|^{2}_{L^{2}})-C(R_{0})\|\hat{\mathbf{u}}\|^{2}_{L^{2}}.$
Finally, putting all estimates above into (3.3) and applying Cauchy-Schwarz’s
and Gronwall’s inequalities, we get (3.5). ∎
In the next lemma we utilize Lemma 3.5 to control
$\big{(}(\epsilon\partial_{t})p,(\epsilon\partial_{t}){\mathbf{u}},(\epsilon\partial_{t})\mathbf{H}\big{)}$.
###### Lemma 3.6.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be the solution to the
Cauchy problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. Set
$(p_{\beta},{\mathbf{u}}_{\beta},\mathbf{H}_{\beta},\theta_{\beta}):=\partial^{\beta}\big{(}(\epsilon\partial_{t})p,(\epsilon\partial_{t}){\mathbf{u}},(\epsilon\partial_{t})\mathbf{H},(\epsilon\partial_{t})\theta\big{)},$
where $1\leq|\beta|\leq s-1$. Then there exist a constant $l_{5}>0$ and an
increasing function $C(\cdot)$ such that, for any $\epsilon\in(0,1]$ and
$t\in[0,T]$, $T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\|(p_{\beta},{\mathbf{u}}_{\beta},\mathbf{H}_{\beta})(\tau)\|^{2}_{L^{2}}+$
$\displaystyle
l_{5}\int^{t}_{0}\|\nabla({\mathbf{u}}_{\beta},\mathbf{H}_{\beta})\|_{L^{2}}^{2}(\tau)d\tau$
$\displaystyle\leq C(\mathcal{O}_{0})\exp\\{(\sqrt{T})C(\mathcal{O}(T))\\}.$
(3.40)
###### Proof.
An application of the operator $\partial^{\beta}(\epsilon\partial_{t})$ to the
system (1.11)–(1.14) leads to
$\displaystyle\partial_{t}p_{\beta}+({\mathbf{u}}\cdot\nabla)p_{\beta}+\frac{1}{\epsilon}{\rm
div}(2{\mathbf{u}}_{\beta}-\bar{\kappa}a(\epsilon
p)b(\theta)\nabla\theta_{\beta})=\epsilon a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm curl\,}\mathbf{H}_{\beta}]$
$\displaystyle\quad\quad\qquad+\epsilon a(\epsilon
p)\Psi({\mathbf{u}}):\nabla{\mathbf{u}}_{\beta}+\bar{\kappa}a(\epsilon
p)b(\theta)\nabla p\cdot\nabla\theta_{\beta}+\tilde{g}_{1},$ (3.41)
$\displaystyle
b(-\theta)[\partial_{t}{\mathbf{u}}_{\beta}+({\mathbf{u}}\cdot\nabla){\mathbf{u}}_{\beta}]+\frac{\nabla
p_{\beta}}{\epsilon}=a(\epsilon p)[({\rm
curl\,}\mathbf{H})\times\mathbf{H}_{\beta}+{\rm
div}\Psi({\mathbf{u}}_{\beta})]+\tilde{g}_{2},$ (3.42)
$\displaystyle\partial_{t}\mathbf{H}_{\beta}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H}_{\beta})-\bar{\nu}\Delta\mathbf{H}_{\beta}=\tilde{g}_{3},\quad{\rm
div}\mathbf{H}_{\beta}=0,$ (3.43)
$\displaystyle\partial_{t}\theta_{\beta}+({\mathbf{u}}\cdot\nabla)\theta_{\beta}+{\rm
div}{\mathbf{u}}_{\beta}=\epsilon^{2}a(\epsilon p)[\bar{\nu}\,{\rm
curl\,}\mathbf{H}:{\rm curl\,}\mathbf{H}_{\beta}]$
$\displaystyle\quad\quad\qquad+\epsilon^{2}a(\epsilon
p)\Psi({\mathbf{u}}):\nabla{\mathbf{u}}_{\beta}+\bar{\kappa}a(\epsilon p){\rm
div}(b(\theta)\nabla\theta_{\beta})+\tilde{g}_{4},$ (3.44)
where
$\displaystyle\tilde{g}_{1}:=$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),{\mathbf{u}}]\cdot\nabla
p+\frac{1}{\epsilon}[\partial^{\beta}(\epsilon\partial_{t}),(\bar{\kappa}a(\epsilon
p)b(\theta))]\Delta\theta$
$\displaystyle+\frac{1}{\epsilon}[\partial^{\beta}(\epsilon\partial_{t}),\nabla(\bar{\kappa}a(\epsilon
p)b(\theta))]\cdot\nabla\theta+\epsilon\bar{\nu}\,[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon
p){\rm curl\,}\mathbf{H}]:{\rm curl\,}\mathbf{H}$
$\displaystyle+\epsilon[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon
p)\Psi({\mathbf{u}})]:\nabla{\mathbf{u}}+[\partial^{\beta}(\epsilon\partial_{t}),\bar{\kappa}a(\epsilon
p)b(\theta)\nabla p]\cdot\nabla\theta,$ $\displaystyle\tilde{g}_{2}:=$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\partial_{t}{\mathbf{u}}-[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta){\mathbf{u}}]\cdot\nabla{\mathbf{u}},$
$\displaystyle+[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon p){\rm
curl\,}\mathbf{H}]\times\mathbf{H}+[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon
p)]{\rm div}\Psi({\mathbf{u}}),$ $\displaystyle\tilde{g}_{3}:=$
$\displaystyle\,\partial^{\beta}(\epsilon\partial_{t})\big{(}{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})\big{)}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H}_{\beta}),$ $\displaystyle\tilde{g}_{4}:=$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),{\mathbf{u}}]\cdot\nabla\theta+\epsilon^{2}\bar{\nu}\,[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon
p){\rm curl\,}\mathbf{H}]:{\rm curl\,}\mathbf{H}$
$\displaystyle+\epsilon^{2}[\partial^{\beta}(\epsilon\partial_{t}),a(\epsilon
p)\Psi({\mathbf{u}})]:\nabla{\mathbf{u}}$
$\displaystyle+\bar{\kappa}\partial^{\beta}(\epsilon\partial_{t})\big{(}a(\epsilon
p){\rm div}(b(\theta)\nabla\theta_{\alpha})\big{)}-\bar{\kappa}a(\epsilon){\rm
div}(b(\theta)\nabla\theta_{\alpha}).$
It follows from the linear estimate (3.5) that for some $l_{4}>0$,
$\displaystyle\sup_{\tau\in[0,t]}\|(p_{\beta},{\mathbf{u}}_{\beta},\mathbf{H}_{\beta})(\tau)\|^{2}_{L^{2}}+l_{4}\int^{t}_{0}\|\nabla({\mathbf{u}}_{\beta},\mathbf{H}_{\beta})\|_{L^{2}}^{2}(\tau)d\tau$
$\displaystyle\leq
e^{TC(R)}\|(p_{\beta},\tilde{\mathbf{u}}_{\beta},\mathbf{H}_{\beta})(0)\|^{2}_{L^{2}}+C(R)e^{TC(R)}\sup_{\tau\in[0,T]}\|\nabla\theta_{\beta}(\tau)\|^{2}_{L^{2}}$
$\displaystyle\quad+TC(R)\sup_{\tau\in[0,T]}\|\nabla(\epsilon{\mathbf{u}}_{\beta},\epsilon\mathbf{H}_{\beta})(\tau)\|^{2}_{L^{2}}+C(R)\int^{T}_{0}\|\nabla\theta_{\beta}\|^{2}_{H^{1}}(\tau)d\tau$
$\displaystyle\quad+C(R)\int^{T}_{0}\left\\{\|\tilde{g}_{1}\|^{2}_{L^{2}}+\|\tilde{g}_{2}\|^{2}_{L^{2}}+\|\tilde{g}_{3}\|^{2}_{L^{2}}+\|\nabla\tilde{g}_{4}\|_{L^{2}}^{2}\right\\}(\tau)d\tau,$
(3.45)
where $R$ is defined as $R_{0}$ in (3.33) with
$(p_{0},{\mathbf{u}}_{0},\mathbf{H}_{0},\theta_{0})$ replaced with
$(p,{\mathbf{u}},\mathbf{H},\theta)$.
It remains to control the terms $\|\tilde{g}_{1}\|_{L^{2}}^{2}$,
$\|\tilde{g}_{2}\|_{L^{2}}^{2}$, $\|\tilde{g}_{3}\|_{L^{2}}^{2}$, and
$\|\nabla\tilde{g}_{4}\|_{L^{2}}^{2}$. The first term of $\tilde{g}_{1}$ can
be bounded as follows.
$\displaystyle\left\|[\partial^{\beta}(\epsilon\partial_{t}),{\mathbf{u}}]\cdot\nabla
p\right\|_{L^{2}}\leq$ $\displaystyle\epsilon
C_{0}(\|{\mathbf{u}}\|_{H^{s-1}}\|(\epsilon\partial_{t})\nabla
p\|_{H^{s-2}}+\|(\epsilon\partial_{t}){\mathbf{u}}\|_{H^{s-1}}\|\nabla
p\|_{H^{s-1}})$ $\displaystyle\leq$ $\displaystyle C(\mathcal{Q}).$
Similarly, the second term of $\tilde{g}_{1}$ admits the following
boundedness:
$\displaystyle\frac{1}{\epsilon}\|[\partial^{\beta}(\epsilon\partial_{t}),(\bar{\kappa}a(\epsilon
p)b(\theta))]\Delta\theta\|$ $\displaystyle\leq C_{0}\big{(}\|a(\epsilon
p)b(\theta)\|_{H^{s-1}}\|\partial_{t}\Delta\theta\|_{H^{s-2}}+\|\partial_{t}(a(\epsilon
p)b(\theta))\|_{H^{s-1}}\|\Delta\theta\|_{H^{s-1}}\big{)}$ $\displaystyle\leq
C(\mathcal{Q})(1+\mathcal{S}).$
The other four terms in $\tilde{g}_{1}$ can be treated similarly and hence can
be bounded from above by $C(\mathcal{Q})(1+\mathcal{S})$.
For the first term of $\tilde{g}_{2}$, one has by the equation (1.12) that
$\displaystyle[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\partial_{t}{\mathbf{u}}=$
$\displaystyle[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\\{({\mathbf{u}}\cdot\nabla){\mathbf{u}}\\}$
$\displaystyle+\frac{1}{\epsilon}[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\\{b^{-1}(-\theta_{0})\nabla
p\\}$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\\{b^{-1}(-\theta)a(\epsilon
p)[({\rm curl\,}\mathbf{H}_{0})\times\mathbf{H}]\\}$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\\{b^{-1}(-\theta)a(\epsilon
p){\rm div}\Psi({\mathbf{u}})\\}.$ (3.46)
Note that the terms on the right-hand side of (3.3) have similar structure as
that of $\tilde{g}_{1}$. Thus, we see that
$\displaystyle\|[\partial^{\beta}(\epsilon\partial_{t}),b(-\theta)]\partial_{t}{\mathbf{u}}\|_{L^{2}}\leq
C(\mathcal{Q})(1+\mathcal{S}).$
Similarly, the other four terms of $\tilde{g}_{2}$ can be bounded from above
by $C(\mathcal{Q})(1+\mathcal{S})$.
Next, by the identity (2.4), one can rewrite $\tilde{g}_{3}$ as
$\displaystyle\tilde{g}_{3}=$
$\displaystyle-[\partial^{\beta}(\epsilon\partial_{t}),{\rm
div}{\mathbf{u}}]\mathbf{H}-[\partial^{\beta}(\epsilon\partial_{t}),{\mathbf{u}}]\cdot\nabla\mathbf{H}+\sum_{i=1}^{3}[\partial^{\beta}(\epsilon\partial_{t}),\nabla{\mathbf{u}}_{i}]\mathbf{H}.$
Following a process similar to that in the estimate of $\tilde{g}_{1}$, one
gets
$\displaystyle\|\tilde{g}_{3}\|_{L^{2}}\leq C(\mathcal{Q})(1+\mathcal{S}).$
And analogously,
$\displaystyle\|\nabla\tilde{g}_{4}\|_{L^{2}}\leq
C(\mathcal{Q})(1+\mathcal{S}).$
We proceed to control the other terms on the right-hand side of (3.3). It
follows from (3.26) that
$C(R)e^{TC(R)}\sup_{\tau\in[0,T]}\|\nabla\theta_{\beta}(\tau)\|^{2}_{L^{2}}\leq
C(\mathcal{O}(T))\exp\\{\sqrt{T}C(\mathcal{O}(T)\\}$
and
$\int^{T}_{0}\|\Delta\theta_{\beta}\|^{2}_{L_{2}}(\tau)d\tau\leq\int^{T}_{0}\|(\epsilon\partial_{t})\theta\|^{2}_{H^{s+1}}(\tau)d\tau\leq
C(\mathcal{O}_{0})\exp\\{\sqrt{T}C(\mathcal{O}(T))\\}.$
Thanks to (3.25), one has
$\displaystyle
TC(R)\sup_{\tau\in[0,T]}\|\nabla(\epsilon{\mathbf{u}}_{\beta},\epsilon\mathbf{H}_{\beta})(\tau)\|^{2}_{L^{2}}$
$\displaystyle\leq
TC(\mathcal{O}(T))\sup_{\tau\in[0,T]}\|(\epsilon\partial_{t})(\epsilon{\mathbf{u}},\epsilon\mathbf{H})(\tau)\|^{2}_{H^{s}}$
$\displaystyle\leq TC(\mathcal{O}(T)).$
Then, the desired inequality (3.6) follows from the above estimates and the
inequality (3.3). ∎
Now we are in a position to estimate the Sobolev norm of $({\rm
div}{\mathbf{u}},\nabla p)$ based on Lemma 3.6.
###### Lemma 3.7.
Let $s\geq 4$ and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be the solution to the
Cauchy problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. Then there exist a
constant $l_{6}>0$ and an increasing function $C(\cdot)$ such that, for any
$\epsilon\in(0,1]$ and $t\in[0,T_{1}]$, $T=\min\\{T_{1},1\\}$, it holds that
$\displaystyle\sup_{\tau\in[0,t]}\left\\{\|p(\tau)\|_{H^{s}}+\|{\rm
div}{\mathbf{u}}(\tau)\|_{H^{s-1}}\right\\}$
$\displaystyle+l_{6}\int^{t}_{0}\left\\{\|\nabla p\|^{2}_{H^{s}}+\|\nabla{\rm
div}{\mathbf{u}}\|^{2}_{H^{s-1}}\right\\}(\tau)d\tau$ $\displaystyle\leq
C(\mathcal{O}_{0})\exp\big{\\{}(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\big{\\}}.$
(3.47)
###### Proof.
Rewrite the equations (1.11) and (1.12) as
$\displaystyle{\rm div}{\mathbf{u}}=$
$\displaystyle-\frac{1}{2}(\epsilon\partial_{t})p-\frac{\epsilon}{2}({\mathbf{u}}\cdot\nabla)p+\frac{1}{2}{\rm
div}(\bar{\kappa}a(\epsilon
p)b(\theta)\nabla\theta)+\frac{\epsilon^{2}\bar{\nu}}{2}a(\epsilon p_{0})|{\rm
curl\,}\mathbf{H}|^{2}$ $\displaystyle+\frac{\epsilon^{2}}{2}a(\epsilon
p)\Psi({\mathbf{u}}):\nabla{\mathbf{u}}+\frac{\epsilon\bar{\kappa}}{2}a(\epsilon
p)b(\theta)\nabla p\cdot\nabla\theta,$ (3.48) $\displaystyle{\nabla p}=$
$\displaystyle-b(-\theta)(\epsilon\partial_{t}){\mathbf{u}}-{\epsilon}b(-\theta)({\mathbf{u}}\cdot\nabla){\mathbf{u}}$
$\displaystyle+\epsilon a(\epsilon p)[({\rm
curl\,}\mathbf{H})\times\mathbf{H}]+\epsilon a(\epsilon p){\rm
div}\Psi({\mathbf{u}}).$ (3.49)
Then,
$\displaystyle\|{\rm div}{\mathbf{u}}\|_{H^{s-1}}\leq$ $\displaystyle
C_{0}\|(\epsilon\partial_{t})p\|_{H^{s-1}}+C_{0}\epsilon\,\|{\mathbf{u}}\|_{H^{s-1}}\|\nabla
p\|_{H^{s-1}}$ $\displaystyle+C_{0}\|{\rm div}(\bar{\kappa}a(\epsilon
p)b(\theta)\nabla\theta)\|_{H^{s-1}}+C_{0}\|a(\epsilon
p_{0})\|_{L^{\infty}}\|\epsilon\,{\rm curl\,}\mathbf{H}\|_{H^{s-1}}^{2}$
$\displaystyle+C_{0}\|a(\epsilon
p)\|_{L^{\infty}}\|\Psi(\epsilon{\mathbf{u}}):(\epsilon\nabla{\mathbf{u}})\|_{H^{s-1}}$
$\displaystyle+C_{0}\|a(\epsilon p)b(\theta)\|_{L^{\infty}}\|(\epsilon\nabla
p)\|_{H^{s-1}}\|\nabla\theta\|_{H^{s-1}}.$ (3.50)
It follows from Lemmas 3.2–3.4 and 3.6, and the inequalities (3.23)–(3.26)
that
$\displaystyle\|(\epsilon\partial_{t})p\|_{H^{s-1}}\leq$ $\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
$\displaystyle\epsilon\,\|{\mathbf{u}}\|_{H^{s-1}}\|\nabla p\|_{H^{s-1}}\leq$
$\displaystyle\epsilon C(\mathcal{O}),$ $\displaystyle\|{\rm
div}(\bar{\kappa}a(\epsilon p)b(\theta)\nabla\theta)\|_{H^{s-1}}\leq$
$\displaystyle
C_{0}\|\Delta\theta\|_{H^{s-1}}+C_{0}\|\nabla\theta\|_{H^{s-1}}$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
$\displaystyle\|\epsilon\,{\rm curl\,}\mathbf{H}\|_{H^{s-1}}\leq$
$\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
$\displaystyle\|\Psi(\epsilon{\mathbf{u}}):(\epsilon\nabla{\mathbf{u}})\|_{H^{s-1}}\leq$
$\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
$\displaystyle\|(\epsilon\nabla p)\|_{H^{s-1}}\|\nabla\theta\|_{H^{s-1}}\leq$
$\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$
These bounds together with (3.50) imply that
$\displaystyle\sup_{\tau\in[0,t]}\|{\rm
div}{\mathbf{u}}\|(\tau)\|_{H^{s-1}}\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$
Similar arguments applying to the equation (3.49) for $\nabla p$ yield
$\displaystyle\sup_{\tau\in[0,t]}\|p\|_{H^{s}}+l_{6}\int^{t}_{0}\|\nabla
p\|^{2}_{H^{s}}(\tau)d\tau\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}$ (3.51)
for some positive constant $l_{6}>0$.
To obtain the desired inequality (3.7), we shall establish the following
estimate
$\displaystyle\int^{T}_{0}\|\nabla{\rm
div}{\mathbf{u}}\|^{2}_{H^{s-1}}(\tau)d\tau\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$ (3.52)
In fact, for any multi-index $\alpha$ satisfying $1\leq|\alpha|\leq s$, one
can apply the operator $\partial^{\alpha}$ to (3.48) and then take the inner
product with $\partial^{\alpha}{\rm div}{\mathbf{u}}$ to obtain
$\displaystyle\int^{T}_{0}\|\partial^{\alpha}{\rm
div}{\mathbf{u}}\|^{2}_{L^{2}}(\tau)d\tau=$
$\displaystyle-\frac{1}{2}\int^{T}_{0}\langle\partial^{\alpha}(\epsilon\partial_{t})p,\partial^{\alpha}{\rm
div}{\mathbf{u}}\rangle(\tau)d\tau$
$\displaystyle+\int^{T}_{0}\langle\Xi,\partial^{\alpha}{\rm
div}{\mathbf{u}}\rangle(\tau)d\tau,$ (3.53)
where
$\displaystyle\Xi:=$
$\displaystyle-\frac{\epsilon}{2}({\mathbf{u}}\cdot\nabla)p+\frac{1}{2}{\rm
div}(\bar{\kappa}a(\epsilon
p)b(\theta)\nabla\theta)+\frac{\epsilon^{2}\bar{\nu}}{2}a(\epsilon p_{0})|{\rm
curl\,}\mathbf{H}|^{2}$ $\displaystyle+\frac{\epsilon^{2}}{2}a(\epsilon
p)\Psi({\mathbf{u}}):\nabla{\mathbf{u}}+\frac{\epsilon\bar{\kappa}}{2}a(\epsilon
p)b(\theta)\nabla p\cdot\nabla\theta.$
It thus follows from (3.6) and similar arguments to those for (3.51), that for
all $1\leq|\alpha|\leq s$,
$\displaystyle\int^{T}_{0}\|\partial^{\alpha}\Xi(\tau)\|^{2}_{L^{2}}d\tau\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
whence,
$\displaystyle\int^{T}_{0}\left|\langle\Xi,\partial^{\alpha}{\rm
div}{\mathbf{u}}\rangle\right|(\tau)d\tau$ $\displaystyle\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}\left\\{\int^{T}_{0}\|\partial^{\alpha}{\rm
div}{\mathbf{u}}\|^{2}_{L^{2}}(\tau)d\tau\right\\}^{1/2}$ $\displaystyle\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}+\frac{1}{4}\int^{T}_{0}\|\partial^{\alpha}{\rm
div}{\mathbf{u}}\|^{2}_{L^{2}}(\tau)d\tau.$
For the first term on the right-hand side of (3.3), one gets by integration by
parts that
$\displaystyle-\frac{1}{2}\int^{T}_{0}\langle\partial^{\alpha}(\epsilon\partial_{t})p,\partial^{\alpha}{\rm
div}{\mathbf{u}}\rangle(\tau)d\tau=$
$\displaystyle-\frac{1}{2}\langle\partial^{\alpha}p,\epsilon\partial^{\alpha}{\rm
div}(\epsilon{\mathbf{u}})\rangle\Big{|}^{T}_{0}$
$\displaystyle+\frac{1}{2}\int^{T}_{0}\langle\partial^{\alpha}\nabla
p,\partial^{\alpha}(\epsilon\partial_{t}){\mathbf{u}}\rangle(\tau)d\tau.$
By virtue of the estimate (3.3) on $(\epsilon
q,\epsilon{\mathbf{u}},\theta-\bar{\theta})$ and (3.51), we find that
$\displaystyle\left|\frac{1}{2}\langle\partial^{\alpha}p,\epsilon\partial^{\alpha}{\rm
div}(\epsilon{\mathbf{u}})\rangle\Big{|}^{T}_{0}\right|\leq$
$\displaystyle\sup_{\tau\in[0,T]}\\{\|p\|_{H^{s}}\|\epsilon{\mathbf{u}}\|_{H^{s+1}}\\}$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\},$
$\displaystyle\frac{1}{2}\bigg{|}\int^{T}_{0}\langle\partial^{\alpha}\nabla
p,\partial^{\alpha}(\epsilon\partial_{t}){\mathbf{u}}\rangle(\tau)d\tau\bigg{|}\leq$
$\displaystyle\frac{1}{2}\bigg{\\{}\int^{T}_{0}\|\partial^{\alpha}\nabla
p\|^{2}_{L^{2}}(\tau)d\tau\bigg{\\}}^{1/2}$
$\displaystyle\times\bigg{\\{}\int^{T}_{0}\|\partial^{\alpha}(\epsilon\partial_{t}){\mathbf{u}}\|^{2}_{L^{2}}(\tau)d\tau\bigg{\\}}^{1/2}$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$
These bounds, together with (3.3), yield the desired estimate (3.52). This
completes the proof. ∎
### 3.4. $H^{s-1}$-estimate on ${\rm curl\,}{\mathbf{u}}$
The another key point to obtain a uniform bound for ${\mathbf{u}}$ is the
following estimate on ${\rm curl\,}{\mathbf{u}}$.
###### Lemma 3.8.
Let $s>5/2$ be an integer and $(p,{\mathbf{u}},\mathbf{H},\theta)$ be the
solution to the Cauchy problem (1.11)–(1.14), (1.19) on $[0,T_{1}]$. Then
there exist a constant $l_{7}>0$ and an increasing function $C(\cdot)$ such
that, for any $\epsilon\in(0,1]$ and $t\in[0,T_{1}]$, $T=\min\\{T_{1},1\\}$,
it holds that
$\displaystyle\sup_{\tau\in[0,t]}\left\\{\|{\rm
curl\,}(b({-\theta}){\mathbf{u}})(\tau)\|^{2}_{H^{s-1}}+\|{\rm
curl\,}\mathbf{H}(\tau)\|^{2}_{H^{s-1}}\right\\}$
$\displaystyle+l_{7}\int^{t}_{0}\left\\{\|\nabla{\rm
curl\,}(b({-\theta}){\mathbf{u}})\|^{2}_{H^{s-1}}+\|\nabla{\rm
curl\,}\mathbf{H}(\tau)\|^{2}_{H^{s-1}}\right\\}(\tau)d\tau$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$ (3.54)
###### Proof.
Applying the operator _curl_ to the equations (1.12) and (1.13), using the
identities (2.1) and (2.2), and the fact that ${\rm curl\,}\nabla=0$, one
infers
$\displaystyle\partial_{t}({\rm
curl\,}(b(-\theta){\mathbf{u}}))+({\mathbf{u}}\cdot\nabla)({\rm
curl\,}(b(-\theta){\mathbf{u}}))$ $\displaystyle\quad\ \ ={\rm
curl\,}\\{a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\mathbf{H}\\}+\bar{\mu}{\rm div}\\{a(\epsilon
p)b(\theta)\nabla({\rm curl\,}(b(-\theta){\mathbf{u}}))\\}+\Upsilon_{1},$
(3.55) $\displaystyle\partial_{t}({\rm curl\,}\mathbf{H})-{\rm curl\,}[{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})]-\bar{\nu}\Delta({\rm
curl\,}\mathbf{H})=0,$ (3.56)
where $\Upsilon_{1}$ is defined by
$\displaystyle\Upsilon_{1}:=$ $\displaystyle\bar{\mu}{\rm
div}\big{(}a(\epsilon p)(\nabla b(\theta))\otimes{\rm
curl\,}(b(-\theta){\mathbf{u}})\big{)}-\bar{\mu}\nabla a(\epsilon
p)\cdot\nabla(b(\theta){\rm curl\,}(b(-\theta){\mathbf{u}}))$
$\displaystyle-\bar{\mu}a(\epsilon p)\Delta((\nabla
b(\theta))\times(b(-\theta){\mathbf{u}}))-\nabla a(\epsilon
p)\times(\bar{\mu}\Delta{\mathbf{u}}+(\bar{\mu}+\bar{\lambda})\nabla{\rm
div}{\mathbf{u}})$ $\displaystyle+{\rm
curl\,}(b(-\theta){\mathbf{u}}\partial_{t}\theta)+[{\rm
curl\,},{\mathbf{u}}]\cdot\nabla(b(-\theta){\mathbf{u}})+{\rm
curl\,}(b(-\theta){\mathbf{u}}({\mathbf{u}}\cdot\nabla\theta)).$
For any multi-index $\alpha$ satisfying $0\leq|\alpha|\leq s-1$, we apply the
operator $\partial^{\alpha}$ to (3.55) and (3.56) to obtain
$\displaystyle\partial_{t}\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))+({\mathbf{u}}\cdot\nabla)[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))]$ $\displaystyle\quad\quad\quad\ \
=\partial^{\alpha}{\rm curl\,}\\{a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\mathbf{H}\\}$ $\displaystyle\quad\quad\quad\quad\ \
+\bar{\mu}{\rm div}\\{a(\epsilon p)b(\theta)\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))]\\}+\partial^{\alpha}\Upsilon_{1}+\Upsilon_{2},$
(3.57) $\displaystyle\partial_{t}\partial^{\alpha}({\rm
curl\,}\mathbf{H})-\partial^{\alpha}{\rm curl\,}[{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})]-\bar{\nu}\Delta({\rm
curl\,}\mathbf{H})=0,$ (3.58)
where
$\displaystyle\Upsilon_{2}:=$
$\displaystyle-[\partial^{\alpha},{\mathbf{u}}]\cdot\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))]$ $\displaystyle-[\partial^{\alpha},{\rm
div}(a(\epsilon p)b(\theta))]\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))].$
Multiplying (3.57) by $\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))$ and (3.58) by $\partial^{\alpha}({\rm
curl\,}\mathbf{H})$ respectively, summing up, and integrating over
$\mathbb{R}^{3}$, we deduce that
$\displaystyle\frac{1}{2}\frac{d}{dt}\\{\|\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\|^{2}_{L^{2}}+\|\partial^{\alpha}({\rm
curl\,}\mathbf{H})\|^{2}_{L^{2}}\\}+\bar{\nu}\|\partial^{\alpha}({\rm
curl\,}\mathbf{H})\|^{2}_{L^{2}}$ $\displaystyle+\langle a(\epsilon
p)b(\theta)\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))],\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))]\rangle$ $\displaystyle=$
$\displaystyle-\langle({\mathbf{u}}\cdot\nabla)[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))],\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\rangle$
$\displaystyle-\langle\partial^{\alpha}{\rm curl\,}\\{a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\mathbf{H}\\},\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\rangle$
$\displaystyle+\langle\partial^{\alpha}{\rm curl\,}[{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})],\partial^{\alpha}({\rm
curl\,}\mathbf{H})\rangle+\langle\partial^{\alpha}\Upsilon_{1}+\Upsilon_{2},\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\rangle$ $\displaystyle:=$
$\displaystyle\mathcal{J}_{1}+\mathcal{J}_{2}+\mathcal{J}_{3}+\mathcal{J}_{4},$
(3.59)
where $\mathcal{J}_{i}$ ($i=1,\cdots,4$) will be bounded as follows.
An integration by parts leads to
$\displaystyle|\mathcal{J}_{1}|\leq\|{\rm
div}{\mathbf{u}}\|_{L^{\infty}}\|\nabla[\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))]\|^{2}_{L^{2}}.$
By virtue of (2.3), the Cauchy-Schwarz and Moser-type inequalities (see [34]),
the term $\mathcal{J}_{2}$ can be bounded as follows.
$\displaystyle|\mathcal{J}_{2}|\leq$
$\displaystyle|\langle\partial^{\alpha}\\{a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\mathbf{H}\\},\partial^{\alpha}{\rm curl\,}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\rangle|$ $\displaystyle\leq$
$\displaystyle|\partial^{\alpha}\\{a(\epsilon p)({\rm
curl\,}\mathbf{H})\times\mathbf{H}\\}\|_{L^{2}}\|\partial^{\alpha}\nabla({\rm
curl\,}(b(-\theta){\mathbf{u}}))\|_{L^{2}}$ $\displaystyle\leq$
$\displaystyle\eta_{2}\|\partial^{\alpha}\nabla({\rm
curl\,}(b(-\theta){\mathbf{u}}))\|_{L^{2}}^{2}$
$\displaystyle+C_{\eta}\\{\|{\rm
curl\,}\mathbf{H}\|^{2}_{L^{\infty}}\|a(\epsilon
p)\mathbf{H}\|_{H^{s-1}}^{2}+\|a(\epsilon
p)\mathbf{H}\|^{2}_{L^{\infty}}\|{\rm curl\,}\mathbf{H}\|^{2}_{H^{s-1}}\\},$
where $\eta_{2}>0$ is a sufficiently small constant independent of $\epsilon$.
If we integrate by parts, make use of (2.3) and the fact that ${\rm
curl\,}{\rm curl\,}\mathbf{a}=\nabla\,{\rm div}\,\mathbf{a}-\Delta\mathbf{a}$
and ${\rm div}\mathbf{H}=0$, we see that the term $\mathcal{J}_{3}$ can be
rewritten as
$\displaystyle\mathcal{J}_{3}=\left\langle\partial^{\alpha}{\rm
curl\,}({\mathbf{u}}\times\mathbf{H}),\partial^{\alpha}\Delta\mathbf{H}\right\rangle,$
which, together with the Moser-type inequality, implies that
$\displaystyle|\mathcal{J}_{3}|\leq
C(\mathcal{S})+\eta_{3}\|\mathbf{H}^{\epsilon}(\tau)\|^{2}_{s+1},$
where $\eta_{3}>0$ is a sufficiently small constant independent of $\epsilon$.
To handle $\mathcal{J}_{4}$, we note that the leading order terms in
$\Upsilon_{1}$ are of third-order in $\theta$ and of second-order in
${\mathbf{u}}$, and the leading order terms in $\Upsilon_{2}$ are of order
$s+1$ in ${\mathbf{u}}$ and of order $s+1$ in $(\epsilon p,\theta)$. Then it
follows that
$\displaystyle|\mathcal{J}_{4}|$ $\displaystyle\leq$ $\displaystyle
C_{0}(\|\partial^{\alpha}\Upsilon_{1}\|_{L^{2}}+\|\Upsilon_{2}\|_{L^{2}})\|\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\|_{L^{2}}$ $\displaystyle\leq$ $\displaystyle
C(\mathcal{S})\|\partial^{\alpha}({\rm
curl\,}(b(-\theta){\mathbf{u}}))\|_{L^{2}}.$
Putting the above estimates into the (3.4), choosing $\eta_{2}$ and $\eta_{3}$
sufficient small, summing over $\alpha$ for $0\leq|\alpha|\leq s-1$, and then
integrating the result on $[0,t]$, we conclude
$\displaystyle\sup_{\tau\in[0,t]}\left\\{\|{\rm
curl\,}(b({-\theta}){\mathbf{u}})(\tau)\|^{2}_{H^{s-1}}+\|{\rm
curl\,}\mathbf{H}(\tau)\|^{2}_{H^{s-1}}\right\\}$
$\displaystyle\quad+l_{7}\int^{t}_{0}\left\\{\|\nabla{\rm
curl\,}(b({-\theta}){\mathbf{u}})\|^{2}_{H^{s-1}}+\|\nabla{\rm
curl\,}\mathbf{H}(\tau)\|^{2}_{H^{s-1}}\right\\}(\tau)d\tau$
$\displaystyle\leq C_{0}\big{\\{}\|{\rm
curl\,}(b({-\theta}){\mathbf{u}})(0)\|^{2}_{H^{s-1}}+\|{\rm
curl\,}\mathbf{H}(0)\|^{2}_{H^{s-1}}\big{\\}}$
$\displaystyle\quad+C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}$
$\displaystyle\leq
C(\mathcal{O}_{0})\exp\\{(\sqrt{T}+\epsilon)C(\mathcal{O}(T))\\}.$
∎
###### Proof of Proposition 2.4.
By virtue of the definition of the norm $\|\cdot\|_{s,\epsilon}$ and the fact
that
$\displaystyle\|{\mathbf{v}}\|_{H^{m+1}}\leq K\big{(}\|{\rm
div}{\mathbf{v}}\|_{H^{m}}+\|{\rm
curl\,}{\mathbf{v}}\|_{H^{m}}+\|{\mathbf{v}}\|_{H^{m}}\big{)},\quad\forall\,\,{\mathbf{v}}\in
H^{m+1}(\mathbb{R}^{3}),$
Proposition 2.4 follows directly from Lemmas 3.1, 3.3, 3.7 and 3.8. ∎
Once Proposition 2.4 is established, the existence part of Theorem 1.1 can be
proved by directly applying the same arguments as in [2, 40], and hence we
omit the details here.
## 4\. Decay of the local energy and zero Mach number limit
In this section, we shall prove the convergence part of Theorem 1.1 by
modifying the arguments developed by Métivier and Schochet [40], see also some
extensions in [1, 2, 38].
###### Proof of the convergence part of Theorem 1.1.
The uniform estimate (1.21) implies that
$\displaystyle\sup_{\tau\in[0,T_{0}]}\|(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(\tau)\|_{H^{s}}+\sup_{\tau\in[0,T_{0}]}\|\theta^{\epsilon}-\bar{\theta}\|_{H^{s+1}}<+\infty.$
Thus, after extracting a subsequence, one has
$\displaystyle(p^{\epsilon},{\mathbf{u}}^{\epsilon})\rightharpoonup(\bar{p},{\mathbf{w}})$
$\displaystyle\text{weakly-}\ast\ \text{in}$ $\displaystyle\quad\quad
L^{\infty}(0,T_{0};H^{s}(\mathbb{R}^{3})),$ (4.1)
$\displaystyle\mathbf{H}^{\epsilon}\rightharpoonup{\mathbf{B}}$
$\displaystyle\text{weakly-}\ast\ \text{in}$ $\displaystyle\qquad
L^{\infty}(0,T_{0};H^{s}(\mathbb{R}^{3})),$ (4.2)
$\displaystyle\theta^{\epsilon}-\bar{\theta}\rightharpoonup\vartheta-\bar{\theta}\\!\\!\\!\\!\\!\\!$
$\displaystyle\text{weakly-}\ast\ \text{in}$ $\displaystyle\qquad
L^{\infty}(0,T_{0};H^{s+1}(\mathbb{R}^{3})).$ (4.3)
It follows from the equations for $\mathbf{H}^{\epsilon}$ and
$\theta^{\epsilon}$ that
$\displaystyle\partial_{t}\mathbf{H}^{\epsilon},\,\partial_{t}\theta^{\epsilon}\in
C([0,T_{0}],H^{s-2}(\mathbb{R}^{3})).$ (4.4)
(4.2)–(4.4) implies, after further extracting a subsequence, that for all
$s^{\prime}<s$,
$\displaystyle\mathbf{H}^{\epsilon}\rightarrow{\mathbf{B}}\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!$
strongly in $\displaystyle\quad
C([0,T_{0}],H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3})),$ (4.5)
$\displaystyle\theta^{\epsilon}-\bar{\theta}\rightarrow\vartheta-\bar{\theta}\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!$
strongly in $\displaystyle\quad
C([0,T_{0}],H^{s^{\prime}+1}_{\mathrm{loc}}(\mathbb{R}^{3})),$ (4.6)
where the limit ${\mathbf{B}}\in
C([0,T_{0}],H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))\cap
L^{\infty}(0,T_{0};H^{s}_{\mathrm{loc}}(\mathbb{R}^{3}))$ and
$(\vartheta-\bar{\theta})\in
C([0,T_{0}],H^{s^{\prime}+1}_{\mathrm{loc}}(\mathbb{R}^{3}))\cap
L^{\infty}(0,T_{0};H^{s+1}_{\mathrm{loc}}(\mathbb{R}^{3}))$.
Similarly, from (3.8) we get
$\displaystyle{\rm
curl\,}\big{(}e^{-\theta^{\epsilon}}{\mathbf{u}}^{\epsilon}\big{)}\rightarrow{\rm
curl\,}\big{(}e^{-\vartheta}{\mathbf{w}}\big{)}\quad\text{strongly in}\quad
C([0,T_{0}],H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3}))$ (4.7)
for all $s^{\prime}<s$.
In order to obtain the limit system, one needs to show that the limits in
(4.1) hold in the strong topology of
$L^{2}(0,T_{0};H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))$ for all
$s^{\prime}<s$. To this end, we first show that $\bar{p}=0$ and ${\rm
div}(2{\mathbf{w}}-\bar{\kappa}e^{\vartheta}\nabla\vartheta)=0$. In fact, the
equations (1.11) and (1.12) can be rewritten as
$\displaystyle\epsilon\,\partial_{t}p^{\epsilon}+{\rm
div}(2{\mathbf{u}}^{\epsilon}-\bar{\kappa}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla\theta^{\epsilon})=\epsilon
f^{\epsilon},$ (4.8)
$\displaystyle\epsilon\,e^{-\theta^{\epsilon}}\partial_{t}{\mathbf{u}}^{\epsilon}+\nabla
p^{\epsilon}=\epsilon\,\mathbf{g}^{\epsilon}.$ (4.9)
By virtue of (1.21), $f^{\epsilon}$ and $\mathbf{g}^{\epsilon}$ are uniformly
bounded in $C([0,T_{0}],H^{s-1}(\mathbb{R}^{3}))$. Passing to the weak limit
in (4.8) and (4.9), respectively, we see that $\nabla\bar{p}=0$ and ${\rm
div}(2{\mathbf{w}}-\bar{\kappa}e^{\vartheta}\nabla\vartheta)=0$. Since
$\bar{p}\in L^{\infty}(0,T_{0};H^{s}(\mathbb{R}^{3}))$, we infer that
$\bar{p}=0$.
Notice that by virtue of (4.7), the strong compactness for the incompressible
component of $e^{-\theta^{\epsilon}}{\mathbf{u}}^{\epsilon}$ holds. So, it is
sufficient to prove the following proposition on the acoustic components in
order to get the strong convergence of ${\mathbf{u}}^{\epsilon}$.
###### Proposition 4.1.
Suppose that the assumptions in Theorem 1.1 hold. Then, $p^{\epsilon}$
converges to $0$ strongly in
$L^{2}(0,T_{0};H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))$ and ${\rm
div}(2{\mathbf{u}}^{\epsilon}-\bar{\kappa}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla\theta^{\epsilon})$ converges to $0$
strongly in $L^{2}(0,T_{0};H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3}))$
for all $s^{\prime}<s$.
The proof of Proposition 4.1 is based on the following dispersive estimates on
the wave equation obtained by Métivier and Schochet [40] and reformulated in
[2].
###### Lemma 4.2.
([40, 2]) Let $T>0$ and $v^{\epsilon}$ be a bounded sequence in
$C([0,T],H^{2}(\mathbb{R}^{3}))$, such that
$\displaystyle\epsilon^{2}\partial_{t}(a^{\epsilon}\partial_{t}v^{\epsilon})-\nabla\cdot(b^{\epsilon}\nabla
v^{\epsilon})=c^{\epsilon},$
where $c^{\epsilon}$ converges to $0$ strongly in
$L^{2}(0,T;L^{2}(\mathbb{R}^{3}))$. Assume further that for some $s>3/2+1$,
the coefficients $(a^{\epsilon},b^{\epsilon})$ are uniformly bounded in
$C([0,T],H^{s}(\mathbb{R}^{3}))$ and converge in
$C([0,T],H^{s}_{\mathrm{loc}}(\mathbb{R}^{3}))$ to a limit $(a,b)$ satisfying
the decay estimates
$\displaystyle|a(x,t)-\hat{a}|\leq
C_{0}|x|^{-1-\zeta},\quad|\nabla_{x}a(x,t)|\leq C_{0}|x|^{-2-\zeta},$
$\displaystyle|b(x,t)-\hat{b}|\leq
C_{0}|x|^{-1-\zeta},\quad|\nabla_{x}b(x,t)|\leq C_{0}|x|^{-2-\zeta},$
for some positive constants $\hat{a}$, $\hat{b}$, $C_{0}$ and $\zeta$. Then
the sequence $v^{\epsilon}$ converges to $0$ strongly in
$L^{2}(0,T;L^{2}_{\mathrm{loc}}(\mathbb{R}^{3}))$.
###### Proof of Proposition 4.1.
We fist show that $p^{\epsilon}$ converges to $0$ strongly in
$L^{2}(0,T_{0};\linebreak H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))$ for
all $s^{\prime}<s$. Applying $\epsilon^{2}\partial_{t}$ to (1.11), we find
that
$\displaystyle\epsilon^{2}\partial_{t}\\{\partial_{t}p^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)p^{\epsilon}\\}+{\epsilon}\partial_{t}\\{{\rm
div}(2{\mathbf{u}}^{\epsilon}-\bar{\kappa}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla\theta^{\epsilon})\\}$
$\displaystyle\quad\quad=\epsilon^{3}\partial_{t}\\{e^{-\epsilon
p^{\epsilon}}[\bar{\nu}|{\rm
curl\,}\mathbf{H}^{\epsilon}|^{2}+\Psi({\mathbf{u}}^{\epsilon}):\nabla{\mathbf{u}}^{\epsilon}]\\}+\epsilon^{2}\partial_{t}\\{\bar{\kappa}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla
p^{\epsilon}\cdot\nabla\theta^{\epsilon}\\}.$ (4.10)
Dividing (1.12) by $e^{-\theta^{\epsilon}}$ and then applying the operator
_div_ to the resulting equations, one gets
$\displaystyle\epsilon\partial_{t}{\rm div}{\mathbf{u}}^{\epsilon}+{\rm
div}\big{(}e^{\theta^{\epsilon}}{\nabla p^{\epsilon}}\big{)}=$
$\displaystyle-\epsilon{\rm
div}\\{({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon}\\}$
$\displaystyle+\epsilon{\rm div}\big{\\{}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}[({\rm curl\,}\mathbf{H})\times\mathbf{H}+{\rm
div}\Psi^{\epsilon}({\mathbf{u}}^{\epsilon})]\big{\\}},$ (4.11)
Subtracting (4) from (4), we have
$\displaystyle\epsilon^{2}\partial_{t}\Big{(}\frac{1}{2}\partial_{t}p^{\epsilon}\Big{)}-{\rm
div}\big{(}e^{\theta^{\epsilon}}{\nabla p^{\epsilon}}\big{)}=\epsilon
F^{\epsilon}(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon}),$
(4.12)
where
$F^{\epsilon}(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$
is a smooth function in its variables with $F(0)=0$. By the uniform
boundedness of
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$
one infers that
$\displaystyle\epsilon
F^{\epsilon}(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})\rightarrow
0\quad\text{strongly in}\quad L^{2}(0,T_{0};L^{2}(\mathbb{R}^{3})).$
By the the strong convergence of $\theta^{\epsilon}$, the initial conditions
(1.22), and the arguments in Section 8.1 in [2], one can easily prove that the
coefficients in (4.12) satisfy the conditions in Lemma 4.2. Therefore, we can
apply Lemma 4.2 to obtain
$\displaystyle p^{\epsilon}\rightarrow 0\quad\text{strongly in}\quad
L^{2}(0,T_{0};L^{2}_{\mathrm{loc}}(\mathbb{R}^{3})).$
Since $p^{\epsilon}$ is bounded uniformly in
$C([0,T_{0}],H^{s}(\mathbb{R}^{3}))$, an interpolation argument gives
$\displaystyle p^{\epsilon}\rightarrow 0\quad\text{strongly in}\quad
L^{2}(0,T_{0};H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))\ \ \text{for
all}\ \ s^{\prime}<s.$
Similarly, we can obtain the strong convergence of ${\rm
div}(2{\mathbf{u}}^{\epsilon}-\kappa^{\epsilon}e^{-\epsilon
p^{\epsilon}+\theta^{\epsilon}}\nabla\theta^{\epsilon})$. This completes the
proof. ∎
We continue our proof of Theorem 1.1. It follows from Proposition 4.1 and
(4.6) that
$\displaystyle{\rm div}\,{\mathbf{u}}^{\epsilon}\rightarrow{\rm
div}\,{\mathbf{w}}\quad\text{strongly in}\quad
L^{2}(0,T_{0};H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3})).$
Thus, using (4.7), one obtains
$\displaystyle{\mathbf{u}}^{\epsilon}\rightarrow{\mathbf{w}}\quad\text{strongly
in}\quad
L^{2}(0,T_{0};H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{3}))\qquad\mbox{for
all }s^{\prime}<s.$
By (4.5), (4.6), and Proposition 4.1, we find that
$\begin{array}[]{lcl}\nabla{\mathbf{u}}^{\epsilon}\rightarrow\nabla{\mathbf{w}}&\text{strongly
in}&L^{2}(0,T_{0};H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3}));\\\
\nabla\mathbf{H}^{\epsilon}\rightarrow\nabla{\mathbf{B}}&\text{strongly
in}&L^{2}(0,T_{0};H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3}));\\\
\nabla\theta^{\epsilon}\rightarrow\nabla\vartheta&\text{strongly
in}&L^{2}(0,T_{0};H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{3})).\end{array}$
Passing to the limits in the equations for $p^{\epsilon}$,
$\mathbf{H}^{\epsilon}$, and $\theta^{\epsilon}$, one sees that the limit
$(0,{\mathbf{w}},{\mathbf{B}},\vartheta)$ satisfies, in the sense of
distributions, that
$\displaystyle{\rm
div}(2{\mathbf{w}}-\bar{\kappa}\,e^{\vartheta}\nabla\vartheta)=0,$ (4.13)
$\displaystyle\partial_{t}{\mathbf{B}}-{\rm
curl\,}({\mathbf{w}}\times{\mathbf{B}})-\bar{\nu}\Delta{\mathbf{B}}=0,\quad{\rm
div}{\mathbf{B}}=0,$ (4.14)
$\displaystyle\partial_{t}\vartheta+({\mathbf{w}}\cdot\nabla)\vartheta+{\rm
div}{\mathbf{w}}=\bar{\kappa}\,{\rm div}(e^{\vartheta}\nabla\vartheta).$
(4.15)
On the other hand, applying _curl_ to the momentum equations (1.12), using the
equations (1.11) and (1.14) on $p^{\epsilon}$ and $\theta^{\epsilon}$, and
then taking to the limit on the resulting equations, we deduce that
$\displaystyle{\rm
curl\,}\left\\{\partial_{t}\big{(}e^{-\vartheta}{\mathbf{w}})+{\rm
div}\big{(}{\mathbf{w}}e^{-\vartheta}\otimes{\mathbf{w}}\big{)}-({\rm
curl\,}{\mathbf{B}})\times{\mathbf{B}}-{\rm div}\Phi({\mathbf{w}})\right\\}=0$
holds in the sense of distributions. Therefore it follows from (4.13)–(4.15)
that
$\displaystyle
e^{-\vartheta}[\partial_{t}{\mathbf{w}}+({\mathbf{w}}\cdot\nabla){\mathbf{w}}]+\nabla\pi=({\rm
curl\,}{\mathbf{B}})\times{\mathbf{B}}+{\rm div}\Phi({\mathbf{w}}),$ (4.16)
for some function $\pi$.
Following the same arguments as those in the proof of Theorem 1.5 in [40], we
conclude that $({\mathbf{w}},{\mathbf{B}},\vartheta)$ satisfies the initial
condition
$\displaystyle({\mathbf{w}},{\mathbf{B}},\vartheta)|_{t=0}=({\mathbf{w}}_{0},{\mathbf{B}}_{0},\vartheta_{0}).$
(4.17)
Moreover, the standard iterative method shows that the system (4.13)–(4.16)
with initial data (4.17) has a unique solution
$({\mathbf{w}}^{*},{\mathbf{B}}^{*},\vartheta^{*}-\bar{\theta})\in
C([0,T_{0}],H^{s}(\mathbb{R}^{3}))$. Thus, the uniqueness of solutions to the
limit system (4.13)–(4.16) implies that the above convergence holds for the
full sequence of
$(p^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon},\theta^{\epsilon})$.
Therefore the proof is completed. ∎
Acknowledgements: This work was partially done when Li was visiting the
Institute of Mathematical Sciences, CUHK. He would like to thank the institute
for hospitality. Jiang was supported by the National Basic Research Program
under the Grant 2011CB309705 and NSFC (Grant No. 11229101). Ju was supported
by NSFC (Grant No. 11171035). Li was supported by NSFC (Grant No. 11271184,
10971094), NCET-11-0227, PAPD, and the Fundamental Research Funds for the
Central Universities. Xin was supported in part by the Zheng Ge Ru Foundation,
Hong Kong RGC Earmarked Research Grants CUHK4042/08P and CUHK4041/11p.
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|
arxiv-papers
| 2011-11-12T13:36:54 |
2024-09-04T02:49:24.247703
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Song Jiang, Qiangchang Ju, Fucai Li, and Zhouping Xin",
"submitter": "Fucai Li",
"url": "https://arxiv.org/abs/1111.2925"
}
|
1111.2926
|
# Incompressible limit of the compressible non-isentropic magnetohydrodynamic
equations with zero magnetic diffusivity
Song Jiang LCP, Institute of Applied Physics and Computational Mathematics,
P.O. Box 8009, Beijing 100088, P.R. China jiang@iapcm.ac.cn , Qiangchang Ju
Institute of Applied Physics and Computational Mathematics, P.O. Box 8009-28,
Beijing 100088, P.R. China qiangchang_ju@yahoo.com and Fucai Li Department
of Mathematics, Nanjing University, Nanjing 210093, P.R. China fli@nju.edu.cn
###### Abstract.
We study the incompressible limit of the compressible non- isentropic
magnetohydrodynamic equations with zero magnetic diffusivity and general
initial data in the whole space $\mathbb{R}^{d}$ ($d=2,3$). We first establish
the existence of classic solutions on a time interval independent of the Mach
number. Then, by deriving uniform a priori estimates, we obtain the
convergence of the solution to that of the incompressible magnetohydrodynamic
equations as the Mach number tends to zero.
###### Key words and phrases:
Compressible MHD equations, non-isentropic, zero magnetic diffusivity,
incompressible limit
###### 2000 Mathematics Subject Classification:
76W05, 35B40
## 1\. Introduction
This paper is concerned with the incompressible limit to the compressible non-
isentropic magnetohydrodynamic (MHD) equations with zero magnetic diffusivity
and general initial data in the whole space $\mathbb{R}^{d}$ ($d=2,3$).
In the study of a highly conducting fluid, for example, the magnetic fusion,
it is rational to ignore the magnetic diffusion term in the MHD equations
since the magnetic diffusion coefficient (resistivity coefficient) is
inversely proportional to the electrical conductivity coefficient, see [10].
In this situation, the system, describing the motion of the fluid in
${\mathbb{R}^{d}}$, can be described by the following compressible non-
isentropic MHD equations with zero magnetic diffusivity:
$\displaystyle\partial_{t}\rho+{\rm div}(\rho{\mathbf{u}})=0,$ (1.1)
$\displaystyle\partial_{t}(\rho{\mathbf{u}})+{\rm
div}\left(\rho{\mathbf{u}}\otimes{\mathbf{u}}\right)+{\nabla
p}=(\nabla\times\mathbf{H})\times\mathbf{H}+{\rm div}\Psi,$ (1.2)
$\displaystyle\partial_{t}\mathbf{H}-\nabla\times({\mathbf{u}}\times\mathbf{H})=0,\quad{\rm
div}\mathbf{H}=0,$ (1.3) $\displaystyle\partial_{t}{\mathcal{E}}+{\rm
div}\left({\mathbf{u}}({\mathcal{E}}^{\prime}+p)\right)={\rm
div}(({\mathbf{u}}\times\mathbf{H})\times\mathbf{H})+{\rm
div}({\mathbf{u}}\Psi+\kappa\nabla\theta).$ (1.4)
Here $\rho$ denotes the density, ${\mathbf{u}}\in{\mathbb{R}}^{d}$ the
velocity, $\mathbf{H}\in{\mathbb{R}}^{d}$ the magnetic field, and $\theta$ the
temperature, respectively; $\Psi$ is the viscous stress tensor given by
$\Psi=2\mu\mathbb{D}({\mathbf{u}})+\lambda{\rm
div}{\mathbf{u}}\;\mathbf{I}_{d}$
with
$\mathbb{D}({\mathbf{u}})=(\nabla{\mathbf{u}}+\nabla{\mathbf{u}}^{\top})/2$,
and $\mathbf{I}_{d}$ being the $d\times d$ identity matrix, and
$\nabla{\mathbf{u}}^{\top}$ the transpose of the matrix $\nabla{\mathbf{u}}$;
${\mathcal{E}}$ is the total energy given by
${\mathcal{E}}={\mathcal{E}}^{\prime}+|\mathbf{H}|^{2}/2$ and
${\mathcal{E}}^{\prime}=\rho\left(e+|{\mathbf{u}}|^{2}/2\right)$ with $e$
being the internal energy, $\rho|{\mathbf{u}}|^{2}/2$ the kinetic energy, and
$|\mathbf{H}|^{2}/2$ the magnetic energy. The viscosity coefficients $\lambda$
and $\mu$ of the fluid satisfy $2\mu+d\lambda>0$ and $\mu>0$; $\kappa>0$ is
the heat conductivity. For simplicity, we assume that $\mu,\lambda$ and
$\kappa$ are constants. The equations of state $p=p(\rho,\theta)$ and
$e=e(\rho,\theta)$ relate the pressure $p$ and the internal energy $e$ to the
density $\rho$ and the temperature $\theta$ of the flow.
For the smooth solution to the system (1.1)–(1.4), we can rewrite the total
energy equation (1.4) in the form of the internal energy. In fact, multiplying
(1.2) by ${\mathbf{u}}$ and (1.3) by $\mathbf{H}$, and summing the resulting
equations together, we obtain
$\displaystyle\frac{d}{dt}\Big{(}\frac{1}{2}\rho|{\mathbf{u}}|^{2}+\frac{1}{2}|\mathbf{H}|^{2}\Big{)}$
$\displaystyle+\frac{1}{2}{\rm
div}\big{(}\rho|{\mathbf{u}}|^{2}{\mathbf{u}}\big{)}+\nabla
p\cdot{\mathbf{u}}$ $\displaystyle={\rm
div}\Psi\cdot{\mathbf{u}}+(\nabla\times\mathbf{H})\times\mathbf{H}\cdot{\mathbf{u}}+\nabla\times({\mathbf{u}}\times\mathbf{H})\cdot\mathbf{H}.$
(1.5)
Using the identities
$\displaystyle{\rm
div}(\mathbf{H}\times(\nabla\times\mathbf{H}))=|\nabla\times\mathbf{H}|^{2}-\nabla\times(\nabla\times\mathbf{H})\cdot\mathbf{H},$
(1.6) $\displaystyle{\rm
div}(({\mathbf{u}}\times\mathbf{H})\times\mathbf{H})=(\nabla\times\mathbf{H})\times\mathbf{H}\cdot{\mathbf{u}}+\nabla\times({\mathbf{u}}\times\mathbf{H})\cdot\mathbf{H},$
(1.7)
and subtracting (1) from (1.4), we thus obtain the internal energy equation
$\partial_{t}(\rho e)+{\rm div}(\rho{\mathbf{u}}e)+({\rm
div}{\mathbf{u}})p=\Psi:\nabla{\mathbf{u}}+\kappa\Delta\theta,$ (1.8)
where $\Psi:\nabla{\mathbf{u}}$ denotes the scalar product of two matrices:
$\Psi:\nabla{\mathbf{u}}=\sum^{3}_{i,j=1}\frac{\mu}{2}\left(\frac{\partial
u^{i}}{\partial x_{j}}+\frac{\partial u^{j}}{\partial
x_{i}}\right)^{2}+\lambda|{\rm
div}{\mathbf{u}}|^{2}=2\mu|\mathbb{D}({\mathbf{u}})|^{2}+\lambda(\mbox{tr}\mathbb{D}({\mathbf{u}}))^{2}.$
Using the Gibbs relation
$\theta\mathrm{d}S=\mathrm{d}e+p\,\mathrm{d}\left(\frac{1}{\rho}\right),$
(1.9)
we can further replace the equation (1.8) by
$\partial_{t}(\rho S)+{\rm div}(\rho
S{\mathbf{u}})=\Psi:\nabla{\mathbf{u}}+\kappa\Delta\theta,$ (1.10)
where $S$ denotes the entropy.
In the present paper, we assume that $\kappa=0$ in (1.10). Now, as in [27], we
reconsider the equations of state as functions of $S$ and $p$, i.e.,
$\rho=R(S,p)$ and $\theta=\Theta(S,p)$ for some positive smooth functions $R$
and $\Theta$ defined for all $S$ and $p>0$, and satisfying $\partial
R/\partial p>0$. For instance, we have $\rho=p^{1/\gamma}e^{-S/\gamma}$ for
ideal fluids. Then, by utilizing (1.1) together with the constraint ${\rm
div}{\mathbf{H}}=0$, the system (1.1), (1.2), (1.4) and (1.10) can be
rewritten as
$\displaystyle A(S,p)(\partial_{t}p+({\mathbf{u}}\cdot\nabla)p)+{\rm
div}{\mathbf{u}}=0,$ (1.11) $\displaystyle
R(S,p)(\partial_{t}{\mathbf{u}}+({\mathbf{u}}\cdot\nabla){\mathbf{u}})+\nabla
p=(\nabla\times\mathbf{H})\times\mathbf{H}+{\rm div}\Psi,$ (1.12)
$\displaystyle\partial_{t}{\mathbf{H}}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})=0,\quad{\rm div}\mathbf{H}=0,$ (1.13)
$\displaystyle
R(S,p)\Theta(S,p)(\partial_{t}S+({\mathbf{u}}\cdot\nabla)S)=\Psi:\nabla{\mathbf{u}},$
(1.14)
where
$\displaystyle A(S,p)=\frac{1}{R(S,p)}\frac{\partial R(S,p)}{\partial p}.$
(1.15)
Considering the physical explanation of the incompressible limit, we introduce
the dimensionless parameter $\epsilon$, the Mach number, and make the
following changes of variables:
$\displaystyle p(x,t)=p^{\epsilon}(x,\epsilon t),\quad
S(x,t)=S^{\epsilon}(x,\epsilon t),$
$\displaystyle{{\mathbf{u}}}(x,t)=\epsilon{\mathbf{u}}^{\epsilon}(x,\epsilon
t),\;\;\;{\mathbf{H}}(x,t)=\epsilon\mathbf{H}^{\epsilon}(x,\epsilon t),$
and
$\displaystyle\mu=\epsilon\,\mu^{\epsilon},\;\;\;\lambda=\epsilon\,\lambda^{\epsilon}.$
As the analysis in [27], we use the transformation $p^{\epsilon}(x,\epsilon
t)=\underline{p}e^{\epsilon q^{\epsilon}(x,\epsilon t)}$ for some positive
constant $\underline{p}$. Under these changes of variables, the system
(1.11)–(1.14) becomes
$\displaystyle a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}q^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)q^{\epsilon})+\frac{1}{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon}=0,$ (1.16) $\displaystyle
r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})+\frac{1}{\epsilon}\nabla
q^{\epsilon}=({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}+{\rm
div}\Psi^{\epsilon},$ (1.17)
$\displaystyle\partial_{t}{\mathbf{H}}^{\epsilon}-{\rm
curl\,}({\mathbf{u}}^{\epsilon}\times\mathbf{H}^{\epsilon})=0,\quad{\rm
div}\mathbf{H}^{\epsilon}=0,$ (1.18) $\displaystyle
b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}S^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)S^{\epsilon})=\epsilon^{2}\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon},$
(1.19)
where we have used the abbreviations
$\Psi^{\epsilon}=2\mu^{\epsilon}\mathbb{D}({\mathbf{u}}^{\epsilon})+\lambda^{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon}\;\mathbf{I}_{d}$ and
$\displaystyle a^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$
$\displaystyle:=A(S^{\epsilon},\underline{p}e^{\epsilon
q^{\epsilon}})\underline{p}e^{\epsilon
q^{\epsilon}}=\frac{\underline{p}e^{\epsilon
q^{\epsilon}}}{R(S^{\epsilon},\underline{p}e^{\epsilon
q^{\epsilon}})}\cdot\frac{\partial R(S^{\epsilon},s)}{\partial
s}\Big{|}_{s=\underline{p}e^{\epsilon q^{\epsilon}}},$ (1.20) $\displaystyle
r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$
$\displaystyle:=\frac{R(S^{\epsilon},\underline{p}e^{\epsilon
q^{\epsilon}})}{\underline{p}e^{\epsilon q^{\epsilon}}},\quad
b^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon}):=R(S^{\epsilon},\epsilon
q^{\epsilon})\Theta(S^{\epsilon},\epsilon q^{\epsilon}).$ (1.21)
Formally, we obtain from (1.16) and (1.17) that $\nabla
q^{\epsilon}\rightarrow 0$ and ${\rm div}{\mathbf{u}}^{\epsilon}=0$ as
$\epsilon\rightarrow 0$. Applying the operator _curl_ to (1.17), using the
fact that ${\rm curl\,}\nabla=0$, and letting $\epsilon\rightarrow 0$ and
$\mu^{\epsilon}\rightarrow\mu>0$, we obtain
$\displaystyle{\rm
curl\,}\big{(}r(\bar{S},0)(\partial_{t}{\mathbf{v}}+{\mathbf{v}}\cdot\nabla{\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}-\mu\Delta{\mathbf{v}}\big{)}=0,$
where we have assumed that
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$
and $r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ converge to
$(\bar{S},0,{\mathbf{v}},\bar{\mathbf{H}})$ and $r(\bar{S},0)$ in some sense,
respectively. Finally, applying the identity
$\displaystyle{\rm
curl\,}({{\mathbf{u}}}\times{\mathbf{H}})={{\mathbf{u}}}({\rm
div}{\mathbf{H}})-{\mathbf{H}}({\rm
div}{{\mathbf{u}}})+({\mathbf{H}}\cdot\nabla){{\mathbf{u}}}-({{\mathbf{u}}}\cdot\nabla){\mathbf{H}},$
(1.22)
we expect to get the following incompressible non-isentropic MHD equations
$\displaystyle
r(\bar{S},0)(\partial_{t}{\mathbf{v}}+({\mathbf{v}}\cdot\nabla){\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}+\nabla\pi=\mu\Delta{\mathbf{v}},$
(1.23)
$\displaystyle\partial_{t}\bar{\mathbf{H}}+({{\mathbf{v}}}\cdot\nabla)\bar{\mathbf{H}}-(\bar{\mathbf{H}}\cdot\nabla){{\mathbf{v}}}=0,$
(1.24) $\displaystyle\partial_{t}\bar{S}+({\mathbf{v}}\cdot\nabla)\bar{S}=0,$
(1.25) $\displaystyle{\rm div}\,{\mathbf{v}}=0,\quad{\rm
div}\bar{\mathbf{H}}=0$ (1.26)
for some function $\pi$.
The aim of this paper is to establish the above limit process rigorously in
the whole space $\mathbb{R}^{d}$.
Before stating our main results, we review the previous related works. We
begin with the results for the Euler and Navier-Stokes equations. For well-
prepared initial data, Schochet [33] obtained the convergence of the
compressible non-isentropic Euler equations to the incompressible non-
isentropic Euler equations in a bounded domain for local smooth solutions. For
general initial data, Métivier and Schochet [27] proved rigorously the
incompressible limit of the compressible non-isentropic Euler equations in the
whole space ${\mathbb{R}}^{d}$. There are two key points in the article [27].
First, they obtained the uniform estimates in Sobolev norms for the acoustic
component of the solutions, which are propagated by a wave equation with
unknown variable coefficients. Second, they proved that the local energy of
the acoustic wave decays to zero in the whole space case. This approach was
extended to the non-isentropic Euler equations in the exterior domain and the
full Navier-Stokes equations in the whole space by Alazard in [1] and [2],
respectively, and to the dispersive Navier-Stokes equations by Levermore, Sun
and Trivisa [26]. For the spatially periodic case, Métivier and Schochet [28]
showed the incompressible limit of the one-dimensional non-isentropic Euler
equations with general data. Compared to the non-isentropic case, the
treatment of the propagation of oscillations in the isentropic case is simpler
and there are many works on this topic. For example, see Ukai [35], Asano [3],
Desjardins and Grenier [7] in the whole space; Isozaki [16, 17] on the
exterior domain; Iguchi [15] on the half space; Schochet [32] and Gallagher
[11] in a periodic domain; and Lions and Masmoudi [30], and Desjardins, et al.
[8] in a bounded domain. Recently, Jiang and Ou [22] investigated the
incompressible limit of the non-isentropic Navier-Stokes equations with zero
heat conductivity and well-prepared initial data in three-dimensional bounded
domains. The justification of the incompressible limit of the non-isentropic
Euler or Navier-Stokes equations with general initial data in a bounded domain
or a multi-dimensional periodic domain is still open. The interested reader
can refer to [5] for formal computations on the case of viscous polytropic
gases and [28, 4] for some analysis on the non-isentropic Euler equations in a
multi-dimensional periodic domain. For more results on the incompressible
limit of the Euler and Navier-Stokes equations, please see the monograph [9]
and the survey articles [6, 31, 34].
For the isentropic compressible MHD equations, the justification of the low
Mach limit was given in several aspects. In [23], Klainerman and Majda first
studied the incompressible limit of the isentropic compressible ideal MHD
equations in the spatially periodic case with well-prepared initial data.
Recently, the incompressible limit of the isentropic viscous (including both
viscosity and magnetic diffusivity) of compressible MHD equations with general
data was studied in [14, 18, 19]. In [14], Hu and Wang obtained the
convergence of weak solutions of the compressible viscous MHD equations in
bounded, spatially periodic domains and the whole space, respectively. In
[18], the authors employed the modulated energy method to verify the limit of
weak solutions of the compressible MHD equations in the torus to the strong
solution of the incompressible viscous or partial viscous MHD equations (the
shear viscosity coefficient is zero but the magnetic diffusion coefficient is
a positive constant). In [19], the authors obtained the convergence of weak
solutions of the viscous compressible MHD equations to the strong solution of
the ideal incompressible MHD equations in the whole space by using the
dispersion property of the wave equation if both the shear viscosity and the
magnetic diffusion coefficients go to zero.
For the full compressible MHD equations, the incompressible limit in the
framework of the so-called variational solutions was studied in [24, 25, 29].
Recently, the authors [20] justified rigourously the low Mach number limit of
classical solutions to the ideal or full compressible non-isentropic MHD
equations with small entropy or temperature variations. When the heat
conductivity and large temperature variations are present, the low Mach number
limit for the full compressible non-isentropic MHD equations justified in
[21]. We emphasize here that the arguments in [21] are different from the
present paper and depend essentially on positivity of fluid viscosity and
magnetic diffusivity coefficients.
As aforementioned, in this paper we want to establish rigorously the limit as
$\epsilon\to 0$ to the system (1.16)–(1.19) for $\mu_{\epsilon}\to\mu>0$. In
this case, the magnetic equation is purely hyperbolic due to the lack of
magnetic diffusivity. The first-order derivatives of $\mathbf{H}^{\epsilon}$
in the momentum equation and magnetic equation cannot be controlled. It is
very hard to study the system (1.16)–(1.19). To the author’s knowledge, there
is very few mathematical analysis on the system (1.16)–(1.19) with fixed or
unfixed $\epsilon$, even for the isentropic case. Our main idea is trying to
make full use of the fluid viscosities to control the higher order derivatives
of the magnetic field.
Now, we supplement the system (1.16)–(1.19) with initial conditions
$\displaystyle(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})|_{t=0}=(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})$
(1.27)
and state the main results as follows.
###### Theorem 1.1.
Let $s>d/2+2$ be an integer. Assume that $\mu^{\epsilon}\rightarrow\mu>0$ and
$\lambda^{\epsilon}\rightarrow\lambda$ as $\epsilon\to 0$. Suppose that the
initial data
$(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})$
satisfy
$\displaystyle\|(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})\|_{H^{s}({\mathbb{R}}^{d})}\leq
M_{0}.$ (1.28)
Then there exists a $T>0$ such that for any $\epsilon\in(0,1]$, the Cauchy
problem (1.16)–(1.19), (1.27) has a unique solution
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\in
C^{0}([0,T],H^{s}({\mathbb{R}}^{d}))$, and there exists a positive constant
$N$, depending only on $T$ and $M_{0}$, such that
$\displaystyle\|\big{(}S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon}\big{)}(t)\|_{H^{s}({\mathbb{R}}^{d})}\leq
N,\quad\forall\,t\in[0,T].$ (1.29)
Furthermore, if there exist positive constants $\underline{S}$, $N_{0}$ and
$\delta$ such that $S^{\epsilon}_{0}(x)$ satisfies
$|S^{\epsilon}_{0}(x)-\underline{S}\,\,|\leq{N}_{0}|x|^{-1-\delta},\quad|\nabla
S^{\epsilon}_{0}(x)|\leq N_{0}|x|^{-2-\delta},$ (1.30)
then the sequence of solutions
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$
converges weakly in $L^{\infty}(0,T;H^{s}({\mathbb{R}}^{d}))$ and strongly in
$L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}({\mathbb{R}}^{d}))$ for all
$s^{\prime}<s$ to a limit $(\bar{S},0,{\mathbf{v}},\bar{\mathbf{H}})$, where
$(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})$ is the unique solution in
$C([0,T],H^{s}({\mathbb{R}}^{d}))$ of (1.23)–(1.26) with initial data
$(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})|_{t=0}=(S_{0},{\mathbf{w}}_{0},\mathbf{H}_{0})$,
where ${\mathbf{w}}_{0}\in H^{s}({\mathbb{R}}^{d})$ is determined by
${\rm div}\,{\mathbf{w}}_{0}=0,\,\;{\rm
curl\,}(r(S_{0},0){\mathbf{w}}_{0})={\rm
curl\,}(r(S_{0},0){\mathbf{v}}_{0}),\;\,r(S_{0},0):=\lim_{\epsilon\rightarrow
0}r^{\epsilon}(S^{\epsilon}_{0},0).$ (1.31)
The function $\pi\in C([0,T]\times{\mathbb{R}}^{d})$ satisfies $\nabla\pi\in
C([0,T],H^{s-1}({\mathbb{R}}^{d}))$.
We briefly describe the strategy of the proof. The proof of Theorems 1.1
includes two main steps: the uniform estimates of the solutions, and the
convergence from the original scaling equations to the limiting ones. Once we
have established the uniform estimates (1.29) of the solutions in Theorems
1.1, the convergence of solutions is easily proved by using the local energy
decay theorem for fast waves in the whole space, which is shown by Métivier
and Schochet in [27]. Thus, the main task in the present paper is to obtain
the uniform estimates (1.29). For this purpose, we shall modify the approach
developed in [27]. In fact, due to the strong coupling of hydrodynamic motion
and magnetic field, and the lack of magnetic diffusivity, new difficulties
arise in obtaining the uniform estimates for the solutions to (1.16)–(1.19),
(1.27). First of all, when we perform the operator
$(\\{E^{\epsilon}\\}^{-1}L(\partial_{x}))^{\sigma}$ to the continuity and
momentum equations, or the operator curl to the momentum equations, one order
more spatial derivatives arise for the magnetic field, and this prevents us
from closing the energy estimates. Second, since the coefficients of the
acoustic wave equations depend on the entropy, we could not get the estimates
of $\|{\mathbf{u}}^{\epsilon}\|_{L^{2}(0,T;H^{s+1})}$ directly from the
system. The ideas to overcome these difficulties here are the following: We
transfer one spatial derivative from the magnetic field to the velocity with
the help of the special coupled way between magnetic field and fluid velocity.
Then, to control $\|{\mathbf{u}}^{\epsilon}\|_{L^{2}(0,T;H^{s+1})}$, we employ
a kind of Helmholtz decomposition of the velocity. Third, we make full use the
special structure of the magnetic field equation and the estimates on
${\mathbf{u}}$ to control $\|\mathbf{H}^{\epsilon}\|_{L^{\infty}(0,T;H^{s})}$.
We point out that our arguments in this paper can be modified slightly to the
case of the the compressible non-isentropic MHD equations with infinite
Reynolds number. We shall give a brief discussion in Section 5.
This paper is arranged as follows. In Section 2, we give notations, recall
basic facts, and present commutators estimates. In Section 3 we establish the
uniform boundeness of the solutions and prove the existence part of Theorem
1.1. In Section 4, we use the decay of the local energy to the acoustic wave
equations to prove the convergent part of Theorem 1.1. In the last section, we
consider the incompressible limit to the compressible non-isentropic MHD
equations with infinite Reynolds number.
## 2\. Preliminary
We give notations and recall basic facts which will be used frequently in the
proofs.
(1) We denote $\langle\cdot,\cdot\rangle$ the standard inner product in
$L^{2}({\mathbb{R}}^{d})$ with $\langle f,f\rangle=\|f\|^{2}$ and $H^{k}$ the
usual Sobolev space $W^{k,2}$ with norm $\|\cdot\|_{k}$, in particular,
$\|\cdot\|_{0}=\|\cdot\|$. The notation $\|(A_{1},\dots,A_{k})\|$ means the
summation of $\|A_{i}\|$ ($i=1,\cdots,k$), and it also applies to other norms.
For a multi-index $\alpha=(\alpha_{1},\dots,\alpha_{d})$, we denote
$D^{\alpha}=\partial^{\alpha_{1}}_{x_{1}}\dots\partial^{\alpha_{d}}_{x_{d}}$
and $|\alpha|=|\alpha_{1}|+\cdots+|\alpha_{d}|$. We also omit the spatial
domain ${\mathbb{R}}^{d}$ in integrals for convenience. We use the symbols $K$
or $C_{0}$ to denote the generic positive constants, and $C(\cdot)$ and
$\tilde{C}(\cdot)$ to denote the smooth functions, which may vary from line to
line.
(2) For a scalar function $f$ and vector functions $\mathbf{a}$, $\mathbf{b}$
and $\mathbf{c}$, we have the following basic vector identities:
$\displaystyle{\rm div}(\mathbf{a}\times\mathbf{b})$
$\displaystyle=\mathbf{b}\cdot{\rm curl\,}\mathbf{a}-\mathbf{a}\cdot{\rm
curl\,}\mathbf{b},$ (2.1) $\displaystyle\nabla(|\mathbf{a}|^{2})$
$\displaystyle=2(\mathbf{a}\cdot\nabla)\mathbf{a}+2\mathbf{a}\times{\rm
curl\,}\mathbf{a},$ (2.2) $\displaystyle{\rm curl\,}(f\mathbf{a})$
$\displaystyle=f\cdot{\rm curl\,}\mathbf{a}-\nabla f\times\mathbf{a},$ (2.3)
$\displaystyle{\rm curl\,}(\mathbf{a}\times\mathbf{b})$
$\displaystyle=(\mathbf{b}\cdot\nabla)\mathbf{a}-(\mathbf{a}\cdot\nabla)\mathbf{b}+\mathbf{a}({\rm
div}\mathbf{b})-\mathbf{b}({\rm div}\mathbf{a}),$ (2.4) $\displaystyle{\rm
div}\big{(}(\mathbf{a}\times\mathbf{b})\times\mathbf{c}\big{)}$
$\displaystyle=\mathbf{c}\cdot{\rm
curl\,}(\mathbf{a}\times\mathbf{b})-(\mathbf{a}\times\mathbf{b})\cdot{\rm
curl\,}\mathbf{c}.$ (2.5)
(3) We have the following well-known nonlinear estimates [12].
(i) Let $\alpha=(\alpha_{1},\alpha_{2},\alpha_{d})$ be a multi-index such that
$|\alpha|=k$. Then, for all $\sigma\geq 0$, and $f,g\in
H^{k+\sigma}({\mathbb{R}}^{d})$, there exists a generic constant $C_{0}$ such
that
$\displaystyle\|[f,\partial^{\alpha}]g\|_{H^{\sigma}}\leq$ $\displaystyle
C_{0}(\|f\|_{W^{1,\infty}}\|g\|_{H^{\sigma+k-1}}+\|f\|_{H^{\sigma+k}}\|g\|_{L^{\infty}}).$
(2.6)
(ii) For integers $k\geq 0$, $l\geq 0$, $k+l\leq\sigma$ and $\sigma>d/2$, the
product maps continuously $H^{\sigma-k}({\mathbb{R}}^{d})\times
H^{\sigma-l}({\mathbb{R}}^{d})$ to $H^{\sigma-k-l}({\mathbb{R}}^{d})$ and
$\displaystyle\|uv\|_{\sigma-k-l}\leq K\|u\|_{\sigma-k}\|v\|_{\sigma-l}.$
(2.7)
(iii) Let $\sigma>d/2$ be an integer. Assume that $F(u)$ is a smooth function
such that $F(0)=0$ and $u\in H^{\sigma}({\mathbb{R}}^{d})$, then $F(u)\in
H^{\sigma}({\mathbb{R}}^{d})$ and its norm is bounded by
$\displaystyle\|F(u)\|_{\sigma}\leq C(\|u\|_{\sigma})\|u\|_{\sigma},$ (2.8)
where $C(\cdot)$ is independent of $u$ and maps $[0,\infty)$ into
$[0,\infty)$.
## 3\. Uniform estimates
In this section and the first part of the next section we assume that
$\mu^{\epsilon}\equiv\mu>0$ and $\lambda^{\epsilon}\equiv\lambda$ for
simplicity of the presentation. The general case can be treated by a slight
modification in the arguments presented here.
In view of [27] and the classical local existence results obtained by Vol’pert
and Khudiaev [36] for hyperbolic-parabolic systems, the key point in the proof
of the existence part of Theorem 1.1 is to establish the uniform estimate
(1.29), which can be deduced from the following _a priori_ estimate.
###### Theorem 3.1.
For any $\epsilon>0$, let
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\in
C([0,T],H^{s}({\mathbb{R}}^{d}))$ be the solution to (1.16)–(1.19). Then there
exists an increasing $C(\cdot)$ from $[0,\infty)$ to $[0,\infty)$, such that
$\displaystyle\mathcal{M}_{\epsilon}(T)\leq
C_{0}+(T+\epsilon)C(\mathcal{M}_{\epsilon}(T)),$ (3.1)
where
$\displaystyle\mathcal{M}_{\epsilon}(T):=\,$
$\displaystyle{\mathcal{N}_{\epsilon}(T)}^{2}+\int_{0}^{T}\|{\mathbf{u}}^{\epsilon}\|_{s+1}^{2}d\tau,$
(3.2)
with
$\displaystyle\mathcal{N}_{\epsilon}(T):=\,$
$\displaystyle\sup_{t\in[0,T]}\|(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(t)\|_{s}.$
(3.3)
The remainder of this section is devoted to establishing (3.1). In the
calculations that follow, we always suppose that the assumptions in Theorem
1.1 hold. We consider a solution
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$ to
the problem (1.16)–(1.19), (1.27) on $C([0,T],H^{s}({\mathbb{R}}^{d}))$ with
initial data satisfying (1.28).
The main idea for proving the uniform estimate (3.1) is motivated by the work
[27] where the operator $(\\{E^{\epsilon}\\}^{-1}L(\partial_{x}))^{m}$ is
introduced to control the acoustic components of velocity for the Euler
equations. When the strong coupling of the fluid and magnetic filed is
present, however, the arguments in [27] cannot be directly applied to get a
uniform estimate of the acoustic parts due to lack of magnetic diffusion in
the magnetic equation. Instead, here we transfer one order spatial derivative
from $\mathbf{H}^{\epsilon}$ to ${\mathbf{u}}^{\epsilon}$, and then employ the
fluid viscosity to control higher derivatives. We remark that the reason that
these techniques work is due to the special structure of coupling between the
fluid and magnetic fields.
We begin with the estimate on the entropy $S^{\epsilon}$.
###### Lemma 3.2.
There exist a constant $C_{0}>0$ and a function $C(\cdot)$, independent of
$\epsilon$, such that for all $t\in(0,T]$,
$\displaystyle\|S^{\epsilon}(t)\|^{2}_{s}\leq
C_{0}+tC(\mathcal{M}_{\epsilon}(T))+\epsilon^{2}C(\mathcal{M}_{\epsilon}(T)).$
(3.4)
###### Proof.
For the multi-index $\alpha$ satisfying $|\alpha|\leq s-1$, denote
$f_{\alpha}=\partial_{x}^{\alpha}S^{\epsilon}$. In view of the positivity of
$b^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$, we deduce from (1.19)
that,
$\displaystyle\partial_{t}f_{\alpha}+({\mathbf{u}}^{\epsilon}\cdot\nabla)f_{\alpha}=g_{\alpha}+\epsilon^{2}h_{\alpha},$
(3.5)
where
$\displaystyle
g_{\alpha}=-[\partial_{x}^{\alpha},{\mathbf{u}}^{\epsilon}]\cdot\nabla
S^{\epsilon},\;\;h_{\alpha}=\epsilon^{2}\partial^{\alpha}_{x}\left(\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right).$
The commutator inequality (2.6) and Sobolev embedding theorem imply that
$\|g_{\alpha}\|\leq C(M_{\epsilon}(T))$. On the other hand, from the Sobolev
embedding theorem and the Moser-type inequality [23] we get
$\displaystyle\|h_{\alpha}\|$ $\displaystyle\leq
K\Big{(}\|(\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon})\|_{L^{\infty}}\Big{\|}D^{s}\Big{(}\frac{1}{b^{\epsilon}}\Big{)}\Big{\|}+\|D^{s}(\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon})\|\Big{\|}\frac{1}{b^{\epsilon}}\Big{\|}_{L^{\infty}}\Big{)}$
$\displaystyle\leq
C(\mathcal{N}_{\epsilon}(T))+C(\mathcal{N}_{\epsilon}(T))\|{\mathbf{u}}^{\epsilon}\|_{s+1}.$
Multiplying (3.5) by $f_{\alpha}$ and integrating over
$[0,t]\times{\mathbb{R}}^{d}$ with $t\leq T$, we obtain
$\displaystyle\|f_{\alpha}(t)\|^{2}\leq$
$\displaystyle\|f_{\alpha}(0)\|^{2}+\|\partial_{x}{\mathbf{u}}^{\epsilon}\|_{L^{\infty}((0,t)\times{\mathbb{R}}^{d})}\int^{t}_{0}\|f_{\alpha}(\tau)\|^{2}d\tau$
$\displaystyle+2\int^{t}_{0}\|g_{\alpha}(\tau)\|\,\|f_{\alpha}(\tau)\|d\tau+2\epsilon^{2}\int^{t}_{0}\|h_{\alpha}(\tau)\|\,\|f_{\alpha}(\tau)\|d\tau$
$\displaystyle\leq$ $\displaystyle
C_{0}+tC(\mathcal{M}_{\epsilon}(T))+\epsilon^{2}C(\mathcal{M}_{\epsilon}(T)),$
where we have used Young’s inequality and the embedding
$H^{\sigma}\hookrightarrow L^{\infty}$ for $\sigma>d/2$. The conclusion then
follows by adding up these estimates for all $|\alpha|\leq s-1$. ∎
The following $L^{2}$-bound of
$(q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$ can be obtained
directly using the energy method due to the skew-symmetry of the singular term
in the system and the special structure of coupling between the magnetic field
and fluid velocity. This $L^{2}$-bound is very important in our arguments,
since the induction analysis will be used to get the desired Sobolev
estimates.
###### Lemma 3.3.
There exist constants $C_{0}>0$ and $0<\xi_{1}<\mu$, and a function $C(\cdot)$
independent of $\epsilon$, such that for all $t\in[0,T]$,
$\displaystyle\|(q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(t)\|^{2}+\xi_{1}\int^{t}_{0}\|\nabla{\mathbf{u}}^{\epsilon}(\tau)\|^{2}d\tau\leq
C_{0}+tC(\mathcal{M}_{\epsilon}(T)).$ (3.6)
###### Proof.
Multiplying (1.16) by $q^{\epsilon}$, (1.17) by ${\mathbf{u}}^{\epsilon}$, and
(1.18) by $\mathbf{H}^{\epsilon}$, respectively, integrating over
${\mathbb{R}}^{d}$, and adding the resulting equations together, we obtain
$\displaystyle\langle
a^{\epsilon}\partial_{t}q^{\epsilon},q^{\epsilon}\rangle+\langle
r^{\epsilon}\partial_{t}{\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle+\langle\partial_{t}\mathbf{H}^{\epsilon},\mathbf{H}^{\epsilon}\rangle+\mu\|\nabla{\mathbf{u}}^{\epsilon}\|^{2}+(\mu+\lambda)\|{\rm
div}{\mathbf{u}}^{\epsilon}\|^{2}$ $\displaystyle\qquad+\langle
a^{\epsilon}({\mathbf{u}}^{\epsilon}\cdot\nabla)q^{\epsilon},q^{\epsilon}\rangle+\langle
r^{\epsilon}({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle$
$\displaystyle=\int[({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}]\cdot{\mathbf{u}}^{\epsilon}\;dx+\int{\rm
curl\,}({\mathbf{u}}^{\epsilon}\times\mathbf{H}^{\epsilon})\cdot\mathbf{H}^{\epsilon}\;dx.$
(3.7)
Here the singular terms involving $1/\epsilon$ are canceled. Using the
identity (1.7) and integrating by parts, we immediately obtain that
$\displaystyle\int[({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}]\cdot{\mathbf{u}}^{\epsilon}\;dx+\int{\rm
curl\,}({\mathbf{u}}^{\epsilon}\times\mathbf{H}^{\epsilon})\cdot\mathbf{H}^{\epsilon}\;dx=0.$
In view of the positivity and smoothness of
$a^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ and
$r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$, we get directly from
(1.16), (1.19) and (2.8) that
$\displaystyle\|\partial_{t}S^{\epsilon}\|_{s-1}\leq
C(\mathcal{N}_{\epsilon}(T)),\quad\|\epsilon\partial_{t}q^{\epsilon}\|_{s-1}\leq
C(\mathcal{N}_{\epsilon}(T)),$ (3.8)
while by the Sobolev embedding theorem, we find that
$\displaystyle\|(\partial_{t}a^{\epsilon},\partial_{t}r^{\epsilon})\|_{L^{\infty}}\leq\|(\partial_{t}a^{\epsilon},\partial_{t}r^{\epsilon})\|_{s-2}\leq
C(\mathcal{N}_{\epsilon}(T)).$ (3.9)
By the definition of $\mathcal{N}_{\epsilon}(T)$ and the Sobolev embedding
theorem, it is easy to see that
$\displaystyle\|(\nabla a^{\epsilon},\nabla r^{\epsilon})\|_{L^{\infty}}\leq
C(\mathcal{N}_{\epsilon}(T)).$
Since $\mu>0,2\mu+d\lambda>0$, there exists a positive constant $\kappa_{1}$
such that
$\displaystyle\mu\|\nabla{\mathbf{u}}^{\epsilon}\|^{2}+(\mu+\lambda)\|{\rm
div}{\mathbf{u}}^{\epsilon}\|^{2}\geq\kappa_{1}\|\nabla{\mathbf{u}}^{\epsilon}\|^{2}.$
Thus, from (3.7) we get that
$\displaystyle\langle a^{\epsilon}q^{\epsilon},q^{\epsilon}\rangle+$
$\displaystyle\langle
r^{\epsilon}{\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle+\langle\mathbf{H}^{\epsilon},\mathbf{H}^{\epsilon}\rangle+\kappa_{1}\int^{t}_{0}\|\nabla{\mathbf{u}}^{\epsilon}(\tau)\|^{2}d\tau$
$\displaystyle\leq$ $\displaystyle\big{\\{}\langle
a^{\epsilon}q^{\epsilon},q^{\epsilon}\rangle+\langle
r^{\epsilon}{\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle+\langle\mathbf{H}^{\epsilon},\mathbf{H}^{\epsilon}\rangle\big{\\}}\big{|}_{t=0}$
$\displaystyle+C(\mathcal{M}_{\epsilon}(T))\int^{t}_{0}\big{\\{}|q^{\epsilon}(\tau)|^{2}+|{\mathbf{u}}^{\epsilon}(\tau)|^{2}+|\mathbf{H}^{\epsilon}(\tau)|^{2}\big{\\}}d\tau.$
(3.10)
Moreover, we have
$\displaystyle\|q^{\epsilon}\|^{2}+\|{\mathbf{u}}^{\epsilon}\|^{2}$
$\displaystyle\leq\|(a^{\epsilon})^{-1}\|_{L^{\infty}}\langle
a^{\epsilon}q^{\epsilon},q^{\epsilon}\rangle+\|(r^{\epsilon})^{-1}\|_{L^{\infty}}\langle
r^{\epsilon}{\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle$
$\displaystyle\leq C_{0}(\langle
a^{\epsilon}q^{\epsilon},q^{\epsilon}\rangle+\langle
r^{\epsilon}{\mathbf{u}}^{\epsilon},{\mathbf{u}}^{\epsilon}\rangle),$ (3.11)
since $a^{\epsilon}$ and $r^{\epsilon}$ are uniformly bounded away from zero.
Applying Gronwall’s Lemma to (3), we conclude
$\displaystyle\|(q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(t)\|^{2}\leq
C_{0}\|(q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})\|^{2}\exp\\{tC(\mathcal{M}_{\epsilon}(T))\\}.$
Therefore, the estimate (3.6) follows from an elementary inequality
$\displaystyle e^{Ct}\leq 1+\tilde{C}t,\quad 0\leq t\leq T_{0},$ (3.12)
where $T_{0}$ is some fixed constant. ∎
Concerning the desired higher order estimates, we cannot directly get them by
differentiating the system as done in [27], since the coefficients
$a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon}),r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ and
$b^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ contain two scales
$S^{\epsilon}$ and $\epsilon q^{\epsilon}$. We shall adapt and modify the
techniques developed in [27] to derive the higher order estimates. Set
$\displaystyle E^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})=\left(\begin{array}[]{cc}a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})&0\\\ 0&r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})\mathbf{I}_{d}\\\
\end{array}\right),\quad{\mathbf{U}}^{\epsilon}=\left(\begin{array}[]{c}q^{\epsilon}\\\
{\mathbf{u}}^{\epsilon}\end{array}\right),$ $\displaystyle
L(\partial_{x})=\left(\begin{array}[]{cc}0&{\rm div}\\\ \nabla&0\\\
\end{array}\right),$
where $\mathbf{I}_{d}$ denotes the $d\times d$ unit matrix.
Let $L_{E^{\epsilon}}(\partial_{x})=\\{E^{\epsilon}\\}^{-1}L(\partial_{x})$
and $r_{0}(S^{\epsilon})=r^{\epsilon}(S^{\epsilon},0)$. Note that
$r_{0}(S^{\epsilon})$ is smooth, positive, and bounded away from zero with
respect to each $\epsilon$. First, using Lemma 3.2 and employing the same
analysis as in [27], we have
###### Lemma 3.4.
There exist constants $C_{1}>0$, $K>0$, and a function $C(\cdot)$, depending
only on $M_{0}$, such that for all $\sigma\in[1,\dots,s]$ and $t\in[0,T]$,
$\displaystyle\|{\mathbf{U}}^{\epsilon}\|_{\sigma}\leq
K\|L(\partial_{x}){{\mathbf{U}}}^{\epsilon}\|_{\sigma-1}+\tilde{C}\big{(}\|{\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}+\|{\mathbf{U}}^{\epsilon}\|_{\sigma-1}\big{)}$
(3.13)
and
$\displaystyle\|{\mathbf{U}}^{\epsilon}\|_{\sigma}\leq\tilde{C}\big{\\{}\|\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{U}}}^{\epsilon}\|_{0}+\|{\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}+\|{\mathbf{U}}^{\epsilon}\|_{\sigma-1}\big{\\}},$
(3.14)
where $\tilde{C}:=C_{1}+tC(\mathcal{M}_{\epsilon}(T))+\epsilon
C(\mathcal{M}_{\epsilon}(T))$.
We remark that the inequalities (3.13) and (3.14) are similar to the well
known Helmholtz decomposition, and the estimate on
$\|S^{\epsilon}(t)\|^{2}_{s}$ in Lemma 3.2 plays a key role in the proof of
Lemma 3.4.
Our next task is to derive a bound on
$\|\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{U}}}^{\epsilon}\|_{0}$
and $\|{\rm curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}$ by induction
arguments. Let
${\mathbf{W}}_{\sigma}^{\epsilon}=\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}(0,{\mathbf{u}}^{\epsilon})^{\top}$.
We first show the following estimate.
###### Lemma 3.5.
There exist a sufficiently small constant $\eta_{1}>0$ and two constants
$C_{0}>0$, $0<\xi_{2}<\mu$, and a function $C(\cdot)$ from $[0,\infty)$ to
$[0,\infty)$, independent of $\epsilon$, such that for all
$\sigma\in[1,\dots,s]$ and $t\in[0,T]$,
$\displaystyle\|\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{U}}}^{\epsilon}(t)\|^{2}$
$\displaystyle+\frac{\xi_{2}}{2}\int_{0}^{t}\|\nabla{\mathbf{W}}_{\sigma}^{\epsilon}\|^{2}(\tau)d\tau$
$\displaystyle\leq
C_{0}+tC(\mathcal{N}_{\epsilon}(T))+\eta_{1}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}(\tau)\|_{s+1}^{2}d\tau.$
(3.15)
###### Proof.
Let
${\mathbf{U}}^{\epsilon}_{\sigma}:=\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{\mathbf{U}}^{\epsilon}$,
$\sigma\in\\{0,\dots,s\\}.$ For simplicity, we set
$\mathcal{M}:=\mathcal{M}_{\epsilon}(T)$,
$\mathcal{N}:=\mathcal{N}_{\epsilon}(T)$, and
$E:=E^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$. The case $k=0$ is an
immediate consequence of Lemma 3.3. It is easy to verify that the operator
$L_{E}(\partial_{x})$ is bounded from $H^{k}$ to $H^{k-1}$ for
$k\in\\{1,\dots,s+1\\}$. Note that the equations (1.16), (1.17) can be written
as
$\displaystyle(\partial_{t}+{\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{U}}^{\epsilon}+\frac{1}{\epsilon}E^{-1}L(\partial_{x}){\mathbf{U}}^{\epsilon}=E^{-1}(\mathbf{J}^{\epsilon}+\mathbf{V}^{\epsilon})$
(3.16)
with
$\displaystyle\mathbf{J}^{\epsilon}=\left(\begin{array}[]{c}0\\\ ({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}\end{array}\right),\quad\mathbf{V}^{\epsilon}=\left(\begin{array}[]{c}0\\\
{\rm div}\Psi^{\epsilon}\end{array}\right).$
For $\sigma\geq 1$, we commute the operator $\\{L_{E}\\}^{\sigma}$ with (3.16)
and multiply the resulting equation by $E$ to infer that
$\displaystyle
E(\partial_{t}+{\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{U}}^{\epsilon}_{\sigma}+\frac{1}{\epsilon}L(\partial_{x}){\mathbf{U}}^{\epsilon}_{\sigma}=E(\mathbf{f}_{\sigma}+\mathbf{g}_{\sigma}+\mathbf{h}_{\sigma}),$
(3.17)
where
$\displaystyle\mathbf{f}_{\sigma}:=$
$\displaystyle[\partial_{t}+{\mathbf{u}}^{\epsilon}\cdot\nabla,\\{L_{E}\\}^{\sigma}]{\mathbf{U}}^{\epsilon},$
$\displaystyle\mathbf{g}_{\sigma}:=$
$\displaystyle\\{L_{E}\\}^{\sigma}(E^{-1}\mathbf{J}^{\epsilon}),$
$\displaystyle\mathbf{h}_{\sigma}:=$
$\displaystyle\\{L_{E}\\}^{\sigma}(E^{-1}\mathbf{V}^{\epsilon}).$
Multiplying (3.17) by ${\mathbf{U}}^{\epsilon}_{\sigma}$ and integrating over
$[0,t]\times{\mathbb{R}}^{d}$ with $t\leq T$, noticing that the singular terms
cancel out since $L(\partial_{x})$ is skew-adjoint, we use the inequalities
(3.8) and (3.9), and Cauchy-Schwarz’s inequality to deduce that
$\displaystyle\langle
E(t){\mathbf{U}}_{\sigma}^{\epsilon}(t),{\mathbf{U}}_{\sigma}^{\epsilon}(t)\rangle\leq$
$\displaystyle\langle
E(0){\mathbf{U}}_{\sigma}^{\epsilon}(0),{\mathbf{U}}_{\sigma}^{\epsilon}(0)\rangle+C(\mathcal{M})\int^{t}_{0}\|{\mathbf{U}}^{\epsilon}_{\sigma}(\tau)\|^{2}d\tau$
$\displaystyle+\int^{t}_{0}\|\mathbf{f}_{\sigma}(\tau)\|^{2}d\tau+2\int^{t}_{0}\int_{{\mathbb{R}}^{d}}(E(\mathbf{g}_{\sigma}+\mathbf{h}_{\sigma}){\mathbf{U}}_{\sigma}^{\epsilon})(\tau)d\tau.$
(3.18)
Following the proof process of Lemma 2.4 in [27], we obtain that
$\displaystyle\|\mathbf{f}_{k}(t)\|\leq C(\mathcal{N}_{\epsilon}(t)).$ (3.19)
Now we estimate the nonlinear term in (3) involving $\mathbf{g}_{\sigma}$. We
expand $\mathbf{g}_{\sigma}$ as follows
$\displaystyle\mathbf{g}_{\sigma}=$
$\displaystyle\,\sum^{d}_{i,j=1}\sum_{|\alpha|=\sigma+1}\\{E^{-1}\\}^{k+1}\partial_{x}^{\alpha}H_{i}^{\epsilon}H^{\epsilon}_{j}$
$\displaystyle+\sum^{d}_{i,j=1}\sum_{\Lambda_{1}}\sum_{\Lambda_{2}}\\{E^{-1}\\}^{l}\partial_{x}^{\beta_{1}}\\{E^{-1}\\}\cdots\partial_{x}^{\beta_{k}}\\{E^{-1}\\}\partial_{x}^{\gamma}H_{i}^{\epsilon}\partial_{x}^{\delta}H^{\epsilon}_{j}$
$\displaystyle:=$ $\displaystyle\,B_{1}+B_{2},$
where
$\displaystyle\Lambda_{1}=\\{(\beta_{1},\cdots,\beta_{k},\gamma,\delta)\big{|}|\beta_{1}|+\cdots+|\beta_{k}|+|\gamma|+|\delta|\leq
k+1,0<|\gamma|\leq k,|\delta|\leq k\\},$
$\displaystyle\Lambda_{2}=\\{l\big{|}l=k+1-(|\beta_{1}|+\cdots+|\beta_{k}|),(\beta_{1},\cdots,\beta_{k},0,0)\in\Lambda_{1}\\}.$
Since there is no magnetic diffusion in the system, we cannot deal with
directly the terms involving $B_{1}$. Instead, we transform one spatial
derivative to ${\mathbf{U}}_{\sigma}^{\epsilon}$. Integrating by parts, we
have
$\displaystyle\int_{{\mathbb{R}}^{d}}(EB_{1}{\mathbf{U}}_{\sigma}^{\epsilon})(\tau)dx=$
$\displaystyle-\sum^{d}_{i,j=1}\sum_{|\alpha|=\sigma}\int_{{\mathbb{R}}^{d}}\\{E^{-1}\\}^{k}\partial_{x}^{\alpha}H_{i}^{\epsilon}\partial_{x}H^{\epsilon}_{j}{\mathbf{U}}_{\sigma}^{\epsilon}dx$
$\displaystyle-\sum^{d}_{i,j=1}\sum_{|\alpha|=\sigma}\int_{{\mathbb{R}}^{d}}\partial_{x}\\{E^{-1}\\}^{k}\partial_{x}^{\alpha}H_{i}^{\epsilon}H^{\epsilon}_{j}{\mathbf{U}}_{\sigma}^{\epsilon}dx$
$\displaystyle-\sum^{d}_{i,j=1}\sum_{|\alpha|=\sigma}\int_{{\mathbb{R}}^{d}}\\{E^{-1}\\}^{k}\partial_{x}^{\alpha}H_{i}^{\epsilon}H^{\epsilon}_{j}\partial_{x}({\mathbf{U}}_{\sigma}^{\epsilon})dx$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{N})+\eta_{1}\|{\mathbf{u}}^{\epsilon}(\tau)\|_{s+1}^{2}$ (3.20)
for sufficiently small constant $\eta_{1}>0$.
By virtue of Cauchy-Schwarz’s and Sobolev’s inequalities, and (3.9), a direct
computation implies that
$\displaystyle\int_{{\mathbb{R}}^{d}}|(EB_{2}{\mathbf{U}}_{\sigma}^{\epsilon})(\tau)|d\tau\leq
C(\mathcal{N}).$ (3.21)
Next, we deal with the term involving the viscosity. Recall that
$L(\partial_{x}){\mathbf{U}}^{\epsilon}=({\rm
div}{\mathbf{u}}^{\epsilon},\nabla q^{\epsilon})$. Denote
$\displaystyle L_{1}:=\\{a^{\epsilon}\\}^{-1}{\rm
div},\;\;\;\;L_{2}:=\\{r^{\epsilon}\\}^{-1}\nabla.$
A straightforward computation implies that
$\displaystyle{\mathbf{U}}^{\epsilon}_{k}=\left\\{\begin{array}[]{ll}\left(\begin{array}[]{c}\\{L_{1}L_{2}\\}^{\frac{k-1}{2}}L_{1}{\mathbf{u}}^{\epsilon}\\\
\\{L_{2}L_{1}\\}^{\frac{k-1}{2}}L_{2}q^{\epsilon}\end{array}\right),&\text{if}\
k\ \text{is odd};\\\\[-4.30554pt] \\\
\left(\begin{array}[]{c}\\{L_{1}L_{2}\\}^{k/2}q^{\epsilon}\\\
\\{L_{2}L_{1}\\}^{k/2}{\mathbf{u}}^{\epsilon}\end{array}\right),&\hbox{\text{if}}\
k\ \text{is even}.\end{array}\right.$
Thus, we induce that
$\displaystyle\int^{t}_{0}\int_{{\mathbb{R}}^{d}}(E\mathbf{h}_{\sigma}{\mathbf{U}}_{\sigma}^{\epsilon})(\tau)dxd\tau=\int^{t}_{0}\int_{{\mathbb{R}}^{d}}E\\{L_{E}\\}^{\sigma}(E^{-1}\mathbf{V}^{\epsilon}){\mathbf{W}}_{\sigma}^{\epsilon}dxd\tau.$
An integration by parts gives
$\displaystyle\int^{t}_{0}\int_{{\mathbb{R}}^{d}}EL^{\sigma}_{E}(E^{-1}\mathbf{V}^{\epsilon}){\mathbf{W}}_{\sigma}^{\epsilon}dxd\tau$
$\displaystyle=$
$\displaystyle-\int^{t}_{0}\int_{{\mathbb{R}}^{d}}\mu|\nabla{\mathbf{W}}_{\sigma}^{\epsilon}|^{2}+(\mu+\lambda)|{\rm
div}{\mathbf{W}}_{\sigma}^{\epsilon}|^{2}dxd\tau$
$\displaystyle+\int^{t}_{0}\int_{{\mathbb{R}}^{d}}E[\mu
E^{-1}\Delta+(\mu+\lambda)E^{-1}\nabla{\rm
div},\\{L_{E}\\}^{\sigma}](0,{\mathbf{u}}^{\epsilon})^{T}{\mathbf{W}}_{\sigma}^{\epsilon}dxd\tau,$
where it is easy to verify that
$\displaystyle[E^{-1}\Delta,\\{L_{E}\\}^{\sigma}]=\sum_{i=0}^{k-1}\\{L_{E}\\}^{i}[E^{-1}\Delta,L_{E}]\\{L_{E}\\}^{\sigma-i-1}.$
(3.22)
Noting that $L_{E}(\partial_{x})={E}^{-1}L(\partial_{x})$, we find that
$\displaystyle[E^{-1}\Delta,L_{E}]=-E^{-1}\Delta
E^{-1}L(\partial_{x})+\sum_{i,j=1}^{d}B_{ij}\partial_{x_{ij}},$
where $B_{ij}$ $(i,j=1,\cdots,d$) are the sums of bilinear functions of
$E^{-1}$ and $\partial_{x}\\{E^{-1}\\}$, and Sobolev’s inequalities imply that
$\displaystyle\|B_{ij}\|_{s-1}\leq C(\mathcal{N}).$
Thus, we integrate by parts to infer that
$\displaystyle\mu\int^{t}_{0}\int_{{\mathbb{R}}^{d}}E[E^{-1}\Delta,\\{L_{E}\\}^{\sigma}](0,{\mathbf{u}}^{\epsilon})^{\top}{\mathbf{W}}_{\sigma}^{\epsilon}dxd\tau$
$\displaystyle\leq$
$\displaystyle\frac{\mu}{2}\int_{0}^{t}\|\nabla{\mathbf{W}}_{\sigma}^{\epsilon}\|^{2}d\tau+C(\mathcal{N})\int_{0}^{t}\|{\mathbf{W}}_{\sigma}^{\epsilon}(\tau)\|^{2}d\tau$
$\displaystyle+\int_{0}^{t}\|\tilde{H}_{2}^{\sigma}(\tau)\|^{2}+\|\tilde{H}_{1}^{\sigma}(\tau)\|^{2}d\tau.$
Here $\tilde{H}_{1}^{\sigma}$ is a finite sum of terms of the form
$(\partial^{\alpha_{1}}_{x}e_{1})\cdots(\partial^{\alpha_{l}}_{x}e_{l})\,(\partial^{\beta}_{x}w)\,(\partial^{\gamma}_{x}u^{\epsilon}_{m})$
with $|\alpha_{1}|+\cdots+|\alpha_{l}|+|\beta|+|\gamma|\leq\sigma\leq s$,
$|\gamma|>0$, and thus $|\beta|\leq k-1\leq s-1$, where $(e_{1},\dots,e_{l})$,
$w$ and $u^{\epsilon}_{m}$ denote the coefficients of $E^{-1}$, $C_{j}$ and
${\mathbf{u}}^{\epsilon}$ respectively, with $C_{j}$ taking a form similar to
that of $B_{ij}$. $\tilde{H}_{2}^{\sigma}$ is a finite sum of terms of the
form
$(\partial^{\alpha_{1}}_{x}e_{1})\cdots(\partial^{\alpha_{l}}_{x}e_{l})\,(\partial^{\beta}_{x}w)\,(\partial^{\gamma}_{x}u^{\epsilon}_{m})$
with $|\alpha_{1}|+\cdots+|\alpha_{l}|+|\beta|+|\gamma|\leq\sigma+1\leq s+1$,
$|\gamma|>1$, and thus $|\beta|\leq\sigma-1\leq s-1$, where
$(e_{1},\dots,e_{l})$, $w$ and $u^{\epsilon}_{m}$ denote the coefficients of
$E^{-1}$, $B_{ij}$ and ${\mathbf{u}}^{\epsilon}$ respectively. Hence, we have
$\displaystyle\|\tilde{H}_{1}^{\sigma}\|^{2}+\|\tilde{H}_{2}^{\sigma}\|^{2}\leq
C(\mathcal{M}).$
Similarly, we can show that
$\displaystyle(\mu+\lambda)\int^{t}_{0}\int_{{\mathbb{R}}^{d}}E[E^{-1}\nabla{\rm
div},\\{L_{E}\\}^{\sigma}](0,{\mathbf{u}}^{\epsilon})^{T}{\mathbf{W}}_{\sigma}^{\epsilon}dxd\tau$
$\displaystyle\qquad\quad\leq\frac{\mu+\lambda}{2}\int_{0}^{t}\|{\rm
div}{\mathbf{W}}_{\sigma}^{\epsilon}\|^{2}d\tau+C(\mathcal{N})\int_{0}^{t}\|{\mathbf{W}}_{\sigma}^{\epsilon}\|^{2}d\tau+tC(\mathcal{M}).$
Finally, the above estimates (3.19)–(3.21) and the positivity of $E$ imply
(3.5). ∎
Next, we derive an estimate for $\|{\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}$. Define
$\displaystyle f^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon}):=1-\frac{r_{0}(S^{\epsilon})}{r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}.$ (3.23)
Hereafter we denote $r_{0}(t):=r_{0}(S^{\epsilon}(t))$ and
$f^{\epsilon}(t):=f^{\epsilon}(S^{\epsilon}(t),\epsilon q^{\epsilon}(t))$ for
notational simplicity.
One can factor out $\epsilon q^{\epsilon}$ in $f^{\epsilon}(t)$. In fact,
using Taylor’s expansion, one obtains that there exists a smooth function
$g^{\epsilon}(t)$, such that
$\displaystyle f^{\epsilon}(t)=\epsilon g^{\epsilon}(t):=\epsilon
g^{\epsilon}(S^{\epsilon}(t),\epsilon
q^{\epsilon}(t)),\quad\|g^{\epsilon}(t)\|_{s}\leq
C(\mathcal{M}_{\epsilon}(T)).$ (3.24)
Since
$\partial_{t}S^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)S^{\epsilon}=\epsilon^{2}\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})},$
the equations for ${\mathbf{u}}^{\epsilon}$ are equivalent to
$\displaystyle[\partial_{t}+({\mathbf{u}}^{\epsilon}\cdot\nabla)](r_{0}{\mathbf{u}}^{\epsilon})+\frac{1}{\epsilon}\nabla
q^{\epsilon}=$ $\displaystyle g^{\epsilon}\nabla q^{\epsilon}+(1-\epsilon
g^{\epsilon})({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}$
$\displaystyle+(1-\epsilon g^{\epsilon}){\rm
div}\Psi^{\epsilon}+\epsilon^{2}\,r^{\prime}_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon}\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}.$ (3.25)
We perform the operator _curl_ to the equation (3) to obtain that
$\displaystyle[\partial_{t}+({\mathbf{u}}^{\epsilon}\cdot\nabla)]({\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon}))$ $\displaystyle=$
$\displaystyle[{\mathbf{u}}^{\epsilon}\cdot\nabla,{\rm
curl\,}](r_{0}{\mathbf{u}}^{\epsilon})+{\rm curl\,}[(1-\epsilon
g^{\epsilon})({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}]$
$\displaystyle+{\rm curl\,}[(1-\epsilon g^{\epsilon}){\rm
div}\Psi^{\epsilon}]+[{\rm curl\,},g^{\epsilon}]\nabla q^{\epsilon}$
$\displaystyle+\epsilon^{2}{\rm
curl\,}\left\\{r^{\prime}_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon}\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right\\}.$ (3.26)
###### Lemma 3.6.
There exist constants $C_{0}>0$, $0<\xi_{3}<\mu$, a function $C(\cdot)$ from
$[0,\infty)$ to $[0,\infty)$ and a sufficiently small constant $\eta_{2}>0$,
such that for all $\epsilon\in(0,1]$ and all $t\in[0,T]$,
$\displaystyle\|\\{{\rm curl\,}(r_{0}{\mathbf{u}}^{\epsilon}),{\rm
curl\,}\mathbf{H}^{\epsilon}\\}(t)\|^{2}_{s-1}$
$\displaystyle+\xi_{3}\int_{0}^{t}\|\nabla{\rm
curl\,}{\mathbf{u}}^{\epsilon}\|_{s-1}^{2}d\tau$ $\displaystyle\leq$
$\displaystyle
C_{0}+tC(\mathcal{N}_{\epsilon}(T))+\eta_{2}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|^{2}_{s+1}d\tau.$
(3.27)
###### Proof.
Set $\mathcal{N}:=\mathcal{N}_{\epsilon}(T)$ and $\omega={\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})$. Taking $\partial^{\alpha}_{x}$
$(|\alpha|\leq s-1)$ on (3), multiplying the resulting equations by
$\partial^{\alpha}_{x}\omega$, and integrating over
$[0,t]\times{\mathbb{R}}^{d}$ with $t\leq T$, we obtain
$\displaystyle\frac{1}{2}\|\partial^{\alpha}\omega(t)\|^{2}=\,$
$\displaystyle\frac{1}{2}\|\partial^{\alpha}\omega(0)\|^{2}-\int^{t}_{0}\langle({\mathbf{u}}^{\epsilon}\cdot\nabla)\partial^{\alpha}\omega,\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\langle[{\mathbf{u}}^{\epsilon}\cdot\nabla,\partial_{x}^{\alpha}]\omega,\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\langle\partial^{\alpha}\\{[{\rm
curl\,},g^{\epsilon}]\nabla
q^{\epsilon}\\},\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\langle\partial^{\alpha}\\{[{\mathbf{u}}^{\epsilon}\cdot\nabla,{\rm
curl\,}](r_{0}{\mathbf{u}}^{\epsilon})\\},\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\left\langle\partial^{\alpha}\left\\{\epsilon^{2}\,{\rm
curl\,}\left\\{r^{\prime}_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon}\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right\\}\right\\},\partial^{\alpha}\omega\right\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\langle\partial^{\alpha}\\{{\rm curl\,}[(1-\epsilon
g^{\epsilon})({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}]\\},\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle+\int^{t}_{0}\langle\partial^{\alpha}\\{{\rm curl\,}[(1-\epsilon
g^{\epsilon}){\rm
div}\Psi^{\epsilon}]\\},\partial^{\alpha}\omega\rangle(\tau)d\tau$
$\displaystyle=$
$\displaystyle:\frac{1}{2}\|\partial^{\alpha}\omega(0)\|^{2}+\int^{t}_{0}\sum_{i=1}^{7}I_{i}(\tau).$
(3.28)
We have to estimate the terms $I_{i}(\tau)$ ($1\leq i\leq 7$) on the right-
hand side of (3). Applying partial integrations, we have
$\displaystyle
I_{1}(\tau)=\int_{\mathbb{R}^{d}}|\partial^{\alpha}\omega|^{2}{\rm
div}{\mathbf{u}}^{\epsilon}\leq\|{\rm
div}{\mathbf{u}}^{\epsilon}(\tau)\|_{L^{\infty}}\|\partial^{\alpha}\omega(\tau)\|^{2}\leq
C(\mathcal{N})\|\partial^{\alpha}\omega(\tau)\|^{2},$ (3.29)
while for the term $I_{2}(\tau)$, an application of Cauchy-Schwarz’s
inequality gives
$\displaystyle|I_{2}(\tau)|\leq\|\partial^{\alpha}\omega\|\,\|\mathbf{h}_{\alpha}(\tau)\|,\quad\mathbf{h}_{\alpha}(\tau):=[{\mathbf{u}}^{\epsilon}\cdot\nabla,\partial_{x}^{\alpha}]\omega.$
The commutator $\mathbf{h}_{\alpha}$ is a sum of terms
$\partial^{\beta}_{x}{\mathbf{u}}^{\epsilon}\partial^{\gamma}_{x}\omega$ with
multi-indices $\beta$ and $\gamma$ satisfying $|\beta|+|\gamma|\leq s$,
$|\beta|>0$, and $|\gamma|>0$. Thus, the inequality (2.7) with
$\sigma=s-1>d/2$ implies that $\|\mathbf{h}_{\alpha}(\tau)\|\leq
C(\mathcal{N})$. Hence, we have
$\displaystyle|I_{2}(\tau)|\leq
C(\mathcal{N})+\|\partial^{\alpha}\omega(\tau)\|^{2}.$ (3.30)
Noting that $([{\rm
curl\,},g^{\epsilon}]\,\mathbf{a}\,)_{i,j}=a_{i}\partial_{x_{j}}g^{\epsilon}-a_{j}\partial_{x_{i}}g^{\epsilon}$
for $\mathbf{a}=(a_{1},\cdots,a_{d})$, the inequality (2.7), and the estimate
(3.24), we can control the term $I_{3}(\tau)$ as follows
$\displaystyle|I_{3}(\tau)|$ $\displaystyle\leq\|\partial^{\alpha}\\{[{\rm
curl\,},g^{\epsilon}]\nabla q^{\epsilon}\\}\|\;\|\partial^{\alpha}\omega\|$
$\displaystyle\leq K\|[{\rm curl\,},g^{\epsilon}]\nabla
q^{\epsilon}\|_{s-1}\,\|\partial^{\alpha}\omega\|$ $\displaystyle\leq
K\|\nabla g^{\epsilon}(\tau)\|_{s-1}\|\nabla
q^{\epsilon}(\tau)\|_{s-1}\,\|\partial^{\alpha}\omega\|$ $\displaystyle\leq
C(\mathcal{N})+\|\partial^{\alpha}\omega(\tau)\|^{2}.$ (3.31)
Similarly, the term $I_{4}(\tau)$ can be bounded as follows.
$\displaystyle|I_{4}(\tau)|$ $\displaystyle\leq
K\|\partial^{\alpha}\\{[{\mathbf{u}}^{\epsilon}\cdot\nabla,{\rm
curl\,}](r_{0}{\mathbf{u}}^{\epsilon})\\}\|\;\|\partial^{\alpha}\omega\|$
$\displaystyle\leq K\|[{\mathbf{u}}^{\epsilon}\cdot\nabla,{\rm
curl\,}](r_{0}{\mathbf{u}}^{\epsilon})\|_{s-1}\,\|\partial^{\alpha}\omega\|$
$\displaystyle\leq K\|[{\mathbf{u}}^{\epsilon}_{j},{\rm
curl\,}]\partial_{x_{j}}(r_{0}{\mathbf{u}}^{\epsilon})\|_{s-1}\,\|\partial^{\alpha}\omega\|$
$\displaystyle\leq C(\mathcal{N})+\|\partial^{\alpha}\omega(\tau)\|^{2}.$
(3.32)
To bound the term $I_{5}(\tau)$, we use the Moser-type inequality (see [23])
to deduce
$\displaystyle|I_{5}(\tau)|\leq$
$\displaystyle\epsilon^{2}K\left\|\partial^{\alpha}\left\\{{\rm
curl\,}\left\\{r^{\prime}_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon}\frac{\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right\\}\right\\}\right\|\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle=$
$\displaystyle\epsilon^{2}K\left\|\partial^{\alpha}\left[\left(\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right){\rm
curl\,}{\mathbf{u}}^{\epsilon}\right]\right\|\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle+\epsilon^{2}K\left\|\partial^{\alpha}\left[\nabla\left(\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\right)\times{\mathbf{u}}^{\epsilon}\right]\right\|\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle\leq$ $\displaystyle\epsilon^{2}K\|{\rm
curl\,}{\mathbf{u}}^{\epsilon}\|_{L^{\infty}}\left\|D^{s-1}\Big{(}\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})}\Big{)}\right\|\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle+\epsilon^{2}K\|D^{s-1}({\rm
curl\,}{\mathbf{u}}^{\epsilon})\|\left\|\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon^{2}q^{\epsilon})}\right\|_{L^{\infty}}\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle+\epsilon^{2}K\|{\mathbf{u}}^{\epsilon}\|_{L^{\infty}}\left\|D^{s}\Big{(}\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon^{2}q^{\epsilon})}\Big{)}\right\|\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle+\epsilon^{2}K\|D^{s-1}{\mathbf{u}}^{\epsilon}\|\left\|\nabla\left(\frac{r^{\prime}_{0}(S^{\epsilon})\Psi^{\epsilon}:\nabla{\mathbf{u}}^{\epsilon}}{b^{\epsilon}(S^{\epsilon},\epsilon^{2}q^{\epsilon})}\right)\right\|_{L^{\infty}}\cdot\|\partial^{\alpha}\omega\|$
$\displaystyle\leq$ $\displaystyle
C(\mathcal{N})+\epsilon^{2}C(\mathcal{N})\|{\mathbf{u}}^{\epsilon}\|_{s+1}^{2}+\|\partial^{\alpha}\omega(\tau)\|^{2},$
(3.33)
where the condition $s>2+d/2$ and the inequality (2.8) have been used. For the
term $I_{6}(\tau)$, by virtue of (2.1), ${\rm curl\,}{\rm
curl\,}\mathbf{a}=\nabla\,{\rm div}\,\mathbf{a}-\Delta\mathbf{a}$. Thus, we
integrate by parts to see that
$\displaystyle I_{6}(\tau)=$
$\displaystyle\langle\partial^{\alpha}\\{(1-\epsilon g^{\epsilon})({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}\\},\partial^{\alpha}\\{{\rm
curl\,}{\rm curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\\}\rangle$ $\displaystyle=$
$\displaystyle\langle\partial^{\alpha}\\{(1-\epsilon g^{\epsilon})({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}}\\},\partial^{\alpha}\\{\nabla{\rm
div}(r_{0}{\mathbf{u}}^{\epsilon})-\Delta(r_{0}{\mathbf{u}}^{\epsilon})\\}\rangle,$
and use Cauchy-Schwarz’s inequality and (2.8) to conclude
$\displaystyle|I_{6}(\tau)|\leq
C(\mathcal{N})+\theta_{1}\|{\mathbf{u}}^{\epsilon}(\tau)\|^{2}_{s+1},$ (3.34)
where $\theta_{1}>0$ is a sufficiently small constant independent of
$\epsilon$. Next, we deal with the term $I_{7}(\tau)$. By the vector
identities and integration by parts, we see that there exists a sufficiently
small $\theta_{2}$, such that
$\displaystyle I_{7}(\tau)\leq$
$\displaystyle-\inf\\{r_{0}(S^{\epsilon})\\}\|\nabla{\rm
curl\,}{\mathbf{u}}^{\epsilon}(\tau)\|_{\sigma-1}+C(\mathcal{N})$
$\displaystyle+\theta_{2}\|{\mathbf{u}}^{\epsilon}(\tau)\|_{s+1}+\epsilon
C(\mathcal{N})\|{\mathbf{u}}^{\epsilon}(\tau)\|_{s+1}.$ (3.35)
Finally, to estimate $\|{\rm curl\,}\mathbf{H}^{\epsilon}\|_{s-1}$, we apply
the operator _curl_ to (1.18) and use the vector identity (1.22) to obtain
$\displaystyle\partial_{t}({\rm
curl\,}\mathbf{H}^{\epsilon})+{\mathbf{u}}^{\epsilon}\cdot\nabla({\rm
curl\,}\mathbf{H}^{\epsilon})$ $\displaystyle=-[{\rm
curl\,},{\mathbf{u}}^{\epsilon}]\cdot\nabla\mathbf{H}^{\epsilon}+{\rm
curl\,}((\mathbf{H}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon}-\mathbf{H}^{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon}).$ (3.36)
By the commutator inequality and Sobolev’s inequalities, we find that
$\displaystyle\|[{\rm
curl\,},{\mathbf{u}}^{\epsilon}]\cdot\nabla\mathbf{H}^{\epsilon}\|\leq
C(\mathcal{N})$
and
$\displaystyle\|{\rm
curl\,}((\mathbf{H}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon}-\mathbf{H}^{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon})\|\leq
C(\mathcal{N})+\theta_{3}\|{\mathbf{u}}^{\epsilon}\|_{s+1}^{2}$
for sufficiently small constant $\theta_{3}>0$. Then, by arguments similar to
those used in Lemma 3.2, we derive that
$\displaystyle\|{\rm curl\,}\mathbf{H}^{\epsilon}\|_{s-1}^{2}\leq
C_{0}+tC(\mathcal{N})+\theta_{3}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}(\tau)\|_{s+1}^{2}d\tau.$
(3.37)
Thus, the lemma follows from adding up the estimates (3)–(3.37) for all
$|\alpha|\leq s-1$ and choosing constants $\theta_{1}$, $\theta_{2}$ and
$\theta_{3}$ appropriately small. ∎
Next we complete the proof of Theorem 3.1 by the following estimate.
###### Lemma 3.7.
There exist constants $C_{0}>0$, $0<\xi_{4}<\mu$, and a function $C(\cdot)$
from $[0,\infty)$ to $[0,\infty)$, such that for all $\epsilon\in(0,1]$ and
$t\in[0,T]$,
$\displaystyle\|(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})(t)\|^{2}_{s}+\xi_{4}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|^{2}d\tau\leq
C_{0}+(t+\epsilon)C(\mathcal{M}_{\epsilon}(T)).$ (3.38)
###### Proof.
First, from (3.14) we get
$\displaystyle\|{\mathbf{u}}^{\epsilon}\|_{\sigma+1}^{2}\leq\tilde{C}\Big{\\{}\|\nabla(\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{W}}}^{\epsilon}\|_{0}^{2}+\|\nabla{\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}+\|{\mathbf{u}}^{\epsilon}\|_{\sigma}^{2}\Big{\\}},$
(3.39)
where $\tilde{C}:=C_{1}+tC(\mathcal{M}_{\epsilon}(T))+\epsilon
C(\mathcal{M}_{\epsilon}(T))$. Moreover, using Lemma 3.2, we obtain
$\displaystyle\|{\mathbf{u}}^{\epsilon}\|_{\sigma+1}^{2}\leq\tilde{C}\Big{\\{}\|\nabla\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{W}}}^{\epsilon}\|_{0}^{2}+\|\nabla({\rm
curl\,}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}+\|{\mathbf{u}}^{\epsilon}\|_{\sigma}^{2}\Big{\\}}.$
(3.40)
In view of (3.40), there exists a constant $\kappa_{2}$ such that
$\displaystyle\frac{\xi_{2}}{2}\|\nabla{\mathbf{W}}_{\sigma}^{\epsilon}\|_{0}^{2}+\xi_{3}\|\nabla{\rm
curl\,}{\mathbf{u}}^{\epsilon}\|_{\sigma-1}^{2}\geq$
$\displaystyle\kappa_{2}\|{\mathbf{u}}^{\epsilon}\|_{\sigma+1}^{2}-\tilde{C}_{1}\Big{\\{}\|\nabla\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{W}}}^{\epsilon}\|_{0}^{2}$
$\displaystyle+\|\nabla({\rm
curl\,}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}\Big{\\}}-\tilde{C}\|{\mathbf{u}}^{\epsilon}\|_{\sigma}^{2},$
(3.41)
where $\tilde{C}_{1}=tC(\mathcal{M}_{\epsilon}(T))+\epsilon
C(\mathcal{M}_{\epsilon}(T))$. Now, we combine the estimates (3.5) and (3.6)
with (3.41), and use the fact that ${\rm div}\mathbf{H}=0$ to conclude that
there exists a positive constant $\kappa_{3}$, such that
$\displaystyle\|(L_{E}(\partial_{x}))^{\sigma}{\mathbf{U}}^{\epsilon}\|_{0}^{2}+\|{\rm
curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}+\kappa_{3}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|_{\sigma+1}^{2}dx$
$\displaystyle\leq$ $\displaystyle
C_{0}+tC(\mathcal{M}_{\epsilon}(T))+\tilde{C}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|_{\sigma}^{2}d\tau$
$\displaystyle+\tilde{C}_{1}\int_{0}^{t}\Big{\\{}\|\nabla\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{{\mathbf{W}}}^{\epsilon}\|_{0}^{2}+\|\nabla({\rm
curl\,}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}\Big{\\}}d\tau$
$\displaystyle\leq$ $\displaystyle
C_{0}+(t+\epsilon)C(\mathcal{M}_{\epsilon}(T))+\tilde{C}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|_{\sigma}^{2}d\tau$
for sufficiently small $\eta_{1}$ and $\eta_{2}$. Thus by induction, we
conclude that
$\displaystyle\|\\{L_{E^{\epsilon}}(\partial_{x})\\}^{\sigma}{\mathbf{U}}^{\epsilon}\|_{0}^{2}$
$\displaystyle+\|{\rm curl\,}(r_{0}{\mathbf{u}}^{\epsilon})\|_{\sigma-1}^{2}$
$\displaystyle+\kappa_{3}\int_{0}^{t}\|{\mathbf{u}}^{\epsilon}\|_{\sigma+1}^{2}d\tau\leq
C_{0}+(t+\epsilon)C(\mathcal{M}_{\epsilon}(T)).$
Using (3.14) again, we obtain the estimate (3.38) by induction on
$\sigma\in\\{0,\dots,s\\}$. ∎
## 4\. Incompressible limit
In this section, we shall prove the convergence part of Theorem 1.1 by
modifying the method developed by Métivier and Schochet [27], see also some
extensions [1, 2, 26].
###### Proof of the convergence part of Theorem 1.1.
The uniform bound (1.29) implies that, after extracting a subsequence, one
gets the following limits:
$\displaystyle(q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\rightharpoonup(q,{\mathbf{v}},\bar{\mathbf{H}})\quad\text{weakly-}\ast\
\text{in}\quad L^{\infty}(0,T;H^{s}(\mathbb{R}^{d})).$ (4.1)
The equations (1.18) and (1.19) imply that $\partial_{t}S^{\epsilon}$ and
$\partial_{t}\mathbf{H}^{\epsilon}\in C([0,T],H^{s-1}(\mathbb{R}^{d}))$. Thus,
after further extracting a subsequence, we obtain that, for all
$s^{\prime}<s$,
$\displaystyle S^{\epsilon}\rightarrow\bar{S}\quad\text{strongly in}\quad
C([0,T],H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d})),$ (4.2)
$\displaystyle\mathbf{H}^{\epsilon}\rightarrow\bar{\mathbf{H}}\quad\text{strongly
in}\quad C([0,T],H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d})),$ (4.3)
where the limit $\bar{\mathbf{H}}\in
C([0,T],H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))\cap
L^{\infty}(0,T;H^{s}_{\mathrm{loc}}(\mathbb{R}^{d}))$. Similarly, by (3) and
the uniform bound (1.29), we have
$\displaystyle{\rm
curl\,}(r_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon})\rightarrow{\rm
curl\,}(r_{0}(\bar{S}){\mathbf{v}})\quad\text{strongly in}\quad
C([0,T],H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{d}))$ (4.4)
for all $s^{\prime}<s$, where $r_{0}(\bar{S})=\lim_{\epsilon\rightarrow
0}r_{0}(S^{\epsilon}):=\lim_{\epsilon\rightarrow
0}r^{\epsilon}(S^{\epsilon},0)$.
In order to obtain the limit system, we need to prove that the convergence in
(4.1) holds in the strong topology of
$L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))$ for all
$s^{\prime}<s$. To this end, we first show that $q=0$ and ${\rm
div}{\mathbf{v}}=0$. In fact, from (3.16) we get
$\displaystyle\epsilon E^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})\partial_{t}{\mathbf{U}}^{\epsilon}+L(\partial_{x}){\mathbf{U}}^{\epsilon}=-\epsilon
E^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon}){\mathbf{u}}^{\epsilon}\cdot\nabla{\mathbf{U}}^{\epsilon}+\epsilon(\mathbf{J}^{\epsilon}+\mathbf{V}^{\epsilon}).$
(4.5)
Since
$\displaystyle E^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})-E^{\epsilon}(S^{\epsilon},0)=O(\epsilon),$
we have
$\displaystyle\epsilon
E^{\epsilon}(S^{\epsilon},0)\partial_{t}{\mathbf{U}}^{\epsilon}+L(\partial_{x}){\mathbf{U}}^{\epsilon}=\epsilon\mathbf{h}^{\epsilon},$
(4.6)
where $\mathbf{h}^{\epsilon}$ is uniformly bounded in
$C([0,T],H^{s-2}(\mathbb{R}^{d}))$ in view of (1.29). Passing to the weak
limit to (4.6), we obtain $\nabla q=0$ and ${\rm div}{\mathbf{v}}=0$. Since
$q\in L^{\infty}(0,T;H^{s}(\mathbb{R}^{d}))$, we infer that $q=0$. Noticing
that the strong compactness for the incompressible components by (4.4), it is
sufficient to prove the following proposition on the acoustic components.
###### Proposition 4.1.
Suppose that the assumptions in Theorem 1.1 hold, then $q^{\epsilon}$
converges strongly to $0$ in
$L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))$ for all
$s^{\prime}<s$, and ${\rm div}{\mathbf{u}}^{\epsilon}$ converges strongly to
$0$ in $L^{2}(0,T;H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{d}))$ for all
$s^{\prime}<s$.
The proof of Proposition 4.1 is built on the the following dispersive
estimates on the wave equations obtained by Métivier and Schochet [27] and
reformulated in [2].
###### Lemma 4.2 ([27, 2]).
Let $T>0$ and $w^{\epsilon}$ be a bounded sequence in
$C([0,T],H^{2}(\mathbb{R}^{d}))$, such that
$\displaystyle\epsilon^{2}\partial_{t}(a^{\epsilon}\partial_{t}w^{\epsilon})-\nabla\cdot(b^{\epsilon}\nabla
w^{\epsilon})=c^{\epsilon},$
where $c^{\epsilon}$ converges to $0$ strongly in
$L^{2}(0,T;L^{2}(\mathbb{R}^{d}))$. Assume further that, for some $s>d/2+1$,
the coefficients $(a^{\epsilon},b^{\epsilon})$ are uniformly bounded in
$C([0,T];H^{s}(\mathbb{R}^{d}))$ and converges in
$C([0,T];H^{s}_{\mathrm{loc}}(\mathbb{R}^{d}))$ to a limit $(a,b)$ satisfying
the decay estimate
$\displaystyle|a(x,t)-\underline{a}|\leq
C_{0}|x|^{-1-\delta},\quad|\nabla_{x}a(x,t)|\leq C_{0}|x|^{-2-\delta},$
$\displaystyle|b(x,t)-\underline{b}|\leq
C_{0}|x|^{-1-\delta},\quad|\nabla_{x}b(x,t)|\leq C_{0}|x|^{-2-\delta},$
for some given positive constants $\underline{a}$, $\underline{b}$, $C_{0}$
and $\delta$. Then the sequence $w^{\epsilon}$ converges to $0$ in
$L^{2}(0,T;L^{2}_{\mathrm{loc}}(\mathbb{R}^{d}))$.
###### Proof of Proposition 4.1.
We first show that $q^{\epsilon}$ converges strongly to $0$ in
$L^{2}(0,T;\linebreak H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))$ for all
$s^{\prime}<s$. An application of the operator $\epsilon^{2}\partial_{t}$ to
(1.16) gives
$\displaystyle\epsilon^{2}\partial_{t}(a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})\partial_{t}q^{\epsilon})+\epsilon\partial_{t}{\rm
div}{\mathbf{u}}^{\epsilon}=-\epsilon^{2}\partial_{t}\\{a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})({\mathbf{u}}^{\epsilon}\cdot\nabla)q^{\epsilon}\\}.$ (4.7)
Dividing (1.17) by $r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ and then
taking the operator _div_ to the resulting equation, one has
$\displaystyle\partial_{t}{\rm
div}{\mathbf{u}}^{\epsilon}+\frac{1}{\epsilon}{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}\nabla
q^{\epsilon}\Big{)}$ $\displaystyle\qquad\qquad=-{\rm
div}(({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})+{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}({\rm
curl\,}\mathbf{H}^{\epsilon})\times\mathbf{H}^{\epsilon}\Big{)}$
$\displaystyle\qquad\qquad\quad+{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}{\rm
div}\Psi({\mathbf{u}}^{\epsilon})\Big{)}.$ (4.8)
Subtracting (4) from (4.7), we obtain
$\displaystyle\epsilon^{2}\partial_{t}(a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})\partial_{t}q^{\epsilon})-{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}\nabla
q^{\epsilon}\Big{)}=F^{\epsilon}(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon}),$
(4.9)
where
$\displaystyle
F^{\epsilon}(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})=\,$
$\displaystyle\epsilon{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}({\rm
curl\,}\mathbf{H}^{\epsilon})\times\mathbf{H}^{\epsilon}\Big{)}$
$\displaystyle+\epsilon{\rm
div}\Big{(}\frac{1}{r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})}{\rm
div}\Psi({\mathbf{u}}^{\epsilon})\Big{)}-\epsilon{\rm
div}(({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})$
$\displaystyle-\epsilon^{2}\partial_{t}\\{a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})({\mathbf{u}}^{\epsilon}\cdot\nabla)q^{\epsilon}\\}.$
In view of the uniform boundedness of
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$,
the smoothness and positivity assumptions on $a^{\epsilon}$ and
$r^{\epsilon}$, and the convergence of $S^{\epsilon}$, we find that
$\displaystyle
F^{\epsilon}(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\rightarrow
0\quad\text{strongly in}\quad L^{2}(0,T;L^{2}(\mathbb{R}^{d})),$
and the coefficients in (4.9) satisfy the requirements in Lemma 4.2.
Therefore, by virtue of Lemma 4.2,
$\displaystyle q^{\epsilon}\rightarrow 0\quad\text{strongly in}\quad
L^{2}(0,T;L^{2}_{\mathrm{loc}}(\mathbb{R}^{d})).$
On the other hand, the uniform boundedness of $q^{\epsilon}$ in
$C([0,T],H^{s}(\mathbb{R}^{d}))$ and an interpolation argument yield that
$\displaystyle q^{\epsilon}\rightarrow 0\quad\text{strongly in}\quad
L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))\ \ \text{for all}\ \
s^{\prime}<s.$
Similarly, we can obtain the convergence of ${\rm div}u^{\epsilon}$. ∎
We continue our proof of Theorem 1.1. From Proposition 4.1, we know that
$\displaystyle{\rm div}\,{\mathbf{u}}^{\epsilon}\rightarrow{\rm
div}\,{\mathbf{v}}\quad\mathrm{in}\quad
L^{2}(0,T;H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{d})).$
Hence, from (4.4) it follows that
$\displaystyle{\mathbf{u}}^{\epsilon}\rightarrow{\mathbf{v}}\quad\mathrm{in}\quad
L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}(\mathbb{R}^{d}))\qquad\mbox{for all
}s^{\prime}<s.$
By (4.2), (4.3) and Proposition 4.2, we obtain
$\begin{array}[]{ccl}r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})\rightarrow
r_{0}(\bar{S})&\mathrm{in}&L^{\infty}(0,T;L^{\infty}(\mathbb{R}^{d}));\\\
\nabla{\mathbf{u}}^{\epsilon}\rightarrow\nabla{\mathbf{v}}&\mathrm{in}&L^{2}(0,T;H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{d}));\\\
\nabla\mathbf{H}^{\epsilon}\rightarrow\nabla\bar{\mathbf{H}}&\mathrm{in}&L^{2}(0,T;H^{s^{\prime}-1}_{\mathrm{loc}}(\mathbb{R}^{d})).\end{array}$
Passing to the limit in the equations for $S^{\epsilon}$ and
$\mathbf{H}^{\epsilon}$, we see that the limits $\bar{S}$ and
$\bar{\mathbf{H}}$ satisfy
$\displaystyle\partial_{t}\bar{S}+({\mathbf{v}}\cdot\nabla)\bar{S}=0,\quad\partial_{t}\bar{\mathbf{H}}+({{\mathbf{v}}}\cdot\nabla)\bar{\mathbf{H}}-(\bar{\mathbf{H}}\cdot\nabla){{\mathbf{v}}}=0$
in the sense of distributions. Since $r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})-r_{0}(S^{\epsilon})=O(\epsilon)$, we have
$\displaystyle(r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})-r_{0}(S^{\epsilon}))(\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})\rightarrow
0,$
whence,
$\displaystyle r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})=\,$
$\displaystyle(r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})-r_{0}(S^{\epsilon}))(\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})$
$\displaystyle+\partial_{t}(r_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon})+({\mathbf{u}}^{\epsilon}\cdot\nabla)(r_{0}(S^{\epsilon}){\mathbf{u}}^{\epsilon})$
$\displaystyle\rightarrow$
$\displaystyle\,r_{0}(\bar{S})(\partial_{t}{\mathbf{v}}+({\mathbf{v}}\cdot\nabla){\mathbf{v}})$
in the sense of distributions.
Applying the operator _curl_ to the momentum equation (1.17) and taking to the
limit, we find that
$\displaystyle{\rm
curl\,}\big{(}r_{0}(\bar{S})(\partial_{t}{\mathbf{v}}+{\mathbf{v}}\cdot\nabla{\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}-\mu\Delta{\mathbf{v}}\big{)}=0.$
Therefore, by the fact that ${\rm curl\,}\nabla=0$, the limit
$(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})$ satisfies
$\displaystyle
r(\bar{S},0)(\partial_{t}{\mathbf{v}}+({\mathbf{v}}\cdot\nabla){\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}-\mu\Delta{\mathbf{v}}+\nabla\pi=0,$
(4.10)
$\displaystyle\partial_{t}\bar{\mathbf{H}}+({{\mathbf{v}}}\cdot\nabla)\bar{\mathbf{H}}-(\bar{\mathbf{H}}\cdot\nabla){{\mathbf{v}}}=0,$
(4.11) $\displaystyle\partial_{t}\bar{S}+({\mathbf{v}}\cdot\nabla)\bar{S}=0,$
(4.12) $\displaystyle{\rm div}{\mathbf{v}}=0,\quad{\rm div}\bar{\mathbf{H}}=0$
(4.13)
for some function $\pi$.
If we employ the same arguments as in the proof of Theorem 1.5 in [27], we
find that $(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})$ satisfies the initial
conditions (1.31). Moreover, the standard iterative method shows that the
system (4.10)–(4.13) with initial data (1.31) has a unique solution
$(S^{*},{\mathbf{v}}^{*},\mathbf{H}^{*})\in C([0,T],H^{s}(\mathbb{R}^{d})).$
Thus, the uniqueness of solutions to the limit system (4.10)–(4.13) implies
that the above convergence results hold for the full sequence of
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$.
Thus, the proof is completed. ∎
## 5\. Compressible non-isentropic MHD equations with infinite Reynolds
number
In the study of magnetohydrodynamics, for some local processes in the cosmic
system, the effect of the magnetic diffusion will become very important, see
[13]. Moreover, when the Reynolds number of a fluid is very high and the
temperature changes very slowly, we can ignore the viscosity and the heat
conductivity of the fluid in the MHD equations. In such situation, the
compressible MHD equations in the non-isentropic case take the form:
$\displaystyle\partial_{t}\rho+{\rm div}(\rho{\mathbf{u}})=0,$ (5.1)
$\displaystyle\partial_{t}(\rho{\mathbf{u}})+{\rm
div}\left(\rho{\mathbf{u}}\otimes{\mathbf{u}}\right)+\nabla p=({\rm
curl\,}\mathbf{H})\times\mathbf{H},$ (5.2)
$\displaystyle\partial_{t}{\mathcal{E}}+{\rm
div}\left({\mathbf{u}}({\mathcal{E}}^{\prime}+p)\right)={\rm
div}\big{[}({\mathbf{u}}\times\mathbf{H})\times\mathbf{H}+\nu\mathbf{H}\times({\rm
curl\,}\mathbf{H})\big{]},$ (5.3) $\displaystyle\partial_{t}\mathbf{H}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})=-{\rm curl\,}(\nu\,{\rm
curl\,}\mathbf{H}),\quad{\rm div}\mathbf{H}=0.$ (5.4)
As before here $\rho$ denotes the density, ${\mathbf{u}}\in{\mathbb{R}}^{d}$
($d=2,3$) the velocity, $\mathbf{H}\in{\mathbb{R}}^{d}$ the magnetic field;
${\mathcal{E}}$ the total energy given by
${\mathcal{E}}={\mathcal{E}}^{\prime}+|\mathbf{H}|^{2}/2$ and
${\mathcal{E}}^{\prime}=\rho\left(e+|{\mathbf{u}}|^{2}/2\right)$ with $e$
being the internal energy, $\rho|{\mathbf{u}}|^{2}/2$ the kinetic energy, and
$|\mathbf{H}|^{2}/2$ the magnetic energy. The equations of state
$p=p(\rho,\theta)$ and $e=e(\rho,\theta)$ relate the pressure $p$ and the
internal energy $e$ to the density $\rho$ and the temperature $\theta$. The
constant $\nu>0$ is the magnetic diffusivity acting as a magnetic diffusion
coefficient of the magnetic field.
Using the Gibbs relation (1.9) and the identities (1.6) and (1.7), the
equation of energy conservation (5.3) can be replaced by
$\partial_{t}(\rho S)+{\rm div}(\rho S{\mathbf{u}})=\frac{\nu}{\theta}|{\rm
curl\,}\mathbf{H}|^{2},$ (5.5)
where $S$ denotes the entropy.
As in Section 1, we reconsider the equations of state as functions of $S$ and
$p$, i.e., $\rho=R(S,p)$ and $\theta=\Theta(S,p)$ for some positive smooth
functions $R$ and $\Theta$ defined for all $S$ and $p>0$, and satisfying
$\partial R/\partial p>0$. Then, by utilizing (1.1) together with the
constraint ${\rm div}{\mathbf{H}}=0$, the system (5.1), (5.2), (5.4) and (5.5)
can be rewritten as
$\displaystyle A(S,p)(\partial_{t}p+({\mathbf{u}}\cdot\nabla)p)+{\rm
div}{\mathbf{u}}=0,$ (5.6) $\displaystyle
R(S,p)(\partial_{t}{\mathbf{u}}+({\mathbf{u}}\cdot\nabla){\mathbf{u}})+\nabla
p=({\rm curl\,}\mathbf{H})\times\mathbf{H},$ (5.7)
$\displaystyle\partial_{t}{\mathbf{H}}-{\rm
curl\,}({\mathbf{u}}\times\mathbf{H})=-{\rm curl\,}(\nu\,{\rm
curl\,}\mathbf{H}),\quad{\rm div}\mathbf{H}=0,$ (5.8) $\displaystyle
R(S,p)\Theta(S,p)(\partial_{t}S+({\mathbf{u}}\cdot\nabla)S)={\nu}|{\rm
curl\,}\mathbf{H}|^{2},$ (5.9)
where $A(S,p)$ is defined by (1.15). By introducing the dimensionless
parameter $\epsilon$, and making the following changes of variables:
$\displaystyle p(x,t)=p^{\epsilon}(x,\epsilon t),\quad
S(x,t)=S^{\epsilon}(x,\epsilon t),$
$\displaystyle{{\mathbf{u}}}(x,t)=\epsilon{\mathbf{u}}^{\epsilon}(x,\epsilon
t),\;\;\;{\mathbf{H}}(x,t)=\epsilon\mathbf{H}^{\epsilon}(x,\epsilon
t),\;\;\;\nu=\epsilon\,\mu^{\epsilon},$
and $p^{\epsilon}(x,\epsilon t)=\underline{p}e^{\epsilon
q^{\epsilon}(x,\epsilon t)}$ for some positive constant $\underline{p}$, the
system (5.6)–(5.9) can be rewritten as
$\displaystyle a^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}q^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)q^{\epsilon})+\frac{1}{\epsilon}{\rm
div}{\mathbf{u}}^{\epsilon}=0,$ (5.10) $\displaystyle
r^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}{\mathbf{u}}^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla){\mathbf{u}}^{\epsilon})+\frac{1}{\epsilon}\nabla
q^{\epsilon}=({\rm
curl\,}{\mathbf{H}^{\epsilon}})\times{\mathbf{H}^{\epsilon}},$ (5.11)
$\displaystyle\partial_{t}{\mathbf{H}}^{\epsilon}-{\rm
curl\,}({\mathbf{u}}^{\epsilon}\times\mathbf{H}^{\epsilon})-\mu^{\epsilon}\Delta\mathbf{H}^{\epsilon}=0,\quad{\rm
div}\mathbf{H}^{\epsilon}=0,$ (5.12) $\displaystyle
b^{\epsilon}(S^{\epsilon},\epsilon
q^{\epsilon})(\partial_{t}S^{\epsilon}+({\mathbf{u}}^{\epsilon}\cdot\nabla)S^{\epsilon})=\epsilon^{2}{\mu}^{\epsilon}|{\rm
curl\,}\mathbf{H}^{\epsilon}|^{2},$ (5.13)
where we have used the identity ${\rm curl\,}{\rm
curl\,}\mathbf{H}^{\epsilon}=\nabla{\rm
div}\mathbf{H}^{\epsilon}-\Delta\mathbf{H}^{\epsilon},$ the constraint ${\rm
div}\mathbf{H}^{\epsilon}=0$, and the abbreviations (1.20) and (1.21).
Formally, we obtain from (5.10) and (5.11) that $\nabla
q^{\epsilon}\rightarrow 0$ and ${\rm div}{\mathbf{u}}^{\epsilon}=0$ as
$\epsilon\rightarrow 0$. Applying the operator _curl_ to (5.11), using the
fact that ${\rm curl\,}\nabla=0$, and letting $\epsilon\rightarrow 0$, we
obtain
$\displaystyle{\rm
curl\,}\big{(}r(\bar{S},0)(\partial_{t}{\mathbf{v}}+{\mathbf{v}}\cdot\nabla{\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}\big{)}=0,$
where we have assumed that
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$
and $r^{\epsilon}(S^{\epsilon},\epsilon q^{\epsilon})$ converge to
$(\bar{S},0,{\mathbf{v}},\bar{\mathbf{H}})$ and $r(\bar{S},0)$ in some sense,
respectively. Finally, Letting $\mu^{\epsilon}\rightarrow\mu>0$ and applying
the identity (1.22), we expect to get the following incompressible MHD
equations:
$\displaystyle
r(\bar{S},0)(\partial_{t}{\mathbf{v}}+({\mathbf{v}}\cdot\nabla){\mathbf{v}})-({\rm
curl\,}\bar{\mathbf{H}})\times\bar{\mathbf{H}}+\nabla\hat{\pi}=0,$ (5.14)
$\displaystyle\partial_{t}\bar{\mathbf{H}}+({{\mathbf{v}}}\cdot\nabla)\bar{\mathbf{H}}-(\bar{\mathbf{H}}\cdot\nabla){{\mathbf{v}}}-\mu\Delta\bar{\mathbf{H}}=0,$
(5.15) $\displaystyle\partial_{t}\bar{S}+({\mathbf{v}}\cdot\nabla)\bar{S}=0,$
(5.16) $\displaystyle{\rm div}\,{\mathbf{v}}=0,\quad{\rm
div}\bar{\mathbf{H}}=0$ (5.17)
for some function $\hat{\pi}$.
We supplement the system (1.16)–(1.19) with initial conditions
$\displaystyle(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})|_{t=0}=(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0}).$
(5.18)
The main result of this section reads as follows.
###### Theorem 5.1.
Let $s>d/2+2$ be an integer and $\mu^{\epsilon}\rightarrow\mu>0$. For any
constant $M_{0}>0$, there is a positive constant $T=T(M_{0})$, such that for
all $\epsilon\in(0,1]$ and any initial data
$(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})$
satisfying
$\displaystyle\|(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})\|_{H^{s}({\mathbb{R}}^{d})}\leq
M_{0},$ (5.19)
the Cauchy problem (5.10)–(5.13), (5.18) has a unique solution
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})\in
C^{0}([0,T],H^{s}({\mathbb{R}}^{d}))$, satisfying that for all
$\epsilon\in(0,1]$ and $t\in[0,T]$,
$\displaystyle\|\big{(}S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon}\big{)}(t)\|_{H^{s}({\mathbb{R}}^{d})}\leq
N\quad\mbox{for some constant }N=N(M_{0})>0.$ (5.20)
Moreover, suppose further that the initial data
$(S^{\epsilon}_{0},q^{\epsilon}_{0},{\mathbf{u}}^{\epsilon}_{0},\mathbf{H}^{\epsilon}_{0})$
converge strongly in $H^{s}({\mathbb{R}}^{d})$ to
$(S_{0},0,{\mathbf{v}}_{0},\mathbf{H}_{0})$ and $S^{\epsilon}_{0}$ decays
sufficiently rapidly at infinity in the sense that
$|S^{\epsilon}_{0}(x)-\underline{S}\,\,|\leq{N}_{0}|x|^{-1-\iota},\quad|\nabla
S^{\epsilon}_{0}(x)|\leq N_{0}|x|^{-2-\iota}$ (5.21)
for all $\epsilon\in(0,1]$ and some positive constants $\underline{S}$,
$N_{0}$ and $\iota$. Then,
$(S^{\epsilon},q^{\epsilon},{\mathbf{u}}^{\epsilon},\mathbf{H}^{\epsilon})$
converges weakly in $L^{\infty}(0,T;H^{s}({\mathbb{R}}^{d}))$ and strongly in
$L^{2}(0,T;H^{s^{\prime}}_{\mathrm{loc}}({\mathbb{R}}^{d}))$ to a limit
$(\bar{S},0,{\mathbf{v}},\bar{\mathbf{H}})$ for all $s^{\prime}<s$. Moreover,
$(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})$ is the unique solution in $C([0,T],$
$H^{s}({\mathbb{R}}^{d}))$ to the system (5.14)–(5.17) with initial data
$(\bar{S},{\mathbf{v}},\bar{\mathbf{H}})|_{t=0}=(S_{0},{\mathbf{w}}_{0},\mathbf{H}_{0})$,
where ${\mathbf{w}}_{0}\in H^{s}({\mathbb{R}}^{d})$ is determined by
${\rm div}\,{\mathbf{w}}_{0}=0,\,\;{\rm
curl\,}(r(S_{0},0){\mathbf{w}}_{0})={\rm
curl\,}(r(S_{0},0){\mathbf{v}}_{0}),\;\,r(S_{0},0):=\lim_{\epsilon\rightarrow
0}r^{\epsilon}(S^{\epsilon}_{0},0).$ (5.22)
The function $\hat{\pi}\in C([0,T]\times{\mathbb{R}}^{d})$ satisfies
$\nabla\hat{\pi}\in C([0,T],H^{s-1}({\mathbb{R}}^{d}))$.
###### Sketch of the proof of Theorem 5.1.
As explained before, the main step is to establish the uniform estimate
(5.20). For this purpose, we define $\mathcal{M}_{\epsilon}(T)$ as follows
$\displaystyle\mathcal{M}_{\epsilon}(T):=\,$
$\displaystyle{\mathcal{N}_{\epsilon}(T)}^{2}+\int_{0}^{T}\|\mathbf{H}^{\epsilon}\|_{s+1}^{2}d\tau,$
(5.23)
where $\mathcal{N}_{\epsilon}(T)$ is defined by (3.3). By arguments similar to
those used in the proof of Theorem 3.1, one can obtain the desired estimate.
Indeed, the arguments are easier since one can use the magnetic diffusion term
to control the terms involving $\mathbf{H}$ in the momentum equations, and
therefore we omit the details here. ∎
Acknowledgements: The authors would like to thank Prof. Fanghua Lin for
suggesting this problem and for helpful discussions. This work was partially
done when Li visited the Institute of Applied Physics and Computational
Mathematics in Beijing. He would like to thank the institute for hospitality.
Jiang was supported by the National Basic Research Program under the Grant
2011CB309705 and NSFC (Grant No. 40890154). Ju was supported by NSFC (Grant
No. 40890154, 11171035). Li was supported by NSFC (Grant No. 10971094), PAPD,
and the Fundamental Research Funds for the Central Universities.
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|
arxiv-papers
| 2011-11-12T13:45:42 |
2024-09-04T02:49:24.264027
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Song Jiang, Qiangchang Ju, and Fucai Li",
"submitter": "Fucai Li",
"url": "https://arxiv.org/abs/1111.2926"
}
|
1111.3083
|
11institutetext: Université de Lyon, Lyon, F-69003, France; Université Lyon 1,
Villeurbanne, F-69622, France; CNRS, UMR 5574, Centre de Recherche
Astrophysique de Lyon; École normale supérieure de Lyon, 46, allée d’Italie,
F-69364 Lyon Cedex 07, France
11email: [Guillaume.Laibe;Jean-Francois.Gonzalez]@ens-lyon.fr 22institutetext:
Centre for Stellar and Planetary Astrophysics, School of Mathematical
Sciences, Monash University, Clayton Vic 3168, Australia
22email: guillaume.laibe@monash.edu 33institutetext: Centre for Astrophysics
and Supercomputing, Swinburne University, PO Box 218, Hawthorn, VIC 3122,
Australia
33email: smaddison@swin.edu.au
# Revisiting the “radial-drift barrier” of planet formation and its relevance
in observed protoplanetary discs
G. Laibe 1122 J.-F. Gonzalez 11 S.T. Maddison 33
(Received 7 July 2010; Accepted ??)
###### Abstract
Context. To form metre-sized pre-planetesimals in protoplanetary discs,
growing grains have to decouple from the gas before they are accreted onto the
central star during their phase of fast radial migration and thus overcome the
so-called “radial-drift barrier” (often inaccurately referred to as the
“metre-size barrier”).
Aims. To predict the outcome of the radial motion of dust grains in
protoplanetary discs whose surface density and temperature follow power-law
profiles, with exponent $p$ and $q$ respectively. We investigate both the
Epstein and the Stokes drag regimes which govern the motion of the dust.
Methods. We analytically integrate the equations of motion obtained from
perturbation analysis. We compare these results with those from direct
numerical integration of the equations of motion. Then, using data from
observed discs, we predict the fate of dust grains in real discs.
Results. When a dust grain reaches the inner regions of the disc, the
acceleration due to the increase of the pressure gradient is counterbalanced
by the increase of the gas drag. We find that most grains in the Epstein
(resp. the Stokes) regime survive their radial migration if
$-p+q+\frac{1}{2}\leq 0$ (resp. if $q\leq\frac{2}{3}$). The majority of
observed discs satisfies both $-p+q+\frac{1}{2}\leq 0$ and $q\leq\frac{2}{3}$:
a large fraction of both their small and large grains remain in the disc, for
them the radial drift barrier does not exist.
###### Key Words.:
planetary systems: protoplanetary discs — methods: analytical
## 1 Introduction
Much of the information about the gas structure of protoplanetary discs is
inferred from the emission by the dust component and an assumed dust-to-gas
ratio. Interpretations of recent observations in the (sub)millimetre domain
(Andrews & Williams 2005, 2007; Lommen et al. 2007) show that observed discs
typically have masses between $10^{-4}$ and $10^{-1}\ M_{\odot}$ and a spatial
extent of a few hundred AU. Their radial surface density and temperature
profiles are approximated by power laws ($\Sigma\propto r^{-p}$,
$\mathcal{T}\propto r^{-q}$), whose respective exponents $p$ and $q$ have
positive values typically of order unity.
Seminal studies describe the dust motion in protoplanetary discs, which
depends strongly on the gas structure. Weidenschilling (1977a, hereafter W77)
and Nakagawa et al. (1986, hereafter NSH86) demonstrated that dust grains from
micron sizes to pre-planetesimals (a few metres in size) experience a radial
motion through protoplanetary discs. This motion is called radial drift or
migration. Due to its pressure gradient, the gas orbits the central star at a
sub-Keplerian velocity. Grains therefore have a differential velocity with
respect to the gas. The ensuing drag transfers linear and angular momentum
from the dust to the gas. Thus, dust particles can not sustain the Keplerian
motion they would have without the presence of gas and as a result migrate
toward the central star.
This migration motion depends strongly on the grain size, which sets the
magnitude of the drag, as well as the nature of the drag regime. Specifically,
as shown by W77 and NSH86, grains of a critical size pass through the disc in
a fraction of the disc lifetime. This catastrophic outcome is called the
“radial-drift barrier” of planet formation. More precisely, we will adopt the
subsequent definition for the “radial-drift barrier” in this study: “the
ability of grains of be accreted onto the central star/depleted from the disc
within its lifetime”. Historically, this process was first studied in a
Minimum Mass Solar Nebula (MMSN, see Weidenschilling 1977b; Hayashi 1981;
Desch 2007; Crida 2009), in which the critical size corresponds to metre-sized
bodies and thus was called the “metre-size barrier”. However, planets are
frequently observed (besides the 8 planets in our solar system, more than 700
extra-solar planets have been discovered to date111http://exoplanet.eu): some
solid material must therefore have overcome this barrier and stayed in the
disc to form larger bodies. Moreover, if the small grains of every disc were
submitted to the radial-drift barrier, we would barely detect them since their
emission via optical/IR scattering and IR thermal radiation is due to small
grains. As discs are frequently observed, the grains from which the emission
is detected cannot be strongly depleted for a substantial fraction of discs.
From a theoretical point of view, such a discrepancy between the observations
and the theoretical predictions imply that the seminal theory has to be
extended (some physical element is lacking) or that it has not been fully
exploited. This second option has been investigated by Youdin & Shu (2002,
hereafter YS02). They highlight the fact that, contrary to the primary
hypothesis of W77, observed dusty discs are drastically different from the
MMSN prototype. As the radial surface density and temperature profiles fix
both the radial pressure gradient and the magnitude of the gas drag, different
values for the power-law exponents $p$ and $q$ affect the optimal grain size
of migration and thus induce different radial motions for the dust through the
disc. Specifically, YS02 showed that for steep surface density profiles and
smooth temperature profiles, the grains radial velocity decreases when the
grains reach the inner discs regions. Grains in such discs therefore
experience a “pile-up”. However, while important, the work of YS02 does not
provide a precise conclusion on the outcome of the grains nor any quantitative
criterion for the “pile-up” process to be efficient enough to avoid the
radial-drift barrier. Furthermore, YS02 restricts their study to the special
case of a gas phase with a low density (e.g. the grain size smaller than the
gas mean free path, called the Epstein regime). This hypothesis is not valid
anymore when considering the radial drift of pre-planetesimals, whose grain
sizes are larger than the gas mean free path and are submitted to the Stokes
drag regime. Although the radial drift of pre-planetesimals has already been
studied in different situations with numerical or semi-analytical methods —
see e.g. Haghighipour & Boss (2003); Birnstiel et al. (2009); Youdin (2011) —
its rigorous theory for the standard case of a simple disc has not yet been
derived.
Within this context, we see that (i) the seminal theory describing the radial
motion of dust grains has been developed within the limits of the Epstein
regime but does not treat the Stokes regime ; (ii) here exists no clear
theoretical criterion to predict the impact of the “pile-up effect” on the
outcome of the dust radial motion ; (iii) there exists no criterion to predict
whether a given disc will be submitted to the “radial-drift barrier”
phenomenon. To answer these three points, we re-visit in this study the work
of W77 and NSH86 and extend the developments of YS02 for both the Epstein and
the Stokes regime. Performing rigorous perturbative expansions, we find two
theoretical criteria (one for each regime) which predict when the “pile-up”
effect is sufficient for the grains not to be accreted onto the central star.
We then test when these theoretical criteria can be applied in real discs.
Additionally, our work is motivated by the recent observational results of
Ricci et al. (2010a, b). From their observations they claim that “a mechanism
halting or slowing down the inward radial drift of solid particles is required
to explain the data”. In this work we aim to show that contrary to what is
usually invoked, local pressure maxima due to turbulent vortices or spiral
density waves may help but are not necessarily required to explain the
observations. Ricci et al. (2010a) also mention that “the observed flux of the
fainter discs are instead typically overpredicted even by more than an order
of magnitude”. Here, we also aim to provide a quantitative criterion to
determine which discs are faint and which one are not. Thus, revisiting the
seminal theory of the radial drift is timely, all the more so than an
important quantity of new data is soon to be provided by ALMA, the Atacama
Large Millimeter/submillimeter Array.
In this paper, we first recall some general properties of grain motion in
protoplanetary discs for both the Epstein and Stokes regime in Sect. 2. We
then focus on the radial motion of non-growing grains in the Epstein regime.
We expand the radial motion equations assuming a weak pressure gradient in
Sect. 3 and detail the two different modes of migration which grains may
experience in Sects. 3.1 and 3.2. This allows us to derive an analytic
criterion which determines the asymptotic dust behaviour in the Epstein regime
in Sect. 3.3. We transpose these derivations for the Stokes regime at low
Reynolds number in Sect. 4 and obtain a similar criterion for this regime. We
also discuss the grains outcome for large Reynolds numbers. In Sect. 5, we
discuss the relevance of these criteria and study their implications for
observed protoplanetary discs and planet formation in Sect. 6. Our conclusions
are presented in Sect. 7.
## 2 Dynamics of dust grains
To reduce the parameter space for this study, we assume the following:
1. 1.
The disc is a thin, non-magnetic, non-self-graviting, inviscid perfect gas
disc which is vertically isothermal. Its radial surface density and
temperature are described by power-law profiles. Notations are described in
Appendix A. The flow is laminar and in stationary equilibrium. Consequently,
the gas velocity and density are described by well-known relations, which we
present in Appendix B.
2. 2.
The grains are compact homogeneous spheres of fixed radius. The collisions
between grains and the collective effects due to large dust concentrations are
neglected. When the grains are small compared to the mean free path of the gas
($\lambda_{\mathrm{g}}>{4s}/{9}$, where $s$ is the grain size), their
interactions with the gas are treated by the Epstein drag force for diluted
media (Epstein 1924; Baines et al. 1965; Stepinski & Valageas 1996). This drag
is caused by the transfer of momentum by individual collisions with gas
molecules at the grains surface. Assuming specular reflections on the grain
and when the differential velocity with the gas is negligible compared to the
gas sound speed, the now common expression of the drag force is
$\left\\{\begin{array}[]{rcl}\mathbf{F}_{\mathrm{D}}&=&-\displaystyle\frac{m_{\mathrm{d}}}{t_{\mathrm{s}}}\,\Delta\mathbf{v}\\\\[10.00002pt]
t_{\mathrm{s}}&=&\displaystyle\frac{\rho_{\mathrm{d}}s}{\rho_{\mathrm{g}}c_{\mathrm{s}}}\,,\end{array}\right.$
(1)
where $m_{\mathrm{d}}$ is the dust grain’s mass, $t_{\mathrm{s}}$ the stopping
time, $\rho_{\mathrm{g}}$ the gas density, $c_{\mathrm{s}}$ the local gas
sound speed, $\rho_{\mathrm{d}}$ the intrinsic dust density, and
$\Delta\mathbf{v}=\mathbf{v}-\mathbf{v}_{\mathrm{g}}$ the differential
velocity between dust and the mean gas motion. In classical T Tauri star
(CTTS) protoplanetary discs, drag forces for particles smaller than $\sim 10$
m are well described by the Epstein regime (Garaud et al. 2004, see also Sect.
6.1). Small grains which produce the emission of observed protoplanetary discs
satisfy this criterion.
The interactions between large dust particles
($\lambda_{\mathrm{g}}<{4s}/{9}$) and the gas are treated by the Stokes drag
force (Whipple 1972; Stepinski & Valageas 1996). In this case, the gas mean
free path is small and the dust particle is locally surrounded by a viscous
fluid. Depending on the local Reynolds number of the flow around the grains
$R_{\mathrm{g}}\leavevmode\nobreak\ =\leavevmode\nobreak\
\frac{2s|\Delta\mathbf{v}|}{\nu}$, where $\nu$ is the microscopic kinematic
viscosity of the gas, the drag force takes the following expression:
$\textbf{F}_{\mathrm{D}}=-\frac{1}{2}C_{\mathrm{D}}\pi
s^{2}\rho_{\mathrm{g}}\left|\Delta\mathbf{v}\right|\Delta\mathbf{v},$ (2)
where the drag coefficient $C_{\mathrm{D}}$ is given by
$C_{\mathrm{D}}=\left\\{\begin{array}[]{ll}24R_{\mathrm{g}}^{-1}&\mathrm{for}\
R_{\mathrm{g}}<1\\\\[10.00002pt] 24R_{\mathrm{g}}^{-0.6}&\mathrm{for}\
1<R_{\mathrm{g}}<800\\\\[10.00002pt] 0.44&\mathrm{for}\
800<R_{\mathrm{g}}\,.\end{array}\right.$ (3)
If $R_{\mathrm{g}}<1$, the drag force remains linear in $\Delta\mathbf{v}$.
In this work, the physical relations are written in cylindrical coordinates
($r$, $\theta$, $z$). The related unit vector system is given by
$\left(\textbf{e}_{r},\textbf{e}_{\theta},\textbf{e}_{z}\right)$. As the
system is invariant by rotation around the vertical axis $\textbf{e}_{z}$, the
physical quantities depend only on $r$ and $z$. The physical quantities of the
gas, designated by subscript g, are first determined in a general way. Then,
the limit $z=0$ is taken to study the restricted radial motion.
Dust dynamics depends on both the magnitude of the drag (driven by the
differential velocity) and on its relative contribution with respect to the
gravity of the central star. Seminal studies of dust dynamics were conducted
by Whipple (1972), W77, Weidenschilling (1980) and NSH86, and extended by
others (YS02; Takeuchi & Lin 2002; Haghighipour & Boss 2003; Garaud et al.
2004; Youdin & Chiang 2004). Here we recall the major points of those studies.
We consider two forces acting on the grain: the gravity of the central star
and gas drag. (We assume that the momentum transferred by drag from a single
grain on the gas phase is negligible.) Thus
$m_{\mathrm{d}}\frac{\mathrm{d}\textbf{v}}{\mathrm{d}t}=-\textbf{F}_{\mathrm{D}}+m_{\mathrm{d}}\textbf{g},$
(4)
where $\textbf{F}_{\mathrm{D}}$ is the drag force. As shown by Eqs. (1)–(2), a
general expression of the ratio
$\frac{\textbf{F}_{\mathrm{D}}}{m_{\mathrm{d}}}$ is of the form
$\frac{\textbf{F}_{\mathrm{D}}}{m_{\mathrm{d}}}=-\frac{\tilde{\mathcal{C}}\left(r,z\right)}{s^{y}}|\textbf{v}-\textbf{v}_{\mathrm{g}}|^{\lambda}\left(\textbf{v}-\textbf{v}_{\mathrm{g}}\right),$
(5)
where the quantities $\tilde{\mathcal{C}}$, $y$ and $\lambda$ are defined for
both the Epstein and the Stokes regime in Appendix C.
## 3 Radial motion in the Epstein regime: perturbation analysis at small
pressure gradients
Considering the Epstein (small grains) regime, Eq. (4) reduces to:
$m_{\mathrm{d}}\frac{\mathrm{d}\textbf{v}}{\mathrm{d}t}=-\frac{m_{\mathrm{d}}}{t_{\mathrm{s}}}\left(\textbf{v}-\textbf{v}_{\mathrm{g}}\right)+m_{\mathrm{d}}\textbf{g}.$
(6)
Writing Eq. (6) in $\left(r,\theta,z\right)$ coordinates leads to
$\left\\{\begin{array}[]{rcl}\displaystyle\frac{\mathrm{d}v_{r}}{\mathrm{d}t}-\frac{v_{\theta}^{2}}{r}+\frac{\left(v_{r}-v_{\mathrm{g}r}\right)}{t_{\mathrm{s}}}+\frac{\mathcal{G}Mr}{\left(z^{2}+r^{2}\right)^{3/2}}&=&\displaystyle
0\\\
\displaystyle\frac{\mathrm{d}v_{\theta}}{\mathrm{d}t}+\frac{v_{r}v_{\theta}}{r}+\frac{\left(v_{\theta}-v_{\mathrm{g}\theta}\right)}{t_{\mathrm{s}}}&=&\displaystyle
0\\\
\displaystyle\frac{\mathrm{d}v_{z}}{\mathrm{d}t}+\frac{\left(v_{z}-v_{\mathrm{g}z}\right)}{t_{\mathrm{s}}}+\frac{\mathcal{G}Mz}{\left(z^{2}+r^{2}\right)^{3/2}}&=&\displaystyle
0.\end{array}\right.$ (7)
To highlight the important parameters involved in the grains dynamics, we
introduce dimensionless quantities (see Appendix C). It is crucial to note
that the ratio $\frac{t_{\mathrm{s}}}{t_{\mathrm{k}}}$ of the two timescales
related to the physical processes acting on the grain is given by
$\frac{t_{\mathrm{s}}}{t_{\mathrm{k}}}=\frac{t_{\mathrm{s}0}}{t_{\mathrm{k}0}}R^{p}\mathrm{e}^{\frac{Z^{2}}{2R^{3-q}}}=S_{0}R^{p}\mathrm{e}^{\frac{Z^{2}}{2R^{3-q}}}.$
(8)
With Eq. (1) and noting $\Omega_{\mathrm{k}}$ the Keplerian angular velocity,
this ratio can be written as
$\frac{t_{\mathrm{s}}}{t_{\mathrm{k}}}=\frac{s}{\left(\frac{\rho_{\mathrm{g}}c_{\mathrm{s}}}{\rho_{\mathrm{d}}\Omega_{\mathrm{k}}}\right)}=\frac{s}{s_{\mathrm{opt}}}=S,$
(9)
where
$s_{\mathrm{opt}}=\frac{\rho_{\mathrm{g}}c_{\mathrm{s}}}{\rho_{\mathrm{d}}\Omega_{\mathrm{k}}}$.
This timescale ratio therefore corresponds to a dimensionless size
$S=S_{0}R^{p}e^{\frac{Z^{2}}{2R^{3-q}}}$ for the grain. If $S\ll 1$ (resp.
$S\gg 1$), the effects of drag will occur much faster (resp. slower) than
gravitational effects. If $S\simeq 1$, both gravity and drag will act on the
same timescale. Interestingly, $s_{\mathrm{opt}}$ varies in the disc midplane
as $r^{-p}$, as does surface density.
Then using Eqs. (8), (51) and (81), we obtain for
($\textbf{e}_{r},\textbf{e}_{\theta},\textbf{e}_{z}$):
$\left\\{\begin{array}[]{rcl}\displaystyle\frac{\mathrm{d}\tilde{v}_{r}}{\mathrm{d}T}-\frac{\tilde{v}_{\theta}^{2}}{R}+\frac{\tilde{v}_{r}}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}\mathrm{e}^{-\frac{Z^{2}}{2R^{3-q}}}+\frac{R}{\left(R^{2}+\phi_{0}Z^{2}\right)^{3/2}}&=&0\\\\[8.61108pt]
\lx@intercol\displaystyle\frac{\mathrm{d}\tilde{v}_{\theta}}{\mathrm{d}T}+\frac{\tilde{v}_{\theta}\tilde{v}_{r}}{R}+\hfil\lx@intercol&&\\\
\frac{\left(\tilde{v}_{\theta}-\sqrt{\frac{1}{R}-\eta_{0}R^{-q}-q\left(\frac{1}{R}-\frac{1}{\sqrt{R^{2}+\phi_{0}Z^{2}}}\right)}\right)}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}\mathrm{e}^{-\frac{Z^{2}}{2R^{3-q}}}&=&0\\\\[8.61108pt]
\displaystyle\frac{\mathrm{d}^{2}Z}{\mathrm{d}T^{2}}+\frac{1}{S_{0}}\frac{\mathrm{d}Z}{\mathrm{d}T}R^{-\left(p+\frac{3}{2}\right)}\mathrm{e}^{-\frac{Z^{2}}{2R^{3-q}}}+\frac{Z}{\left(R^{2}+\phi_{0}Z^{2}\right)^{3/2}}&=&0.\end{array}\right.$
(10)
These equations depend on five control parameters: the initial dimensionless
grain size, $S_{0}$, the radial surface density and temperature exponents, $p$
and $q$, the square of the disc aspect ratio,
$\phi_{0}=\left(H_{0}/r_{0}\right)^{2}$, and the subkeplerian parameter,
$\eta_{0}$, given by Eq. (82). The equations can be simplified in some cases,
e.g. if the vertical motion is considered to occur faster than the radial
motion, $R\simeq 1$ and $\frac{\mathrm{d}^{2}Z}{\mathrm{d}T^{2}}$ simplifies
to the damped harmonic oscillator equation. If we consider only the radial
motion (for a 2D disc), we have $Z=0$, and
$\left\\{\begin{array}[]{l}\displaystyle\frac{\mathrm{d}\tilde{v}_{r}}{\mathrm{d}T}=\displaystyle\frac{\tilde{v}_{\theta}^{2}}{R}-\frac{\tilde{v}_{r}}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}-\frac{1}{R^{2}}\\\
\displaystyle\frac{\mathrm{d}\tilde{v}_{\theta}}{\mathrm{d}T}=\displaystyle-\frac{\tilde{v}_{\theta}\tilde{v}_{r}}{R}-\frac{\left(\tilde{v}_{\theta}-\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}\right)}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}.\end{array}\right.$
(11)
Even for discs in two dimensions, Eq. (11) is not analytically tractable.
However, as some of the parameters involved in the equation are small,
approximations of the solution can be found by performing perturbative
expansions. Some of the classical results detailed below have been studied in
W77 and NSH86, but are here properly justified. The principle of those
expansions is described on Fig. 1.
Figure 1: Principle of the various perturbative expansions for the grain
radial motion. Expanding first with respect to the small pressure gradient
($\eta_{0}R^{-q+1}\ll 1$) leads to NSH86 equations. Expanding first with
respect to the grain sizes ($S_{0}R^{p}\ll 1$ or $S_{0}R^{p}\gg 1$) leads to
W77 expressions for the particular case $p=0$. Combining both leads to the A-
and B- mode, respectively for small and large grains.
At $t=0$, $r=r_{0}$ which implies that $R\left(T=0\right)=1$. Because of gas
drag, a grain dissipates both its energy and angular momentum and therefore,
experiences a radial inward motion, i.e. $R<1$. The first parameter with
respect to which a perturbative expansion can be performed is $\eta_{0}$
(linked to the pressure gradient by Eq. (83)) as it takes values of
approximately $10^{-3}$–$10^{-2}$ in real protoplanetary discs (see NSH86),
and thus $\eta_{0}\ll 1$. We consider that this inequality also implies that
$\eta_{0}R^{-q}\ll\frac{1}{R}.$ (12)
This inequality is always verified when $q\leq 1$ and thus applies to observed
discs (see Sect. 6). For $q>1$, there is a region where this inequality is not
verified. However, in this case, the pressure gradient has the same order of
magnitude as the gravity of the central star and the model of a power-law
profile for the radial temperature is not accurate enough to model realistic
discs. We thus consider that for real discs Eq. (12) is always justified.
Then, following NSH86, we consider the system of equations given by Eq. (11).
We set
$\left\\{\begin{array}[]{l}\displaystyle\tilde{v}_{r}=\displaystyle\tilde{v}_{r0}+\eta_{0}\tilde{v}_{r1}+\mathcal{O}\left(\eta_{0}^{2}\right)\\\
\displaystyle\tilde{v}_{\theta}=\displaystyle\tilde{v}_{\theta
0}+\eta_{0}\tilde{v}_{\theta
1}+\mathcal{O}\left(\eta_{0}^{2}\right)\end{array}\right.$ (13)
and look at the orders $\mathcal{O}\left(1\right)$,
$\mathcal{O}\left(\eta_{0}\right)$,… of the expansion – see Appendix E. We
find that:
$\tilde{v}_{r}=\eta_{0}\tilde{v}_{r1}+\mathcal{O}\left(\eta_{0}^{2}\right)=-\frac{2S_{0}R^{p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)}{1+R^{2p}S_{0}^{2}}+\mathcal{O}\left(\eta_{0}^{2}\right).$
(14)
The pressure gradient term has been retained to keep the generality, however
since we assume that $\eta_{0}\ll 1$, we equivalently have
$\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}-\sqrt{\frac{1}{R}}=-\frac{\eta_{0}}{2}R^{-q+\frac{1}{2}}+\mathcal{O}\left(\eta_{0}^{2}\right).$
(15)
Thus, to order $\mathcal{O}\left(\eta_{0}\right)$,
$\tilde{v}_{r}=-\frac{\eta_{0}S_{0}R^{p-q+\frac{1}{2}}}{1+R^{2p}S_{0}^{2}}+\mathcal{O}\left(\eta_{0}^{2}\right),$
(16)
or equivalently, using Eqs. (84) and (81),
$v_{r}=\frac{rc_{\mathrm{s}}^{2}}{v_{\mathrm{k}}}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}\frac{\left(t_{\mathrm{s}}/t_{\mathrm{k}}\right)}{1+\left(t_{\mathrm{s}}/t_{\mathrm{k}}\right)^{2}}.$
(17)
$S_{0}R^{p}$ is the dimensionless expression of the ratio
$t_{\mathrm{s}}/t_{\mathrm{k}}$. Eq. (16) shows that $R^{2p}S_{0}^{2}\ll 1$ or
$R^{2p}S_{0}^{2}\gg 1$, and thus $R^{p}S_{0}\ll 1$ or $R^{p}S_{0}\gg 1$,
resulting in asymptotic behaviours for the radial grain motion. These
asymptotic regimes were first described by W77 for the particular case $p=0$.
They correspond physically to two limiting cases: where the gas drag
dominates, which we call the A-mode, and where gravity dominates, which we
call the B-mode. In the next sections, we study and describe these two so-
called “regimes of migration” or “modes of migration” before treating the
global evolution of grains given by Eq. (16).
### 3.1 A-mode (Radial differential migration)
The A-mode corresponds to the regime $R^{p}S_{0}\ll 1$ (or equivalently
$t_{\mathrm{s}}/t_{\mathrm{k}}\ll 1$). In the A-mode, Eq. (16) reduces to
$\tilde{v}_{r}=\frac{\mathrm{d}r}{\mathrm{d}t}=\frac{\mathrm{d}R}{\mathrm{d}T}=-\eta_{0}S_{0}R^{p-q+\frac{1}{2}},$
(18)
or equivalently
$v_{r}=\frac{rc_{\mathrm{s}}^{2}}{v_{\mathrm{k}}}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}\frac{t_{\mathrm{s}}}{t_{\mathrm{k}}},$
(19)
where the $\mathcal{O}\left(R^{p}S_{0}\right)$ has been neglected. In this
mode of migration, the stopping time is much smaller than the Keplerian time
scale. Considering one grain’s orbit around the central star, its orbital
velocity is forced by the gas drag to become sub-Keplerian in just a few
stopping times, i.e. almost instantaneously. Thus, the centrifugal
acceleration is not efficient enough to counterbalance the gravitational
attraction of the central star, and the grain feels an inward radial
differential acceleration. The gas drag counterbalances this radial motion and
the grain reaches a local limit velocity in a few stopping times. We call the
physical process of the A-mode of migration “Radial Differential Migration”.
The A-mode of migration originates first from a perturbative expansion for
$\eta_{0}\ll 1$ (rigorously for $\eta_{0}R^{-q+1}\ll 1$) and second from a
perturbative expansion for $S_{0}\ll 1$ (rigorously for $S_{0}R^{p}\ll 1$).
Formally, we have performed: $\lim\limits_{\begin{subarray}{c}S_{0}\ll
1\end{subarray}}\lim\limits_{\begin{subarray}{c}\eta_{0}\ll
1\end{subarray}}\left[...\right]$. Historically, the A-mode had been derived
by W77 to explain the radial motion of small grains. In his study, he
neglected the radial dependence of the stopping time and assumed $S_{0}\ll 1$
(this approximation also implies that $R^{2p}S_{0}^{2}\ll 1$, as $R<1$, see
Appendix F).
It is straightforward to integrate the differential equation Eq. (18) by
separating the $R$ and $T$ variables. Noting that $R\left(T=0\right)=1$, we
have:
* •
If $-p+q+\frac{1}{2}\neq 0$:
$\left\\{\begin{array}[]{l}\displaystyle
R=\displaystyle\left[1-\left(-p+q+\frac{1}{2}\right)\eta_{0}S_{0}T\right]^{\frac{1}{-p+q+\frac{1}{2}}}\\\
\displaystyle
T=\displaystyle\frac{1-R^{-p+q+\frac{1}{2}}}{\left(-p+q+\frac{1}{2}\right)\eta_{0}S_{0}}.\end{array}\right.$
(20)
* •
If $-p+q+\frac{1}{2}=0$:
$\left\\{\begin{array}[]{l}\displaystyle
R=\displaystyle\mathrm{e}^{-\eta_{0}S_{0}T}\\\ \displaystyle
T=\displaystyle-\frac{\mathrm{ln}\left(R\right)}{\eta_{0}S_{0}}.\end{array}\right.$
(21)
The outcome of the dust radial motion comes from a competition between two
effects. As the grain reaches smaller radii, (1) gas drag increases, which
slows down the radial drift, and (2) the differential acceleration due to the
pressure gradient increases which enhances the migration efficiency. Point (1)
is related to $s_{\mathrm{opt}}$, which scales as the surface density profile,
while the acceleration due to the pressure gradient in (2) is related to the
temperature profile (see Eq. (81)). Depending on which process is dominant,
the grain’s dynamics can lead to two regimes:
* •
If $-p+q+\frac{1}{2}\leq 0$: then as
$T(R)\sim\displaystyle\frac{-R^{-p+q+\frac{1}{2}}}{\left(-p+q+\frac{1}{2}\right)\eta_{0}S_{0}}=\mathcal{O}\left(R^{-p+q+\frac{1}{2}}\right),$
(22)
the time it takes the grain to reach smaller and smaller radii increases
drastically, according to the diverging power-law. Importantly, this behaviour
constitutes our definition of the grain “pile-up”. Mathematically speaking,
accretion onto the central star occurs in an infinite time, i.e.
$\lim\limits_{\begin{subarray}{c}T\to+\infty\end{subarray}}R=0.$ (23)
* •
If $-p+q+\frac{1}{2}>0$: the grain is accreted onto the central star in a
finite migration time given by
$T_{\mathrm{m}}=\frac{1}{\eta_{0}S_{0}\left(-p+q+\frac{1}{2}\right)},$ (24)
which increases as $S_{0}$ and $\eta_{0}$ decrease, so that
$\lim\limits_{\begin{subarray}{c}T\to T_{\mathrm{m}}\end{subarray}}R=0.$ (25)
The presence or absence of a physical grain pile-up is therefore demonstrated
considering the asymptotic behaviour of $R(T)$ at large times. It is important
to realise that the pile-up is a cumulative effect that can not arise from
velocities only (which however provide qualitative information on the grain’s
motion) but can only be found by integrating the equation of motion. This
rigorously allows us to distinguish two different behaviours for the outcome
of the grain’s radial motion, and thus two classes of discs with respect to
the A-mode.
### 3.2 B-mode (Drift forced by a resistive torque)
Returning to Eq. (16), the B-mode corresponds to the other asymptotic regime,
where $R^{p}S_{0}\gg 1$ (or equivalently $t_{\mathrm{s}}/t_{\mathrm{k}}\gg
1$). In this case,
$\tilde{v}_{r}=\frac{\mathrm{d}R}{\mathrm{d}T}=-\frac{\eta_{0}}{S_{0}}R^{-p-q+\frac{1}{2}},$
(26)
or equivalently
$v_{r}=\frac{rc_{\mathrm{s}}^{2}}{v_{\mathrm{k}}}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}\frac{t_{\mathrm{k}}}{t_{\mathrm{s}}}.$
(27)
In this mode of migration, the stopping time is much larger than the Keplerian
time scale. Hence, the orbital velocity of a grain around the central star is
almost the Keplerian velocity. However, because of the pressure gradient, the
gas orbits around the central star at a sub-Keplerian velocity. Thus, the
azimuthal differential velocity between the gas and the grain generates an
azimuthal drag force whose torque extracts angular momentum from the grain.
Given that the Keplerian angular momentum increases with radius
($l\propto\sqrt{r}$), this torque results in the inward migration of the
grain. We call the physical process of this B-mode of migration “Drift Forced
by a Resistive Torque”.
As for the A-mode, the B-mode of migration can also be derived first from an
expansion in $\left(S_{0}R^{p}\right)^{-2}$ and then from an expansion in
$\eta_{0}$. Historically, W77 found an expression while only assuming that
$S_{0}\gg 1$ since he considered a flat density profile. To find the
expression derived by W77 for large grains, we must assume that $S_{0}R^{p}\gg
1$. It is crucial to see that this expression does not imply that $S_{0}\gg 1$
when $p>0$ and $R\to 0$.
It is straightforward to integrate the differential equation Eq. (26) by
separating the $R$ and $T$ variables. Noting that $R\left(T=0\right)=1$, we
have:
* •
If $p+q+\frac{1}{2}\neq 0$:
$\left\\{\begin{array}[]{l}\displaystyle
R=\displaystyle\left[1-\left(p+q+\frac{1}{2}\right)\frac{\eta_{0}}{S_{0}}T\right]^{\frac{1}{p+q+\frac{1}{2}}}\\\
\displaystyle
T=\displaystyle\frac{S_{0}}{\eta_{0}}\frac{1-R^{p+q+\frac{1}{2}}}{\left(p+q+\frac{1}{2}\right)}.\end{array}\right.$
(28)
* •
If $p+q+\frac{1}{2}=0$:
$\left\\{\begin{array}[]{l}\displaystyle
R=\displaystyle\mathrm{e}^{-\frac{\eta_{0}}{S_{0}}T}\\\ \displaystyle
T=\displaystyle-\frac{S_{0}}{\eta_{0}}\mathrm{ln}\left(R\right).\end{array}\right.$
(29)
As for the A-mode, two kinds of behaviours appear, depending on the $p$ and
$q$ exponents:
* •
If $p+q+\frac{1}{2}\leq 0$: The grain migrates inwards, piles-up in the disc’s
inner regions and falls onto the star in an infinite time:
$\lim\limits_{\begin{subarray}{c}T\to+\infty\end{subarray}}R=0.$ (30)
However, the negative exponents required to be in this regime do not
correspond to physical discs. Therefore, the grain dynamics in the B-mode in
real discs belong to the second case:
* •
If $p+q+\frac{1}{2}>0$: The grain is accreted onto the central star in a
finite time
$T_{\mathrm{m}}=\frac{S_{0}}{\eta_{0}\left(p+q+\frac{1}{2}\right)},$ (31)
which increases as $S_{0}$ increases and $\eta_{0}$ decreases and so that
$\lim\limits_{\begin{subarray}{c}T\to T_{\mathrm{m}}\end{subarray}}R=0.$ (32)
As for the A-mode, considering the limit of $R(T)$ at large times also proves
the existence of two classes of discs with respect to the B-mode of migration.
The radial motion of a grain in the B-mode of migration is also driven by a
competition between the increase of both the drag and the acceleration due to
the pressure gradient. However, for real discs, $p>0$ and $q>0$, and
therefore, $p+q+\frac{1}{2}>0$. Grains migrating in B-mode in such discs fall
in a finite time onto the central star.
### 3.3 Radial evolution and asymptotic behaviour of single grains
As we have seen, the grains behaviour is divided into two asymptotic regimes,
called the A-mode and the B-mode, which come from two different physical
origins. However, the two criteria determining if the grains are accreted onto
the central object in a finite or infinite time differ for the A- and the
B-mode. It is thus crucial to determine in which mode a dust grain ends its
motion to predict if the grain is ultimately accreted or not. Returning again
to Eq. (16), we have
$\frac{\mathrm{d}R}{\mathrm{d}T}=\frac{-\eta_{0}S_{0}R^{p-q+\frac{1}{2}}}{1+R^{2p}S_{0}^{2}}.$
(33)
We can separate the $R$ and $T$ variables and integrate to obtain an
expression for $T(R)$:
$T=\frac{1}{\eta_{0}S_{0}}\left[T_{\mathrm{A}}\left(R\right)+T_{\mathrm{B}}\left(R\right)\right],$
(34)
where
$T_{\mathrm{A}}\left(R\right)=\left\\{\begin{array}[]{ll}\displaystyle\frac{1-R^{-p+q+\frac{1}{2}}}{-p+q+\frac{1}{2}}&\mathrm{if}\
-p+q+\frac{1}{2}\neq 0\\\\[10.00002pt] -\mathrm{ln}\left(R\right)&\mathrm{if}\
-p+q+\frac{1}{2}=0.\end{array}\right.$ (35)
and
$T_{\mathrm{B}}\left(R\right)=\left\\{\begin{array}[]{ll}\displaystyle
S_{0}^{2}\frac{1-R^{p+q+\frac{1}{2}}}{p+q+\frac{1}{2}}&\mathrm{if}\
p+q+\frac{1}{2}\neq 0\\\\[10.00002pt] \displaystyle-
S_{0}^{2}\mathrm{ln}\left(R\right)&\mathrm{if}\
p+q+\frac{1}{2}=0.\end{array}\right.$ (36)
Eq. (34) provides the asymptotic behaviour of the grains at large times.
Interestingly, as $p\geq 0$ for realistic discs, the contribution of the
B-mode becomes negligible when $R\ll S_{0}^{-1/p}$. Hence, grains initially
migrating in the A-mode stay in the A-mode, but grains initially migrating in
the B-mode end their radial motion in the A-mode. This behaviour is summarized
on Fig. 2 and detailed in Appendix G. This result was not predicted by W77, as
he neglected the radial dependence of the stopping time. Mathematically
speaking, it comes from the fact that the perturbative expansion of W77 has
been performed with respect to powers of $S_{0}$ and not powers of
$S_{0}R^{p}$.
Figure 2: Radial evolution of grains in the ($R$,$S_{0}$) plane showing that a
grain in the Epstein drag regime ends its radial motion in the A-mode. The
solid curves represent $R^{-p}$ for various values of $p$, they separate the
A-mode (below) from the B-mode (above) regions. The horizontal dashed lines
show trajectories of grains as they migrate inwards from $R=1$. The shaded
area is a forbidden zone.
To illustrate the radial motion of dust grains in protoplanetary discs, we
numerically integrate the equations of motion for different values of the
parameters $\eta_{0}$, $S_{0}$, $p$, $q$. We set $\eta_{0}=10^{-2}$ to mimic a
realistic disc and vary the order of magnitude of $S_{0}$ from $10^{-4}$ to
$10^{2}$ for two sets of ($p$,$q$) values. First, we choose ($p=0$,
$q=\frac{3}{4}$); according to the NSH86 expansion, the grains are accreted by
the central star in a time $T_{\mathrm{m}}$. Second, we set ($p=\frac{3}{2}$,
$q=\frac{3}{4}$); the grains fall onto the central star in an infinite time
from the same approximation. This set of ($p$,$q$) values is taken to mimic
discs profiles that are commonly used and for which $-p+q+\frac{1}{2}$ can
take a positive or a negative value. Consequently, we interpret the radial
motion $R\left(T\right)$ of dust grains plotted in Fig. 3.
* •
The top panel of Fig. 3 shows the results for ($p=0$, $q=\frac{3}{4}$): Grains
fall onto the central star, initially in the A-mode for the small grains and
in the B-mode for the large ones. The radial-drift process is long for small
and large grains but is optimal for grains with $S=S_{\mathrm{m}}=1$ for which
the accretion time is $T_{\mathrm{m}}=1.6/\eta_{0}$, or $T_{\mathrm{m}}=160$
with $\eta_{0}=10^{-2}$.
Figure 3: Radial motion $R\left(T\right)$ of dust grains in the Epstein regime
for $\eta_{0}=10^{-2}$. $S_{0}$ varies from $10^{-4}$ to $10^{2}$. Top: $p=0$,
$q=\frac{3}{4}$, here $-p+q+\frac{1}{2}>0$ and the grain is accreted onto the
central star in a finite time. Bottom: $p=\frac{3}{2}$, $q=\frac{3}{4}$, here
$-p+q+\frac{1}{2}<0$ and the grain piles up and is consequently accreted onto
the central star in an infinite time.
* •
The bottom panel of Fig. 3 shows the results for ($p=\frac{3}{2}$,
$q=\frac{3}{4}$): In this case, the radial density profile is steep enough to
ensure that the grains are not accreted onto the central star. To reach a
given radius (for example $R_{\mathrm{f}}=0.1)$, the optimal size is
$S_{\mathrm{m,f}}\simeq 2.9=\mathcal{O}\left(1\right)$ (see Eq. (112)). Hence,
in this case grains efficiently reach the disc inner regions without ever
being accreted onto the central star. The transition from the B-mode to the
A-mode (for which $R\propto T^{-4}$ in this case) for the large grains is
visible in this plot.
## 4 Radial motion in the Stokes regime
Radial migration of large particles occurs in the Stokes drag regime, which
depends on the dynamical viscosity $\mu$ of the gas. For hydrogen molecules:
$\mu=\frac{5m\sqrt{\pi}}{64\sigma_{\mathrm{s}}}\sqrt{\frac{k_{\mathrm{B}}T}{m}},$
(37)
where $m=2m_{\mathrm{H}}=3.347446922\times 10^{-27}$ kg and
$\sigma_{\mathrm{s}}=2.367\times 10^{-19}$ m2 is the molecular cross section
of the molecule (Chapman & Cowling 1970). The kinematic viscosity $\nu$ is
then defined by $\mu=\rho_{\mathrm{g}}\nu$ and the gas collisional mean free
path is given by
$\lambda_{g}=\sqrt{\frac{\pi}{2}}\frac{\nu}{c_{\mathrm{s}}}.$ (38)
We now generalise the procedure outlined in Sects. 2 and 3 to the three Stokes
regimes of Eq. (3). Using the dimensionless coordinates described above, we
have
$\mu=\mu_{0}R^{-\frac{q}{2}},$ (39)
$\mu_{0}=\frac{5m\sqrt{\pi}}{64\sigma_{\mathrm{s}}}c_{\mathrm{s}0},$ (40)
where $c_{\mathrm{s}0}$ is given by Eq. (72). Thus, the expression of the
kinematic viscosity $\nu$ is
$\nu=\displaystyle\nu_{0}R^{\frac{3}{2}+p-q}e^{\frac{Z^{2}}{2R^{3-q}}},$ (41)
$\nu_{0}=\displaystyle\frac{\mu_{0}}{\rho_{\mathrm{g}0}}\,.$ (42)
First, if $R_{\mathrm{g}}<1$, the drag force is linear in
$\textbf{v}-\textbf{v}_{\mathrm{g}}$ and thus has the same structure as for
the Epstein regime. Comparing the expressions of $\mathcal{C}\left(R,0\right)$
for the Epstein and the linear Stokes regime (see Appendix C), all the results
found for the radial motion in the Epstein regime can therefore be directly
transposed by setting $q^{\prime}=q$ and $p^{\prime}=\frac{q-3}{2}$. In this
case, the grain radial motion does not depend on $p$ anymore and the NSH86
expansion of the radial motion for small pressure gradients provides (see Eq.
(33))
$\frac{\mathrm{d}R}{\mathrm{d}T}=\frac{-\eta_{0}S_{0}^{2}R^{-\frac{q}{2}-1}}{1+R^{q-3}S_{0}^{4}}.$
(43)
These crucial results follow:
* •
In the A-mode ($S_{0}\ll R^{\frac{3-q}{4}}$), grains experience a pile-up and
migrate onto the central star in an infinite time if
$-p^{\prime}+q^{\prime}+\frac{1}{2}\leq 0$, i.e. if $q\leq-4$ (which never
occurs in real discs).
* •
In the B-mode ($S_{0}\gg R^{\frac{3-q}{4}}$), grains migrate onto the central
star in an infinite time if $p^{\prime}+q^{\prime}+\frac{1}{2}\leq 0$, i.e. if
$q\leq\frac{2}{3}$.
Figure 4: Radial evolution of grains in the ($R$,$S_{0}$) plane showing that a
grain in the linear Stokes drag regime ends its radial motion in the B-mode.
The solid curves represent $R^{\frac{3-q}{4}}$ for various values of $q$, they
separate the A-mode (below) from the B-mode (above) regions. The horizontal
dashed lines show trajectories of grains as they migrate inwards from $R=1$.
Shaded area: forbidden.
Thus, similar to the Epstein regime, we derived one criterion for each mode
and need to determine in which mode the grain ends its motion. For observed
discs, $q-3<0$ (see Sect. 6), and as the particle migrates inward, $R$ becomes
smaller than $S_{0}^{\frac{4}{3-q}}$ and grains end their radial motion in the
B-mode (see Fig. 4). This result is fundamentally different to the one we
obtained for the Epstein regime. Indeed, for grains migrating in the A-mode in
the Stokes regime at low Reynolds numbers, the criterion obtained for a pile-
up in the A-mode is never satisfied for real discs. However, after migrating
inside a critical radius, grains switch to the B-mode, for which the pile-up
can potentially occur, depending on the value of $q$. The corollary is that in
discs having $q\leq\frac{2}{3}$, i.e. a shallow enough temperature profile,
large grains in the Stokes regime at small Reynolds numbers remain in the
disc. Such a criterion is applicable for real protoplanetary discs.
Second, if $R_{\mathrm{g}}>800$, the drag force is quadratic in
$\textbf{v}-\textbf{v}_{\mathrm{g}}$. Assuming that the radial motion is
decoupled from the vertical motion, we perform the NSH expansion at small
pressure gradient (cf. Eq. (13)). We find that whatever the integer $j$,
$\left(\eta_{0}R^{-p}\right)^{j}\left(\textbf{v}-\textbf{v}_{\mathrm{g}}\right)\to
0$ at the limit $\eta_{0}\to 0$. This means that both $v_{r}$ and
$v_{\theta}-\sqrt{1/R}$ are flat functions as their Taylor series expansion
equals zero at each order. Consequently, they can not be determined by
perturbation analysis. This property comes from the quadratic dependency of
the drag with respect to the differential velocity and thus is not related to
the grain size. Consequently, in this drag regime, the drag force is extremely
efficient and the corrections to the Keplerian motion are negligible at every
order of the perturbative expansion. The particles are very well coupled to
the gas and do not migrate significantly.
Third, for the intermediate case, we could not manage to perform the expansion
at small pressure gradients. However, we expect an intermediate behaviour
between the two Stokes regime at small and large Reynolds numbers.
Consequently, if $R_{\mathrm{g}}>1$, the migration motion becomes less
efficient as the drag force is no longer linear with respect to the
differential velocity between the gas and the dust particles. Thus, the main
constraint for the radial-drift barrier due to the Stokes drag comes from the
low Reynolds number regime for which the migration motion is the most
efficient.
Finally, confusion often arises when defining the “radial-drift barrier” as
the difficulty a grain has of “overcoming $s=s_{\mathrm{opt}}$” (i.e.)
reaching the B-mode. Indeed, as we have shown, grains can survive their
migration motion in the Epstein regime when they are in the A-mode whenever
$-p+q+1/2<0$, and grains can start their migration motion in the B-mode but be
accreted in a finite time if $-p+q+1/2>0$. This study also shows for the
Stokes regime that a grain ends its migration motion in the B-mode. However,
as demonstrated in this work, the ability of the grain to overcome the radial
drift barrier is only linked to the value of $q$. If $q>2/3$, the grain will
be accreted onto the central star in a finite time, even if it has
$s>s_{\mathrm{opt}}$. Thus, we would argue that the definition of the radial-
drift barrier has to remain the ability of the grain to be accreted onto the
central star or depleted from the disc within its lifetime .
## 5 Limitations of the model
We have demonstrated that the time it takes for grains to reach smaller and
smaller radii increases dramatically under certain conditions. Specifically,
if $-p+q+\frac{1}{2}<0$ (resp. $q<\frac{2}{3}$), grains experience a pile-up
in the Epstein (resp. Stokes) regime. However, the model developed for the
radial evolution of dust grains in this paper remains simple in that we
neglect several important physical processes: turbulence, grain growth,
collective motion of dust grains, dust feedback on the gas surface density and
temperature profiles. We now discuss how those processes can modify the
criteria derived above.
1. 1.
The local pressure maxima created by turbulence (Cuzzi et al. 2001, 2008) and
the collective effects due to the dust drag onto the gas phase (Youdin &
Goodman 2005) are known to slow down the dust particles. However, the
efficiency of these processes — such as the non linearity of the streaming
instability in global disc models and the life time of the pressure maxima —
in real discs remains difficult to quantify. Omitting these phenomena
constitutes therefore an upper limit for the grain migration efficiency, which
will be slowed by these additional processes.
2. 2.
In this study, we assume that changing the dust distribution does not change
the thermal profile of the disc. We also neglect the viscous evolution of the
disc, assuming that the viscous timescales are larger than the characteristic
timescales of the initial dust evolution. It implies that we assume that $p$
is constant during the whole grains evolution. We can expect that for long
term evolution, the surface density profile will flatten, leading to a smaller
value of $p$. This makes the Epstein criterion harder to be met, while the
Stokes criterion is not affected.
3. 3.
We have shown that even if the velocity of the grain’s inward motion depends
on their sizes, their outcome only depends on the surface density and
temperature profiles. Thus, if we now consider growing (or fragmenting)
grains, we expect that (i) the intensity of the inward motion depends on the
growth efficiency (this point will be discussed in detail in a forthcoming
paper), but that (ii) the grains outcome remains determined by our criteria,
regardless of the growth regime.
Following this discussion, our simple model deals with processes that are
optimized for the grains to be depleted on the central object. Consequently,
our Epstein and Stokes criteria for the radial-drift barrier constitute the
least favourable limit for grain survival. Thus, we are confident when
claiming that the radial-drift barrier does not occur in some classes of
discs. One may however expect that more discs are retaining their grains due
to the complementary processes mentioned above.
It should be noted that the criterion $-p+q+1/2$ quantifies the outcome of the
radial-drift motion of the grains, but not their kinematics (which depends on
the grain size, the grain density, etc…). Thus, we provide predictions for
which discs will retain the largest mass of solid particles, but do not
predict for which discs the radial migration to the inner disc regions is the
fastest. Full simulations like those developed in Brauer et al. (2008) are
required to make predictions of the dust kinetics, even more so when the grain
size evolution is driven by a complex model of growth and fragmentation.
However, this study suggests that while complex simulations are useful to
study the details of the dust dynamics, they are not required to determine the
grains outcome.
As a conclusion of this section, we have mentioned that the physics treated by
our model is not exhaustive. In real discs, the limits $-p+q+\frac{1}{2}=0$
and $q=\frac{2}{3}$ may be softened by the effects of additional physical
phenomena. However, these neglected processes (such as turbulence and grain
growth) tend to decrease the efficiency of the dust radial motion. Our
predictions of when the “radial-draft barrier” does not occur therefore remain
valid. Our model represents a powerful indicator for predicting the dust
behaviour in discs with given power law profiles: we expect that (i) discs
satisfying $-p+q+\frac{1}{2}\leq 0$ retain their small grain population and
that (ii) discs satisfying $q\leq\frac{2}{3}$ keep their large solids. On the
contrary, discs for which $-p+q+\frac{1}{2}>0$ (resp. $q>\frac{2}{3}$) likely
lose their small (resp. large) particles.
## 6 Application to observed discs and planet formation
### 6.1 Validity of the criteria in real protoplanetary discs
We now study how the criteria we derived can be applied when considering the
physical evolution of grains in observed discs, with finite inner radii and
finite lifetimes. The analytic expressions of the previous sections have been
derived using dimensionless quantities. We now provide the physical timescales
of the radial dust motion estimating the parameters involved in real
protoplanetary discs. We consider in this section a typical CTTS disc, of mass
$M_{\mathrm{disc}}=10^{-2}\ M_{\odot}$ around a 1 $M_{\odot}$ star, extending
from $r_{\mathrm{in}}=10^{-2}$ AU to $r_{\mathrm{out}}=10^{3}$ AU. The disc
inner edge is chosen to correspond to the dust sublimation radius for a 1
$M_{\odot}$ star, whereas its outer boundary is representative of the largest
observed discs. Its vertical extent is set by the choice of the temperature
scale. We take $T(1$ AU$)=150$ K, a typical value obtained by Andrews &
Williams (2007) in their disc observations.
The transition from the Epstein to the Stokes regime occurs when
$\lambda_{\mathrm{g}}=\frac{4s}{9}$, or
$s=\frac{9}{4}\lambda_{\mathrm{g}}=\frac{45\pi^{\frac{3}{2}}}{256}\frac{mH^{\prime}_{0}}{\sigma_{\mathrm{s}}\Sigma^{\prime}_{0}}\,r^{p-\frac{q}{2}+\frac{3}{2}},$
(44)
and is represented in the $(r,s)$ plane in Fig. 5 for this typical disc for
different values of the surface density and temperature power-law exponents
$p$ and $q$. The Stokes regime is seen to apply to large bodies in the disc
inner regions and for large values of both $p$ and $q$.
Figure 5: Transition between the Epstein and the Stokes regimes in a
protoplanetary disc of $M_{\rm{disc}}=10^{-2}M_{\odot}$ extending from
$r_{\rm{in}}=10^{-2}$ AU to $r_{\rm{out}}=10^{3}$ AU for several values of $p$
and $q$. Grains with sizes below (resp. above) than the curve experience the
Epstein (resp. Stokes) drag regime.
Disc lifetimes are generally thought to be a few Myr (Haisch et al. 2001;
Carpenter et al. 2005), and thus we take $t_{\mathrm{disc}}\sim 10^{6}$ yr.
For a grain starting at a distance $r_{0}$ from a 1 $M_{\odot}$ star, the
dimensionless value $T_{\mathrm{disc}}$ is therefore
$T_{\mathrm{disc}}=\frac{t_{\mathrm{disc}}}{t_{\mathrm{k}0}}=\frac{\sqrt{\mathcal{G}M_{\odot}}\,t_{\mathrm{disc}}}{r_{0}^{3/2}}\sim\frac{6\times
10^{6}}{[r_{0}\ \mathrm{(AU)}]^{3/2}}.$ (45)
The dimensionless value of the dust disc inner radius ($r_{\mathrm{in}}\sim
0.01$ AU) for a grain starting at $r_{0}$ is
$R_{\mathrm{in}}=\frac{r_{\mathrm{in}}}{r_{0}}\sim\frac{10^{-2}}{r_{0}\
\mathrm{(AU)}}.$ (46)
The link between dimensionless and real grain sizes is made through the
optimal size for radial migration. Considering first the Epstein regime, its
midplane value has a radial dependence given by
$s_{\mathrm{opt}}(r,0)=\frac{\Sigma(r)}{\sqrt{2\pi}\,\rho_{\mathrm{d}}}=\frac{\Sigma^{\prime}_{0}\,r^{-p}}{\sqrt{2\pi}\,\rho_{\mathrm{d}}}.$
(47)
The dimensionless size of a particle of size $s$ starting its migration at
position $r_{0}$ is therefore
$S_{0}=\frac{s}{s_{\mathrm{opt},0}}=\frac{\sqrt{2\pi}\,\rho_{\mathrm{d}}}{\Sigma^{\prime}_{0}}\,s\,r_{0}^{p}.$
(48)
Values of $S_{0}$ are plotted in the $(r_{0},s)$ plane in Fig. 6 for a disc of
total mass $M_{\mathrm{disc}}=0.01\ M_{\odot}$ extending from
$r_{\mathrm{in}}=10^{-2}$ AU to $r_{\mathrm{out}}=10^{3}$ AU, with grains of
intrinsic density $\rho_{\mathrm{d}}=1000$ kg m-3 and $\eta_{0}=10^{-2}$, for
both $p=0$ and $p=\frac{3}{2}$.
Figure 6: Values of $S_{0}$ as a function of grain size $s$ and initial
position $r_{0}$ for a disc of mass $M_{\mathrm{disc}}=0.01\ M_{\odot}$
extending from $r_{\mathrm{in}}=10^{-2}$ AU to $r_{\mathrm{out}}=10^{3}$ AU,
with grains of intrinsic density $\rho_{\mathrm{d}}=1000$ kg m-3 and
$\eta_{0}=10^{-2}$, for $p=0$ and $q=\frac{3}{4}$ (top) and $p=\frac{3}{2}$
and $q=\frac{3}{4}$ (bottom). The thick line shows the limit between grains
that are accreted onto the star ($t_{\mathrm{in}}<t_{\mathrm{disc}}$) and
those that survive in the disc ($t_{\mathrm{in}}>t_{\mathrm{disc}}$), i.e. the
survival limit, in the Epstein regime.
The dimensionless time $T_{\mathrm{in}}$ for a grain to reach
$R_{\mathrm{in}}$ is given by Eq. (34). In combination with Eq. (46), this
gives an expression of $T_{\mathrm{in}}$ as a function of $S_{0}$ and $r_{0}$,
whereas Eq. (45) gives an expression of $T_{\mathrm{disc}}$ as a function of
$r_{0}$. Equating them yields a second order equation in $S_{0}$ as a function
of $r_{0}$, which can be solved to determine under which conditions a grain
reaches the disc inner edge at the end of its lifetime. Using Eq. (48) gives
the corresponding relationship between the grain size and its initial
position:
$\begin{array}[]{l}s=\displaystyle\frac{p+q+\frac{1}{2}}{2\sqrt{2\pi}\,\rho_{\mathrm{d}}}\frac{\Sigma^{\prime}_{0}\,r_{0}^{-p}}{1-\left(\frac{r_{\mathrm{in}}}{r_{0}}\right)^{p+q+\frac{1}{2}}}\left[\frac{\sqrt{\mathcal{G}M}\,t_{\mathrm{disc}}\,\eta_{0}}{r_{0}^{3/2}}\right.\\\\[21.52771pt]
\displaystyle\left.\pm\sqrt{\frac{\mathcal{G}M\,t_{\mathrm{disc}}^{2}\,\eta_{0}^{2}}{r_{0}^{3}}-4\frac{\left(1-\left(\frac{r_{\mathrm{in}}}{r_{0}}\right)^{-p+q+\frac{1}{2}}\right)\left(1-\left(\frac{r_{\mathrm{in}}}{r_{0}}\right)^{p+q+\frac{1}{2}}\right)}{\left(-p+q+\frac{1}{2}\right)\left(p+q+\frac{1}{2}\right)}}\right],\end{array}$
(49)
which is plotted as a thick line in Fig. 6. It separates the $(r_{0},s)$ plane
into regions in which grains reach $r_{\mathrm{in}}$ and leave the disc before
it dissipates ($t_{\mathrm{in}}<t_{\mathrm{disc}}$) or survive in the disc
throughout its lifetime ($t_{\mathrm{in}}>t_{\mathrm{disc}}$). We call this
curve the survival limit.
For $-p+q+\frac{1}{2}>0$, illustrated by the case $(p=0,q=\frac{3}{4})$, Fig.
6 shows as expected that most of the grains are lost during the disc lifetime.
However, small and large grains initially in the outer disc survive, therefore
even with this profile, the disc retains a fraction of its grain population
before it dissipates. Moreover, one may expect growing grains that reach
$S_{0}=1$ to be inevitably accreted onto the star (unless the growth process
is fast enough for grains to outgrow the fast migrating sizes before they
leave the disc, see Laibe et al. 2008). Such discs may not form planets, but
their remaining dust content may still make them observable, although they
would likely be faint.
For $-p+q+\frac{1}{2}\leq 0$, illustrated by the case
$(p=\frac{3}{2},q=\frac{3}{4})$, even though all grains fall on the star in an
infinite time, some of them reach the disc’s inner edge before it dissipates.
On the contrary, for $r_{0}>350$ AU, grains of all sizes remain in the disc.
This is also the case for all (sub)micron-sized grains, whatever their initial
location, as well as most of the grains up to 0.1 mm. These grains likely make
the disc bright and easy to observe, since they are the grains contributing
most to the disc emission at IR and submm wavelengths. A large reservoir of
grains is available to participate in the planet formation process, however a
firm conclusion on their survival as their size evolves would require
incorporating a treatment of grain growth, as discussed in Sect. 5.
It should be noted that the disc used in these examples represents a lower
limit, as it is low mass and very extended. A more massive disc with a smaller
outer radius would have a larger surface density, and the corresponding
survival limit in Fig. 6 would be shifted vertically towards larger sizes and
more and more grains of larger sizes would survive.
Figure 7: Survival limits of grains for different values of $p$ and $q$. Left:
Epstein regime, right: Stokes regime. Grains to the left of the curves
($t_{\mathrm{in}}<t_{\mathrm{disc}}$) are accreted onto the star whereas those
to the right ($t_{\mathrm{in}}>t_{\mathrm{disc}}$) survive in the disc.
The left panels of Fig. 7 show the influence of the surface density and
temperature profiles on the location and shape of the survival limit curve in
the $(r_{0},s)$ plane in the Epstein regime. Increasing $p$ from 0 to 2 shifts
the curve towards smaller radii and larger grain sizes, as well as slightly
tilts it clockwise. The outer disc region in which grains of all sizes survive
extends inwards, as well as the surviving population of small grains as the
curve’s lower branch shifts upwards and becomes flatter. On the contrary, the
steepening of the curve’s upper branch, confining the population of surviving
large solids to the disc outer regions, is less dramatic. Increasing $q$ from
$\frac{1}{4}$ to $\frac{3}{4}$ also tilts the curve clockwards, but shifts it
towards larger radii and smaller grain sizes. However, its effect is more
limited than that of changing $p$. A disc with a steeper surface density
profile and a shallower temperature profile is therefore more efficient at
retaining a larger quantity of small grains and up to larger sizes. Indeed,
large $p$ and small $q$ values are required to meet the $-p+q+\frac{1}{2}<0$
criterion introduced in Sect. 3.
An equation very similar to Eq. (49) can be obtained for the linear Stokes
regime by replacing $p$ and $q$ by $p^{\prime}=\frac{q-3}{2}$ and
$q^{\prime}=q$ (since the equation of motion has the same structure for both
drag regimes, see Sect. 4), and using the expression of $s_{\mathrm{opt},0}$
for that regime, given in Table 2. Here $s_{\mathrm{opt},0}\propto
T^{\frac{1}{4}}$: the weak temperature dependence results in very little
change for a large range of temperatures below or above our adopted value of
$T(1$ AU$)=150$ K. The right panels of Fig. 7 show the survival limit in the
Stokes regime for the different values of $q$ (note that it no longer depends
on $p$). When $q$ increases, the curve’s lower branch slides towards larger
radii, making the survival of particles in the Stokes regime less and less
favourable.
Given the form of Eq. (49) and its different expressions for each drag regime,
it is not possible to compute analytically the survival limit for a grain
transitioning from the Epstein to the Stokes regime as it migrates inwards.
However, large values of $p$ and $q$, for which the Stokes region is the
largest, are not observed in protoplanetary discs (see Sect. 6.2), and in
practical cases the Stokes regime only applies to a small area of the
$(r_{0},s)$ plane. Small grains, which are detected in disc observations at IR
and sub-millimetre (submm) wavelengths, are mostly subject to the Epstein
drag. We therefore focus on that regime in the following.
Figure 8: Isocontours of the survival time (i.e. the time needed to reach the
disc inner edge at $r_{\mathrm{in}}=0.01$ AU) of grains of size $s$ and
initial position $r_{0}$ for $q=\frac{3}{4}$ and different values of $p$.
Equation (49) can give quantitative information about the outcome of the grain
population. Replacing $t_{\mathrm{disc}}$ by any time $t$ gives the location
in the $(r_{0},s)$ plane of grains reaching the disc inner edge (at
$r=r_{\mathrm{in}}$) in that time $t$, which is therefore the survival time of
those grains. Its isocontours are shown in Fig. 8 for different values of $p$
and for $q=\frac{3}{4}$. Only one value of $q$ is shown as the $q$ dependence
is moderate, as can be seen from the left panels of Fig. 7. The fate of
particular dust grains can easily be obtained from these figures. For example,
in the context of disc observations, 1 mm grains initially at 100 AU fall on
the $10^{5}$ yr contour for $p=0$. Their survival time decreases to a few
$10^{4}$ yr for $p\sim 0.8$, and increases again to values larger than
$10^{6}$ yr as $p$ increases. At an initial position of a few hundred AU, 1 mm
grains survive longer than $10^{6}$ yr for any $p$, therefore long enough to
contribute to the disc emission over its entire lifetime. As noted above, such
grains have longer survival times for higher $p$ and lower $q$ values. As
another example, in the context of planetesimal formation, the survival time
of a 1 m particle initially at 1 AU is $\sim 10^{5}$ yr for $p=0$, decreases
to $\sim 10^{2}$ yr for $p\sim 1$, and increases again to $\sim 10^{3}-10^{4}$
yr for $p=2$. The ability of such particles to remain for long enough in the
disc to grow to larger sizes therefore strongly depends on the surface density
profile. As a general rule, the survival of pre-planetesimals in the inner
disc is favoured by small values of $p$.
Figure 9: Time evolution (in nine snapshots from $t=10^{-2}$ to $10^{6}$ yr)
of isocontours of the initial position of grains in the $(r,s)$ plane for the
disc with $p=0$ and $q=\frac{3}{4}$. The label for each contour can be deduced
from its abscissa in the upper left panel at $t=10^{-2}$ yr.
Figure 10: Same as Fig. 9 for $p=\frac{3}{2}$ and $q=\frac{3}{4}$.
Similarly, replacing now $r_{\mathrm{in}}$ with any radius $r$ in Eq. (49)
gives the locus in the $(r_{0},s)$ plane of grains reaching that radius $r$ at
any time $t$. Alternatively, one can plot isocontours of the initial position
$r_{0}$ of grains in the $(r,s)$ plane at various times, thus showing the
radial evolution of grains with the same initial position but different sizes.
This is shown in Fig. 9 for ($p=0$, $q=\frac{3}{4}$) and Fig. 10 for
($p=\frac{3}{2}$, $q=\frac{3}{4}$). These plots make it easy to compare the
radial evolution of any particle to any physical timescale of interest in the
disc. In particular, they show that the disc still contains particles at all
radii at the end of its evolution ($t=10^{6}$ yr). No grains are found to the
right of the $r_{0}=10^{3}$ AU contour, since this was the initial outer disc
radius. In the disc with $(p=0,q=\frac{3}{4})$, no grains between $\sim 0.06$
and $\sim 0.2$ mm remain, and grains of other sizes still present were
initially in the outer disc. Given that the grains of sizes which contribute
to IR and submm emission have come from a small fraction of the initial disc,
this disc is likely faint. In the disc with $(p=\frac{3}{2},q=\frac{3}{4})$,
only grains with $s\sim 0.1$ mm are absent from the very outer regions, and
the observable grains come from a larger portion of the disc, likely making
the disc brighter than in the previous case.
As a conclusion, the analytic criteria derived above apply even when taking
into account the finite lifetime (or inner radius) of the disc. For most CTTS
discs, the dust is in the Epstein drag regime (except for some extreme values
for the grains sizes and discs profiles). Therefore, the grain’s radial
outcome is given by the value of $-p+q+\frac{1}{2}$. However, the transition
between discs for which the radial-drift barrier occurs or not consists more
of a continuum around the value $-p+q+\frac{1}{2}=0$ than in the sharp
transition predicted by the analytic model. Therefore, the radial motion of
the grains has to be studied on a case-by-case basis for discs close to the
transition $-p+q+\frac{1}{2}=0$, using the figures shown above in this
section.
### 6.2 Constraining physical systems
We now turn to observed discs and check if they meet our Epstein and Stokes
criteria to determine whether the radial-drift barrier is constraining for
planet formation. To estimate the values of the $p$ and $q$ exponents for real
discs, we use the results of disc modeling obtained by Andrews & Williams
(2005, 2007) from data on 63 discs in $\rho$ Ophiuchi, Taurus and Aurigae.
Using sub-millimetre fluxes measured at several wavelengths, they fit a range
of disc parameters assuming a geometrically thin irradiated disc with
opacities from Beckwith et al. (1990), a gas-to-dust ratio of 100, a disc
radius of 100 AU and zero disc inclination. The temperature exponent $q$ is
well constrained by the observational data set: the histogram of most probable
$q$ values is shown in Fig. 11.
Figure 11: Histogramm of the $q$ parameters obtained from Andrews & Williams
(2005, 2007) data of 63 observed discs. The distribution is roughly comprised
between 0.4 and 0.8 and, centred around 0.55. Approximately 90 % of the discs
satisfy $q<2/3$.
However, $p$ is not well constrained and is usually assumed to be
$\frac{3}{2}$. Very flat profiles with $p<\frac{1}{2}$ and very steep profiles
with $p>\frac{3}{2}$ seem to be excluded (Dutrey et al. 1996; Wilner et al.
2000; Kitamura et al. 2002; Testi et al. 2003; Isella et al. 2009; Andrews &
Williams 2007; Andrews et al. 2009). We represent the disc distribution
modeled by Andrews & Williams from observations in the $\left(p,q\right)$
diagram of Fig. 12: the histogram of Fig. 11 is represented by the gray-shaded
area and spread over a range of $p$ values, taking into account that extreme
values of $p$ are less probable. The dashed line ($-p+q+1/2=0$) represents the
border between migration in an infinite time and accretion onto the central
star for the A-mode of migration in the Epstein regime, while the thick dotted
line ($q=3/2$) represents that same border for the B-mode of migration in the
Stokes regime at low Reynolds number. The two black circles indicate the discs
used as examples in Sect. 3.3.
Figure 12: Location of the different outcomes of radial migration in the
($p$,$q$) plane. Dashed (resp. dotted) line: limit between accretion without
or with grains pile-up resulting in a finite or infinite time in the A-mode of
the Epstein regime (resp. B-mode of the Stokes regime at small Reynolds
numbers). Shaded area: location of observed discs. Black dots: discs used as
examples in Sect. 3.3.
We have split the disc distribution in four regions in the $\left(p,q\right)$
plane:
1. 1.
region 1: $-p+q+\frac{1}{2}\leq 0$ and $q\leq\frac{2}{3}$: both small and
large grains experience the pile-up effect. Those discs are potentially
observable and may favour planet formation.
2. 2.
region 2: $-p+q+\frac{1}{2}\leq 0$ and $q>\frac{2}{3}$: only small grains
experience the pile-up effect: even though such discs retain their small
grains, the population of pre-planetesimals in the disc inner regions may
efficiently be accreted onto the central star (at least until they reach the
high-$R_{\mathrm{g}}$ Stokes regime).
3. 3.
region 3: $-p+q+\frac{1}{2}>0$ and $q\leq\frac{2}{3}$: if the pre-
planetesimals can form before the entire distribution of small grains has been
accreted onto the central object, they will remain in the disc and may
constitute planet embryos.
4. 4.
region 4: $-p+q+\frac{1}{2}>0$ and $q>\frac{2}{3}$: both small and large
grains are accreted onto the central star.
The Epstein criterion indicates that for $q$ values in the range constrained
by observations, discs which keep their small grain population, and are
therefore likely to be bright in the IR and submm, should have $p$ values
approximately in the $[1;\frac{3}{2}]$ range. This is indeed what is found in
most disc surveys (Ricci et al. 2010a, b). On the contrary, smaller $p$ values
should correspond to discs which lose most of their small grains, and are
therefore more difficult to detect. This is what is found by Andrews et al.
(2010), who pushed their previous observations of the Ophiuchus star forming
region (Andrews et al. 2009) down to fainter discs, finding for this new
sample a median $p$ value of 0.9, lower than for brighter discs. The criterion
we derive in this paper for small grains in the Epstein regime provides
therefore the correct behaviour for explaining the range of $p$ values of
observed discs. However, this result has to be considered carefully for two
reasons. Firstly, the $p$ and the $q$ exponents determined from the
observations have to be considered with their respective errors. Given these
uncertainties, one may not be able to distinguish between a strict negative or
positive value for $-p+q+\frac{1}{2}$. Second, the boundary between the
different zones of the $(p,q)$ plane consists more of a continuum rather than
a strict limit due to the finite lifetime/inner radii of the discs. The
outcome of the grains may thus not be predicted when the value of
$-p+q+\frac{1}{2}$ is close to zero.
Now turning to the Stokes criterion for large solids, Figs. 11 and 12 show
that the vast majority of observed protoplanetary discs have shallow
temperature profiles ($q\leq\frac{2}{3}$) and are thus able to retain their
population of pre-planetesimals. These discs are therefore relevant places to
find evidence of planet formation, provided small grains can efficiently grow
to form pre-planetesimals. For the remainder of the disc population, the
outcome of pre-planetesimals will likely depend on their ability to reach the
high Reynolds number Stokes regime. However, the case of a steep radial
temperature profile can be encountered in at least one particular situation:
circumplanetary discs which typically have temperature profiles with $q=1$
(Ayliffe & Bate 2009). In this environment, we predict from our Stokes
criterion that planetesimals will be accreted onto the planet. The timescale
of the planet formation by the core-accretion process, which usually
corresponds to the time required to release the gravitational energy of the
accreted bodies (Pollack et al. 1996), is thus increased as the drag from the
gas onto the planetesimals releases an additional thermal contribution.
## 7 Conclusion and perspectives
In this study, we have generalised the radial grain motion studies of W77,
NSH86 and YS02 for both the Epstein and the Stokes regimes, taking into
account the effects of both the surface density and temperature profiles in
the disc. As observations do not provide direct information about the three
dimensional structure of discs, radial profiles of surface density and
temperature are often described by power laws:
$\Sigma\left(r\right)=\Sigma_{0}^{\prime}r^{-p}$ and
$\mathcal{T}\left(r\right)=\mathcal{T}_{0}^{\prime}r^{-q}$, where both $p$ and
$q$ take positive values. The radial dust behaviour in those discs is governed
by the competition between gravity and gas drag. The final outcome of the
radial motion is set by two counterbalancing effects. First, the temperature
increases when the radius decreases. Consequently, the deviation from the
Keplerian velocity increases, which accelerates the dust’s radial inward
motion. At the same time, the surface density also increases, which increases
the gas drag efficiency and slows down the dust motion. The competition
between these two effects fixes the ultimate mode of migration of the grain
(A-mode, where the drag dominates or B-mode, where the gravity dominates) and
thus the final outcome for the dust motion. In this work, we have shown that
it can be represented by an analytical criterion which depends on the drag
regime. For the Epstein drag regime (in which the ultimate radial motion is in
the A-mode), if $-p+q+\frac{1}{2}>0$, the dust particle is accreted onto the
central star in a finite time, and if $-p+q+\frac{1}{2}\leq 0$, the grain
pile-up results in an infinite accretion time and small dust grains remain in
the disc. We have shown that, as expected, these conclusions are somewhat
mitigated when taking into account the finite disc lifetime and finite disc
size. However, the outcomes still remain similar: the bulk of the small grain
population is lost to the star in the first case, whereas in the second case
the disc keeps most of its small grains. A similar criterion is found for the
Stokes regime at low Reynolds number: if $q\leq\frac{2}{3}$, the accretion
time is infinite and large pre-planetesimals remain in the disc and can
constitute the primary material for planet formation. However, the Stokes
radial motion differs from the Epstein regime as the ultimate radial motion
occurs in the B-mode.
The observational consequence is that discs with a large population of small
grains should be strong emitters in the infrared and sub-millimetre and should
be easier to observe, and that those having lost most of their small dust
should be fainter and harder to detect. This is indeed what is found: a large
fraction of the observed discs have large $p$ values whereas fainter discs
tend to have lower $p$ values (Andrews et al. 2010), in agreement with this
Epstein criterion. In addition, most of the observed discs have
$q\leq\frac{2}{3}$, allowing them to retain also their large pre-
planetesimals. As noted by Ricci et al. (2010a, b), explaining the data
requires a mechanism halting or slowing down the radial migration of dust
grains. We show here that local pressure maxima need not be invoked, but
rather that the combination of adequate surface density and temperature
profiles is sufficient. The $p$ and $q$ exponents used to reach our
conclusions are of course strongly dependent on the model used to fit the
data. However, even varying the fitting models, a large majority of discs
still satisfy both conditions $-p+q+\frac{1}{2}\leq 0$ and $q\leq\frac{2}{3}$.
Consequently, the radial-drift barrier (or the so-called metre-size barrier
when considering an MMSN disc) does not appear to constitute a problem for
planet formation for the discs that we do observe.
Our conclusions presented in this study assumed that the grain size remains
constant during its motion. However, observations tell us that grains do grow
(Testi et al. 2003; Wilner et al. 2003; Apai et al. 2005; Lommen et al. 2007,
2009). Grain growth is studied in various theoretical studies (Schmitt et al.
1997; Stepinski & Valageas 1997; Suttner et al. 1999; Tanaka et al. 2005;
Dullemond & Dominik 2005; Klahr & Bodenheimer 2006; Garaud 2007; Brauer et al.
2008; Laibe et al. 2008; Birnstiel et al. 2009). In a forthcoming paper, we
will generalise the formalism developed here to explain the radial and
vertical behaviour of growing dust grains.
###### Acknowledgements.
This research was partially supported by the Programme National de Physique
Stellaire and the Programme National de Planétologie of CNRS/INSU, France, and
the Agence Nationale de la Recherche (ANR) of France through contract
ANR-07-BLAN-0221. The authors want to thank C. Terquem, L. Fouchet, S. Arena
and E. Crespe for useful comments and discussions. We also thank the referee
for greatly improving the quality of this work by suggesting we include the
important Stokes regime and present real numbers (time scales and grain sizes)
for our criteria.
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## Appendix A Notations
The notations and conventions used throughout this paper are summarized in
Table 1.
Symbol | Meaning
---|---
$M$ | Mass of the central star
g | Gravity field of the central star
$r_{0}$ | Initial distance to the central star
$\rho_{\mathrm{g}}$ | Gas density
$\bar{\rho}_{\mathrm{g}}\left(r\right)$ | $\rho_{\mathrm{g}}\left(r,z=0\right)$
$c_{\mathrm{s}}$ | Gas sound speed
$\bar{c}_{\mathrm{s}}\left(r\right)$ | $c_{\mathrm{s}}\left(r,z=0\right)$
$c_{\mathrm{s}0}$ | Gas sound speed at $r_{0}$
$T$ | Dimensionless time
$\mathcal{T}$ | Gas temperature ($\mathcal{T}_{0}$: value at $r_{0}$)
$\Sigma_{0}$ | Gas surface density at $r_{0}$
$p$ | Radial surface density exponent
$q$ | Radial temperature exponent
$P$ | Gas pressure
$v_{\mathrm{k}}$ | Keplerian velocity at $r$
$v_{\mathrm{k}0}$ | Keplerian velocity at $r_{0}$
$H_{0}$ | Gas scale height at $r_{0}$
$\phi_{0}$ | Square of the aspect ratio $H_{0}/r_{0}$ at $r_{0}$
$\eta_{0}$ | Sub-Keplerian parameter at $r_{0}$
$s$ | Grain size
$S$ | Dimensionless grain size
$S_{0}$ | Initial dimensionless grain size
$y$ | Grain size exponent in the drag force
$\textbf{v}_{\mathrm{g}}$ | Gas velocity
v | Grain velocity
$\rho_{\mathrm{d}}$ | Dust intrinsic density
$m_{\mathrm{d}}$ | Mass of a dust grain
$t_{\mathrm{s}}$ | Drag stopping time
$t_{\mathrm{s}0}$ | Drag stopping time at $r_{0}$
Table 1: Notations used in the article.
## Appendix B disc structure
### B.1 Hydrostatic equilibrium
At stationary equilibrium ($\frac{\partial}{\partial t}=0$), gas velocities
$\textbf{v}_{\mathrm{g}r}$, $\textbf{v}_{\mathrm{g}\theta}$,
$\textbf{v}_{\mathrm{g}z}$ and the gas density $\rho_{\mathrm{g}}$ obey mass
conservation and the Euler equation:
$\left\\{\begin{array}[]{rcl}\displaystyle\frac{1}{r}\frac{\rho_{\mathrm{g}}\partial
rv_{\mathrm{g}r}}{\partial r}+\frac{\rho_{\mathrm{g}}\partial
v_{\mathrm{g}z}}{\partial z}&=&\displaystyle 0,\\\
\displaystyle\rho_{\mathrm{g}}\textbf{v}_{\mathrm{g}}.\nabla\textbf{v}_{\mathrm{g}}&=&\displaystyle-\nabla
P+\rho_{\mathrm{g}}\textbf{g}.\end{array}\right.$ (50)
The solution of Eq. (50) requires:
$v_{\mathrm{g}r}=v_{\mathrm{g}z}=0,$ (51)
which ensures mass conservation. Projecting the Euler equation on
$\textbf{e}_{z}$:
$\frac{1}{\rho_{\mathrm{g}}}\frac{\partial P}{\partial
z}=-\frac{\mathcal{G}Mz}{\left(z^{2}+r^{2}\right)^{3/2}}.$ (52)
Assuming that:
$P=c_{\mathrm{s}}^{2}\left(r,z\right)\rho_{\mathrm{g}},$ (53)
and dividing both sides of Eq. (52) by $c_{\mathrm{s}}^{2}$, we have:
$\frac{\partial\,\mathrm{ln}\left(c_{\mathrm{s}}^{2}\rho_{\mathrm{g}}\right)}{\partial
z}=-\frac{\mathcal{G}Mz}{\left(z^{2}+r^{2}\right)^{3/2}c_{\mathrm{s}}^{2}}.$
(54)
Integrating Eq. (54) between 0 and $z$ provides:
$\rho_{\mathrm{g}}\left(r,z\right)=\frac{\bar{P}\left(r\right)}{c_{\mathrm{s}}^{2}\left(r,z\right)}\mathrm{e}^{\displaystyle-\int_{0}^{z}\frac{\mathcal{G}Mz^{\prime}\mathrm{d}z^{\prime}}{\left(r^{2}+z^{\prime
2}\right)^{3/2}c_{\mathrm{s}}^{2}\left(r,z^{\prime}\right)}}.$ (55)
This expression can be simplified by the following approximations:
* •
In the special vertically isothermal case, where the sound speed depends only
on the radial coordinate, Eq. (55) simplifies to:
$\rho_{\mathrm{g}}\left(r,z\right)=\bar{\rho}_{\mathrm{g}}\left(r\right)\mathrm{e}^{\displaystyle-\frac{\mathcal{G}M}{\bar{c}_{\mathrm{s}}^{2}\left(r\right)}\left[\frac{1}{r}-\frac{1}{\sqrt{r^{2}+z^{2}}}\right]}.$
(56)
* •
Further, assuming a thin disc ($\left(\frac{z}{r}\right)^{2}\ll 1$), a Taylor
series expansion of Eq. (56) leads to:
$\rho_{\mathrm{g}}\left(r,z\right)=\bar{\rho}_{\mathrm{g}}\left(r\right)\mathrm{e}^{-\frac{z^{2}}{2H\left(r\right)^{2}}},$
(57)
with:
$H\left(r\right)=\frac{r\bar{c}_{\mathrm{s}}\left(r\right)}{v_{\mathrm{k}}\left(r\right)},$
(58)
which is the classical scale height for vertically isothermal thin discs.
### B.2 Azimuthal velocity
The radial component of the Euler equations is given by:
$-\frac{v_{\mathrm{g}\theta}^{2}}{r}=-\frac{1}{\rho_{\mathrm{g}}}\frac{\partial
P}{\partial r}-\frac{\mathcal{G}Mr}{\left(z^{2}+r^{2}\right)^{3/2}},$ (59)
where $\rho_{\mathrm{g}}$ is given by Eq. (55). Thus:
$\rho_{\mathrm{g}}\left(r,z\right)=\frac{\bar{P}\left(r\right)}{c_{\mathrm{s}}^{2}\left(r,z\right)}\mathrm{e}^{-I_{1}\left(r,z\right)},$
(60)
with:
$\left\\{\begin{array}[]{rcl}I_{1}\left(r,z\right)&=&\displaystyle\int_{0}^{z}\mathcal{G}Mc_{\mathrm{s}}^{-2}\left(r,z^{\prime}\right)\partial_{z^{\prime}}f\left(r,z^{\prime}\right)\mathrm{d}z^{\prime},\\\\[8.61108pt]
f\left(r,z\right)&=&\displaystyle\frac{1}{r}-\frac{1}{\sqrt{r^{2}+z^{2}}}.\end{array}\right.$
(61)
To simplify Eq. (59), we first use the following identity:
$\frac{1}{\rho_{\mathrm{g}}}\frac{\,\partial P}{\partial
r}=c_{\mathrm{s}}^{2}\frac{\partial\,\mathrm{ln}\left(c_{\mathrm{s}}^{2}\rho_{\mathrm{g}}\right)}{\partial
r},$ (62)
which becomes with Eq. (60):
$\frac{1}{\rho_{\mathrm{g}}}\frac{\partial P}{\partial
r}=c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\left(\bar{P}\right)}{\mathrm{d}r}-c_{\mathrm{s}}^{2}\frac{\partial\,I_{1}}{\partial
r}.$ (63)
Noting that $f\left(r,z=0\right)=0$ and integrating $I_{1}$ by parts provides:
$I_{1}\left(r,z\right)=\mathcal{G}Mc_{\mathrm{s}}^{-2}f\left(r,z\right)-\int_{0}^{z}\mathcal{G}Mf\left(r,z^{\prime}\right)\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\left(r,z^{\prime}\right)\mathrm{d}z^{\prime},$
(64)
and:
$\displaystyle\frac{1}{\rho_{\mathrm{g}}}\frac{\mathrm{d}P}{\mathrm{d}r}$
$\displaystyle=$ $\displaystyle
c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\left(\bar{P}\right)}{\mathrm{d}r}-\left(-\frac{\mathcal{G}M}{r^{2}}+\frac{\mathcal{G}Mr}{\left(z^{2}+r^{2}\right)^{3/2}}\right.$
$\displaystyle\left.+c_{\mathrm{s}}^{2}\mathcal{G}Mf\partial_{r}c_{\mathrm{s}}^{-2}-c_{\mathrm{s}}^{2}\frac{\partial}{\partial
r}\int_{0}^{z}\mathcal{G}Mf\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}\right).$
Then, Eq. (59) becomes:
$\frac{v_{\mathrm{g}\theta}^{2}}{r}=c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\left(\bar{P}\right)}{\mathrm{d}r}+\frac{\mathcal{G}M}{r^{2}}+c_{\mathrm{s}}^{2}\frac{\partial}{\partial
r}\int_{0}^{z}\mathcal{G}Mf\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}-c_{\mathrm{s}}^{2}\mathcal{G}Mf\partial_{r}c_{\mathrm{s}}^{-2}.$
(66)
Noting that :
$\frac{\partial}{\partial
r}\int_{0}^{z}f\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}=\int_{0}^{z}\frac{\partial
f}{\partial
r}\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}+\int_{0}^{z}f\partial_{r}\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime},$
(67)
and integrating the last term of the right hand side of Eq. (67) by parts
provides:
$\int_{0}^{z}f\partial_{r}\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}=f\partial_{r}c_{\mathrm{s}}^{-2}-\int_{0}^{z}\partial_{z^{\prime}}f\partial_{r}c_{\mathrm{s}}^{-2}\mathrm{d}z^{\prime}.$
(68)
Therefore, Eq. (66) reduces to:
$\frac{v_{\mathrm{g}\theta}^{2}}{r}=\frac{\mathcal{G}M}{r^{2}}+c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}+\mathcal{G}Mc_{\mathrm{s}}^{2}\int_{0}^{z}\left(\partial_{r}f\partial_{z^{\prime}}c_{\mathrm{s}}^{-2}-\partial_{z^{\prime}}f\partial_{r}c_{\mathrm{s}}^{-2}\right)\mathrm{d}z^{\prime},$
(69)
which can be more elegantly written as:
$\frac{v_{\mathrm{g}\theta}^{2}}{r}=\frac{\mathcal{G}M}{r^{2}}+c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}+\mathcal{G}Mc_{\mathrm{s}}^{2}\int_{0}^{z}\left[\nabla
f\times\nabla
c_{\mathrm{s}}^{-2}\right].\textbf{e}_{\theta}\mathrm{d}z^{\prime}.$ (70)
Thus, the expression of the azimuthal velocity of such a disc can be separated
in three terms, called the Keplerian, the pressure gradient and the baroclinic
terms respectively. This last term is neglected in most studies. For a three
dimensional disc, this term rigorously cancels for $c_{\mathrm{s}}=$ constant.
In this case, the flow is inviscid and derives from a potential, and isobars
and isodensity surfaces coincide: thus, there is no source of vorticity and
the azimuthal velocity depends only on the radial coordinate. This terms also
cancels out for flat discs in two dimensions. If the disc is vertically
isothermal, Eq. (70) becomes:
$\frac{v_{\mathrm{g}\theta}^{2}}{r}=\frac{\mathcal{G}M}{r^{2}}+c_{\mathrm{s}}^{2}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}-\mathcal{G}M\left[\frac{1}{r}-\frac{1}{\sqrt{r^{2}+z^{2}}}\right]\partial_{r}\mathrm{ln}c_{\mathrm{s}}^{-2}.$
(71)
### B.3 Radial profiles of surface density and temperature
In this section, we consider that the disc surface density and the temperature
(and thus the sound speed) depend only on the radial coordinate and are given
by the following power-law profiles:
$\left\\{\begin{array}[]{l}\Sigma\left(r\right)=\Sigma_{0}\left(\frac{r}{r_{0}}\right)^{-p}=\Sigma_{0}^{\prime}r^{-p},\\\
T\left(r\right)=T_{0}\left(\frac{r}{r_{0}}\right)^{-q}=T_{0}^{\prime}r^{-q},\\\
c_{\mathrm{s}}\left(r\right)=c_{\mathrm{s}0}\left(\frac{r}{r_{0}}\right)^{-q/2}=c^{\prime}_{\mathrm{s}0}r^{-q/2}.\end{array}\right.$
(72)
For vertically isothermal thin discs, the vertical density is therefore given
by Eq. (57) with the scale height given by Eq. (58), which can be expressed
as:
$H\left(r\right)=H_{0}\left(\frac{r}{r_{0}}\right)^{\frac{3}{2}-\frac{q}{2}}=H_{0}^{\prime}r^{\frac{3}{2}-\frac{q}{2}},$
(73)
with:
$H_{0}^{\prime}=\frac{c^{\prime}_{\mathrm{s}0}}{\sqrt{\mathcal{G}M}}.$ (74)
The expression of $\rho_{\mathrm{g}}$ compatible with the vertical hydrostatic
equilibrium and providing the power-law profile set by Eq. (72) is written as:
$\rho_{\mathrm{g}}=\rho_{\mathrm{g}0}^{\prime}r^{-x}\mathrm{e}^{-\frac{z^{2}}{2H^{2}\left(r\right)}}.$
(75)
Indeed:
$\displaystyle\int_{-\infty}^{+\infty}\rho_{\mathrm{g0}}^{\prime}r^{-x}\mathrm{e}^{-\frac{z^{2}}{2H^{2}\left(r\right)}}\mathrm{d}z$
$\displaystyle=$
$\displaystyle\rho_{\mathrm{g}0}^{\prime}\sqrt{2\pi}H\left(r\right)r^{-x},$
$\displaystyle=$ $\displaystyle\Sigma_{0}^{\prime}r^{-p}.$
Hence, with
$\Sigma_{0}^{\prime}=\sqrt{2\pi}\rho_{\rm{g}0}^{\prime}H_{0}^{\prime}$ and
$x=p-\frac{q}{2}+\frac{3}{2}$,
$\rho_{\mathrm{g}}=\frac{\Sigma_{0}^{\prime}}{\sqrt{2\pi}H_{0}^{\prime}}r^{-\left(p-\frac{q}{2}+\frac{3}{2}\right)}\mathrm{e}^{-\left[\frac{z^{2}}{2H_{0}^{\prime
2}r^{3-q}}\right]},$ (77)
which gives the correct surface density profile when integrated with respect
to $z$. With this expression of $\rho_{\mathrm{g}}$, $\bar{P}\left(r\right)$
is given by:
$\bar{P}\left(r\right)=c_{\mathrm{s}}^{2}\left(r\right)\rho_{\mathrm{g}}\left(r,z=0\right)=c_{\mathrm{s}0}^{\prime
2}\frac{\Sigma_{0}^{\prime}}{\sqrt{2\pi}H_{0}^{\prime}}r^{-\left(p+\frac{q}{2}+\frac{3}{2}\right)},$
(78)
which ensures that:
$\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}=-\frac{\left(p+\frac{q}{2}+\frac{3}{2}\right)}{r},$
(79)
and, using Eq. (71), we find:
$v_{\mathrm{g}\theta}=\sqrt{\frac{\mathcal{G}M}{r}-\left(p+\frac{q}{2}+\frac{3}{2}\right)c_{\mathrm{s}}^{\prime
2}r^{-q}-\mathcal{G}Mq\left(\frac{1}{r}-\frac{1}{\sqrt{r^{2}+z^{2}}}\right)}.$
(80)
## Appendix C Dimensionless quantities and equations of motion
To highlight the important physical parameters involved, we set
$v_{\mathrm{k}0}=\sqrt{\frac{\mathcal{G}M}{r_{0}}}$ and introduce
dimensionless quantities given by the following expressions:
$\left\\{\begin{array}[]{rcl}r/r_{0}&=&R,\\\
\mathcal{T}/\mathcal{T}_{0}&=&R^{-q},\\\ \Sigma/\Sigma_{0}&=&R^{-p},\\\
v_{\mathrm{k}}/v_{\mathrm{k}0}&=&R^{-\frac{1}{2}},\\\
H/H_{0}&=&R^{\frac{3}{2}-\frac{q}{2}},\\\ z/H_{0}&=&Z,\\\
v_{\mathrm{g}\theta}/v_{\mathrm{k}0}&=&\displaystyle\sqrt{\frac{1}{R}-\eta_{0}R^{-q}-q\left(\frac{1}{R}-\frac{1}{\sqrt{R^{2}+\phi_{0}Z^{2}}}\right)},\\\
\rho_{\mathrm{g}}/\rho_{\mathrm{g}0}&=&R^{-\left(p-\frac{q}{2}+\frac{3}{2}\right)}\mathrm{e}^{-\frac{Z^{2}}{2R^{3-q}}}.\\\
\end{array}\right.$ (81)
with:
$\eta_{0}=\left(p+\frac{q}{2}+\frac{3}{2}\right)\left(\frac{c_{\mathrm{s}0}}{v_{\mathrm{k}0}}\right)^{2}.$
(82)
The dimensionless parameter $\eta_{0}$ gives the order of magnitude of the
relative discrepancy between the Keplerian motion and the gas azimuthal
velocity. We note that:
$\left.\frac{rc_{\mathrm{s}}^{2}}{v_{\mathrm{k}}}\frac{\mathrm{d}\,\mathrm{ln}\bar{P}}{\mathrm{d}r}\right/v_{\mathrm{k}0}=-\eta_{0}R^{-q+1/2}.$
(83)
Then, we set $t_{\mathrm{k}0}=\sqrt{\frac{r_{0}^{3}}{\mathcal{G}M}}$ and
define
$\left\\{\begin{array}[]{rcl}\displaystyle\frac{t}{t_{\mathrm{k}0}}&=&\displaystyle
T\\\\[8.61108pt]
\displaystyle\frac{v_{r}}{v_{\mathrm{k}0}}&=&\displaystyle\frac{\mathrm{d}R}{\mathrm{d}T}=\tilde{v}_{r}\\\\[8.61108pt]
\displaystyle\frac{v_{\theta}}{v_{\mathrm{k}0}}&=&\displaystyle
R\frac{\mathrm{d}\theta}{\mathrm{d}T}=\tilde{v}_{\theta}\\\\[8.61108pt]
\displaystyle\frac{z}{H_{0}}&=&\displaystyle Z\\\\[8.61108pt]
\displaystyle\frac{v_{z}}{H_{0}/t_{\mathrm{k}0}}&=&\displaystyle\frac{\mathrm{d}Z}{\mathrm{d}T}.\end{array}\right.$
(84)
Writing the coefficient $\tilde{\mathcal{C}}\left(r,z\right)$ of the drag
force of Eq. (5) as:
$\tilde{\mathcal{C}}\left(r,z\right)=\mathcal{C}_{0}\mathcal{C}\left(R,Z\right),$
(85)
and using dimensionless coordinates, we have:
$\frac{\textbf{F}_{\mathrm{D}}/m_{\mathrm{d}}}{v_{\mathrm{k}0}/t_{\mathrm{k}0}}=-\frac{\mathcal{C}\left(R,Z\right)}{\left[\frac{s}{\left(v_{\mathrm{k}0}^{\lambda}t_{\mathrm{k}0}\mathcal{C}_{0}\right)^{\frac{1}{y}}}\right]^{y}}|\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}|^{\lambda}\left(\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}\right).$
(86)
We also introduce:
$s_{\mathrm{opt},0}=\left(v_{\mathrm{k}0}^{\lambda}t_{\mathrm{k}0}\mathcal{C}_{0}\right)^{\frac{1}{y}},$
(87)
and
$S_{0}=\frac{s}{s_{\mathrm{opt},0}},$ (88)
so that Eq. (86) becomes:
$\frac{\textbf{F}_{\mathrm{D}}/m_{\mathrm{d}}}{v_{\mathrm{k}0}/t_{\mathrm{k}0}}=-\frac{\mathcal{C}\left(R,Z\right)}{S_{0}^{y}}|\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}|^{\lambda}\left(\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}\right).$
(89)
Physically, $s_{\mathrm{opt},0}$ corresponds to the grain size at which the
drag stopping time equals to the Keplerian time at $r_{0}$. In Table 2, we
give the expressions of $y$, $\lambda$, $s_{\mathrm{opt},0}$,
$\mathcal{C}\left(R,Z\right)$ for the Epstein and the three Stokes drag
regimes. The dimensionless equations of motion for a dust grain are then:
$\left\\{\begin{array}[]{rcl}\displaystyle\frac{\mathrm{d}\tilde{v}_{r}}{\mathrm{d}T}-\frac{\tilde{v}_{\theta}^{2}}{R}+\frac{\tilde{v}_{r}}{S_{0}^{y}}\mathcal{C}\left(R,Z\right)|\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}|^{\lambda}+\frac{R}{\left(R^{2}+\phi_{0}Z^{2}\right)^{3/2}}&=&0\\\\[8.61108pt]
\lx@intercol\displaystyle\frac{\mathrm{d}\tilde{v}_{\theta}}{\mathrm{d}T}+\frac{\tilde{v}_{\theta}\tilde{v}_{r}}{R}+\hfil\lx@intercol&&\\\
\frac{\left(\tilde{v}_{\theta}-\sqrt{\frac{1}{R}-\eta_{0}R^{-q}-q\left(\frac{1}{R}-\frac{1}{\sqrt{R^{2}+\phi_{0}Z^{2}}}\right)}\right)}{S_{0}^{y}}\mathcal{C}\left(R,Z\right)|\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}|^{\lambda}&=&0\\\\[8.61108pt]
\displaystyle\frac{\mathrm{d}^{2}Z}{\mathrm{d}T^{2}}+\frac{1}{S_{0}^{y}}\frac{\mathrm{d}Z}{\mathrm{d}T}\mathcal{C}\left(R,Z\right)|\tilde{\textbf{v}}-\tilde{\textbf{v}}_{\mathrm{g}}|^{\lambda}+\frac{Z}{\left(R^{2}+\phi_{0}Z^{2}\right)^{3/2}}&=&0.\end{array}\right.$
(90)
Table 2: Expressions of the coefficients $y$, $\lambda$, $s_{\mathrm{opt},0}$ and $\mathcal{C}\left(R,Z\right)$ for different drag regimes. Drag regime | $y$ | $\lambda$ | $s_{\mathrm{opt},0}$ | $\mathcal{C}\left(R,Z\right)$
---|---|---|---|---
Epstein | 1 | 0 | $\displaystyle\frac{\Sigma_{0}}{\sqrt{2\pi}\rho_{\mathrm{d}}}$ | $\displaystyle R^{-\left(p+\frac{3}{2}\right)}e^{-\frac{Z^{2}}{2R^{3-q}}}$
Stokes ($R_{\mathrm{g}}<1$) | 2 | 0 | $\displaystyle\sqrt{\frac{9t_{\mathrm{k}0}\mu_{0}}{2\rho_{\mathrm{d}}}}$ | $\displaystyle R^{-\frac{q}{2}}$
Stokes ($1<R_{\mathrm{g}}<800$) | 1.6 | 0.4 | $\displaystyle\left(\frac{18r_{0}\nu_{0}^{0.6}\Sigma_{0}}{v_{\mathrm{k}0}^{0.6}\sqrt{2\pi}2^{1.6}\rho_{\mathrm{d}}H_{0}}\right)^{\frac{1}{1.6}}$ | $\displaystyle R^{-\left(\frac{2p}{5}+\frac{q}{10}+\frac{3}{5}\right)}e^{-\frac{2}{5}\frac{Z^{2}}{2R^{3-q}}}$
Stokes ($800<R_{\mathrm{g}}$) | 1 | 1 | $\displaystyle\frac{1.32r_{0}\Sigma_{0}}{8\rho_{\mathrm{d}}\sqrt{2\pi}H_{0}}$ | $\displaystyle R^{-\left(p-\frac{q}{2}+\frac{3}{2}\right)}e^{-\frac{Z^{2}}{2R^{3-q}}}$
## Appendix D Lemma for the different expansions
Lemma: Let $x$ be either $r$ or $\theta$ and $i$ the order of the perturbative
expansion. If:
* •
$\tilde{v}_{r0}=0$, and
* •
$\tilde{v}_{xi}$ can be written as a function of $R$
($\tilde{v}_{xi}=f\left(R\right)$) with $f=\mathcal{O}\left(1\right)$ of the
expansion in $\eta_{0}$,
then, $\displaystyle\frac{\mathrm{d}\tilde{v}_{xi}}{\mathrm{d}T}$ is of order
$\mathcal{O}\left(\eta_{0}\right)$.
Proof:
$\frac{\mathrm{d}\tilde{v}_{xi}}{\mathrm{d}T}=\frac{\mathrm{d}\tilde{v}_{xi}}{\mathrm{d}R}\frac{\mathrm{d}R}{\mathrm{d}T}=\tilde{v}_{r}f^{\prime}\left(R\right)=\eta_{0}\tilde{v}_{r1}f^{\prime}\left(R\right)+\mathcal{O}\left(\eta_{0}^{2}\right)=\mathcal{O}\left(\eta_{0}\right).$
(91)
## Appendix E Epstein regime: perturbation analysis
* •
Order $\mathcal{O}\left(1\right)$: At this order of expansion,
$\eta_{0}R^{-q}$ is negligible compared to $\frac{1}{R}$. Thus, substituting
Eq. (13) into Eq. (11) provides
$\left\\{\begin{array}[]{l}\displaystyle\frac{\mathrm{d}\tilde{v}_{r0}}{\mathrm{d}T}=\displaystyle\frac{\tilde{v}_{\theta
0}^{2}}{R}-\frac{\tilde{v}_{r0}}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}-\frac{1}{R^{2}}\\\
\displaystyle\frac{\mathrm{d}\tilde{v}_{\theta
0}}{\mathrm{d}T}=\displaystyle-\frac{\tilde{v}_{\theta
0}\tilde{v}_{r0}}{R}-\frac{\left(\tilde{v}_{\theta
0}-\sqrt{\frac{1}{R}}\right)}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}.\end{array}\right.$
(92)
At this stage, we do not know the order of $\frac{\mathrm{d}\tilde{v}_{\theta
0}}{\mathrm{d}T}$. We show in the lemma of Appendix D that
$\frac{\mathrm{d}\tilde{v}_{\theta
0}}{\mathrm{d}T}=\mathcal{O}\left(\eta_{0}\right)$. Applying this lemma, we
see that at the order $\mathcal{O}\left(1\right)$, taking $\tilde{v}_{r0}=0$
and $\tilde{v}_{\theta 0}=\sqrt{\frac{1}{R}}$ (which ensures that
$\frac{\mathrm{d}\tilde{v}_{\theta
0}}{\mathrm{d}T}=\mathcal{O}\left(\eta_{0}\right)$) is a relevant solution for
the equations of motion (which corresponds to circular Keplerian motion).
Thus,
$\left\\{\begin{array}[]{l}\displaystyle\tilde{v}_{r0}=\displaystyle 0\\\
\displaystyle\tilde{v}_{\theta
0}=\displaystyle\sqrt{\frac{1}{R}}.\end{array}\right.$ (93)
* •
Order $\mathcal{O}\left(\eta_{0}\right)$: Applying the lemma in this order of
expansion and noting that
$\frac{\mathrm{d}\tilde{v}_{\theta
0}}{\mathrm{d}T}=-\frac{1}{2}R^{-3/2}\tilde{v}_{r}=-\frac{\eta_{0}}{2}R^{-3/2}\tilde{v}_{r1}+\mathcal{O}\left(\eta_{0}^{2}\right),$
(94)
Eq. (11) becomes
$\left\\{\begin{array}[]{l}\displaystyle
0=-\frac{R^{-\left(p+\frac{3}{2}\right)}}{S_{0}}\tilde{v}_{r1}+\displaystyle
2R^{-3/2}\tilde{v}_{\theta 1}\\\ \displaystyle
0=\displaystyle-\frac{1}{2}R^{-3/2}\tilde{v}_{r1}-\frac{\left(\tilde{v}_{\theta
1}-\frac{1}{\eta_{0}}\left(\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}-\sqrt{\frac{1}{R}}\right)\right)}{S_{0}}R^{-\left(p+\frac{3}{2}\right)}.\end{array}\right.$
(95)
Solving the linear system Eq. (95) for $\left(\tilde{v}_{r1},\tilde{v}_{\theta
1}\right)$ provides
$\left\\{\begin{array}[]{l}\displaystyle\tilde{v}_{r1}=\displaystyle-\frac{1}{\eta_{0}}\frac{2S_{0}R^{p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)}{1+R^{2p}S_{0}^{2}}\\\
\displaystyle\tilde{v}_{\theta
1}=\displaystyle-\frac{1}{\eta_{0}}\frac{R^{-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)}{1+R^{2p}S_{0}^{2}}.\end{array}\right.$
(96)
In addition to the expression of $\tilde{v}_{r}$ given in Sect. 3, we also
note that
$\begin{array}[]{rcl}\tilde{v}_{\theta}&=&\displaystyle\sqrt{\frac{1}{R}}+\eta_{0}\tilde{v}_{\theta
1}+\mathcal{O}\left(\eta_{0}^{2}\right)\\\
&=&\displaystyle\sqrt{\frac{1}{R}}-\frac{R^{-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)}{1+R^{2p}S_{0}^{2}}+\mathcal{O}\left(\eta_{0}^{2}\right),\end{array}$
(97)
which provides with Eq. (15)
$\tilde{v}_{\theta}=\sqrt{\frac{1}{R}}-\frac{\eta_{0}R^{-q+\frac{1}{2}}}{2\left(1+R^{2p}S_{0}^{2}\right)}.$
(98)
## Appendix F Link with W77’s original derivation
Following W77’s historic reasoning for small grains (see Sect. 3.1), we
perform a perturbative expansion of the radial equation of motion with the
$S_{0}$ variable. We verify that taking the limit at small $\eta_{0}$ provides
the expression found for the A-mode in the NSH86 expansion. (Formally, we will
show that $\lim\limits_{\begin{subarray}{c}S_{0}\ll
1\end{subarray}}\lim\limits_{\begin{subarray}{c}\eta_{0}\ll
1\end{subarray}}\left[...\right]=\lim\limits_{\begin{subarray}{c}\eta_{0}\ll
1\end{subarray}}\lim\limits_{\begin{subarray}{c}S_{0}\ll
1\end{subarray}}\left[...\right]$ ). Hence, we set:
$\left\\{\begin{array}[]{l}\displaystyle\tilde{v}_{r}=\displaystyle\tilde{v}_{r0}+S_{0}\tilde{v}_{r1}+S_{0}^{2}\tilde{v}_{r2}+\mathcal{O}\left(S_{0}^{3}\right),\\\
\displaystyle\tilde{v}_{\theta}=\displaystyle\tilde{v}_{\theta
0}+S_{0}\tilde{v}_{\theta 1}+S_{0}^{2}\tilde{v}_{\theta
2}+\mathcal{O}\left(S_{0}^{3}\right),\end{array}\right.$ (99)
where we have used for convenience the same formalism as for the expansion in
$\eta_{0}$ – see Eq. (13) (noting of course that $\tilde{v}$ represents
different functions). An important point is that the lemma of Appendix D holds
when substituting $S_{0}$ to $\eta_{0}$. Therefore, substituting Eq. (99) into
Eq. (11) provides the equations of motion for different orders of
$\mathcal{O}\left(S_{0}\right)$:
* •
Order $\mathcal{O}\left(\frac{1}{S_{0}}\right)$: Eq. (11) provides
$\mathcal{O}\left(1\right)$ expressions for the velocities:
$\left\\{\begin{array}[]{r@{\ }l}\displaystyle\tilde{v}_{r0}=&\displaystyle
0,\\\ \displaystyle\tilde{v}_{\theta
0}=&\displaystyle\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}.\end{array}\right.$ (100)
In this order of expansion, the azimuthal velocity corresponds to the sub-
Keplerian velocity of the gas. There is no radial motion.
* •
Order $\mathcal{O}\left(1\right)$:
$\left\\{\begin{array}[]{r@{\
}l}\displaystyle\tilde{v}_{r1}=&\displaystyle-2R^{p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right),\\\
\displaystyle\tilde{v}_{\theta 1}=&\displaystyle 0.\end{array}\right.$ (101)
* •
Order $\mathcal{O}\left(S_{0}\right)$:
$\left\\{\begin{array}[]{r@{\ }l}\displaystyle\tilde{v}_{r2}=&\displaystyle
0,\\\ \displaystyle\tilde{v}_{\theta 2}=&\displaystyle
R^{p+\frac{3}{2}}\left[\eta_{0}\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}R^{p-q-\frac{1}{2}}\right.\\\
&\qquad\quad\left.\displaystyle+\frac{1}{2}\frac{-\frac{1}{R^{2}}+\eta_{0}qR^{-q-1}}{\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}\eta_{0}R^{p-q+\frac{1}{2}}}\right].\end{array}\right.$
(102)
Finally, we obtain expressions for $\tilde{v}_{r}$ and $\tilde{v}_{\theta}$:
$\left\\{\begin{array}[]{r@{\
}l}\displaystyle\tilde{v}_{r}=&\displaystyle-2S_{0}R^{p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)+\mathcal{O}\left(S_{0}^{3}\right),\\\\[8.61108pt]
\displaystyle\tilde{v}_{\theta}=&\displaystyle\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}+S_{0}^{2}R^{p+\frac{3}{2}}\left[\eta_{0}\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}R^{p-q-\frac{1}{2}}\right.\\\
&\qquad\quad\left.\displaystyle+\frac{1}{2}\frac{-\frac{1}{R^{2}}+\eta_{0}qR^{-q-1}}{\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}\eta_{0}R^{p-q+\frac{1}{2}}}\right]+\mathcal{O}\left(S_{0}^{3}\right).\end{array}\right.$
(103)
We now compare the NSH86 expansion at $R^{p}S_{0}\ll 1$ (A-mode) and the W77
expansion at $\eta_{0}\ll 1$.
* •
NSH86: From Eqs. (16) and Eq. (98):
$\left\\{\begin{array}[]{r@{\
}l}\displaystyle\tilde{v}_{r}=&\displaystyle-\frac{\eta_{0}S_{0}R^{p-q+\frac{1}{2}}}{1+R^{2p}S_{0}^{2}}+\mathcal{O}\left(\eta_{0}^{2}\right)\\\
=&\displaystyle-\eta_{0}S_{0}R^{p-q+\frac{1}{2}}+\mathcal{O}\left(\eta_{0}^{2}\right)+\mathcal{O}\left(S_{0}^{2}\right),\\\\[8.61108pt]
\displaystyle\tilde{v}_{\theta}=&\displaystyle\sqrt{\frac{1}{R}}-\frac{\eta_{0}}{2}\frac{R^{-q+\frac{1}{2}}}{1+S_{0}^{2}R^{2p}}+\mathcal{O}\left(\eta_{0}^{2}\right)\\\
=&\displaystyle\sqrt{\frac{1}{R}}-\frac{\eta_{0}}{2}R^{-q+\frac{1}{2}}+\frac{\eta_{0}S_{0}^{2}}{2}R^{2p-q+\frac{1}{2}}+\mathcal{O}\left(\eta_{0}^{2}\right)+\mathcal{O}\left(S_{0}^{3}\right).\end{array}\right.$
(104)
* •
W77 small grains: From Eq. (103):
$\left\\{\begin{array}[]{r@{\
}l}\displaystyle\tilde{v}_{r}=&\displaystyle-2S_{0}R^{p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)+\mathcal{O}\left(S_{0}^{2}\right)\\\
=&\displaystyle-\eta_{0}S_{0}R^{p-q+\frac{1}{2}}+\mathcal{O}\left(S_{0}^{2}\right)+\mathcal{O}\left(\eta_{0}^{2}\right),\\\\[8.61108pt]
\displaystyle\tilde{v}_{\theta}=&\displaystyle\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}+S_{0}^{2}R^{p+\frac{3}{2}}\Big{[}\eta_{0}\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}R^{p-q-\frac{1}{2}}\\\
&\displaystyle+\frac{1}{2}\frac{\eta_{0}R^{p-q+\frac{1}{2}}\left(-\frac{1}{R^{2}}+\eta_{0}qR^{-q-1}\right)}{\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}}\Big{]}+\mathcal{O}\left(S_{0}^{3}\right)\\\
=&\displaystyle\sqrt{\frac{1}{R}}-\frac{\eta_{0}}{2}R^{-q+\frac{1}{2}}+\frac{\eta_{0}S_{0}^{2}}{2}R^{2p-q+\frac{1}{2}}+\mathcal{O}\left(S_{0}^{3}\right)+\mathcal{O}\left(\eta_{0}^{2}\right).\end{array}\right.$
(105)
Clearly, Eqs. (104) and (105) are identical, demonstrating that the theories
of W77 and NSH86 are consistent. We also note that if the simplification of
Eq. (15) is not performed, the two W77 expansions directly appear as the
expansion of NSH86 in $\mathcal{O}\left(S_{0}\right)$ or
$\mathcal{O}\left(S_{0}^{-1}\right)$.
Now, in the case of large grains, we perform a perturbative expansion of the
radial equation of motion with respect to $\frac{1}{S_{0}}$ while assuming
that $S_{0}R^{p}\gg 1$, and verify that taking the limit at small $\eta_{0}$
provides the expression found for the B-mode in the NSH86 expansion.
Taking the same precautions as for the previous expansions, we write:
$\left\\{\begin{array}[]{r@{\
}l}\displaystyle\tilde{v}_{r}=&\displaystyle\tilde{v}_{r0}+\frac{1}{S_{0}}\tilde{v}_{r1}+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right),\\\
\displaystyle\tilde{v}_{\theta}=&\displaystyle\tilde{v}_{\theta
0}+\frac{1}{S_{0}}\tilde{v}_{\theta
1}+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right).\end{array}\right.$ (106)
Following the same method as for the small grain sizes expansion, we obtain:
* •
Order $\mathcal{O}\left(1\right)$:
$\left\\{\begin{array}[]{l}\tilde{v}_{r0}=0,\\\ \tilde{v}_{\theta
0}=\displaystyle\sqrt{\frac{1}{R}}.\end{array}\right.$ (107)
It this order of expansion, the azimuthal velocity of the grain is the
standard Keplerian velocity.
* •
Order $\mathcal{O}\left(\frac{1}{S_{0}}\right)$:
$\left\\{\begin{array}[]{l}\tilde{v}_{r1}=\displaystyle-2\left(\sqrt{\frac{1}{R}}-\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}\right)R^{-p},\\\
\tilde{v}_{\theta 1}=0.\end{array}\right.$ (108)
The expansion at higher order is more complicated and will not be used for
further developments. At the order
$\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right)$, we have for $\tilde{v}_{r}$ and
$\tilde{v}_{\theta}$:
$\left\\{\begin{array}[]{l}\displaystyle\tilde{v}_{r}=\displaystyle-\frac{2}{S_{0}}\left(\sqrt{\frac{1}{R}}-\sqrt{\frac{1}{R}-\eta_{0}R^{-q}}\right)R^{-p}+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right),\\\
\displaystyle\tilde{v}_{\theta}=\displaystyle\sqrt{\frac{1}{R}}+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right).\end{array}\right.$
(109)
We now compare the expressions provided by the NSH86 expansion at
$R^{2p}S_{0}^{2}\gg 1$ (B-mode) and the W77 expansion at $\eta_{0}\ll 1$ for
the radial velocity:
* •
NSH86:
$\displaystyle\tilde{v}_{r}$ $\displaystyle=$
$\displaystyle-\frac{\eta_{0}S_{0}R^{p-q+\frac{1}{2}}}{1+R^{2p}S_{0}^{2}}+\mathcal{O}\left(\eta_{0}^{2}\right),$
$\displaystyle=$
$\displaystyle-\frac{\eta_{0}}{S_{0}}R^{-p-q+\frac{1}{2}}+\mathcal{O}\left(\eta_{0}^{2}\right)+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right).$
* •
W77 large grains:
$\displaystyle\tilde{v}_{r}$ $\displaystyle=$
$\displaystyle-\frac{2}{S_{0}}R^{-p-\frac{1}{2}}\left(1-\sqrt{1-\eta_{0}R^{-q+1}}\right)+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right),$
$\displaystyle=$
$\displaystyle-\frac{\eta_{0}}{S_{0}}R^{-p-q+\frac{1}{2}}+\mathcal{O}\left(\frac{1}{S_{0}^{2}}\right)+\mathcal{O}\left(\eta_{0}^{2}\right).$
Once again, we show that the W77 and NSH86 theories are consistent.
## Appendix G Asymptotic radial behaviour of a single grain
Figure 13: The discrepancy (bottom panel) between the exact motion (top panel)
and its NSH86 approximation (central panel) is negligible. This is illustrated
plotting the radial motion of dust grains for $S_{0}=10^{-2}$,
$\eta_{0}=10^{-2}$ and for $p=0,q=\frac{3}{4}$ (solid) and
$p=\frac{3}{2},q=\frac{3}{4}$ (dashed). Top: $R_{1}\left(T\right)$, middle:
$R_{2}\left(T\right)$, bottom: relative difference
$\left(R_{2}\left(T\right)-R_{1}\left(T\right)\right)/R_{1}\left(T\right)$.
Noting $R_{1}\left(T\right)$ the position of a grain integrated directly from
the equation of motion (Eq. (11)) and $R_{2}\left(T\right)$ the position
integrated from the NSH86 approximation (Eq. (34)), we highlight (Fig. 13)
that the discrepancy between the motion from the exact equations and its NSH86
approximation is negligible (the relative error is lower than $10^{-3}$ for
all the considered sizes). It is therefore justified to use the analytical
results derived in Sect. 3 to interpret the grain behaviour.
Figure 14: Values of $S_{\mathrm{m}}$ (left) and
$\eta_{0}T_{\mathrm{m}}\left(S_{\mathrm{m}}\right)$ (right) in the ($p$,$q$)
plane.
Thus, from Eq. (34), we see that the time for a grain starting at $R=1$ to
reach some final radius $R_{\mathrm{f}}$ is minimized for an optimal grain
size $S_{\mathrm{m,f}}$ given by
$S_{\mathrm{m,f}}=\sqrt{\frac{p+q+\frac{1}{2}}{-p+q+\frac{1}{2}}\times
k_{\mathrm{f}}\left(R_{\mathrm{f}}\right)}\,,$ (112)
with
$k_{\mathrm{f}}\left(R_{\mathrm{f}}\right)=\left\\{\begin{array}[]{ll}\frac{\left(1-R_{\mathrm{f}}^{-p+q+\frac{1}{2}}\right)}{\left(1-R_{\mathrm{f}}^{p+q+\frac{1}{2}}\right)}&\mathrm{if}\leavevmode\nobreak\
p+q+\frac{1}{2}\neq 0\\\\[4.30554pt]
-\mathrm{ln}\left(R_{\mathrm{f}}\right)&\mathrm{if}\leavevmode\nobreak\
p+q+\frac{1}{2}=0.\\\ \end{array}\right.$ (113)
As shown in Eq. (34), the outcome of the grain radial motion depends on the
value of $-p+q+\frac{1}{2}$:
* •
If $-p+q+\frac{1}{2}\leq 0$:
$\lim\limits_{\begin{subarray}{c}T\to+\infty\end{subarray}}R=0.$ (114)
For such disc profiles, all grains pile-up and fall onto the central star in
an infinite time. Indeed, the surface density profile given by $p\geq
q+\frac{1}{2}$ is steep enough to counterbalance the increase of the
acceleration due to the pressure gradient. Therefore, grains fall onto the
central star in an infinite time, whatever their initial size. Such an
evolution happens because the grains always end migrating in the A-mode when
they reach the disc’s inner regions. A crucial consequence is that grains are
not depleted on the central star and therefore stay in the disc where they can
potentially form planet embryos.
* •
If $-p+q+\frac{1}{2}>0$:
$\lim\limits_{\begin{subarray}{c}T\to T_{\mathrm{m}}\end{subarray}}R=0,$ (115)
where
$T_{\mathrm{m}}=\frac{1}{\eta_{0}S_{0}}\left(\frac{1}{-p+q+\frac{1}{2}}+\frac{S_{0}^{2}}{p+q+\frac{1}{2}}\right).$
(116)
In this case, grains fall onto the central star in a finite time. The surface
density profile given by $p<q+\frac{1}{2}$ is now too flat to counterbalance
the increasing acceleration due to pressure gradient. We note that:
* –
For small sizes ($S_{0}\ll 1$),
$T_{\mathrm{m}}=\mathcal{O}\left(S_{0}\eta_{0}\right)$.
* –
For large sizes ($S_{0}\gg 1$),
$T_{\mathrm{m}}=\mathcal{O}\left(\frac{S_{0}}{\eta_{0}}\right)$.
* –
$T_{\mathrm{m}}$ reaches a minimal value for a size $S_{\mathrm{m}}$ given by
$S_{\mathrm{m}}=\sqrt{\frac{p+q+\frac{1}{2}}{-p+q+\frac{1}{2}}}.$ (117)
Therefore
$T_{\mathrm{m}}\left(S_{\mathrm{m}}\right)=\frac{2}{\eta_{0}\sqrt{\left(p+q+\frac{1}{2}\right)\left(-p+q+\frac{1}{2}\right)}}.$
(118)
$S_{\mathrm{m}}$ is of order unity and corresponds to an optimal size of
migration. Values of $S_{\mathrm{m}}$ and
$\eta_{0}T_{\mathrm{m}}\left(S_{\mathrm{m}}\right)$ in the $(p,q)$ plane are
shown in Fig. 14. When $S\simeq S_{\mathrm{m}}$, both the A- and B-modes
contribute in an optimal way to the grains radial motion.
In this case, grains can be efficiently accreted by the central star if
$S\simeq S_{\mathrm{m}}=\mathcal{O}\left(1\right)$. Thus, they can not
contribute to the formation of pre-planetesimals. This process is called the
radial-drift barrier for planet formation.
|
arxiv-papers
| 2011-11-14T01:49:47 |
2024-09-04T02:49:24.281474
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Guillaume Laibe (Monash), Jean-Fran\\c{c}ois Gonzalez (CRAL), and Sarah\n T. Maddison (Swinburne)",
"submitter": "Guillaume Laibe",
"url": "https://arxiv.org/abs/1111.3083"
}
|
1111.3089
|
# Dusty gas with SPH — II. Implicit timestepping and astrophysical drag
regimes
Guillaume Laibe, Daniel J. Price
Centre for Stellar and Planetary Astrophysics and School of Mathematical
Sciences, Monash University, Clayton, Vic 3800, Australia
###### Abstract
In a companion paper (Laibe & Price, 2011b), we have presented an algorithm
for simulating two-fluid gas and dust mixtures in Smoothed Particle
Hydrodynamics (SPH). In this paper, we develop an implicit timestepping method
that preserves the exact conservation of the both linear and angular momentum
in the underlying SPH algorithm, but unlike previous schemes, allows the
iterations to converge to arbitrary accuracy and is suited to the treatment of
non-linear drag regimes. The algorithm presented in Paper I is also extended
to deal with realistic astrophysical drag regimes, including both linear and
non-linear Epstein and Stokes drag. The scheme is benchmarked against the test
suite presented in Paper I, including i) the analytic solutions of the
dustybox problem and ii) solutions of the dustywave, dustyshock, dustysedov
and dustydisc obtained with explicit timestepping. We find that the implicit
method is 1–10 times faster than the explicit temporal integration when the
ratio $r$ between the the timestep and the drag stopping time is $1\lesssim
r\lesssim 1000$.
###### keywords:
hydrodynamics — methods: numerical — ISM: dust, extinction — protoplanetary
discs — planets and satellites: formation
††pagerange: Dusty gas with SPH — II. Implicit timestepping and astrophysical
drag regimes–References††pubyear: 2011
## 1 Introduction
Dust in cold astrophysical systems spans a huge range of sizes from sub-micron
sized grains in the interstellar medium to kilometre sized planetesimals
involved in planet formation. Moreover, the ratio of dust to gas as well as
the density and temperature of the gaseous environment in which dust is
embedded can also vary strongly. Handling this full range of physical
parameters presents a challenge to numerical schemes designed to simulate
dusty gas in astrophysics. The main challenges are i) that at high drag (e.g.
small grains), the small timestep required means that purely explicit
timestepping methods become prohibitive and ii) that a wide range of physical
drag prescriptions, including non-linear drag regimes, need to be handled by
the code.
Two main prescriptions for drag between gas and solid particles are applicable
to astrophysics: The Epstein regime — where the gas surrounding a grain can be
treated as a dilute medium — and the Stokes regime — where the grains can be
treated as solid bodies surrounded by a fluid —(see e.g. Baines et al. 1965;
Stepinski & Valageas 1996). The dependance of the drag term on local
parameters of the gas (density, temperature) and the dust (typical grain size,
mass) differ between the two regimes, in turn implying very different dynamics
for the dust grains. For example, in a protoplanetary disc, both of these
regimes may be applicable in different regions of the disc.
In a companion paper (Laibe & Price 2011b, hereafter Paper I), we have
developed a new algorithm for treating two-fluid gas-dust astrophysical
mixtures in Smoothed Particle Hydrodynamics (SPH). Benchmarking of the method
demonstrated that the algorithm gives accurate solutions on a range of test
problems relevant to astrophysics and substantially improve previous
algorithms (Monaghan & Kocharyan, 1995). However, in Paper I, we used only a
simple explicit time stepping and considered only linear drag regimes with a
constant drag coefficient. In this paper, we present an implicit timestepping
method that can be applied to both linear and non-linear drag regimes, which
is both more accurate and more general than the scheme proposed by Monaghan
(1997). We also discuss the SPH implementation of both the Epstein and the
Stokes regimes in their full generality.
The paper is organised as follows: The equations of motion and the
characteristics of the different astrophysical drag regimes for gas and dust
mixtures are given in Sec. 2. We summarise the SPH formalism used for
integrating these equations (derived in detail in Paper I) in Sec. 3 and
extend it to deal with the drag regimes encountered in astrophysics.
Particular attention is paid in Sec. 4 to improving the Monaghan (1997)
implicit timestepping scheme, including its generalisation for non-linear drag
regimes. Finally, the algorithm for non-linear drag regimes is tested against
the analytic solutions of the dustybox problem, as well as the dustywave,
dustysedov, dustyshock and dustydisc tests, in Sec. 5.
## 2 Gas and dust evolution in astrophysical systems
### 2.1 Evolution equations
The equations describing the evolution of astrophysical gas and dust mixtures,
where dust is treated as a pressureless, inviscid, continuous fluid have been
described in detail in Paper I. The equations in the continuum limit are given
by:
$\displaystyle\frac{\partial\hat{\rho}_{\mathrm{g}}}{\partial
t}+\nabla.\left(\hat{\rho}_{\mathrm{g}}\textbf{v}_{\mathrm{g}}\right)$
$\displaystyle=$ $\displaystyle 0,$ (1)
$\displaystyle\frac{\partial\hat{\rho}_{\mathrm{d}}}{\partial
t}+\nabla.\left(\hat{\rho}_{\mathrm{d}}\textbf{v}_{\mathrm{d}}\right)$
$\displaystyle=$ $\displaystyle 0,$ (2)
$\displaystyle\hat{\rho}_{\mathrm{g}}\left(\frac{\partial\textbf{v}_{\mathrm{g}}}{\partial
t}+\textbf{v}_{\mathrm{g}}.\nabla\textbf{v}_{\mathrm{g}}\right)$
$\displaystyle=$ $\displaystyle-\theta\phantom{.}\nabla P_{\rm
g}+\hat{\rho}_{\mathrm{g}}\textbf{f}-F_{\rm drag}^{\rm V},$ (3)
$\displaystyle\hat{\rho}_{\mathrm{d}}\left(\frac{\partial\textbf{v}_{\mathrm{d}}}{\partial
t}+\textbf{v}_{\mathrm{d}}.\nabla\textbf{v}_{\mathrm{d}}\right)$
$\displaystyle=$ $\displaystyle-\left(1-\theta\right)\nabla P_{\rm
g}+\hat{\rho}_{\mathrm{d}}\textbf{f}+F_{\rm drag}^{\rm V},$ (4)
$\displaystyle\frac{{\rm d}u_{{\rm g}}}{{\rm d}t}$ $\displaystyle=$
$\displaystyle-\frac{P_{\mathrm{g}}}{\hat{\rho}_{\mathrm{g}}}\left[\theta\nabla\cdot{\bf
v}_{\rm g}+(1-\theta)\nabla\cdot{\bf v}_{\rm d}\right]+\Lambda_{\rm drag}.$
(5)
where the subscripts ${\rm g}$ and ${\rm d}$ refer to the gas and dust,
respectively such that $P_{\rm g}$ is the gas pressure,
$\textbf{v}_{\mathrm{g}}$ and $\textbf{v}_{\mathrm{d}}$ are the fluid
velocities and $u$ is the specific internal energy of gas. The volume
densities of gas and dust ($\hat{\rho}_{\mathrm{g}}$ and
$\hat{\rho}_{\mathrm{d}}$, respectively) are related to the corresponding
intrinsic densities ($\rho_{\rm{g}}$ and $\rho_{\rm{d}}$, respectively)
according to
$\displaystyle\hat{\rho}_{\mathrm{d}}$ $\displaystyle=$
$\displaystyle(1-\theta)\rho_{\mathrm{d}},$ (6)
$\displaystyle\hat{\rho}_{\mathrm{g}}$ $\displaystyle=$
$\displaystyle\theta\rho_{\mathrm{g}},$ (7)
where $\theta$ is the volume filling fraction of the dust. Finally, the drag
force and heating terms are given by:
$F_{\rm drag}^{\rm V}=K(\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}}),$ (8)
and
$\Lambda_{\rm drag}=K(\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}})^{2}.$
(9)
The drag coefficient K has dimensions of mass per unit volume per unit time
and is generally a function of the relative velocity between the two fluids
$\Delta v\equiv|\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}}|$, implying a
non-linear drag regime with respect to the differential velocity between the
gas and the dust. In most of the astrophysical systems, the dust is diluted
enough into the gas so that the gas filling fraction is $\theta\simeq 1$ to a
very good level of approximation, such that the dust buoyancy term
$\left(1-\theta\right)\nabla P_{\rm g}$ is negligible.
### 2.2 Astrophysical drag regimes
Microscopic collisions of gas molecules on a single dust grain result in a net
exchange of momentum which is equivalent to a drag force $\bf{F}_{\rm{drag}}$
between the two phases. Two limiting cases occur when comparing the typical
geometrical size $s$ of a dust grain to the mean free path $\lambda_{\rm{g}}$
of the gas.
When the typical grain size is negligible compared to the collisional mean
free path of the gas particles ($s\ll\lambda_{\rm{g}}$), the grains are
surrounded by a dilute gas phase and may be treated using the Epstein drag
prescription. In this limit, the analytic expression of the resulting drag
force has been derived, assuming spherical, compact grains with homogeneous
composition, for both specular and diffuse reflections on the grain surface
(see Baines et al. 1965 for the complete derivation). These expressions have
been widely used in astrophysical studies (see e.g. Chiang & Youdin 2010 for
references), sometimes incorrectly where the grains are known to be porous and
have fractal structures (Blum & Wurm, 2008) (thus breaking the assumptions of
the Epstein prescription).
For grain sizes larger than the collisional mean free path
($s\gg\lambda_{\rm{g}}$), grains experience a local differential velocity with
respect to a uniform viscous flow and should be treated using the Stokes drag
prescription (Fan & Zhu, 1998). In this case, the momentum is diffused by
viscosity into the fluid, which implies that the drag expression strongly
depends on the local Reynolds number defined according to
$R_{\mathrm{d}}=\frac{2s\left|\textbf{v}_{\mathrm{d}}-\textbf{v}_{\mathrm{g}}\right|}{\nu},$
(10)
where $\nu$ is the kinematic viscosity of the gas. Analytic expressions for
the drag force can be derived at small Reynolds numbers. At higher Reynolds
number, the drag law is inferred from experiments. Rigorously, additional
contributions to the drag should arise from the grain acceleration (carried
mass and Basset contribution), grain rotation (Magnus) and in the presence of
strong local shear, pressure and temperature gradients (e.g. Fan & Zhu 1998).
These corrections are negligible in nearly all astrophysical contexts.
No current analytic theory describes how both the gas and the dust fluid
exchange momentum in the intermediate regime (i.e. $s\simeq\lambda_{\rm{g}}$).
Generally, an asymptotic continuous interpolation between the two limiting
Epstein and Stokes regimes is used. Stepinski & Valageas (1996) suggest
adopting $s=\frac{4}{9}\lambda_{\rm{g}}$ as a means of obtaining a smooth
transition. It should be noted that although this approach is convenient,
there is no clear measure of the physical accuracy of this assumption.
#### 2.2.1 Epstein regime for dilute media
In a dilute medium ($\lambda_{\mathrm{g}}>{4s}/{9}$), grains are small enough
not to disturb the Maxwellian distribution of the gas velocity. Assuming
grains are spherical, that the mass of a gas molecule is negligible compared
to the mass of a dust grain and that the reflection of gas particles from
collisions with dust grains are specular, the expression of the drag force on
a single grain $\bf{F}_{\rm{drag}}$ (which differs from the volume force
$F_{\rm drag}^{\rm V}$ by a factor $\hat{\rho}_{\mathrm{d}}/m_{\rm d}$, see
Paper I) for the Epstein regime is given by
$\displaystyle\mathbf{F}_{\rm{drag}}=\displaystyle-2\pi
s^{2}\rho_{\mathrm{g}}\Delta v^{2}$
$\displaystyle\left[\frac{1}{2\sqrt{\pi}}\left\\{\left(\frac{1}{\psi}+\frac{1}{2\psi^{3}}\right)e^{-\psi^{2}}+\right.\right.$
(11)
$\displaystyle\left.\left.\left(1+\frac{1}{\psi^{2}}-\frac{1}{4\psi^{4}}\right)\sqrt{\pi}\,\mathrm{erf}\left(\psi\right)\right\\}\right]\textbf{x},$
where $s$ corresponds to the grain radius and $m$, $\rho_{\mathrm{g}}$, $T$
denote the mass of the gas molecules, the intrinsic gas density and the local
temperature of the mixture (the gas and the dust are supposed to have the same
temperature). The thermal sound speed of the gas is thus
$c_{\mathrm{s}}=\sqrt{\gamma k_{\mathrm{B}}T/m}$ and the mean thermal velocity
of the gas, $c_{\mathrm{s}}\sqrt{8/\pi\gamma}$. The dimensionless quantity
$\psi$ is defined according to
$\psi\equiv\displaystyle\sqrt{\frac{\gamma}{2}}\frac{\Delta
v}{c_{\mathrm{s}}},$ (12)
where $\Delta\textbf{v}=\textbf{v}_{\mathrm{d}}-\textbf{v}_{\mathrm{g}}=\Delta
v$ x is the differential velocity (x being a unit vector). However, depending
on the characteristics of the problem (i.e. low or high Mach numbers, or
both), simpler and computationally less expensive approximations may be used.
For $\psi\ll 1$, i.e. low Mach numbers, Eq. 11 can be expanded to third order
in $\psi$, giving
$\mathbf{F}_{\rm{drag}}=\displaystyle-\frac{4\pi}{3}\rho_{\mathrm{g}}s^{2}\sqrt{\frac{8}{\pi\gamma}}c_{\mathrm{s}}\Delta
v\left[1+\frac{\psi^{2}}{5}+\mathcal{O}\left(\psi^{4}\right)\right]\textbf{x},$
(13)
which is usually simplified to its linear term,
$\mathbf{F}_{\rm{drag}}=-\frac{4\pi}{3}\rho_{\mathrm{g}}s^{2}\sqrt{\frac{8}{\pi\gamma}}c_{\mathrm{s}}\Delta\mathbf{v}.$
(14)
For $\psi\gg 1$, i.e. high Mach numbers, the Taylor expansion in $1/\psi$ of
Eq. 11 gives
$\mathbf{F}_{\rm{drag}}=\displaystyle-\left[\pi\rho_{\mathrm{g}}s^{2}\Delta
v^{2}\left(1+\frac{1}{\psi^{2}}-\frac{1}{4\psi^{4}}\right)+\mathcal{O}\left(e^{-\psi^{2}}\right)\right]\textbf{x},$
(15)
which is usually reduced to its quadratic term,
$\mathbf{F}_{\rm{drag}}=-\pi\rho_{\mathrm{g}}s^{2}\Delta v\Delta\mathbf{v}.$
(16)
A convenient way to handle Epstein drag at both low and high Mach numbers is
to use an interpolation between the two asymptotic regimes given by Eqs. 14
and 16 as derived in Kwok (1975) (cf. Paardekooper & Mellema 2006), giving
$\mathbf{F}_{\rm{drag}}=-\frac{4\pi}{3}\rho_{\mathrm{g}}s^{2}\sqrt{\frac{8}{\pi\gamma}}c_{\mathrm{s}}\sqrt{1+\frac{9\pi}{128}\frac{\Delta
v^{2}}{c_{\rm{s}}^{2}}}\Delta\mathbf{v}.$ (17)
The deviation of Eq. 17 from the full expression (Eq. 11) is $\lesssim 1\%$
(Kwok, 1975). Thus, in general, we adopt Eq. 17 for the Epstein regime. We
compare the differences between the various Epstein expressions in Sec. 5.
#### 2.2.2 Stokes regime for dense media
In a dense medium ($\lambda_{\mathrm{g}}>{4s}/{9}$), grains should be treated
with the Stokes drag regime, for which the expression of the drag force
$\bf{F}_{\rm{drag}}$ is:
$\mathbf{F}_{\rm{drag}}=-\frac{1}{2}C_{\mathrm{D}}\pi
s^{2}\rho_{\mathrm{g}}\Delta v\Delta\textbf{v},$ (18)
where the coefficient $C_{\mathrm{D}}$ is a piecewise function of the local
Reynolds number:
$C_{\mathrm{D}}=\begin{cases}24R_{\mathrm{d}}^{-1},&R_{\mathrm{d}}<1;\\\
24R_{\mathrm{d}}^{-0.6},&1<R_{\mathrm{d}}<800;\\\
0.44,&800<R_{\mathrm{d}},\end{cases}$ (19)
where $R_{\mathrm{d}}$ is defined in Eq. 10. Equation 19 indicates that at
small Reynolds numbers ($R_{\mathrm{d}}<1$), the drag force is linear with
respect to the local differential velocity between the grain and the gas. The
relation transitions to a power-law regime
($\mathbf{F}_{\rm{drag}}\propto\Delta v^{0.4}\Delta\textbf{v}$) at
intermediate Reynolds numbers ($1<R_{\mathrm{d}}<800$) and becomes quadratic
at large Reynolds numbers ($R_{\mathrm{d}}>800$). When the local concentration
of dust grains becomes very large (i.e., average distance between the
particles comparable to the grain size), the coefficient $C_{\mathrm{D}}$
should also depend on the local concentration of particles. However, this
extreme situation is not encountered in astrophysical situations.
Assuming gas molecules interact as hard spheres, the dynamic viscosity of the
gas can be computed according to (Chapman & Cowling, 1970):
$\mu=\frac{5m}{64\sigma_{\mathrm{s}}}\sqrt{\frac{\pi}{\gamma}}c_{\mathrm{s}},$
(20)
where $m=2m_{\mathrm{H}}$ and $\sigma_{\mathrm{s}}$ is the geometric cross
section of the molecule ($\sigma_{\mathrm{s}}=2.367\times 10^{-15}$ cm2 for
H2). The gas mean free path $\lambda_{\mathrm{g}}$ and the kinematic viscosity
$\nu$ of the gas are deduced from $\mu$ using
$\lambda_{\mathrm{g}}=\displaystyle\sqrt{\frac{\pi\gamma}{2}}\frac{\mu}{\rho_{\mathrm{g}}c_{\mathrm{s}}},$
(21)
and
$\nu=\displaystyle\frac{\mu}{\rho_{\mathrm{g}}}.$ (22)
## 3 Asytrophysical dust and gas mixtures in SPH
### 3.1 SPH evolution equations
The SPH version of the continuity equations Eqs. 1 – 2 are given by the
density summations for both the gas and the dust phase, computed according to:
$\displaystyle\hat{\rho}_{a}=\sum_{b}m_{b}W_{ab}(h_{a});$ $\displaystyle
h_{a}=\eta\left(\frac{m_{a}}{\hat{\rho}_{a}}\right)^{1/\nu},$ (23)
$\displaystyle\hat{\rho}_{i}=\sum_{j}m_{j}W_{ij}(h_{i});$ $\displaystyle
h_{i}=\eta\left(\frac{m_{j}}{\hat{\rho}_{i}}\right)^{1/\nu},$ (24)
where as in Paper I, the indices $a,b,c$ refer to quantities computed on gas
particles and $i,j,k$ refer to quantities computed on dust particles. The
volume filling fraction $\theta$, is defined on a _gas_ particle, $a$,
according to
$\theta_{a}=1-\frac{\hat{\rho}_{{\rm d},a}}{\rho_{\rm d}},$ (25)
where $\hat{\rho}_{{\rm d},a}$ is the density of _dust_ at the _gas_ particle
location, calculated using
$\hat{\rho}_{{\rm d},a}=\sum^{N_{neigh,dust}}_{j=1}m_{j}W_{aj}(h_{a}),$ (26)
where $h_{a}$ is the smoothing length of the _gas_ particle computed using gas
neighbours. The local density of dust at the gas location can thus be zero
(giving $\theta=1$) if no dust particles are found within the kernel radius
computed with the gas smoothing length. Importantly, as $\hat{\rho}$ and $h$
are mutually dependent, they have to be simultaneously calculated for each
type of particle, e.g. by the iterative procedure described in Price &
Monaghan (2007).
The SPH equations of motion for the gas and the dust particles, corresponding
to the SPH translation of Eqs. 3 and 4, are given by
$\displaystyle\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}=$
$\displaystyle-\sum_{b}m_{b}\left[\frac{P_{a}\tilde{\theta}_{a}}{\Omega_{a}\hat{\rho}_{a}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\
b}\left(h_{a}\right)+\frac{P_{b}\tilde{\theta}_{b}}{\Omega_{b}\hat{\rho}_{b}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\ b}\left(h_{b}\right)\right]$
$\displaystyle-\sum_{j}m_{j}\frac{P_{a}\left(1-\theta_{a}\right)}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\nabla_{\\!\\!\ a}W_{\\!\\!\ a\\!\\!\ j}\left(h_{a}\right)$
$\displaystyle+\nu\sum_{j}m_{j}\frac{K_{aj}}{\hat{\rho}_{a}\hat{\rho}_{j}}\left({\bf
v}_{aj}\cdot\hat{\textbf{r}}_{aj}\right)\hat{\textbf{r}}_{aj}D_{aj}(h_{a}),$
(27)
for an SPH gas particle and
$\displaystyle\frac{\mathrm{d}{\bf v}_{i}}{\mathrm{d}t}=$
$\displaystyle\sum_{b}m_{b}\frac{P_{b}\left(1-\theta_{b}\right)}{\hat{\rho}_{b}\hat{\rho}_{{\mathrm{d},b}}}\nabla_{\\!\\!\
i}W_{\\!\\!\ b\\!\\!\ i}\left(h_{b}\right)$ (28)
$\displaystyle-\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left({\bf
v}_{bi}\cdot\hat{\textbf{r}}_{bi}\right)\hat{\textbf{r}}_{bi}D_{ib}(h_{i}),$
for an SPH dust particle. $\Omega$ is the usual variable smoothing length term
$\Omega_{b}\equiv 1-\frac{\partial
h_{b}}{\partial\hat{\rho}_{b}}\sum_{c}m_{c}\frac{\partial W_{\\!\\!\ b\\!\\!\
c}\left(h_{b}\right)}{\partial h_{b}}.$ (29)
It should be noted that $\Omega_{\rm{d}}$ is computed only using dust particle
neighbours according to:
$\Omega_{\mathrm{d},b}=1-\frac{\partial
h_{b}}{\partial\hat{\rho}_{{\mathrm{d},b}}}\sum_{j}m_{j}\frac{\partial
W_{\\!\\!\ b\\!\\!\ j}\left(h_{b}\right)}{\partial h_{b}}.$ (30)
$\tilde{\theta}$ is defined according to
$\tilde{\theta}\equiv\theta+\frac{\hat{\rho}_{\rm
g}}{\hat{\rho}_{\mathrm{d}}}(1-\theta)(1-\Omega_{\rm d}).$ (31)
At this stage, no assumptions are made with respect to the functional form of
the drag coefficient $K$. The evolution of the internal energy for an SPH gas
particle is given by
$\displaystyle\frac{{\rm d}u_{a}}{{\rm d}t}$
$\displaystyle=\frac{\tilde{\theta}_{a}P_{a}}{\Omega_{a}\hat{\rho}_{a}^{2}}\sum_{b}m_{b}\left(\textbf{v}_{a}-\textbf{v}_{b}\right)\cdot\nabla_{a}W_{ab}(h_{a})$
(32) $\displaystyle+\frac{(1-\theta_{a})P_{a}}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\sum^{N_{neigh,dust}}_{j=1}m_{j}\left({\bf v}_{a}-{\bf
v}_{j}\right)\cdot\nabla_{a}W_{aj}(h_{a})$
$\displaystyle+\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\textbf{v}_{ak}\cdot\hat{\textbf{r}}_{ak}\right)^{2}D_{ak}(h_{a}).$
In Paper I, we showed that the total linear and angular momentum as well as
the total energy are exactly conserved. Thermal coupling terms have been
neglected in this paper.
### 3.2 Kernel functions
Two different kernels are employed to perform the SPH interpolations. First, a
standard bell-shaped kernel $W$:
$W\left(r,h\right)=\frac{\sigma}{h^{\nu}}f\left(q\right),$ (33)
where $h$ denotes the smoothing lengths of each phases, $\nu$ the number of
spatial dimensions and $q\equiv|{\bf r}-{\bf r}^{\prime}|/h$ is the
dimensionless variable used to calculate the densities and the buoyancy terms.
The function $f$ is usually the $M_{4}$ cubic spline kernel (Monaghan, 2005).
The drag interpolation is performed using a second kernel $D$. As shown in
Paper I, double-hump shaped kernels given by
$D\left(r,h\right)=\frac{\tilde{\sigma}}{h^{\nu}}q^{2}f(q),$ (34)
significantly improve the accuracy of the drag interpolation — for the same
computational cost — compared to bell-shaped kernels. The normalisation
constants $\tilde{\sigma}$ for various double hump kernels are given in Paper
I. We adopt the double hump cubic for the drag terms in this paper.
### 3.3 Astrophysical drag regimes in SPH
#### 3.3.1 Gas viscosity and mean free path
The drag coefficients $K_{ak}$ involved in Eqs. 3.1 – 28 and 32 are computed
independently for each pair of any gas particle $a$ and dust particle $k$. We
first use the sound speed $c_{\rm{s},a}$ to estimate the viscosity $\mu_{a}$
on the gas particle $a$ using (see Eq. 20)
$\mu_{a}=\frac{5m}{64\sigma_{\mathrm{s}}}\sqrt{\frac{\pi}{\gamma}}c_{\rm{s},a}.$
(35)
The mean free path is then computed according to Eq. 21, giving
$\lambda_{\mathrm{g},a}=\displaystyle\sqrt{\frac{\pi\gamma}{2}}\frac{\mu_{\rm{a}}}{\hat{\rho}_{a}c_{\rm{s},a}}.$
(36)
Finally, $\lambda_{\mathrm{g},a}$ is compared to the quantity $4s_{k}/9$ —
$s_{k}$ being the grain size of the dust particle — to determine whether the
drag coefficient of the SPH pair $K_{ak}$ is calculated using the Epstein or
the Stokes drag regimes.
#### 3.3.2 Epstein regime
If $4s_{k}/9\leq\lambda_{\mathrm{g},a}$, the drag coefficient $K_{ak}$ is
calculated using the Epstein prescription. Introducing the SPH quantity
$\psi_{ak}$ calculated on a pair of gas and dust SPH particles and defined by
$\psi_{ak}\equiv\sqrt{\frac{\gamma}{2}}\frac{|{\bf
v}_{ak}|}{c_{\mathrm{s},a}},$ (37)
Eq. 11 can be straightforwardly translated to get the drag coefficient
$K_{ak}$ involved in the SPH drag force
$\displaystyle
K_{ak}=\displaystyle-\sqrt{\pi}s^{2}\rho_{\mathrm{g}}\frac{\hat{\rho}_{\mathrm{d}}}{m_{\rm{d}}}|\textbf{v}_{ak}|$
$\displaystyle\left[\frac{1}{2\sqrt{\pi}}\left\\{\left(\frac{1}{\psi_{ak}}+\frac{1}{2\psi_{ak}^{3}}\right)e^{-\psi_{ak}^{2}}+\right.\right.$
(38)
$\displaystyle\left.\left.\left(1+\frac{1}{\psi_{ak}^{2}}-\frac{1}{4\psi_{ak}^{4}}\right)\sqrt{\pi}\,\mathrm{erf}\left(\psi_{ak}\right)\right\\}\right]\textbf{x},$
where $s$ is the grain radius, $m_{\rm{d}}$ is the grain mass and $\gamma$ is
the adiabatic index. Eq. 38 is computationally expensive as it involves
exponential and error functions. The SPH equivalent of Eq.17 is given by
$K_{ak}=\frac{4}{3}\pi\sqrt{\frac{8}{\pi\gamma}}\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\frac{\hat{\rho}_{a}}{\theta_{a}}s^{2}c_{\mathrm{s},a}\sqrt{1+\frac{9\pi}{128}\frac{v_{ak}^{2}}{c_{\mathrm{s},a}^{2}}}.$
(39)
Both Eqs. 38 and 39 reduce to the linear Epstein regime at low Mach numbers
(equivalent of Eq. 14) for which the coefficient $K_{ak}$ is
$K_{ak}=\frac{4}{3}\pi\sqrt{\frac{8}{\pi\gamma}}\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\frac{\hat{\rho}_{a}}{\theta_{a}}s^{2}c_{\mathrm{s},a},$
(40)
and to the quadratic drag regime at high Mach numbers (equivalent of Eq. 16),
for which the coefficient $K_{ak}$ is
$K_{ak}=\pi\rho_{\mathrm{g}}s^{2}\frac{\hat{\rho}_{\mathrm{d}}}{m_{\rm{d}}}|{\bf
v}_{ak}|.$ (41)
#### 3.3.3 Stokes regime
If $4s_{k}/9>\lambda_{\mathrm{g},a}$, the drag coefficient $K_{ak}$ is
calculated using the Stokes prescription (see Eqs. 18–19). The local Reynolds
number $R_{\mathrm{d},ak}$ is computed for each pair of gas and dust particles
using
$R_{\mathrm{d},ak}\equiv\frac{2s\hat{\rho}_{a}\left|\textbf{v}_{ak}\right|}{\mu_{a}\theta_{a}},$
(42)
such that the drag coefficient $K_{ak}$ can be computed according to
$K_{ak}=\begin{cases}\displaystyle
6\pi\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\mu_{a}s&R_{\mathrm{d},ak}<1,\\\
\displaystyle\frac{12\pi}{2^{0.6}}\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\frac{\mu_{a}^{0.6}}{\theta_{a}^{0.4}\hat{\rho}_{a}^{0.6}}s^{1.4}\left|\textbf{v}_{ak}\right|^{0.4}&1<R_{\mathrm{d},ak}<800,\\\
\displaystyle
0.22\pi\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\frac{\hat{\rho}_{a}}{\theta_{a}}s^{2}\left|\textbf{v}_{ak}\right|&R_{\mathrm{d},ak}>800.\\\
\end{cases}$ (43)
These expressions have been used by Ayliffe et al. (2011) to compute the drag
on planetesimals in a protoplanetary disc.
## 4 Timesteping
### 4.1 Explicit timesteping
The simplest method to evolve the evolution equations for the SPH particles is
to use an explicit integrator (e.g. the standard Leapfrog). The stability of
the system is guaranteed provided the timestep remains smaller than a critical
value $\Delta t_{\rm{c}}$. In Paper I, we performed a Von Neumann analysis of
the continuous equations, deriving the explicit timestepping criterion
$\Delta
t_{\rm{c},a}=\min_{k}\left[\frac{\hat{\rho}_{a}\hat{\rho}_{k}}{K_{ak}(\hat{\rho}_{a}+\hat{\rho}_{k})}\right];\hskip
14.22636pt\Delta
t_{\rm{c},i}=\min_{b}\left[\frac{\hat{\rho}_{b}\hat{\rho}_{i}}{K_{bi}(\hat{\rho}_{b}+\hat{\rho}_{i})}\right];$
(44)
for gas and dust particles, respectively, with the minimum being taken over
all the particle’s neighbours. Although this criterion was derived in Paper I
for linear drag regimes only, it remains valid even for non-linear drag
regimes where the drag coefficients depend on the differential velocity
between the particles, i.e. $K_{ak}=K_{ak}\left(|\mathbf{v}_{ak}|\right)$.
### 4.2 Implicit timestepping
When the drag timescale becomes smaller than other time scales in the system
(e.g. the Courant condition or the orbital timescale), the timestep
restriction of the explicit methods may become prohibitive and implicit
methods are required. Monaghan (1997) considered the application of two
implicit schemes (the first-order Backward-Euler and second-order Tischer
scheme) to SPH dust-gas mixtures. Both schemes are unconditionally stable, but
a higher accuracy is achieved with second-order schemes.
#### 4.2.1 Backward-Euler method
The Backward-Euler scheme applied to the drag interaction between SPH dust and
gas particles is given by
$\displaystyle\frac{\mathbf{v}^{n+1}_{a}-\mathbf{v}^{n}_{a}}{\Delta t}$
$\displaystyle=$
$\displaystyle-\nu\sum_{k}m_{k}\frac{K_{ak}^{n+1}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\mathbf{v}^{n+1}_{ak}\cdot\hat{\mathbf{r}}_{ak}\right)\hat{\mathbf{r}}_{ak}D_{ak},$
(45) $\displaystyle\frac{\mathbf{v}^{n+1}_{i}-\mathbf{v}^{n}_{i}}{\Delta t}$
$\displaystyle=$
$\displaystyle+\nu\sum_{b}m_{b}\frac{K_{bi}^{n+1}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left(\mathbf{v}^{n+1}_{bi}\cdot\hat{\mathbf{r}}_{bi}\right)\hat{\mathbf{r}}_{bi}D_{bi}.$
(46)
Although the scheme is unconditionally stable, the implicit equation with
respect to the velocities $\textbf{v}^{n+1}$ must be solved at each time step.
Direct numerical inversion of this linear system would be prohibitive given
the typical number of neighbour interactions for each SPH particle. Thus,
approximate or iterative solutions to Eqs. 45 – 46 are required.
#### 4.2.2 Monaghan (1997) scheme
Monaghan (1997) suggested approximating the velocities $\textbf{v}^{n+1}$ of
Eqs. 45 – 46 using a pairwise treatment in order to preserve the exact
conservation of linear and angular momentum in the SPH formalism. Considering
the interaction between the SPH gas particle $a$ the dust particle $i$,
Monaghan (1997) introduced pairwise auxiliary velocities $\tilde{\bf{v}}$
defined by:
$\displaystyle\tilde{\textbf{v}}_{a}$ $\displaystyle=$ $\displaystyle{\bf
v}_{a}^{n}-m_{i}~{}\Delta t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left(\tilde{\textbf{v}}_{ai}\cdot\hat{\bf{r}}_{ai}\right)\hat{\bf{r}}_{ai},$
(47) $\displaystyle\tilde{\textbf{v}}_{i}$ $\displaystyle=$ $\displaystyle{\bf
v}_{i}^{n}+m_{a}~{}\Delta t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left(\tilde{\textbf{v}}_{ai}\cdot\hat{\bf{r}}_{ai}\right)\hat{\bf{r}}_{ai},$
(48)
Eqs. 47 and 48 are solved, for a given pair of particles, by taking the scalar
product by $\hat{\bf{r}}_{ai}$ of the difference of the two equations, giving
$\tilde{\textbf{v}}_{ai}\cdot\hat{\bf{r}}_{ai}=\frac{\textbf{v}_{ai}^{n}\cdot\hat{\bf{r}}_{ai}}{1+\Delta
t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left(m_{a}+m_{i}\right)}.$ (49)
Substituting this expression into Eq. 47 and 48 gives expressions for
$\tilde{\textbf{v}}_{a}$ and $\tilde{\textbf{v}}_{i}$. Iterating this pairwise
process by looping over all the SPH particles provides an approximate solution
for the velocities $\bf{v}^{n+1}$, i.e.
$\displaystyle\frac{\mathbf{v}^{n+1}_{a}-\mathbf{v}^{n}_{a}}{\Delta t}$
$\displaystyle\simeq$
$\displaystyle-\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\tilde{\mathbf{v}}_{ak}\cdot\hat{\mathbf{r}}_{ak}\right)\hat{\mathbf{r}}_{ak}D_{ak},$
(50) $\displaystyle\frac{\mathbf{v}^{n+1}_{i}-\mathbf{v}^{n}_{i}}{\Delta t}$
$\displaystyle\simeq$
$\displaystyle+\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left(\tilde{\mathbf{v}}_{bi}\cdot\hat{\mathbf{r}}_{bi}\right)\hat{\mathbf{r}}_{bi}D_{bi}.$
(51)
The main drawback of this method is that the approximation given by Eqs. 50
and 51 is inexact – that is, it provides only an approximate solution to Eqs.
45 and 46. Furthermore the accuracy of the approximation is not known _a
priori_ and there is no possibility of performing repeated sweeps in order to
converge to a more accurate solution. In practice, we find that the velocities
obtained by this scheme (for example on the dustybox test) can be
significantly in error, with no possibility of improving the convergence (for
example, by doing several iterations/sweeps).
#### 4.2.3 Alternative pairwise treatment for linear drag regimes
We propose a more consistent method for solving Eqs. 45–46 on a given gas or
dust particle ($a$ and $i$, respectively) by sweeping over all particle pairs
and updating the velocities iteratively according to
$\displaystyle{\textbf{v}}_{a}^{**}$ $\displaystyle=$ $\displaystyle{\bf
v}_{a}^{n}+\Delta t\textbf{F}_{a,\mathrm{drag}}^{*}$ (52) $\displaystyle-$
$\displaystyle m_{i}~{}\Delta t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left[\left({\textbf{v}}_{ai}^{**}-\textbf{v}_{ai}^{*}\right)\cdot\hat{\bf{r}}_{ai}\right]\hat{\bf{r}}_{ai},$
$\displaystyle{\textbf{v}}_{i}^{**}$ $\displaystyle=$ $\displaystyle{\bf
v}_{i}^{n}+\Delta t\textbf{F}_{i,\mathrm{drag}}^{*}$ (53) $\displaystyle+$
$\displaystyle m_{a}~{}\Delta t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left[\left({\textbf{v}}_{ai}^{**}-\textbf{v}_{ai}^{*}\right)\cdot\hat{\bf{r}}_{ai}\right]\hat{\bf{r}}_{ai},$
where ${\bf v}^{**}$ refers to the improved approximation to ${\bf v}^{n+1}$
obtained after updating each pair and ${\bf v}^{*}$ to the previous iteration
value of ${\bf v}^{**}$. Eqs. 52 and 53 are solved for each pair of particles
by taking the dot product of $\hat{\bf{r}}_{ai}$ with the difference of the
two equations, giving
$\left({\textbf{v}}_{ai}\cdot\hat{\bf{r}}_{ai}\right)^{**}=\frac{\left(\textbf{v}_{ai}^{n}+\Delta
t\textbf{F}_{ai,\mathrm{drag}}^{*}+\Delta t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left(m_{a}+m_{i}\right)\textbf{v}_{ai}^{*}\right)\cdot\hat{\bf{r}}_{ai}}{1+\Delta
t\frac{\nu
K_{ai}D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}\left(m_{a}+m_{i}\right)}.$ (54)
Substituting Eq. 54 in Eqs. 52 and 53 gives the updated velocities for the
pair. Note that during the global sweep over particle pairs ${\bf v}^{*}$
begins as ${\bf v}^{n}$ at the first iteration but is updated as soon as new
values become available.
This pairwise correction ensures that i) both the linear and the angular
momentum are exactly conserved and ii) the velocities converge to the correct
solution of the implicit scheme given by Eqs. 45 and 46, since the last term
of Eqs. 52 and 53 tends to zero as the number of iterations increases. We thus
refine our approximation to the solution by performing as many successive
iterations as are required to reach a suitable convergence criterion.
#### 4.2.4 Convergence criterion
We consider that the approximation we obtain from the implicit scheme
described above is accurate enough when
$\frac{|\bf{v}^{k+1}-\bf{v}^{k}|}{\min c_{\rm{s}}}<\varepsilon,$ (55)
is satisfied for each particle. Typically, we adopt $\varepsilon=10^{-4}$,
which ensures that the approximation we make on the time stepping is
negligible compared to the $\mathcal{O}(h^{2})$ truncation error of the
underlying SPH scheme.
#### 4.2.5 Implementation into Leapfrog
The Leapfrog scheme is well suited to the evolution of particle methods
because, for position-dependant forces, it preserves geometric properties of
particle orbits and requires only one evaluation per timestep to give second
order accuracy. In the standard formulation, the evolution is computed
according to
$\begin{array}[]{lrrcll}{\rm Kick}&\left[\right.&{\bf v}^{1/2}&=&{\bf
v}^{0}+\frac{\Delta t}{2}{\bf f}^{0}\left({\bf x}^{0},{\bf
v}^{0}\right),&\left.\right]\\\\[5.0pt] {\rm Drift}&\left[\right.&{\bf
x}^{1}&=&{\bf x}^{0}+\Delta t{\bf v}^{1/2},&\left.\right]\\\\[5.0pt] {\rm
Kick}&\left[\right.&{\bf v}^{1}&=&{\bf v}^{1/2}+\frac{\Delta t}{2}{\bf
f}^{1}\left({\bf x}^{1},{\bf v}^{1}\right),&\left.\right]\end{array}$ (56)
corresponding to Kick, Drift and Kick steps respectively. Adapting Leapfrog to
deal with velocity dependent forces (e.g. drag) is a priori more difficult
since for velocity-dependent forces, the last Kick is implicit in ${\bf
v}^{1}$. For our present purposes, this does not present a major problem since
the drag is already computed implicitly. Splitting the forces into position-
dependent (${\bf f}_{\rm SPH}$) and drag (${\bf f}_{\rm drag}$) contributions,
the scheme becomes
$\begin{array}[]{lrrcll}{\rm Kick}&\left[\right.&{\bf\tilde{v}}^{1/2}&=&{\bf
v}^{0}+\frac{\Delta t}{2}{\bf f}^{0}_{\rm SPH}\left({\bf
x}^{0}\right),&\left.\right]\\\\[5.0pt] {\rm Drift}&\left[\right.&{\bf
x}^{1/2}&=&{\bf x}^{0}+\frac{\Delta t}{2}\tilde{{\bf
v}}^{1/2},&\left.\right]\\\\[5.0pt] {\rm Drag}&\left[\right.&{\bf
v}^{1/2}&=&{\bf\tilde{v}}^{1/2}+\frac{\Delta t}{2}{\bf f}^{1/2}_{\rm
drag}\left({\bf x}^{1/2},{\bf v}^{1/2}\right)&\left.\right]\\\\[5.0pt] {\rm
Drift}&\left[\right.&{\bf x}^{1}&=&{\bf x}^{0}+\Delta t{\bf
v}^{1/2},&\left.\right]\\\\[5.0pt] {\rm
Kick}&\left[\right.&{\bf\tilde{v}}^{1}&=&{\bf v}^{1/2}+\frac{\Delta t}{2}{\bf
f}^{1}_{\rm SPH}\left({\bf x}^{1}\right),&\left.\right]\\\\[5.0pt] {\rm
Drag}&\left[\right.&{\bf v}^{1}&=&{\bf\tilde{v}}^{1}+\frac{\Delta t}{2}{\bf
f}^{1}_{\rm drag}\left({\bf x}^{1},{\bf
v}^{1}\right)&\left.\right]\end{array}$ (57)
where the Drag steps represent the implicit updates computed as described in
Sec. 4.2.3. The disadvantage of Eq. 57 is that two drag force evaluations are
required, removing one of the advantages of the Leapfrog integrator.
Inspection of 57 reveals that an alternative version that requires only one
Drag step can be constructed according to
$\begin{array}[]{lrlr}{\rm Kick}&\bigl{[}&\begin{array}[]{rcl}{\bf
v}^{1/2}&=&{\bf v}^{0}+\frac{\Delta
t_{0}}{2}{\bf\tilde{f}},\end{array}&\bigr{]}\\\\[15.00002pt] {\rm
Drift}&\bigl{[}&\begin{array}[]{rcl}{\bf x}^{1}&=&{\bf x}^{0}+\Delta t_{0}{\bf
v}^{1/2},\end{array}&\bigr{]}\\\\[10.00002pt] {\rm
Drag}&&\left\\{\begin{array}[]{rcl}{\tilde{\bf v}}^{3/2}&=&{\bf
v}^{1/2}+\frac{\Delta t_{0}+\Delta t_{1}}{2}{\bf f}^{1}_{\rm SPH}\left({\bf
x}^{1}\right),\\\\[5.0pt] {\bf v}^{3/2}&=&{\tilde{\bf v}}^{3/2}+\frac{\Delta
t_{0}+\Delta t_{1}}{2}{\bf f}^{1}_{\rm drag}\left({\bf x}^{1},{\bf
v}^{3/2}\right),\\\\[5.0pt] {\tilde{\bf f}}&=&2\left({\tilde{\bf
v}}^{3/2}-{\tilde{\bf v}}^{1/2}\right)/\left(\Delta t_{0}+\Delta
t_{1}\right),\end{array}\right.&\\\\[30.00005pt] {\rm
Kick}&\bigl{[}&\begin{array}[]{rcl}{\bf v}^{1}&=&{\bf v}^{1/2}+\frac{\Delta
t_{0}}{2}{\bf\tilde{f}}.\end{array}&\bigr{]}\end{array}$ (58)
where we have combined the drag steps by predicting the velocity ${\bf
v}^{3/2}$. Note that strictly, the Drag step in this method is semi-implicit
since the force is evaluated using ${\bf x}^{1}$ rather than ${\bf x}^{3/2}$.
However, we expect this approximation to be reasonable as at high drag (for
which the implicit method is designed), the drag mainly changes the
differential velocity between the fluids and has less of an effect on the
positions. Care is also required when the timestep changes between the steps.
We have indicated the correct procedure by specifying $\Delta t_{0}$ and
$\Delta t_{1}$ where $\Delta t_{1}$ is the timestep computed based on ${\bf
x}^{1}$. Finally, Eq. 58 requires that ${\tilde{\bf f}}$ is known at the
beginning of the integration. This can be easily achieved by performing the
Drag step in Eq. 58 with ${\bf v}^{1/2}={\bf v}^{0}$, $\Delta t_{0}=0$ and
$\Delta t_{1}$ equal to the timestep calculated using the initial particle
positions.
#### 4.2.6 Generalisation to non-linear drag regimes
To extend this alternative pairwise treatment to any non-linear drag regime,
two additional points have to be considered. Firstly, although in principle
six quantities ($v_{x,y,z}$ for each particle) have to be determined for each
pair, this can be reduced to a single unknown quantity since the drag
coefficient depends only on the modulus of the differential velocity and the
exchange of momentum is directed along the line of sight joining the
particles. The system of equations for a single pair thus reduces to
$\displaystyle{\textbf{v}}_{a}^{**}$ $\displaystyle=$ $\displaystyle{\bf
v}_{a}^{n}+\Delta t\textbf{F}_{a,\mathrm{drag}}^{*}$ (59) $\displaystyle-$
$\displaystyle m_{i}~{}\Delta t\frac{\nu
D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}K_{ai}\sqrt{\left[\left({\textbf{v}}_{ai}^{**}-\textbf{v}_{ai}^{*}\right)\cdot\hat{\bf{r}}_{ai}\right]^{2}+V^{2,n}_{\rm
orth}},$ $\displaystyle{\textbf{v}}_{i}^{**}$ $\displaystyle=$
$\displaystyle{\bf v}_{i}^{n}+\Delta t\textbf{F}_{i,\mathrm{drag}}^{*}$ (60)
$\displaystyle+$ $\displaystyle m_{a}~{}\Delta t\frac{\nu
D_{ai}}{\hat{\rho}_{a}\hat{\rho}_{i}}K_{ai}\sqrt{\left[\left({\textbf{v}}_{ai}^{**}-\textbf{v}_{ai}^{*}\right)\cdot\hat{\bf{r}}_{ai}\right]^{2}+V^{2,n}_{\rm
orth}},$
where
$V^{2,n}_{\rm
orth}={\textbf{v}}_{ai}^{n}\cdot{\textbf{v}}_{ai}^{n}-\left({\textbf{v}}_{ai}^{n}\cdot\hat{\bf{r}}_{ai}\right)^{2}.$
(61)
Secondly, taking the dot product of $\hat{\bf{r}}_{ai}$ with the difference of
the two equations Eqns. 59 and 60 does not lead in general to an equation
which can be solved analytically. The values of
${\textbf{v}}_{ai}\cdot\hat{\bf{r}}_{ai}$ must therefore be determined using a
numerical rootfinding procedure (we use a Newton-Raphson scheme) before being
substituted in Eqns. 59 and 60 to determine the velocities for both the gas
and the dust particles.
#### 4.2.7 Performance of the implicit scheme
The computational cost of a timestep with the implicit pairwise treatment is
more expensive than an explicit timestep since at least two iterations have to
be performed to ensure that the scheme is converged. However, the implicit
pairwise treatment will be more efficient provided that the number of
iterations is much smaller than the number of explicit timesteps that would
otherwise be required.
We find in practice that the efficiency of the pairwise treatment is mainly
determined by the number of iterations required to satisfy Eq. 55 (this aspect
was not addressed in the Monaghan (1997) scheme where only one iteration is
ever taken in the hope that the approximation is sufficiently accurate). The
rapidity of the convergence depends primarily on the ratio $r=\Delta t/t_{\rm
s}$ of the timestep over the drag stopping time (defined in Eq. (96) of Paper
I) and on the value of $\varepsilon$. Empirically, we have found that, for
$\varepsilon=10^{-4}$ and $1\lesssim r\lesssim 10$, the implicit pairwise
treatment converges efficiently, the ratio $|\bf{v}^{k+1}-\bf{v}^{k}|/(\min
c_{\rm{s}})$ decreasing by $\sim$ two orders of magnitude at each iterations.
Thus, the implicit pairwise treatment improves the computational time by a
factor of $\sim 1$–$10$. However, this rapidity of convergence decreases as
$r$ increases. At very high drag ($r\gtrsim 1000$), we find that the implicit
scheme becomes less efficient than explicit timestepping due to the large
number of iterations required. A similar behaviour has been found using the
Gauss-Seidel iterative scheme developed by Whitehouse et al. (2005) to treat
SPH radiative transfer in the flux-limited diffusion approximation (Bate 2011,
private communication), so this issue is not specific to the pairwise
treatment.
It is important to note that the computational gain obtained with the pairwise
scheme does not solve the resolution issue at high drag extensively discussed
in Paper I. Both of these problems suggest that a more efficient method for
handling high drag regimes is required. Such a method is beyond the scope of
the present paper.
#### 4.2.8 Higher order implicit schemes
Higher temporal accuracy may be achieved by using second instead of first
order implicit schemes. The gain in accuracy is obtained by dividing the drag
timestep $\Delta t$ into two half timesteps. Monaghan (1997) suggested the
‘Tischer’ scheme, where the two half timesteps are given by
$\displaystyle\frac{\mathbf{v}^{n+\frac{1}{2}}_{a}-\mathbf{v}^{n}_{a}}{\Delta
t/2}$ $\displaystyle=$
$\displaystyle-0.6\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\mathbf{v}^{n+\frac{1}{2}}_{ak}\cdot\hat{\mathbf{r}}_{ak}\right)\hat{\mathbf{r}}_{ak}D_{ak},$
(62)
$\displaystyle-0.4\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\mathbf{v}^{n}_{ak}\cdot\hat{\mathbf{r}}_{ak}\right)\hat{\mathbf{r}}_{ak}D_{ak},$
$\displaystyle\frac{\mathbf{v}^{n+\frac{1}{2}}_{i}-\mathbf{v}^{n}_{i}}{\Delta
t/2}$ $\displaystyle=$
$\displaystyle+0.6\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left(\mathbf{v}^{n+\frac{1}{2}}_{bi}\cdot\hat{\mathbf{r}}_{bi}\right)\hat{\mathbf{r}}_{bi}D_{bi},$
(63)
$\displaystyle+0.4\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left(\mathbf{v}^{n}_{bi}\cdot\hat{\mathbf{r}}_{bi}\right)\hat{\mathbf{r}}_{bi}D_{bi},$
and then
$\displaystyle\frac{\mathbf{v}^{n+1}_{a}-\left[1.4\mathbf{v}^{n+\frac{1}{2}}_{a}-0.4\mathbf{v}^{n}_{a}\right]}{\Delta
t/2}$ $\displaystyle=$
$\displaystyle-0.6\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\mathbf{v}^{n+1}_{ak}\cdot\hat{\mathbf{r}}_{ak}\right)\hat{\mathbf{r}}_{ak}D_{ak},$
$\displaystyle\frac{\mathbf{v}^{n+1}_{i}-\left[1.4\mathbf{v}^{n+\frac{1}{2}}_{i}-0.4\mathbf{v}^{n}_{i}\right]}{\Delta
t/2}$ $\displaystyle=$
$\displaystyle+0.6\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left(\mathbf{v}^{n+1}_{bi}\cdot\hat{\mathbf{r}}_{bi}\right)\hat{\mathbf{r}}_{bi}D_{bi}.$
The last terms of Eqs. 62–63 correspond to the explicit drag force involved in
the Forward-Euler scheme. These quantities are computed form the velocities at
the timestep $n$ at the beginning of the scheme, as described in Sec. 4.1.
Then, the successive determination of $\bf{v}^{n+\frac{1}{2}}$ and
$\bf{v}^{n}$ as given by Eqs.
LABEL:eq:Tischer2_scheme1–LABEL:eq:Tischer2_scheme2 consists of two Backward-
Euler steps with step size $\frac{\Delta t}{2}$. They are therefore computed
using our alternative pairwise scheme, until the iterations for each half time
step have converged. Eqs. 62 and 63 concerns the specific case of a linear
drag regime, but this scheme can easily be extended to non-linear drag regimes
as in Sec. 4.2.6.
## 5 Numerical tests
### 5.1 dustybox: Two fluid drag in a periodic box
The dustybox problem presented by Laibe & Price (2011a) and described in
detail in Paper I consists of two fluids in a periodic box moving with a
differential velocity ($\Delta v_{0}=v_{d,0}-v_{g,0}$). This is the only test
where analytic solutions are known for several functional forms corresponding
to non-linear drag regimes (see Laibe & Price 2011a). These represent the
functional forms of the Epstein and Stokes prescription. We thus use the
dustybox problem to benchmark the accuracy of our algorithm for non-linear
drag regimes using both explicit and implicit timestepping.
#### 5.1.1 dustybox: setup
We set up constant densities $\hat{\rho}_{\mathrm{g}}$ and
$\hat{\rho}_{\mathrm{d}}$ and gas pressure $P_{\mathrm{g}}$ in a 3D periodic
domain $x,y,z\in[0,1]$ filled by $20^{3}$ gas particles set up on a regular
cubic lattice and $20^{3}$ dust particles shifted by half of the lattice
spacing in each direction (as in Paper I, we verified that the results are
independent of the offset of the dust lattice). The gas sound speed, the gas
and the dust densities are set to unity in code units and no artificial
viscosity is applied. The intrinsic dust volume is neglected by assuming
$\theta=1$.
Simulations have been performed using both the explicit timestepping presented
in Paper I and the implicit pairwise timestepping described in Sec. 4. For the
latter, we verified that both the total linear and angular momentum are
exactly conserved as expected.
#### 5.1.2 dustybox: different drag regimes
Figure 1: Dust velocity (solid lines) as a function of time in the dustybox
test, using $2\times 20^{3}$ particles, a dust-to-gas ratio of unity and five
different linear and non-linear drag regimes — quadratic, power-law, linear,
third order expansion and mixed, from top to bottom — compared to the exact
solution for each case (long dashed/red lines). The initial velocities are set
to $v_{\mathrm{d},i}=1$, $v_{\mathrm{g},i}=0$ and the time integration is
performed using using the pairwise implicit treatment described in Sec. 4. The
accuracy ($\lesssim 0.1\%$) of the SPH treatment for dust-gas mixtures is
obtained by using the double-hump cubic kernel.
Fig. 1 shows the results of the dustybox test for the five different regimes
given in Table 1 of Laibe & Price (2011a): linear ($K=K_{0}$), quadratic
($K=K_{0}|\Delta v|$), power-law ($K=K_{0}|\Delta v|^{a}$, with $a=0.4$),
third order expansion ($K=K_{0}[1+a_{3}|\Delta v|^{2}]$, with $a_{3}=0.5$) and
mixed ($K=K_{0}\sqrt{1+a_{2}|\Delta v|^{2}}$, with $a_{2}=5$) where we have
used $K_{0}=1$ in each case. The analytic solutions are reproduced within an
accuracy comprised between $0.1\%$ and $1\%$ in every case — both linear and
non-linear. The implicit scheme was find to converge quickly for this problem,
requiring no more than two iterations at every stage of the evolution in each
case.
The efficiency of the damping for the dustybox problem decreases when the
exponent of the drag regime increases, since $||\Delta v||<1$. On the
contrary, additional non-linear terms give an additional contribution to the
drag for the mixed and third order drag regimes, leading to a differential
velocity that is more efficiently damped compared to the linear case.
### 5.2 dustywave: Sound waves in a dust-gas mixture
The exact solution for linear waves propagating in a dust-gas mixture
(dustywave) was derived in Laibe & Price (2011a) assuming a linear drag
regime. Unfortunately, exact solutions prove difficult to obtain for the case
of non-linear drag. Instead, we have verified simply that our simulations of
the dustywave problem for non-linear drag regimes are converged in both space
and time with an explicit timestepping scheme. We have then used these results
to benchmark our simulations using implicit timestepping.
We run the dustyywave problem for the same five non-linear drag regimes used
above. Strictly speaking, It should be noted that assuming $K_{0}$ is constant
(which we assume for this test problem) corresponds to Epstein and Stokes drag
only to first order for the dustywave problem.
#### 5.2.1 dustywave: Setup
The dustywave test is performed in a 1D periodic box, placing equally spaced
particles in the periodic domain $x\in[0,1]$ such that the gas and dust
densities are unity in code units. We do not apply any form of viscosity and
the gas sound speed is set to unity. To remain in the linear acoustic regime,
the relative amplitude of the perturbation of both velocity and density are
set to $10^{-4}$.
#### 5.2.2 dustywave: Different drag regimes
Figure 2: Solution of the 1D dustywave problem showing the SPH gas (solid
black lines) and dust (dashed black) velocities using three different drag
regimes: linear (top panel), power-law with an exponent of $0.4$ (center
panel) and quadratic (bottom panel). The damping is strongly reduced for non-
linear drag regimes compared to the linear case. Results are shown after 5
wave periods and the linear case (top panel) may be compared to the exact
analytic solution (red lines). The dust to gas ratio and $K_{0}$ are set to 1
in code units.
Fig. 2 shows the velocity profiles after 5 periods for three drag regimes —
linear, power law and quadratic — in the 1D dustywave problem using $K_{0}=1$
and a dust-to-gas ratio of unity. The solution obtained for the linear drag
regime shows an efficiently damped perturbation at $t=5$, consistent with a
stopping time of order unity. By comparison, the perturbation is only weakly
damped for the case of the power-law drag regime at the same time. The drag is
weaker still in the quadratic regime, for which both the dust and the gas are
mostly decoupled. Indeed, since the drag stopping time is a decreasing
function of the differential velocity for non-linear drag, and the
differential velocity is small with respect to the sound speed, the damping is
inefficient.
We have also performed dustywave simulations using the third order expansion
and mixed drag regimes. However, for these cases, non-linear terms represent
only negligible corrections compared to the linear term, thus giving the same
results as for the linear case.
### 5.3 dustyshock: Two fluid dust-gas shocks
The dustyshock problem (Paper I) is a two fluid version of the standard Sod
(1978) shock tube problem. Results were presented in Paper I using a constant
drag coefficient $K$ and no heat transfer between the gas and the dust phase.
Here we extend the test to non-linear drag regimes. While the evolution during
the transient stage is dependent on the drag regime, the solution during the
stationary stage remains unchanged, being a fixed function of the gas sound
speed and the dust-to-gas ratio. To facilitate the comparison between a
physical Epstein drag and the case of a constant drag coefficient (Paper I),
we fix the ratio $s^{2}/m_{\rm{d}}$ — $s$ being the grain size and
$m_{\rm{d}}$ the grain mass— to unity in code units, such that the Epstein
drag coefficient is unity in regions where $\hat{\rho}_{\mathrm{g}}=1$.
#### 5.3.1 dustyshock: Setup
Equal mass particles are placed in the 1D domain $x\in[-0.5,0.5]$, where for
$x<0$ we use $\rho_{\mathrm{g}}=\rho_{\mathrm{d}}=1$,
$v_{\rm{g}}=v_{\rm{d}}=0$ and $P_{\rm{g}}=1$, while for $x>0$
$\rho_{\mathrm{g}}=\rho_{\mathrm{d}}=0.125$, $v_{\rm{g}}=v_{\rm{d}}=0$ and
$P_{\rm{g}}=0.1$. We use an ideal gas equation of state $P=(\gamma-1)\rho u$
with $\gamma=5/3$. Initial particle spacing to the left of the shock in both
fluids is $\Delta x=0.001$ while to the right it is $\Delta x=0.008$, giving
$569$ equal mass particles in each phase. Standard SPH artificial viscosity
and conductivity terms are applied as in Paper I.
#### 5.3.2 dustyshock: Different drag regimes
Figure 3: Results of the dustyshock problem with a linear constant drag
coefficient ($K=1$) (left panel) and a non-linear Epstein drag (right panel)
where the drag coefficients are initially the same ahead of the shock. The
dust-to-gas ratio is set to unity. At $t=0.2$, the solutions are in the
transient stage where the analytic solution is not known. As an indication,
the solution for the later stationary stage is shown by the dotted/red lines.
The profiles differ as the damping is less efficient using the non-linear
Epstein regime.
Fig. 3 illustrates how the transient regime of the dustyshock is affected when
treating the drag with an astrophysical prescription where the drag
coefficient depends on the local density and the gas sound speed (right panel)
rather than a constant coefficient (left panel). Specifically, we use the non-
linear Epstein drag regime given by Eq. 39. In this dustyshock test, the non-
linear terms constitute a small correction to the linear Epstein drag regime
given by Eq. 40. Fig. 3 shows that in the case of the Epstein regime, the drag
is less efficient than for the constant coefficient case, leading to a larger
($\sim$ by a factor $7$) differential velocity between the gas and the dust
after $t=0.2$ in code units. The dust velocity profile is also smoother than
in the constant coefficient case. As a result, the kinetic energy is less
efficiently dissipated by the drag, leading to a less sharp peak in the
internal energy of the gas. The density profile of the dust is also closer to
its initial profile behind and ahead of the shock.
### 5.4 dustysedov: Two fluid dust-gas blast wave
The dustysedov problem (Paper I) involves the propagation of a blast wave in
an astrophysical mixture of dust and gas. We adopt physical units for this
problem, assuming a box size of $1$ pc, an ambient sound speed of $2\times
10^{4}$ cm/s and a gas density of $\rho_{0}=6\times 10^{-23}$ g/cm3 the energy
of the blast is $2\times 10^{51}$ erg and time is measured in units of $100$
years, roughly corresponding to a supernova blast wave propagating into the
interstellar medium. We these units, we choose the grain size, $0.1\mu$m and
the dust-to-gas ratio, $0.01$, to be typical of the interstellar medium. In
code units, this corresponds to an initial drag coefficient of $K=1$ outside
the blast radius. As for the dustyshock, we compare the results using a non-
linear Epstein drag prescription with with the constant coefficient case
described in Paper I.
#### 5.4.1 dustysedov: Setup
The problem is set up in a 3D periodic box ($x,y,z\in[-0.5,0.5]$), filled by
$50^{3}$ particles for both the gas and the dust. Gas particles are set up on
a regular cubic lattice, with the dust particles also on a cubic lattice but
shifted by half of the lattice step in each direction. For shock-capturing, we
set $\alpha_{\rm SPH}=1$ and $\beta_{\rm SPH}=2$ for the artificial viscosity
terms and $\alpha_{u}=1$ for the artificial conductivity term. An ideal gas
equation of state $P=(\gamma-1)\rho u$ is adopted with $\gamma=5/3$.
The internal energy is distributed of the gas over the particles located
inside a radius $r<r_{\rm{b}}$ where $r_{\rm b}$ is set to 2h (i.e., the
radius of the smoothing kernel which for $50^{3}$ particles and $\eta=1.2$ is
$0.048$). In code units the total blast energy is $E=1$, with
$\hat{\rho}_{\mathrm{g}}=1$ and $\hat{\rho}_{\mathrm{d}}=0.01$. For
$r>r_{\rm{b}}$, the gas sound speed is set to be $2\times 10^{-5}$ in code
units.
#### 5.4.2 dustysedov: Different drag regimes
Figure 4: Results of the 3D dustysedov test, showing the density in the gas
(left figure) and dust (right figure) from a Sedov blast wave propagating in
an astrophysical ($1\%$ dust-to-gas ratio) mixture of gas and $0.1\mu$m dust
grains in a $1$ pc box. The drag coefficient is constant ($K=1$, top panels)
or given by the Epstein regime (bottom panels). The low dust-to-gas ratio
means that the gas is only weakly affected by the drag from the dust, and is
thus close to the self-similar Sedov solution (dotted/red line). In the
Epstein case, the drag is much higher inside the blast radius and the dust
particles are efficiently piled up by the passage of the gas over-density.
Figure 5: Cross-section slice showing density in the midplane in the 3D
dustysedov problem, for both the gas (left panel) and the dust (right panel)
at $t=0.1$. Initially, the dust-to-gas ratio is $0.01$ and the drag
coefficient is given by the Epstein regime for grains of $0.1\mu$m in size.
$50^{3}$ SPH particles have been used in each phase.
Figs. 4 and 5 show the evolution of the gas and dust mixture where a constant
drag coefficient is used (top panels of Fig. 4) compared to a drag prescribed
by the Epstein regime (bottom panels of Fig. 4, Fig. 5). The gas profiles are
similar in both cases since the gas is poorly affected by the dust given the
low dust-to-gas ratio. However, the dust density profiles differ, essentially
due to the fact that the drag coefficient scales with the sound speed and is
thus higher in the inner blast region for the Epstein case. Thus, the dust is
efficiently piled up and accumulates in the gas over-density. As a result, the
dust is cleaned up by the gas in the inner regions of the blast, but is more
concentrated (by $\sim 10\%$) close to the gas over-density than for the
constant drag coefficient case.
The results using either explicit or implicit timestepping were found to be
indistinguishable. For the Epstein case, we found that roughly ten iterations
were required for the implicit scheme to converge on this problem.
### 5.5 dustydisc
The dustydisc problem concerns the evolution of a dusty gas mixture in a
protoplanetary disc (see Paper I for details). For our test case, we study how
the dust distribution is affected when considering a general non-linear
Epstein drag instead of the standard linear regime. The results obtained when
implicitly integrating the non-linear drag regime have been found to be
similar to benchmark tests performed with explicit integration.
#### 5.5.1 dustydisc: Setup
We setup $10^{5}$ gas particles and $10^{5}$ dust particles in a
$0.01M_{\odot}$ gas disc (with $0.0001M_{\odot}$ of dust) surrounding a
$1M_{\odot}$ star. The disc extends from 10 to 400 AU. Both gas and dust
particles are placed using a Monte-Carlo setup such that the surface density
profiles of both phases are $\Sigma\left(r\right)\propto r^{-1}$. The radial
profile of the gas temperature is taken to be $T\left(r\right)\propto
r^{-0.6}$ with a flaring $H/r=0.05$ at 100 AU. One code unit of time
corresponds to $10^{3}$ yrs.
#### 5.5.2 dustydisc: Evolution of the particles
Figure 6: Rendering of the density for the dust of a typical T-Tauri Star
protoplanerary disc using $2\times 10^{5}$ SPH particles, using an explicit
time integration in the linear Epstein regime (left panel) and an implicit
integrator in the full non-linear Epstein drag regimes (right panel).
Fig. 6 shows a face-on view of a protoplanetary disc, integrating the linear
Epstein regime (left panel) and the full non-linear Epstein drag (right
panel). The dust distributions are not found to exhibit significant
discrepancies. In the non-linear drag regime case however, the dust
distribution is slightly smoother since the drag (and thus, the coupling with
the gas phase) is more efficient.
Figure 7: Vertical settling of a dust grain ($1$cm in size) initially located
at $r_{0}=100$ AU and $z_{0}=2$ AU (solid/black), integrating implicitly the
non-linear Epstein regime. SPH results are compared to the explicit
integration of the linear Epstein regime (dashed/red) and the estimation given
by the damped harmonic oscillator approximation (pointed/red). In the full
non-linear drag regime, the settling is more efficient than for the linear
case since the vertical oscillations in the dust motion reaches a fraction
$z_{0}/H$ of the sound speed.
Fig. 7 compares the vertical motion of a dust grain initially located at
$z=z_{0}$ using the linear (explicit integration) and the full non linear
(implicit integration) Epstein regimes. In the full non-linear case, the
settling is more efficient since the vertical differential velocity between
the dust grains and the gas in the mid plane of the disc reaches a fraction
$z_{0}/H$ of the sound speed, meaning that the non-linear terms are no longer
negligible.
## 6 Conclusions
We have extended the SPH formalism for two-fluid dust and gas mixtures
developed in Paper I to handle the drag regimes usually encountered in a large
range of astrophysical contexts. Specifically, our algorithm is now designed
to treat the dynamics of grains surrounded by a dilute medium (Epstein regime)
or dense fluid (Stokes regime), for which the drag force can be either linear
or non-linear with respect to the differential velocity between the gas and
the dust.
Particular attention has been paid to developing an implicit timestepping
scheme to efficiently simulate the case of high drag, extending the scheme
proposed by Monaghan (1997) which we found to be unsatisfactory. We have
presented a new pairwise implicit scheme that, like the Monaghan (1997)
scheme, preserves the exact conservation of linear and angular momentum but,
unlike the Monaghan (1997) scheme, i) provides control over the accuracy of
the iterative procedure and ii) can incorporate non-linear terms for both
Epstein and Stokes drag. We found that when the ratio $r$ between the the
timestep and the drag stopping time is $1\lesssim r\lesssim 1000$, the
implicit timestepping is faster than a standard explicit integration. However,
at higher values of $r$, the algorithm is less efficient.
The accuracy of the generalised algorithm is benchmarked against the suite of
test problems presented in Paper I. In particular, the solutions obtained for
the dustybox problem are compared to their known analytic solutions for a
large range of non-linear drag regimes and the solutions of the dustywave,
dustyshock, dustysedov and dustydisc problems are benchmarked against
converged results obtained with explicit timestepping.
The two key issues addressed in this paper complete the study of our algorithm
developed in Paper I. Our intention is to apply it to various astrophysical
problems involving gas and dust mixtures in star and planet formation. A first
application is given in Ayliffe et al. (2011).
## Acknowledgments
We thank Ben Ayliffe and Matthew Bate and Joe Monaghan for useful discussions
and comments. Figures have been produced using splash (Price, 2007) with the
new giza backend by DP and James Wetter. We are grateful to the Australian
Research Council for funding via Discovery project grant DP1094585.
## References
* Ayliffe et al. (2011) Ayliffe B., Laibe G., Price D. J., Bate M. R., 2011, MNRAS, submitted
* Baines et al. (1965) Baines M. J., Williams I. P., Asebiomo A. S., 1965, MNRAS, 130, 63
* Blum & Wurm (2008) Blum J., Wurm G., 2008, ARA&A, 46, 21
* Chapman & Cowling (1970) Chapman C., Cowling T., 1970, The mathematical theory of non-uniform gases. Cambridge at the university press
* Chiang & Youdin (2010) Chiang E., Youdin A. N., 2010, Annual Review of Earth and Planetary Sciences, 38, 493
* Fan & Zhu (1998) Fan L.-S., Zhu C., 1998, Principles of Gas-Solid Flows. Cambridge University Press
* Kwok (1975) Kwok S., 1975, ApJ, 198, 583
* Laibe & Price (2011b) Laibe G., Price D. J., 2011b, MNRAS, submitted
* Laibe & Price (2011a) Laibe G., Price D. J., 2011a, MNRAS, in press
* Monaghan (1997) Monaghan J. J., 1997, Journal of Computational Physics, 138, 801
* Monaghan (2005) Monaghan J. J., 2005, Reports on Progress in Physics, 68, 1703
* Monaghan & Kocharyan (1995) Monaghan J. J., Kocharyan A., 1995, Computer Physics Communications, 87, 225
* Paardekooper & Mellema (2006) Paardekooper S.-J., Mellema G., 2006, A&A, 453, 1129
* Price (2007) Price D. J., 2007, Publications of the Astronomical Society of Australia, 24, 159
* Price & Monaghan (2007) Price D. J., Monaghan J. J., 2007, MNRAS, 374, 1347
* Sod (1978) Sod G. A., 1978, Journal of Computational Physics, 27, 1
* Stepinski & Valageas (1996) Stepinski T. F., Valageas P., 1996, A&A, 309, 301
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|
arxiv-papers
| 2011-11-14T02:28:49 |
2024-09-04T02:49:24.297384
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Guillaume Laibe (Monash), Daniel J. Price (Monash)",
"submitter": "Guillaume Laibe",
"url": "https://arxiv.org/abs/1111.3089"
}
|
1111.3090
|
# Dusty gas with SPH — I. Algorithm and test suite
Guillaume Laibe, Daniel J. Price
Monash Centre for Astrophysics (MoCA) and School of Mathematical Sciences,
Monash University, Clayton, Vic 3800, Australia
###### Abstract
We present a new algorithm for simulating two-fluid gas and dust mixtures in
Smoothed Particle Hydrodynamics (SPH), systematically addressing a number of
key issues including the generalised SPH density estimate in multi-fluid
systems, the consistent treatment of variable smoothing length terms, finite
particle size, time step stability, thermal coupling terms and the choice of
kernel and smoothing length used in the drag operator. We find that using
double-hump shaped kernels improves the accuracy of the drag interpolation by
a factor of several hundred compared to the use of standard SPH bell-shaped
kernels, at no additional computational expense. In order to benchmark our
algorithm, we have developed a comprehensive suite of standardised, simple
test problems for gas and dust mixtures: dustybox, dustywave, dustyshock,
dustysedov and dustydisc, the first three of which have known analytic
solutions.
We present the validation of our algorithm against all of these tests. In
doing so, we show that the spatial resolution criterion $\Delta\lesssim c_{\rm
s}t_{\rm s}$ is a necessary condition in all gas+dust codes that becomes
critical at high drag (i.e. small stopping time $t_{\rm s}$) in order to
correctly predict the dynamics. Implicit timestepping and the implementation
of realistic astrophysical drag regimes are addressed in a companion paper.
###### keywords:
hydrodynamics — methods: numerical — ISM: dust, extinction — protoplanetary
discs — planets and satellites: formation
††pagerange: Dusty gas with SPH — I. Algorithm and test
suite–References††pubyear: 2011
## 1 Introduction
Most of our observational information regarding the interstellar medium comes
to us via dust. Over the last few years, observations using the _Spitzer_ and
_Herschel_ space telescopes have substantially improved our observational
sensitivity to and resolution of dust emission in a wide range of
astrophysical environments. Dust grains provide the materials from which the
solid cores required for the planet formation process are built (see e.g.
Chiang & Youdin 2010a). They also modify the dynamical evolution of the
surrounding gas by exchanging momentum and energy via microscopic collisions
(Epstein, 1924; Baines et al., 1965). Dust grains are also the main sources of
the opacities in star-forming molecular clouds, thus determining their
evolution by controlling the thermodynamics. Accurate determination of both
the dynamics of the star and planet formation process and its observational
signature thus require modelling the coupled evolution of gas and dust.
Given that the N-body evolution of solid particles in a mixture of gas and
solid material would be prohibitive in terms of both physical complexity and
computational cost, the usual approach is to regard the solid phase as a
continuum and the mixture as a two-fluid system coupled by a drag term. This
requires averaging physical quantities over a control volume $V$ that is large
enough to be statistically meaningful but sufficiently small compared to the
macroscopic scale to allow a continuum description. In diluted astrophysical
media, the frequency of collisions between dust particles are infrequent
enough that the intrinsic pressure of the dust phase can be regarded as
negligible to a very good level of approximation, leaving the dust as a free-
streaming collisionless fluid whose motion is controlled solely by
gravitational forces and the drag-term interaction with the gas.
However, even with a continuous description of the mixture, the equations can
be solved analytically only for a few simple cases (the solutions to two
specific problems, dustybox and dustywave, corresponding to mutually
interpenetrating fluids and acoustic waves propagating in a dusty gas,
respectively, are derived in Laibe & Price 2011a). As a result, numerical
codes have been developed in order to model more realistic systems based
either on $N$-body dust particles in Eulerian grid-based hydrodynamics (e.g.
Fromang & Papaloizou 2006; Paardekooper & Mellema 2006; Johansen et al. 2007;
Miniati 2010; Bai & Stone 2010) or with a two-fluid Smoothed Particle
Hydrodynamics (SPH) approach.
SPH methods for simulating two-fluid mixtures were first developed by Monaghan
& Kocharyan (1995), improved (via an implicit treatment of the drag terms) in
Monaghan (1997) and applied in an astrophysical context to the dynamics of
dust grains in protoplanetary disks (Maddison et al., 2003; Rice et al., 2004;
Barrière-Fouchet et al., 2005). The particle-based nature of the SPH formalism
means that the addition of a dusty fluid is natural. More importantly, the
drag term that couples the two phases can be implemented such that the total
linear and angular momentum of the system are exactly (and simultaneously)
conserved, in line with the Hamiltonian and exactly conservative nature of the
core SPH method (e.g. Price, 2011).
However, the standard methods for treating dusty gas in SPH were developed
over 15 years ago and our initial attempts to simply apply the existing
formulations uncovered several issues that needed to be addressed.
Specifically: 1) the original formulations assumed a spatially constant SPH
smoothing length; 2) the SPH terms for the conservative part of the equations
should be derived from a Lagrangian; 3) we found that the use of the standard
cubic spline kernel for drag terms could be significantly inaccurate; 4) we
encountered several previously unexplored resolution issues in simulating two-
fluid mixtures; 5) aspects of the implicit timestepping scheme suggested by
Monaghan (1997) were found to be problematic; 6) that treatments of drag have
to date generally limited to linear drag regimes; and finally 7) that the
existing schemes — having been developed with both astrophysical and
geophysical dust problems in mind — have not been widely benchmarked on
problems appropriate to astrophysics; Indeed there is a general lack of
standardised test problems for two-fluid dust/gas codes, a problem partially
addressed by our first paper (Laibe & Price, 2011a).
In this and a companion paper (Laibe & Price 2011c, hereafter Paper II), we
set out to systematically address issues 1)–7) in order to develop a robust
and accurate code for simulating the dynamics of dust in star and planet
formation. The importance of modelling the dust-gas interaction has been
highlighted by recent studies showing that instabilities in dust-gas mixtures
are good candidates for triggering the concentration of dust during
planetesimal formation (Goodman & Pindor, 2000; Youdin & Goodman, 2005).
The continuum equations and the relevant parameters describing the evolution
of dust-gas mixtures are given in Section 2.1. Section 2 describes the two-
fluid SPH algorithm, addressing issues 1)-3). The code is benchmarked against
a suite of test problems that we have specifically designed in order to
provide standardised benchmarks for other two-fluid gas/dust codes, addressing
issues 4) and 7) (Sec. 4). The implicit timestepping scheme and treatment of
non-linear drag (issues 5 and 6) are discussed in Paper II.
## 2 Two-fluid mixtures in SPH
### 2.1 Two-fluid gas and dust mixtures
#### 2.1.1 Densities
The fact that dust grains of finite size occupy a finite volume is accounted
for by defining the volume fraction available to the gas according to (e.g.
Marble, 1970; Harlow & Amsden, 1975)
$\theta=1-\frac{\hat{\rho}_{\mathrm{d}}}{\rho_{\mathrm{d}}}.$ (1)
This means that the volume densities of gas and dust $\hat{\rho}_{\mathrm{g}}$
and $\hat{\rho}_{\mathrm{d}}$, respectively, are distinguished from the
intrinsic densities denoted $\rho_{\mathrm{g}}$ and $\rho_{\mathrm{d}}$,
respectively, according to
$\displaystyle\hat{\rho}_{\mathrm{d}}$ $\displaystyle=$
$\displaystyle(1-\theta)\rho_{\mathrm{d}},$ (2)
$\displaystyle\hat{\rho}_{\mathrm{g}}$ $\displaystyle=$
$\displaystyle\theta\rho_{\mathrm{g}}.$ (3)
The effects associated with finite dust particle size are mostly negligible in
astrophysical problems since typically the intrinsic dust density
$\rho_{\mathrm{d}}$ is much higher than the volume density
$\hat{\rho}_{\mathrm{d}}$ and thus $\theta\approx 1$. We retain these terms,
as in earlier SPH formulations (c.f. Monaghan & Kocharyan, 1995) in order to
retain a general algorithm that can be applied both within and outside of
astrophysics.
The conservation of mass in a two-fluid mixture is thus expressed by the
continuity equations
$\displaystyle\frac{\partial\hat{\rho}_{\mathrm{g}}}{\partial
t}+\nabla.\left(\hat{\rho}_{\mathrm{g}}\textbf{v}_{\mathrm{g}}\right)$
$\displaystyle=$ $\displaystyle 0,$ (4)
$\displaystyle\frac{\partial\hat{\rho}_{\mathrm{d}}}{\partial
t}+\nabla.\left(\hat{\rho}_{\mathrm{d}}\textbf{v}_{\mathrm{d}}\right)$
$\displaystyle=$ $\displaystyle 0,$ (5)
where ${\bf v}_{\rm g}$ and ${\bf v}_{\rm d}$ are the gas and dust fluid
velocities, respectively.
#### 2.1.2 Equations of motion
The equations of motion, expressing momentum conservation in a continuous,
inviscid, two-fluid mixture of gas and dust are given by
$\displaystyle\hat{\rho}_{\mathrm{g}}\left(\frac{\partial\textbf{v}_{\mathrm{g}}}{\partial
t}+\textbf{v}_{\mathrm{g}}.\nabla\textbf{v}_{\mathrm{g}}\right)$
$\displaystyle=$ $\displaystyle-\theta\phantom{.}\nabla P_{\rm
g}+\hat{\rho}_{\mathrm{g}}\textbf{f}-\textbf{F}^{\rm V}_{\mathrm{drag}},$ (6)
$\displaystyle\hat{\rho}_{\mathrm{d}}\left(\frac{\partial\textbf{v}_{\mathrm{d}}}{\partial
t}+\textbf{v}_{\mathrm{d}}.\nabla\textbf{v}_{\mathrm{d}}\right)$
$\displaystyle=$ $\displaystyle-\nabla P_{\rm d}-\left(1-\theta\right)\nabla
P_{\rm g}+\hat{\rho}_{\mathrm{d}}\textbf{f}+\textbf{F}^{\rm
V}_{\mathrm{drag}},$ (7)
where $P_{\rm g}$ and $P_{\rm d}$ are the intrinsic pressures. Any intrinsic
viscosities have been neglected. For astrophysical purposes it may be assumed
that the dust is pressureless, i.e. $P_{\rm d}=0$. Similarly, the term
$\left(1-\theta\right)\nabla P_{\rm g}$ in the momentum equation for the dust
phase — a buoyancy term related to the finite size of the dust particles — is
in general negligibly small. The reader should note that the definitions of
physical quantities in a two fluid medium require the local fluid volume over
which the averaging is performed to be defined (see, e.g. Marble, 1970; Fan &
Zhu, 1998).
The two fluids exchange momentum $\textbf{F}^{\rm V}_{\mathrm{drag}}$, the
drag force per unit volume, the expression for which is obtained by averaging
the local drag stress tensor (denoted $\mathbf{\epsilon}^{ij}_{\rm drag}$)
over the surface area of the dust grains:
$F_{\rm drag}^{{\rm
V},i}=\frac{1}{V}\int_{A_{\mathrm{d}}}\mathbf{\epsilon}^{ij}_{\rm
drag}\mathrm{d}A^{j}.$ (8)
In the case where the local distribution of dust particles is homogeneous
(i.e., dust particles have the same mass, size and intrinsic density), Eq. 8
simplifies to
$\textbf{F}^{\rm
V}_{\mathrm{drag}}=K(\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}}).$ (9)
Note that since $\textbf{F}^{\rm V}_{\mathrm{drag}}$ is a force per unit
volume, the drag coefficient $K$, has dimensions of mass per unit volume per
unit time. This coefficient is related to the drag coefficient on a single
grain (denoted $K_{\rm s}$) by
$K=\frac{\hat{\rho}_{\mathrm{d}}}{m_{\mathrm{d}}}K_{\rm s}$ (10)
where $m_{\mathrm{d}}$ is the mass per grain. The drag force (not per unit
volume) on a single grain is given by
${\bf F}_{\rm drag}=K_{\rm
s}(\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}}).$ (11)
In general $K$ (or equivalently $K_{s}$) can itself be a function of the
relative velocity between the two fluids $\Delta
v\equiv|\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}}|$, resulting in a non-
linear drag regime. In this, Paper I we consider the simplest case of linear
drag, where $K$ is constant with respect to $\Delta v$. Extension of our
scheme to the main non-linear regimes applicable to astrophysics are
considered in Paper II.
Finally, it should be noted that in general additional forces (e.g. the
carried mass, Basset and Saffman forces, Fan & Zhu 1998) may be present in
two-fluid systems. We have assumed in adopting Eqs. 6–7 that these forces can
be neglected for astrophysical applications.
Figure 1: Computing density in SPH gas (solid points) and dust (hollow
circles) mixtures. Standard bell-shaped, Gaussian-like, kernels are adopted
(weighting indicated by the shading), with a single smoothing length on each
particle related to the local number density of particles of the _same type_.
This provides good density estimates in both extremes — where dust is
concentrated below the gas scale (left panel) and where gas is concentrated
below the dust scale (right panel). The density of another fluid at the
position of a reference fluid (e.g. dust density at the location of a gas
particle) is computed using the same smoothing length but only neighbours of
the desired type. This density is thus allowed to be identically zero, as
would be the case for the density of gas-at-dust in the left panel (top), or
dust-at-gas in the right panel.
#### 2.1.3 Energy equation
The evolution equation for the specific internal energy of the gas, $u_{\rm
g}$, is given by
$\hat{\rho}_{\mathrm{g}}\frac{{\rm d}u_{\rm g}}{{\rm
d}t}=-P_{\mathrm{g}}\left[\theta\nabla\cdot{\bf v}_{\rm
g}+(1-\theta)\nabla\cdot{\bf v}_{\rm d}\right]+\Lambda_{\rm drag}+\Lambda_{\rm
therm},$ (12)
where the first term corresponds to the usual compressive ($P{\rm d}V$) term
with the volume reduced by the dust filling factor $\theta$. The second term
is the work done by the gas in triggering buoyancy effects. The third term is
the frictional heating due to the drag force, given by
$\Lambda_{\rm
drag}=\hat{\rho}_{\mathrm{g}}K(\textbf{v}_{\mathrm{g}}-\textbf{v}_{\mathrm{d}})^{2}.$
(13)
The fourth, thermal coupling, term arises when the internal temperature of the
grains differs from the gas temperature (c.f. Marble, 1970; Harlow & Amsden,
1975), and in general consists of terms related to heat transfer due to
conduction ($\Lambda_{\rm cond}$) and radiation ($\Lambda_{\rm rad}$), given
by
$\Lambda_{\rm therm}\equiv\Lambda_{\rm cond}+\Lambda_{\rm rad}=Q(T_{\rm
g}-T_{\rm d})+R(aT_{\rm g}^{4}-aT_{\rm d}^{4}),$ (14)
where $T_{\rm g}$ and $T_{\rm d}$ are the temperatures of the gas and dust,
respectively, $a$ is the radiation constant and $Q$ and $R$ are coefficients,
dependent on gas and dust properties, that characterise the heat transfer. The
thermal energy of the dust evolves according to
$\hat{\rho}_{\mathrm{d}}\frac{{\rm d}u_{\rm d}}{{\rm d}t}=-\Lambda_{\rm
therm}.$ (15)
### 2.2 Densities for two-fluid mixtures in SPH
#### 2.2.1 Computing densities in two-fluid SPH
For two-fluid mixtures, we require a density estimate _for each phase_ ,
corresponding to the exact solution of Eqs. 4 and 5 in SPH. The main
complication arises from the fact that the local particle spacing can be
different for each fluid, implying that the two fluids should have different
resolution lengths calculated based on the local particle number density of
their own type. Figure 1 illustrates the two limiting cases, i.e. a high
concentration of dust in a diluted gas (left panel) and conversely a high
concentration of gas in a low density fluid of dust (right panel). In each
case the smoothing length for each type is determined by the local number
density of particles of _the same type_. That is, the SPH translation of Eqs.
4 and 5 correspond to
$\displaystyle\hat{\rho}_{a}=\sum_{b}m_{b}W_{ab}(h_{a});$ $\displaystyle
h_{a}=\eta\left(\frac{m_{a}}{\hat{\rho}_{a}}\right)^{1/\nu},$ (16)
$\displaystyle\hat{\rho}_{i}=\sum_{j}m_{j}W_{ij}(h_{i});$ $\displaystyle
h_{i}=\eta\left(\frac{m_{j}}{\hat{\rho}_{i}}\right)^{1/\nu},$ (17)
where $\nu$ is the number of spatial dimensions and $\eta$ is a constant
determining the resolution length as a function of the local particle spacing
(typically $\eta=1.2$ is a good choice for the standard cubic spline kernel,
see Price 2011). We adopt the convention that the indices $a,b,c$ refer to
quantities computed on gas particles while $i,j,k$ refer to quantities
computed on dust particles. Note that the densities and smoothing lengths are
independently computed for each fluid and are thus — so far — only defined on
particles of the same type. The numerical solution of Eqs. 16 and 17 involves
determining both $\hat{\rho}$ and $h$ for each type simultaneously, since they
are mutually dependent, thus requiring an iterative procedure. The procedure
is identical to that adopted in standard variable smoothing length SPH
formulations (see e.g. Price & Monaghan 2007 for details).
An additional complication arises from the need to compute the volume filling
fraction $\theta$ (Eq. 1), defined on a _gas_ particle, $a$, according to
$\theta_{a}=1-\frac{\hat{\rho}_{{\rm d},a}}{\rho_{\rm d}},$ (18)
which depends on the density of _dust_ at the gas particle location. Initially
we considered computing this density using a second smoothing length for each
particle based on neighbours of the _other_ type (this would in turn lead to
multiple smoothing lengths on each particle if more types were present).
However, the key point, illustrated by the right hand panel of Fig. 1, is that
it should be possible for the local density of dust at the gas location to be
identically zero (giving $\theta=1$) if no dust particles are found within the
kernel radius computed with the gas smoothing length. Thus, the density of
dust-at-gas should be calculated according to
$\hat{\rho}_{{\rm d},a}=\sum^{N_{neigh,dust}}_{j=1}m_{j}W_{aj}(h_{a}),$ (19)
where $h_{a}$ is the smoothing length of the _gas_ particle computed using gas
neighbours as in Eq. 16. In the case where no dust neighbours fall within the
kernel radius, $\rho_{{\rm d},a}=0$. This is a very simple and efficient
method that can easily be generalised to multiple fluids, requires only one
smoothing length per particle and does not require any significant additional
computational expense.
The discussion above resolves the first issue highlighted in Sec. 1, namely
how to deal with variable resolution in multi-fluid SPH, generalising the
earlier fixed-smoothing-length formulation of Monaghan & Kocharyan (1995). A
similar discussion to the above applies to gravitational force softening on
multiple fluids in $N$-body/SPH codes where the softening formulation is
derived from a kernel density estimate (Price & Monaghan, 2007), in particular
for the case of a mixture of dark matter and baryonic gas (e.g. Merlin et al.,
2010; Iannuzzi & Dolag, 2011).
#### 2.2.2 Kernel function
The kernel function itself can be written as a function of the smoothing
length $h$ and the dimensionless variable $q=|{\bf r}-{\bf r}^{\prime}|/h$ in
the form
$W\left(r,h\right)=\frac{\sigma}{h^{\nu}}f\left(q\right),$ (20)
where $\sigma$ is a normalisation constant. The standard Gaussian kernel is
given by
$f(q)=e^{-q^{2}},$ (21)
where $\sigma=\pi^{-\nu/2}$. The Gaussian is infinitely smooth
(differentiable) but has the practical disadvantage of infinite range. A
standard alternative (providing a Gaussian-like kernel but truncated at $2h$)
is the $M_{4}$ cubic spline kernel (Monaghan, 1992)
$f(q)=\begin{cases}1-\frac{3}{2}q^{2}+\frac{3}{4}q^{3},&0\leq q<1;\\\
\frac{1}{4}\left(2-q\right)^{3},&1\leq q<2;\\\ 0,&q\geq 2,\end{cases}$ (22)
where $\sigma=\left[2/3,10/\left(7\pi\right),1/\pi\right]$ in
$\left[1,2,3\right]$ dimensions. An error analysis of the SPH density estimate
(e.g. Price, 2011) shows that in general the measure of a good density kernel
is that the normalisation condition
$\sum_{b}\frac{m_{b}}{\rho_{b}}W_{ab}\approx 1,$ (23)
is well satisfied for typical SPH particle distributions, corresponding to
$\int W{\rm d}V=1$ in the continuum limit. In general most bell-shaped
(Gaussian-like) kernels, such as the cubic spline, fulfil this criterion (Fulk
& Quinn, 1996).
More accurate density estimates can be obtained — at the price of additional
computational expense — by using kernels with extended range that form a
better approximation to the Gaussian (see Price, 2011). In particular the
$M_{6}$ quintic kernel, truncated at $3h$, gives results that are in practice
largely indistinguishable from the Gaussian, with the functional form
$f(q)=\begin{cases}(3-q)^{5}-6(2-q)^{5}+15(1-q)^{5},&\text{$0\leq q<1$;}\\\
(3-q)^{5}-6(2-q)^{5},&\text{$1\leq q<2$;}\\\ (3-q)^{5},&\text{$2\leq q<3$;}\\\
0,&\text{$q\geq 3$}.\end{cases}$ (24)
where $\sigma=[1/24,96/(1199\pi),1/(20\pi)]$. Use of the quintic is a factor
of $(3/2)^{3}\approx 3.4$ times more expensive than the cubic spline (or other
$2h$-truncated kernels) in three dimensions.
The functional form of the $M_{4}$ cubic and $M_{6}$ quintic spline kernels
are shown in the top row of Fig. 2, showing the kernel function $f(q)$
(solid/black lines) and its first (dashed/red lines) and second (short
dashed/green line) derivatives.
### 2.3 Equations of motion
As discussed by Price (2011), specifying the manner in which the density is
calculated in SPH can be used to self-consistently determine the equations of
motion and energy from a variational principle, using only the additional
constraint of the first law of thermodynamics. For a two-fluid system, only
the dissipationless part of the algorithm can be derived in this manner — that
is, not including the drag terms.
#### 2.3.1 Lagrangian
For a system consisting of gas and dust, the Lagrangian is given by
$L=\sum_{b}m_{b}\left[\frac{1}{2}\textbf{v}_{b}^{2}-u_{b}\left(\rho_{b},s_{b}\right)\right]+\sum_{k}m_{k}\left(\frac{1}{2}\textbf{v}_{k}^{2}\right)$
(25)
where $u_{b}$ is the thermal energy per unit mass of the gas (in general a
function of the entropy $s$ and _intrinsic_ density $\rho$). The equations of
motion can be derived from the Euler-Lagrange equations,
$\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{\partial L}{\partial{\bf
v}}\right)=\frac{\partial L}{\partial{\bf r}}.$ (26)
#### 2.3.2 Equations of motion for the gas
We first consider the evolution of the gas particles. The partial derivative
of the Lagrangian with respect to the velocity $\textbf{v}_{a}$ of a given gas
particle $a$ provides:
$\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{\partial
L}{\partial\textbf{v}_{a}}\right)=m_{a}\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}.$
(27)
The partial derivative of the Lagrangian with respect to the position
$\textbf{r}_{a}$ of the gas particle $a$ is given by
$\frac{\partial L}{\partial\textbf{r}_{a}}=-\sum_{b}m_{b}\left.\frac{\partial
u_{b}}{\partial\rho_{b}}\right|_{s}\frac{\partial\rho_{b}}{\partial\textbf{r}_{a}},$
(28)
where the entropy is constant for a non-dissipative system. Eq. 28 differs
from the usual expression for single fluids because of the distinction between
the intrinsic and volume density of the gas caused by the finite volume
occupied by dust. That is, the thermal energy depends on the _intrinsic_
density rather than the volume density, giving
$\left.\frac{\partial
u_{b}}{\partial\rho_{b}}\right|_{s}=\frac{P_{b}}{\rho_{b}^{2}}=\frac{\theta_{b}^{2}P_{b}}{\hat{\rho}_{b}^{2}}.$
(29)
The derivative of the intrinsic density with respect to the particle
coordinates is given by
$\frac{\partial\rho_{b}}{\partial\textbf{r}_{a}}=\left.\frac{\partial\rho_{b}}{\partial\hat{\rho}_{b}}\right|_{\theta_{b}}\frac{\partial\hat{\rho}_{b}}{\partial\textbf{r}_{a}}+\left.\frac{\partial\rho_{b}}{\partial\theta_{b}}\right|_{\hat{\rho}_{b}}\frac{\partial\theta_{b}}{\partial\textbf{r}_{a}},$
(30)
where, from 18, we have
$\left.\frac{\partial\rho_{b}}{\partial\hat{\rho}_{b}}\right|_{\theta_{b}}=\displaystyle\frac{1}{\theta_{b}};\hskip
28.45274pt\left.\frac{\partial\rho_{b}}{\partial\theta_{b}}\right|_{\hat{\rho}_{b}}=-\frac{\hat{\rho}_{b}}{\theta_{b}^{2}}.$
(31)
The spatial derivative of the density sum for the gas (Eq. 16) is given by
$\frac{\partial\hat{\rho}_{b}}{\partial\textbf{r}_{a}}=\frac{1}{\Omega_{b}}\sum_{c}m_{c}\left(\delta_{ba}-\delta_{ca}\right)\nabla_{a}W_{bc}(h_{b}),$
(32)
where $\Omega$ is the usual variable smoothing length term
$\Omega_{b}\equiv 1-\frac{\partial
h_{b}}{\partial\hat{\rho}_{b}}\sum_{c}m_{c}\frac{\partial W_{\\!\\!\ b\\!\\!\
c}\left(h_{b}\right)}{\partial h_{b}}.$ (33)
The spatial derivative of the volume filling fraction $\theta$ is given, from
Eq. 18 by
$\frac{\partial\theta_{b}}{\partial\textbf{r}_{a}}=-\frac{1}{\rho_{\rm
d}}\frac{\partial\hat{\rho}_{{\rm d},b}}{\partial\textbf{r}_{a}},$ (34)
where
$\displaystyle\frac{\partial\hat{\rho}_{{\rm d},b}}{\partial\textbf{r}_{a}}=$
$\displaystyle\sum_{j}m_{j}\left(\delta_{ba}-\delta_{ja}\right)\nabla_{a}W_{bj}(h_{b})$
$\displaystyle+\frac{1-\Omega_{\mathrm{d},b}}{\Omega_{b}}\sum_{c}m_{c}\left(\delta_{ba}-\delta_{ca}\right)\nabla_{a}W_{bc}(h_{b}),$
(35)
where $\Omega_{\rm d}$ is $\Omega$ computed only using dust particle
neighbours, i.e.
$\Omega_{\mathrm{d},b}=1-\frac{\partial
h_{b}}{\partial\hat{\rho}_{{\mathrm{d},b}}}\sum_{j}m_{j}\frac{\partial
W_{\\!\\!\ b\\!\\!\ j}\left(h_{b}\right)}{\partial h_{b}}.$ (36)
Collecting Eqs. 28–36, noting that $\delta_{ja}=0$ (since a gas and dust index
can never refer to the same particle) and using the fact that $\nabla_{\\!\\!\
a}W_{\\!\\!\ b\\!\\!\ c}=-\nabla_{\\!\\!\ a}W_{\\!\\!\ c\\!\\!\ b}$, gives
$\displaystyle\frac{\partial L}{\partial\textbf{r}_{a}}=$ $\displaystyle-
m_{a}\sum_{b}m_{b}\left[\frac{P_{a}\tilde{\theta}_{a}}{\Omega_{a}\hat{\rho}_{a}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\
b}\left(h_{a}\right)+\frac{P_{b}\tilde{\theta}_{b}}{\Omega_{b}\hat{\rho}_{b}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\ b}\left(h_{b}\right)\right]$ $\displaystyle-
m_{a}\sum_{j}m_{j}\frac{P_{a}\left(1-\theta_{a}\right)}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\nabla_{\\!\\!\ a}W_{\\!\\!\ a\\!\\!\ j}\left(h_{a}\right),$ (37)
where we have defined $\tilde{\theta}$ to include the correction terms for a
variable smoothing length, i.e.
$\tilde{\theta}\equiv\theta+\frac{\hat{\rho}_{\rm
g}}{\hat{\rho}_{\mathrm{d}}}(1-\theta)(1-\Omega_{\rm d}).$ (38)
Although this correction is necessary for strict energy conservation, it is
expected to be negligibly small in practice, since $(1-\theta)$ is negligible
for small grains and $(1-\Omega_{\rm d})$ is $\mathcal{O}(h^{2})$. Finally,
the equations of motion for a gas particle, from the Euler-Lagrange equations,
are given by
$\displaystyle\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}=$
$\displaystyle-\sum_{b}m_{b}\left[\frac{P_{a}\tilde{\theta}_{a}}{\Omega_{a}\hat{\rho}_{a}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\
b}\left(h_{a}\right)+\frac{P_{b}\tilde{\theta}_{b}}{\Omega_{b}\hat{\rho}_{b}^{2}}\nabla_{\\!\\!\
a}W_{\\!\\!\ a\\!\\!\ b}\left(h_{b}\right)\right]$
$\displaystyle-\sum_{j}m_{j}\frac{P_{a}\left(1-\theta_{a}\right)}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\nabla_{\\!\\!\ a}W_{\\!\\!\ a\\!\\!\ j}\left(h_{a}\right).$ (39)
The reader should note that while the first term is a summation over gas
particle neighbours, the second is summed over dust particle neighbours. Eq.
39 may be straightforwardly shown to be a direct translation of Eq. 6 into SPH
form. Note that the summation over dust particles (the buoyancy term) does not
involve $\Omega$ since the smoothing length is independent of the dust
particle positions.
#### 2.3.3 Equations of motion for the dust
The partial derivative of the Lagrangian with respect to the velocity ${\bf
v}_{i}$ of a given dust particle $i$ gives
$\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{\partial L}{\partial{\bf
v}_{i}}\right)=m_{i}\frac{\mathrm{d}{\bf v}_{i}}{\mathrm{d}t}.$ (40)
A buoyancy term arises in the dust because of the dependence of the gas
internal energy on $\theta$, which in turn depends on the positions of dust
particles. That is,
$\frac{\partial L}{\partial{\bf r}_{i}}=-\sum_{b}m_{b}\left.\frac{\partial
u_{b}}{\partial\rho_{b}}\right|_{s}\frac{\partial\rho_{b}}{\partial{\bf
r}_{i}},$ (41)
where
$\frac{\partial\rho_{b}}{\partial{\bf
r}_{i}}=\frac{\partial\rho_{b}}{\partial\theta_{b}}\frac{\partial\theta_{b}}{\partial{\bf
r}_{i}},$ (42)
and in turn,
$\displaystyle\frac{\partial\theta_{b}}{\partial{\bf r}_{i}}=$
$\displaystyle-\frac{1}{\rho_{\rm d}}\frac{\partial\hat{\rho}_{{\rm
d},b}}{\partial{\bf r}_{i}}$ $\displaystyle=$
$\displaystyle-\frac{1}{\rho_{\rm
d}}\sum_{j}m_{j}\left(\delta_{bi}-\delta_{ji}\right)\nabla_{i}W_{bj}(h_{b})$
$\displaystyle-\frac{1-\Omega_{\mathrm{d},b}}{\Omega_{b}}\sum_{c}m_{c}\left(\delta_{bi}-\delta_{ci}\right)\nabla_{i}W_{bc}(h_{b})$
(43)
Collecting Eqs. 27–43 and noting that $\delta_{bi}=\delta_{ci}=0$, we obtain
the equations of motion for a dust particle in the form
$\frac{\mathrm{d}{\bf
v}_{i}}{\mathrm{d}t}=\sum_{b}m_{b}\frac{P_{b}\left(1-\theta_{b}\right)}{\hat{\rho}_{b}\hat{\rho}_{{\mathrm{d},b}}}\nabla_{\\!\\!\
i}W_{\\!\\!\ b\\!\\!\ i}\left(h_{b}\right).$ (44)
where we have written the kernel using $\nabla_{\\!\\!\ i}W_{\\!\\!\ i\\!\\!\
b}=-\nabla_{\\!\\!\ i}W_{\\!\\!\ b\\!\\!\ i}$ to show that the force is equal
and opposite to that in the gas (Eq. 39). It may be straightforwardly verified
that Eq. 44 is indeed a direct translation of Eq. 7 in SPH form.
Equations 39 and 44 may be combined to show that the total momentum is exactly
conserved, i.e.
$\frac{\rm d}{{\rm d}t}\left(\sum_{a}m_{a}{\bf v}_{a}+\sum_{i}m_{i}{\bf
v}_{i}\right)=0.$ (45)
#### 2.3.4 Internal energy equation for the gas
The SPH form of the non-dissipative terms in the internal energy equation for
the gas (Eq. 12) can similarly be derived from the SPH density estimates. In
the absence of dissipation the evolution equation for a given gas particle $a$
is given by
$\frac{{\rm d}u_{a}}{{\rm d}t}=\frac{P_{a}}{\rho_{a}^{2}}\frac{{\rm
d}\rho_{a}}{{\rm
d}t}=\frac{P_{a}}{\rho_{a}^{2}}\left[\left.\frac{\partial\rho_{a}}{\partial\hat{\rho}_{a}}\right|_{\theta_{a}}\frac{{\rm
d}\hat{\rho}_{a}}{{\rm
d}t}+\left.\frac{\partial\rho_{a}}{\partial\theta_{a}}\right|_{\hat{\rho}_{a}}\frac{{\rm
d}\theta_{a}}{{\rm d}t}\right].$ (46)
Using the expressions (31) and simplifying using (18), we have
$\frac{{\rm d}u_{a}}{{\rm
d}t}=\frac{\theta_{a}P_{a}}{\hat{\rho}_{a}^{2}}\frac{{\rm
d}\hat{\rho}_{a}}{{\rm
d}t}+\frac{(1-\theta_{a})P_{a}}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\frac{{\rm d}\hat{\rho}_{{\rm d},a}}{{\rm d}t}.$ (47)
Taking the time derivative of the density sums (16) and (19) we have
$\displaystyle\frac{{\rm d}\hat{\rho}_{a}}{{\rm d}t}$ $\displaystyle=$
$\displaystyle\frac{1}{\Omega_{a}}\sum_{b}m_{b}\left(\textbf{v}_{a}-\textbf{v}_{b}\right)\cdot\nabla_{a}W_{ab}(h_{a}),$
(48) $\displaystyle\frac{{\rm d}\hat{\rho}_{{\rm d},a}}{{\rm d}t}$
$\displaystyle=$ $\displaystyle\sum^{N_{neigh,dust}}_{j=1}m_{j}\left({\bf
v}_{a}-{\bf v}_{j}\right)\cdot\nabla_{a}W_{aj}(h_{a}),$ (49)
$\displaystyle+\frac{(1-\Omega_{{\rm
d},a})}{\Omega_{a}}\sum_{b}m_{b}\left(\textbf{v}_{a}-\textbf{v}_{b}\right)\cdot\nabla_{a}W_{ab}(h_{a})$
giving the SPH internal energy equation in the form
$\displaystyle\frac{{\rm d}u_{a}}{{\rm d}t}$
$\displaystyle=\frac{\tilde{\theta}_{a}P_{a}}{\Omega_{a}\hat{\rho}_{a}^{2}}\sum_{b}m_{b}\left(\textbf{v}_{a}-\textbf{v}_{b}\right)\cdot\nabla_{a}W_{ab}(h_{a})$
$\displaystyle+\frac{(1-\theta_{a})P_{a}}{\hat{\rho}_{a}\hat{\rho}_{{\rm
d},a}}\sum^{N_{neigh,dust}}_{j=1}m_{j}\left({\bf v}_{a}-{\bf
v}_{j}\right)\cdot\nabla_{a}W_{aj}(h_{a}),$ (50)
which indeed can be shown to be an SPH translation of the first two terms in
Eq. 12.
### 2.4 SPH representation of drag terms
#### 2.4.1 Drag interpolation
The remaining aspect is to provide an SPH representation of Eq. 9, specifying
the drag term $\textbf{F}^{\rm V}_{\mathrm{drag}}$ involved in Eqs. 6–7.
Monaghan & Kocharyan (1995) proposed an SPH interpolation of the drag term
given by
$\left<K\Delta\mathbf{v}\right>=\nu\int
K\left(\mathbf{x},\mathbf{x}^{\prime}\right)\left\\{\left[\textbf{v}_{\mathrm{g}}(\mathbf{x})-\textbf{v}_{\mathrm{d}}(\mathbf{x}^{\prime})\right]\cdot\hat{\mathbf{r}}\right\\}\hat{\mathbf{r}}D\left(\mathbf{x}-\mathbf{x}^{\prime},h\right)\mathrm{d}\mathbf{x}^{\prime},$
(51)
where $\hat{\mathbf{r}}$ is the unit vector defined by:
$\hat{\mathbf{r}}=\frac{\mathbf{x}-\mathbf{x}^{\prime}}{|\mathbf{x}-\mathbf{x}^{\prime}|},$
(52)
and $\nu$ is the number of spatial dimensions of the system (and _not_ the
inverse of the number of spatial dimensions as one might intuitively guess —
see below). Monaghan & Kocharyan (1995) proposed this formulation — with
velocity difference projected along the line of sight joining the particles —
mainly because it gives exact conservation of both linear and angular momentum
in the resulting drag terms.
As the SPH interpolation of the drag term does not come from the Euler-
Lagrange equations derived for non-dissipative term form the SPH Lagrangian,
the kernel function used in the drag term is not constrained to be the same
function $W$ used for the density (as assumed by Monaghan & Kocharyan 1995).
Indeed one of our findings from this paper (discussed below) is that use of a
standard (bell-shaped) density kernel for drag computations can be
significantly inaccurate. We thus use $D$ to denote the kernel employed for
the drag interpolation.
#### 2.4.2 Choice of smoothing length in the drag terms
A key issue is the choice of smoothing length involved in the interpolation
term (51) when the gas and dust have different spatial resolutions, as
illustrated in Fig. 1. We have found from experiment that it is very important
to smooth the drag term using the maximum smoothing length of the two fluids,
rather than using an average (c.f. Sec. 4.6 and also Ayliffe et al. 2011).
Otherwise, unphysical resolution-dependent clumping of one fluid below the
scale of the other can occur. For gas this presents less of a problem because
there remain pressure gradients that prevent such clumping. However, for dust
it is crucial since there are no forces that can otherwise counterbalance any
artificial over-concentration. Since most astrophysical problems involve the
concentration of dust in a flow of gas, a straightforward approach is to
simply use the gas smoothing length when computing the drag interaction.
Unless otherwise specified (Sec. 4.6) this is the approach we adopt in this
paper.
#### 2.4.3 Errors in the integral drag interpolant
The origin of Eq. 51 can be understood by considering the projection of
$\left<\Delta\mathbf{v}\right>$ onto
$\hat{\mathbf{r}}_{\alpha\alpha^{\prime}}$, the projection of
$\hat{\mathbf{r}}$ onto the coordinate $\alpha$ (which equivalently denotes
the coordinates $x$, $y$ or $z$ as the system is invariant by rotation) and
use a Taylor expansion of $K$ and
$\textbf{v}_{\mathrm{d}}(\mathbf{x}^{\prime})$ around their values on
$\mathbf{x}$:
$\displaystyle\left<K\Delta\mathbf{v}\right>^{\alpha}=$
$\displaystyle\nu\int\mathrm{d}\mathbf{x}^{\prime}\hat{\bf r}^{\alpha}$
$\displaystyle\left\\{K\Delta\mathbf{v}(\mathbf{x})+\frac{\partial(K\Delta\mathbf{v})}{\partial\mathbf{x}}\cdot(\mathbf{x}-\mathbf{x}^{\prime})+\mathcal{O}\left((\mathbf{x}-\mathbf{x}^{\prime})^{2}\right)\right\\}$
$\displaystyle\cdot\hat{\mathbf{r}}D\left(\mathbf{x}-\mathbf{x}^{\prime},h\right).$
(53)
giving
$\left<K\Delta\mathbf{v}\right>^{\alpha}=\nu
K\Delta\mathbf{v}(\mathbf{x})^{\beta}I^{\alpha\beta}+\nu\frac{\partial(K\Delta\mathbf{v}^{\alpha})}{\partial\mathbf{x}^{\gamma}}J^{\alpha\beta\gamma}+\mathcal{O}\left(h^{2}\right),$
(54)
where
$\displaystyle I^{\alpha\beta}$ $\displaystyle\equiv$
$\displaystyle\int\mathrm{d}\mathbf{x}^{\prime}\hat{r}^{\alpha}~{}\hat{r}^{\beta}D\left(\mathbf{x}-\mathbf{x}^{\prime},h\right),$
(55) $\displaystyle J^{\alpha\beta\gamma}$ $\displaystyle\equiv$
$\displaystyle\int\mathrm{d}\mathbf{x}^{\prime}\hat{r}^{\alpha}\hat{r}^{\beta}\hat{r}^{\gamma}D\left(\mathbf{x}-\mathbf{x}^{\prime},h\right).$
(56)
This shows that Eq. (51) is a second-order approximation to the drag term,
that is,
$\left<K\Delta\mathbf{v}\right>^{\alpha}=K\Delta\mathbf{v}^{\alpha}+\mathcal{O}\left(h^{2}\right),$
(57)
provided the normalisation conditions
$\displaystyle I^{\alpha\beta}$ $\displaystyle=$
$\displaystyle\frac{\delta^{\alpha\beta}}{\nu},$ (58) $\displaystyle
J^{\alpha\beta\gamma}$ $\displaystyle=$ $\displaystyle 0,$ (59)
hold. Condition (59) and the zeroing of the off-diagonal terms in Eq. 58 may
be proved straightforwardly by the fact that the integrals in (55)–(56) are
odd. The normalisation condition of the diagonal terms in Eq. 58 arises
because in 3D we have
$\displaystyle I^{xx}+I^{yy}+I^{zz}$ $\displaystyle=$ $\displaystyle 1,$ (60)
$\displaystyle I^{xx}=I^{yy}=I^{zz},$ (61)
giving $I^{xx}=I^{yy}=I^{zz}=1/\nu$. This explains the factor of $\nu$ in
front of the drag summation term.
#### 2.4.4 Discretisation of drag term
Discretising Eq. 51 provides the SPH translation of the acceleration due to
the drag term for both the gas and the dust. Replacing the integral by a
summation (over particles of the opposing type) and $\rho{\rm d}V$ with the
particle mass, we have
$\left(\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}\right)_{\mathrm{drag}}=\frac{1}{\hat{\rho}_{\mathrm{g}}}\left<K\Delta\mathbf{v}\right>=\nu\sum_{j}m_{j}\frac{K_{aj}}{\hat{\rho}_{a}\hat{\rho}_{j}}\left({\bf
v}_{aj}\cdot\hat{\textbf{r}}_{aj}\right)\hat{\textbf{r}}_{aj}D_{aj}(h_{a}),$
(62)
for a gas particle and
$\left(\frac{\mathrm{d}{\bf
v}_{i}}{\mathrm{d}t}\right)_{\mathrm{drag}}=\frac{1}{\hat{\rho}_{\mathrm{g}}}\left<K\Delta\mathbf{v}\right>=-\nu\sum_{b}m_{b}\frac{K_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}\left({\bf
v}_{bi}\cdot\hat{\textbf{r}}_{bi}\right)\hat{\textbf{r}}_{bi}D_{ib}(h_{b}),$
(63)
for a dust particle, where we have defined ${\bf v}_{aj}\equiv{\bf v}^{\rm
g}_{a}-{\bf v}^{\rm d}_{j}$ and $\hat{\bf r}_{aj}\equiv({\bf r}_{a}-{\bf
r}_{j})/|{\bf r}_{a}-{\bf r}_{j}|$. Importantly, from Eqs. 62–63, we have:
$\sum_{a}m_{a}\left(\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}\right)_{\mathrm{drag}}+\sum_{i}m_{i}\left(\frac{\mathrm{d}{\bf
v}_{i}}{\mathrm{d}t}\right)_{\mathrm{drag}}=0,$ (64)
which ensures that the momentum is exactly exchanged between the gas and the
dust phase by the SPH formalism. Similarly
$\sum_{a}m_{a}{\bf
r}_{a}\times\left(\frac{\mathrm{d}\textbf{v}_{a}}{\mathrm{d}t}\right)_{\mathrm{drag}}+\sum_{i}m_{i}{\bf
r}_{i}\times\left(\frac{\mathrm{d}{\bf
v}_{i}}{\mathrm{d}t}\right)_{\mathrm{drag}}=0,$ (65)
showing that the total angular momentum is conserved.
#### 2.4.5 Errors in the SPH drag interpolation
A key point to note in the formulation of drag terms is that the criterion for
an accurate kernel drag estimate is different from that required for an
accurate density estimate (Eq. 23). Taking Eq. 62 and expanding the velocities
and the drag coefficient $K$ around the position of the gas particle ${\bf
r}_{a}$, to lowest order, i.e.
$K_{aj}\left({\bf v}^{g}_{a}-{\bf v}^{d}_{j}\right)=K_{a}\left({\bf
v}^{g}_{a}-{\bf v}^{d}_{a}\right)+\mathcal{O}(h),$ (66)
we find
$-\nu\frac{K_{a}}{\hat{\rho}_{a}}\left({\bf v}^{g}_{a}-{\bf
v}^{d}_{a}\right)\cdot\sum_{j}\frac{m_{j}}{\hat{\rho}_{j}}\hat{\textbf{r}}_{aj}\hat{\textbf{r}}_{aj}D_{aj}+\mathcal{O}(h),$
(67)
implying a discrete normalisation condition on the drag kernel of the form
$\nu\sum_{j}\frac{m_{j}}{\hat{\rho}_{j}}\hat{\textbf{r}}_{aj}^{\alpha}\hat{\textbf{r}}_{aj}^{\beta}D_{aj}\approx\delta^{\alpha\beta}.$
(68)
The condition (68) implies that the summation on the diagonal terms ($xx$,
$yy$, $zz$) are equal to unity, while the summations on off-diagonal terms
($xy$, $xz$, $yz$) should be zero. The accuracy with which this normalisation
condition is satisfied depends on the particle arrangement. While we find that
the diagonal terms are well computed using standard (bell-shaped) kernels, we
find that — apart from the special case where the dust particles lie on top of
the gas particles — the off-diagonal terms can be very poorly normalised. Fig.
3 shows the $xy$ component of Eq. 68 as a function of the smoothing length (in
units of the particle spacing, $\Delta x$) computed for a dust particle offset
by $\Delta x/4$ in the x-direction from a cubic lattice of gas particles in
3D. Using the cubic spline kernel (top left) results in errors of order 5-10%
of the diagonal terms for reasonable neighbour numbers ($h/\Delta x\approx
1.1-1.5$). Furthermore, improving the smoothness of the kernel by using the
$M_{6}$ quintic or even the Gaussian (top right) does _not_ significantly
reduce the error. In numerical tests (Sec. 4.2) this manifests as a large
error in the drag between the two fluids, implying that a more suitable kernel
is highly desirable.
Figure 2: Functional form of the standard bell-shaped cubic spline (top left)
and quintic (top right) kernels, compared to the double-hump versions of these
kernels (bottom row). Kernel functions are shown by the solid/black lines,
while the long-dashed/red and short-dashed/green lines correspond to the first
and second derivatives, respectively. We find that double-hump kernels are
significantly more accurate than bell-shaped kernels when computing SPH drag
terms (see Fig. 3). Figure 3: Accuracy with which the normalisation condition
for the drag force is computed using standard bell-shaped kernels (top row)
and double-hump kernels (bottom row). The plots show the $xy$ component of Eq.
68 computed on a dust particle offset from a regular cubic lattice of gas
particles, as a function of the smoothing length in units of the particle
spacing ($h/\Delta x$). With the bell-shaped kernels (top row) the errors are
of order $5-10\%$ (the off-diagonal terms should sum to zero). Changing to
double-hump shaped kernels (bottom row) gives errors $\lesssim 0.5\%$.
#### 2.4.6 Drag kernel function
After conducting a search for suitable alternative kernels, we found that the
so-called “double-hump” shaped kernels (Fulk & Quinn, 1996) gave a substantial
improvement in accuracy — that is, giving errors in the computation of Eq. 68
of similar order to the bell-shaped kernels in computing Eq. 23. Defining the
kernel function as previously
$D\left(r,h\right)=\frac{\sigma}{h^{\nu}}g\left(q\right),$ (69)
we construct double-hump kernels from the $M_{4}$ cubic and $M_{6}$ quintic
kernels using
$g(q)=q^{2}f(q),$ (70)
giving, for example, the “double $M_{4}$ cubic” (bottom left panel of Fig. 2),
the “double $M_{6}$ quintic” (bottom right panel of Fig. 2) and similarly the
“double Gaussian”. The normalisation constants are found in the usual manner
by enforcing $\int D{\rm d}V=1$, i.e.,
$\sigma\int g(q){\rm d}V=1,$ (71)
where ${\rm d}V$ corresponds to ${\rm d}q$, $2\pi q{\rm d}q$ and $4\pi
q^{2}{\rm d}q$ in one, two and three dimensions, respectively. The
normalisation constants for the double cubic, double quintic and double
Gaussian are given by
$\displaystyle\sigma_{\textrm{double M}_{4}}$ $\displaystyle=$
$\displaystyle\left[2,\frac{70}{31\pi},\frac{10}{9\pi}\right];$ (72)
$\displaystyle\sigma_{\textrm{double M}_{6}}$ $\displaystyle=$
$\displaystyle\left[\frac{1}{60},\frac{42}{2771\pi},\frac{1}{168\pi}\right];$
(73) $\displaystyle\sigma_{\textrm{double Gaussian}}$ $\displaystyle=$
$\displaystyle\frac{2}{\nu}\pi^{-\nu/2},$ (74)
in [1,2,3] dimensions. The computation of the off-diagonal term in (68) for
the double-cubic and double-Gaussian kernels are shown in the bottom row of
Fig. 3 and indeed show a substantial improvement, giving errors $\lesssim
0.5\%$ compared to the $5-10\%$ errors obtained using the standard kernels
(top row). This improvement in accuracy is also reflected in our numerical
tests (c.f. Sec. 4.2).
It is also possible to physically understand the reason why the double hump
kernel is suited to deal with drag computation. In treating multi-fluid
interactions, one requires the information of one type of particle at the
location of a particle of the opposing type. Assuming that the number of
dimensions of the space is three, the SPH smoothing of the physical quantity
$A$ corresponds approximately to
$\begin{array}[]{rcl}\displaystyle
4\pi^{2}\\!\\!\int_{0}^{1}\\!\\!\\!A\left(q\right)D(q)\mathrm{d}q&\simeq&\displaystyle
4\pi^{2}\\!\\!\int_{0}^{1}\\!\\!\\!A\left(q\right)\frac{\delta\left(q+q_{\mathrm{M}}\right)+\delta\left(q-q_{\mathrm{M}}\right)}{2}\mathrm{d}q\\\\[13.00005pt]
&&=\displaystyle\frac{A\left(-q_{\mathrm{M}}\right)+A\left(q_{\mathrm{M}}\right)}{2},\end{array}$
(75)
showing that — for a given particle — the double hump kernel provides an
average value of a physical quantity stored in the neighbours of the other
species and located at a distance $q=q_{\mathrm{M}}$ of the particle. For the
same reason, the poor accuracy of the bell-shaped kernels can be understood
because the maximum weight corresponds to $q=0$, where in general, no particle
of the other type is present.
#### 2.4.7 Frictional heating terms due to drag
When the system is made of a single gas fluid, the _specific_ thermal energy
$u$ is a function of state whose total derivative is expressed by:
$\rm{d}u=T\rm{d}s+\frac{P}{\rho^{2}}\rm{d}\rho.$ (76)
To generalise this relation with two-fluids interacting with a drag term, we
derive an additional term for Eq. 76 arising from exchange of momentum (all
the other quantities fixed) considering a closed thermally isolated system
made of gas and dust SPH particles, whose energy exchange arises only because
of momentum exchange (i.e. drag) between two states (denoted _i_ and _f_ ,
respectively). Applying the first law of thermodynamics to an infinitesimal
transformation of the system, we have:
$\rm{d}u+\rm{d}e_{\rm{k}}=\delta w_{\rm{i}\to\rm{f}}+\delta
q_{\rm{i}\to\rm{f}},$ (77)
where $u$ is the total specific internal energy, $e_{\rm{k}}$ is the
macroscopic kinetic energy of the system and $w$ is the total work and $q$ is
the total heat exchanged during the transformation. Assuming that the
transformation occurs slowly enough for the gas to remain in thermodynamic
equilibrium, Eq. 77 reduces to:
$\rm{d}u|_{s,\rho}+\rm{d}e_{\rm{k}}=T\rm{d}s+\frac{P}{\rho^{2}}\rm{d}\rho=0.$
(78)
Consequently,
$\displaystyle\rm{d}u|_{s,\rho}=-\rm{d}e_{\rm{k}}=$
$\displaystyle-\sum_{a}\frac{\left(\bf{v}_{a}+\left.\rm{d}\bf{v}_{a}\right|_{s_{a},\rho_{a}}\right)^{2}}{2}+\sum_{a}\frac{\bf{v}_{a}^{2}}{2}$
(79)
$\displaystyle-\sum_{k}\frac{\left(\bf{v}_{k}+\left.\rm{d}\bf{v}_{k}\right|_{s_{k},\rho_{k}}\right)^{2}}{2}+\sum_{k}\frac{\bf{v}_{k}^{2}}{2}.$
As the conservation of the momentum during the transformation ensures that:
$\sum_{k}\left.\rm{d}\bf{v}_{k}\right|_{s_{k},\rho_{k}}=-\sum_{a}\left.\rm{d}\bf{v}_{a}\right|_{s_{a},\rho_{a}},$
(80)
we obtain:
$\rm{d}u|_{s,\rho}=\sum_{a}\bf{v}_{ka}\cdot\left.\rm{d}\bf{v}_{a}\right|_{s_{a},\rho_{a}}.$
(81)
Using the expression Eq. 62 which gives the evolution of the velocity of a gas
particle due to drag term (i.e. at constant specific entropy and density):
$\frac{\rm{d}u|_{s,\rho}}{\rm{d}t}=\sum_{a}\frac{\rm{d}u_{a}}{\rm{d}t}=\sum_{a}\mathbf{v}_{ka}\cdot\left[\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\textbf{v}_{ak}\cdot\hat{\textbf{r}}_{ak}\right)\hat{\textbf{r}}_{ak}D_{ak}(h_{a})\right].$
(82)
which implies that the evolution of the specific internal energy for each gas
particle is given by:
$\left(\frac{\rm{d}u_{a}}{\rm{d}t}\right)_{\rm drag}=\frac{\Lambda_{\rm
drag}}{\hat{\rho}_{a}}=\nu\sum_{k}m_{k}\frac{K_{ak}}{\hat{\rho}_{a}\hat{\rho}_{k}}\left(\textbf{v}_{ak}\cdot\hat{\textbf{r}}_{ak}\right)^{2}D_{ak}(h_{a}).$
(83)
The positive value of ${\rm d}u_{a}/{\rm{d}t}$ (and the fact that the kernel
function is always positive) ensures a positive definite contribution to the
specific internal energy of each SPH gas particle from frictional drag
heating. Eq. 83 provides the SPH translation of the drag heating term (Eq.
13). It is straightforward to show that with the above expression the total
energy is exactly conserved, i.e.
$\sum_{a}m_{a}\frac{{\rm d}u_{a}}{{\rm d}t}+\sum_{a}m_{a}{\bf
v}_{a}\cdot\frac{{\rm d}{\bf v}_{a}}{{\rm d}t}+\sum_{i}m_{i}{\bf
v}_{i}\cdot\frac{{\rm d}{\bf v}_{i}}{{\rm d}t}=0.$ (84)
#### 2.4.8 Thermal coupling terms
The thermal coupling terms can be expressed in SPH form using (Monaghan &
Kocharyan, 1995)
$\Lambda_{{\rm
therm},a}=\sum_{j}m_{a}\frac{Q_{aj}}{\hat{\rho}_{a}\hat{\rho}_{j}}(T_{a}-T_{j})W_{aj}+\sum_{j}m_{a}\frac{R_{aj}}{\hat{\rho}_{a}\hat{\rho}_{j}}a(T_{a}^{4}-T_{j}^{4})W_{aj},$
(85)
for a gas particle, and
$\Lambda_{{\rm
therm},i}=-\sum_{b}m_{i}\frac{Q_{bi}}{\hat{\rho}_{i}\hat{\rho}_{b}}(T_{b}-T_{j})W_{bi}-\sum_{b}m_{b}\frac{R_{bi}}{\hat{\rho}_{b}\hat{\rho}_{i}}a(T_{b}^{4}-T_{i}^{4})W_{bi},$
(86)
for a dust particle. Note that we use the standard SPH kernel for the thermal
coupling terms. Detailed study of the effect of the thermal coupling terms, or
specific expressions for $Q$ and $R$, are beyond the scope of this paper, and
thus for the tests in Sec. 4 we simply set $T_{\rm d}=T_{\rm g}$.
#### 2.4.9 Drag coefficient
In general the drag coefficient $K$ is a function of the properties of both
the gas and dust. For example in the linear Epstein regime relevant to dilute
gases in the limit of low Mach numbers, the coefficient $K_{ak}$ is given by
$K_{ak}=\frac{4}{3}\pi\sqrt{\frac{8}{\pi\gamma}}\frac{\hat{\rho}_{k}}{m_{\rm{d}}}\frac{\hat{\rho}_{a}}{\theta_{a}}s^{2}c_{\mathrm{s},a},$
(87)
where $s$ is the grain radius, $m_{\rm{d}}$ is the grain mass, $\gamma$ is the
adiabatic index, $c_{\mathrm{s},a}$ is the gas sound speed. Since the basic
dust-gas algorithm described above is insensitive to the specific form of the
drag, we consider only constant drag coefficients in this paper in order to
benchmark the method. The detailed implementation of a full range of both
linear and non-linear physical drag formulations is considered in Paper II.
## 3 Timestepping
### 3.1 Empirical timestep criterion
The drag terms impose an additional constraint on the timestep $\Delta t$,
such that it has to be smaller than a critical value $\Delta t_{\rm{c}}$ for
an explicit scheme (e.g. Leapfrog) to remain stable. Empirically, Monaghan &
Kocharyan (1995) use the criterion:
$\Delta t<\min\left(\frac{\rho}{K}\right),$ (88)
which is essentially the minimum of the drag stopping time taken over all of
the SPH particles.
### 3.2 Von Neumann stability analysis
A more precise criterion can be derived by considering the stability of a
simple explicit scheme such as the Forward Euler method. We consider the
evolution of the drag terms over a single drag timetstep $\Delta t$,
calculating the velocities at the timestep $n+1$ from the velocities at the
timestep $n$. Considering only the time-discretisation of the equations, we
have
$\displaystyle\frac{{\bf v}^{n+1}_{\rm g}-\mathbf{v}^{n}_{\rm g}}{\Delta t}$
$\displaystyle=$ $\displaystyle-\frac{K}{\hat{\rho}_{\mathrm{g}}}\left({\bf
v}_{\rm g}-{\bf v}_{\rm d}\right),$ (89) $\displaystyle\frac{{\bf
v}^{n+1}_{\rm d}-\mathbf{v}^{n}_{\rm d}}{\Delta t}$ $\displaystyle=$
$\displaystyle+\frac{K}{\hat{\rho}_{\mathrm{d}}}\left({\bf v}_{\rm g}-{\bf
v}_{\rm d}\right),$ (90)
We then perform a standard Von Neumann analysis, considering a perturbation of
the velocity field with respect to equilibrium at the timestep $m$
corresponding to a monochromatic plane wave, i.e.
$\displaystyle{\bf v}^{m}_{\rm g}$ $\displaystyle=$ $\displaystyle{\bf
V}^{m}_{\rm{g}}e^{ikx},$ (91) $\displaystyle{\bf v}^{m}_{\rm d}$
$\displaystyle=$ $\displaystyle{\bf V}^{m}_{\rm{d}}e^{ikx},$ (92)
where ${\bf v}^{m}_{\rm{g}}$ and ${\bf v}^{m}_{\rm{d}}$ are complex constants
and $k$ is the wavenumber. Substituting Eqs. 91 and 92 into Eqs. 89 and 90
leads to the linear system
$\left(\begin{array}[]{c}{\bf V}_{\rm{g}}\\\ {\bf
V}_{\rm{d}}\end{array}\right)^{n+1}=\left(\begin{array}[]{cc}1-\Delta
t~{}\frac{K}{\hat{\rho}_{\mathrm{g}}}&\Delta
t~{}\frac{K}{\hat{\rho}_{\mathrm{g}}}\\\ \Delta
t~{}\frac{K}{\hat{\rho}_{\mathrm{d}}}&1-\Delta
t~{}\frac{K}{\hat{\rho}_{\mathrm{d}}}\end{array}\right)\left(\begin{array}[]{c}{\bf
V}_{\rm{g}}\\\ {\bf V}_{\rm{d}}\end{array}\right)^{n}.$ (93)
The two complex eigenvalues $\Lambda_{\pm,aj}$ of the matrix $\mathcal{M}$ are
given by:
$\Lambda_{\pm,ai}=1-\frac{\Delta
t}{2}\left(\frac{K}{\hat{\rho}_{\mathrm{g}}}+\frac{K}{\hat{\rho}_{\mathrm{d}}}\right)\pm\frac{\Delta
t}{2}\left(\frac{K}{\hat{\rho}_{\mathrm{g}}}+\frac{K}{\hat{\rho}_{\mathrm{d}}}\right).$
(94)
The condition for the numerical scheme to remain stable ($|\Lambda_{-}|<1$)
implies a minimum timestep given by
$\Delta t<\Delta
t_{\rm{c}}=\frac{\hat{\rho}_{\mathrm{g}}\hat{\rho}_{\mathrm{d}}}{K(\hat{\rho}_{\mathrm{g}}+\hat{\rho}_{\mathrm{d}})}(\equiv
t_{\rm s}).$ (95)
We note that this expression differs slightly to the one suggested by Monaghan
& Kocharyan (1995) (Eq. 88) as it involves the _physical_ drag stopping time
$t_{\rm s}$ — i.e. the typical time to damp the differential velocity between
the gas and the dust fluids — rather than $\frac{K}{\rho}$. The timestep of
Monaghan & Kocharyan (1995) is thus correct in the limit where the density of
one phase is negligible compared to the density of the other phase, but
becomes erroneous in the case of two fluids having densities of the same order
of magnitude. Note this would apply to grid codes also.
### 3.3 SPH explicit timestep
The stability criterion for the full SPH system (and also for other explicit
schemes) is expected to be similar to that derived for the continuum case (Eq.
95). The main difference is that the drag coefficient $K$ is in general only
defined on particle _pairs_ rather than individual particles. We thus take the
minimum of Eq. 95 over a particle’s neighbours, i.e.
$\Delta
t_{\rm{c},a}=\min_{k}\left[\frac{\hat{\rho}_{a}\hat{\rho}_{k}}{K_{ak}(\hat{\rho}_{a}+\hat{\rho}_{k})}\right];\hskip
14.22636pt\Delta
t_{\rm{c},i}=\min_{b}\left[\frac{\hat{\rho}_{b}\hat{\rho}_{i}}{K_{bi}(\hat{\rho}_{b}+\hat{\rho}_{i})}\right];$
(96)
for gas and dust particles, respectively.
### 3.4 Implicit timestepping
For strong drag regimes, the timestep restriction imposed by Eq. (96) becomes
prohibitive, and an implicit timestepping algorithm is required, as proposed
by Monaghan (1997). We use only explicit timestepping for the tests shown in
this paper, with implicit timestepping methods discussed in detail in Paper
II.
## 4 Numerical tests
Despite a number of codes having already been developed for simulating
astrophysical gas-dust mixtures, none have been benchmarked against a wide
range of test problems relevant to astrophysics. For example, while Monaghan &
Kocharyan (1995), Maddison et al. (2003) consider drag on a single dust
particle in a box of gas (similar to our dustybox test below), no waves or
shocks are considered. Similarly Paardekooper & Mellema (2006) benchmark their
algorithm against a single dust-gas shock problem with only a qualitative
solution. Other authors simply check that the timescale for settling in an
accretion disc is roughly consistent (Barrière-Fouchet et al., 2005) or
provide no tests at all (Rice et al., 2004; Fromang & Nelson, 2005; Fromang &
Papaloizou, 2006). In the absence of known analytic solutions for simple
problems, Johansen et al. (2007), Miniati (2010) and Bai & Stone (2010) use
the linear growth rates for the streaming instability (Youdin & Goodman, 2005)
as a test problem, though this is already a complicated problem.
In this paper we present a comprehensive suite of test problems designed to
investigate all aspects of our algorithm relevant to astrophysics. These we
refer to as dustybox, dustywave, dustyshock, dustysedov and dustydisc.
Analytic solutions for the dustybox and dustywave problems have been derived
in Laibe & Price (2011a), while the solution for dustyshock is known in the
limit of high drag. The solutions for dustysedov and dustydisc are more
qualitative but are important reference problems for astrophysical dust-gas
mixtures. We consider simulation of the streaming instability to be of
sufficient importance and complexity to be covered in detail in a separate
paper. Note that all of the tests considered in this paper are performed, for
simplicity, using a constant drag coefficient $K$. Realistic drag regimes are
considered in Paper II.
Figure 4: Dust velocity as a function of time in the dustybox problem, using
$2\times 20^{3}$ particles, a dust-to-gas ratio of unity and a constant drag
coefficient $K=1$. Use of the standard cubic spline kernel for the drag terms
(short dashed/green line) results in errors of order $10\%$ in the velocities
compared to the exact solution (long dashed/red line). Using a double-hump
kernel (solid black) improves the accuracy — for the same computational cost —
by a factor of several hundred, to $\lesssim 0.1\%$.
### 4.1 Code implementation
Implementation of the algorithm into a standard SPH code with explicit
timestepping is relatively straightforward. The main changes are i) to store a
particle type allowing setup and simulation of multiple fluids; ii) to compute
the densities and smoothing lengths on each fluid as described in Sec. 2.2.1;
iii) compute the drag term between each fluid according to Eqs. 62–63 and the
heating term given by Eq. 83 and iv) (optional for astrophysics) to implement
the modifications to the equations of motion due to the volume-filling
fraction of the dust. We have implemented the two-fluid algorithm into both
the $N-$dimensional ndspmhd test code (Price, 2011) and into the parallel
phantom code for 3D problems (Price & Federrath, 2010; Lodato & Price, 2010).
### 4.2 dustybox: Two fluid drag in a periodic box
The dustybox problem presented by Laibe & Price (2011a) involves two fluids in
a periodic box moving with a differential velocity ($\Delta
v_{0}=v_{d,0}-v_{g,0}$). It is similar to the test performed by Monaghan &
Kocharyan (1995) showing the drag on a single dust grain in a box of gas,
except that here we consider the dust as a fluid, meaning that the densities
and smoothing lengths of both phases are computed self-consistently.
#### 4.2.1 dustybox: Setup
We setup the particles in a 3D periodic domain $x,y,z\in[0,1]$ such that the
densities $\hat{\rho}_{\mathrm{g}}$ and $\hat{\rho}_{\mathrm{d}}$ and the gas
pressure $P_{\mathrm{g}}$ are constant, and neglect the dust intrinsic volume
by fixing the volume fraction $\theta=1$. The box is filled by $20^{3}$ SPH
gas particles set up on a regular cubic lattice and $20^{3}$ dust particles
set up on a cubic lattice shifted by half of the lattice step in each
direction. The gas sound speed, the gas and the dust densities are set to
unity in code units and no artificial viscosity terms are applied. We give the
fluids initial velocities $v_{\rm d}=1$ and $v_{\rm g}=0$.
During the simulation, we verified that both the total linear and angular
momentum are exactly conserved as expected (Eqs. 64–65). We have also verified
that 1) the offset of the dust lattice with respect to the gas lattice and 2)
the timestepping scheme do not affect the results.
#### 4.2.2 dustybox: Choice of drag kernel
Fig. 4 shows the dust velocity as a function of time in the dustybox test
using $\hat{\rho}_{\mathrm{g}}=\hat{\rho}_{\mathrm{d}}=1$ (i.e., a dust to gas
ratio of unity) and $K=1$, with the exact solution from Laibe & Price (2011a)
shown by the long-dashed/red line. Using the cubic spline $M_{4}$ kernel
(short-dashed/green line), the errors are of order 10%. Since these errors are
due to intrinsic bias in the kernel interpolation of the drag terms (Fig. 3),
they are independent of resolution, though can be improved – at considerable
cost – by increasing the ratio of smoothing length to particle spacing (i.e.,
the neighbour number). By comparison, use of the double-hump cubic spline
kernel gives errors $\lesssim 0.1\%$ (solid/black line) with no additional
overhead in terms of cost.
#### 4.2.3 dustybox: Effect of drag coefficient and dust-to-gas ratio
Fig. 5 is identical to Fig. 4 but for a range of drag coefficients
$K=0.01,0.1,1,10,100$, compared to the exact solution in each case given by a
solid/black line. Irrespective of the value of $K$, both gas and dust
velocities relax to the barycentric velocity
($\textbf{v}_{\mathrm{g}}=\textbf{v}_{\mathrm{d}}=0.5$) in a few stopping
times
$t_{\rm{s}}=(\hat{\rho}_{\mathrm{g}}\hat{\rho}_{\mathrm{d}})/[K(\hat{\rho}_{\mathrm{g}}+\hat{\rho}_{\mathrm{d}})]$.
Using the double-hump cubic, an accuracy between 0.1 and 1% is achieved in all
cases (long dashed/red lines).
Fig. 6 is similar, but varying the dust-to-gas ratio using
$\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{d}}=0.01,0.1,1,10,100$ (achieved
by varying $\hat{\rho}_{\mathrm{d}}$ with $\hat{\rho}_{\mathrm{g}}=1$) and
using $K=1$. This changes both the drag stopping time and the barycentric
velocity towards which the system relaxes. Here again, an accuracy between 0.1
and 1% is achieved in all cases.
Figure 5: As in Fig. 4 but using only the double hump cubic kernel with a
range of drag coefficients $K=0.01,0.1,1,10,100$ (top-to-bottom, solid/black
lines), compared with the exact solution in each case given by the long-
dashed/red lines. Figure 6: As in Figs. 4 and 5 but varying the dust-to-gas
ratio $\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{d}}=0.01,0.1,1,10,100$
(top-to-bottom, solid/black lines) and a fixed drag coefficient $K=1$ using
the double hump cubic kernel. Exact solutions for each case are given by the
long-dashed/red lines.
### 4.3 dustywave: Sound waves in a dust-gas mixture
The exact solution for linear waves propagating in a dust-gas mixture
(dustywave) has been presented by Laibe & Price (2011a). We have performed a
series of tests involving the propagation of a sound wave along the $x-$axis
in both one and three dimensions in a periodic box, adopting the setup
described in Table 2 of Laibe & Price (2011a). The dustywave problem is more
complex than the dustybox problem as the motion of the mixture is driven by
both the drag and the gas pressure.
Specifically, Laibe & Price (2011a) derive the dispersion relation:
$\omega^{3}+iK\left(\frac{1}{\hat{\rho}_{\mathrm{g}}}+\frac{1}{\hat{\rho}_{\mathrm{d}}}\right)\omega^{2}-k^{2}c_{\rm
s}^{2}\omega-iK\frac{k^{2}c_{\rm s}^{2}}{\hat{\rho}_{\mathrm{d}}}=0,$ (97)
for solutions in the form $e^{i\left(kx-\omega t\right)}$. At high drag, Eq.
97 can be expanded in a Taylor series, which to first order gives:
$\omega=\pm k\tilde{c}_{\rm
s}-i\frac{\hat{\rho}_{\mathrm{g}}\hat{\rho}_{\mathrm{d}}}{K\left(\hat{\rho}_{\mathrm{g}}+\hat{\rho}_{\mathrm{d}}\right)}k^{2}c_{\mathrm{s}}^{2}\left(\frac{1-A^{2}}{2}\right)$
(98)
where the effective sound speed is defined according to
$\tilde{c}_{\rm s}\equiv c_{\rm s}A=c_{\rm
s}\left(1+\frac{\hat{\rho}_{\mathrm{d}}}{\hat{\rho}_{\mathrm{g}}}\right)^{-\frac{1}{2}}.$
(99)
The first term of Eq. 98 gives the propagation of the centre of mass of the
mixture at the effective sound speed $\tilde{c}_{\rm s}$. The second term
corresponds to a corrective dissipative term since $A\in\left[0,1\right]$.
#### 4.3.1 dustywave: Setup
The equilibrium state is characterised by the two phases at rest where the gas
sound speed, and both gas and dust densities are set to unity in code units.
In 1D this is achieved by placing equally spaced particles in the periodic
domain $x\in[0,1]$. For the 3D simulations, the tests are run in a periodic
box $x,y,z\in[0,1]$ with gas particles set up on a regular cubic lattice and
dust particles set up on a cubic lattice shifted by half of the lattice step
in each directions. As previously, no artificial viscosity is applied. We set
the relative amplitude of the perturbation to $10^{-4}$ in both velocity and
density in order to remain in the linear acoustic regime for which the
solution in Laibe & Price (2011a) is derived (we have verified that running
the same simulations setting the relative amplitudes to $10^{-8}$ gives the
same results). The density perturbation is applied to the particles as
described in Appendix B of Price & Monaghan (2004). We adopt an isothermal
equation of state $P=c_{\rm s}^{2}\rho$ with $c_{\rm s}=1$.
Figure 7: Results of the dustywave test in 3D at $t=0$ (top row) and after 1
and 2 wave periods (middle and bottom rows) using $2\times 32^{3}$ particles,
$K=1$ and a dust-to-gas ratio of unity. The analytic solution is given by the
solid/red (gas) and long-dashed/red (dust) lines. The standard cubic spline
kernel (green points, left panel) performs poorly on this test. Using the
double hump cubic kernel (black points, left panel) both the amplitude and
frequency are correct but there remains a small phase error due to the kernel
bias which can be corrected by using the smoother double-hump quintic kernel
(black points, right panel).
#### 4.3.2 dustywave: Effect of the smoothing kernel
The results of the dustywave test in 3D using $2\times 32^{3}$ particles with
$K=1$ and a dust-to-gas ratio of unity is shown in Fig. 7, using the standard
bell-shaped cubic (green points, lower amplitude) and double-hump cubic kernel
(black points, correct amplitudes) for the drag (left panel) and the double-
hump quintic kernel (c.f. Sec. 2.4.6) (right panel), at $t=0$ (top) and after
1 and 2 periods (middle and bottom panels). The numerical solutions (green and
black markers) may be compared to the analytic solutions given by the
solid/red (gas) and long-dashed/red (dust) lines. The amplitude and frequency
of the solution are only correctly captured using double-hump kernels. With
the double-hump cubic employed (left panel, black points) there remains a
slight (few %) phase error in the numerical solution caused by the remaining
kernel bias, independent of resolution. A similar error is found generically
in multidimensional SPH simulations of linear waves (see e.g. Fig. 6 of Price
& Monaghan 2005) and in standard SPH is improved by using a smoother kernel
such as the quintic spline. Indeed, using the double-hump version of the
quintic kernel for the drag terms and the standard quintic for the SPH terms
(right panel) we find the phase error is smaller by a factor of $\sim 5$.
Longer multidimensional simulations of linear waves are more complicated in
SPH because placement of the gas particles on regular lattices are unstable to
low-amplitude transverse modes that cause the particles to rearrange towards a
“glass-like” configuration (Morris, 1996a, b). Furthermore, we find that with
large drag coefficients we require extremely high resolutions to match the
analytic solutions (see below), which becomes prohibitive in 3D. In one
dimension however, the numerical stability of a sound wave is achieved simply
by satisfying the courant condition (i.e. $\Delta t\leq 0.3c_{\rm{s}}/h$) and
the timestep constraint from the drag (Eq. 95). We thus turn to 1D to
investigate the full parameter range of the dustywave solution.
Figure 8: Resolution study for the dustywave test in 1D using a high drag
coefficient ($K=100$) and a dust-to-gas ratio of unity using 32, 64, 128, 256,
512 and 1024 particles from bottom to top. At large drag high resolution is
required to resolve the small differential motions between the fluids and thus
prevent over-damping of the numerical solution, corresponding to the criterion
$h\lesssim c_{\rm s}t_{\rm s}$, here implying $\gtrsim 240$ particles. See
also Fig. 9. Figure 9: As in Fig. 8 but showing the kinetic energy as a
function of time in the numerical solution at progressively increasing
resolution, compared to the analytic solution given by the solid black line.
The kinetic energy decay converges to the analytic solution at $\sim 256-512$
particles per wavelength, implying a demanding resolution criterion
($h\lesssim c_{\rm s}t_{\rm s}$) for high drag.
Figure 10: Parameter study of the 1D dustywave problem, varying the drag
coefficient from $0.01$ to $100$ (top to bottom) and using two different dust-
to-gas ratios ($\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{g}}=1$, left
figure and $\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{g}}=0.01$, right
figure). Results are shown after 5 wave periods using 128 particles except for
the $K=100$ case where $512$ particles are used to satisfy the resolution
criterion $h\lesssim c_{\rm s}t_{\rm s}$. The double hump cubic kernel is used
for the drag terms. The results obtained at these resolutions are
indistinguishable from the analytic solutions (red solid [gas] and dashed
[dust] lines) for both the gas (dots) and the dust (circle) particles.
#### 4.3.3 dustywave: Resolution requirements at high drag
Fig. 8 shows the velocity profiles after 10 periods in the 1D dustywave
problem for a large drag coefficient ($K=100$) and a dust-to-gas ratio of
unity, with numerical resolution as indicated. At low resolutions ($\lesssim
256$ particles per wavelength) and high drag, the amplitude of the wave in the
numerical simulations (black solid [gas] and open [dust] circles, on top of
each other) is severely overdamped compared to the analytic solution (red
solid and long-dashed lines, also on top of each other). This is further
illustrated in Fig. 9 which shows the kinetic energy as a function of time for
simulations at different resolutions, compared to the analytic solution given
by the solid red line.
Figs. 8 and 9 illustrate a key difficulty that arises when considering high
drag coefficients, i.e., where the drag stopping time $t_{\rm{s}}$ defined in
Eq. 95 is much smaller than the period $T$ of the wave. In this case, the drag
term efficiently damps the initial differential velocity between the gas and
the dust in a few $t_{\rm{s}}$. However, as the pressure continues to drive
the propagation of the wave in the gas, a small residual de-phasing of order
$\sim c_{\rm s}t_{\rm s}$ occurs, which is simply the distance travelled by
the gas before it is damped by the dust. This de-phasing induces a small
differential velocity which in turn be damped by the drag. This small
differential effect dissipates the kinetic energy on a timescale $\sim
t_{\rm{s}}$.
The spatial de-phasing between the gas and the dust represents the smallest
length of the problem that must be resolved numerically in order to capture
the physics of the process. If the spatial de-phasing between the gas and the
dust is under-resolved, the differential velocity between the gas and the dust
is artificially larger than the theoretical one, leading to a non-physical
over-dissipation of the kinetic energy of the system, as observed in Figs. 8
and 9.
We thus propose a resolution criterion for resolving the differential drag of
the form
$\Delta\lesssim c_{\rm s}t_{\rm s},$ (100)
where $\Delta$ is the resolution length. For SPH, this becomes
$h\lesssim c_{\rm s}t_{\rm s}.$ (101)
For $K=100$ and $c_{\rm s}=\hat{\rho}_{\mathrm{g}}=\hat{\rho}_{\mathrm{d}}=1$
in code units this implies $h<0.02$, i.e. a minimum of $\sim 240$ particles
(assuming $\eta=1.2$ in Eqs. 16 and 17), which is consistent with Figs. 8 and
9.
Simulating dust-gas interactions at high drag therefore requires a high
spatial resolution in order to accurately resolve the propagation without
over-dissipating the energy of the system. This can lead to a prohibitive
computational cost, somewhat counterintuitively since the drag simply tends to
make the dust stick to the gas. Most importantly, this requirement is not
unique to SPH and is a critical issue for any numerical method. Indeed, Bai &
Stone (2010) find similarly high resolution requirements at short stopping
times in their simulations of the streaming instability.
#### 4.3.4 dustywave: Parameter study
Fig. 10 shows the results of 1D simulations of the dustywave problem for 5
drag coefficients (from $K=10^{-2}$ to $K=10^{2}$) and two different dust to
gas ratio relevant for astrophysical systems ($1$ and $0.01$ for left and
right figures, respectively), showing the velocities after five periods
compared with the analytic solutions in each case. The simulations employ 128
particles except for the $K=100$ case where $512$ particles have been used in
order to satisfy the criterion (101). For this set of parameters, our method
provides results with an excellent accuracy (better than one per cent) on the
frequencies, the amplitudes and the phases of both the gas and the dust
velocities (and consequently the energy of the system).
For equal dust to gas ratios
($\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{g}}=1$, left figure), both
phases are equally affected by the drag. At low drag ($K=0.01$, top panel of
left figure), the damping is not efficient enough for gas or the dust to be
damped as the stopping time is $\sim$ hundreds of periods. At intermediate
drag ($K=1$, middle panel of left figure), the damping is the most efficient
for the two phases. At large drag regimes ($K=100$), the damping of the
differential velocity occurs quickly, but the dust density is large enough to
distort the gas propagation: the wave is de-phased by a half-period compared
to the gas-only solution.
With more typical astrophysical dust to gas ratios
($\hat{\rho}_{\mathrm{d}}/\hat{\rho}_{\mathrm{g}}=0.01$, right figure), the
gas remains essentially unaffected by the dust. It thus propagates almost
freely in the box at a velocity close to the sound speed. By contrast, the
dust phase is strongly affected by the drag as shown by the $K=0.01$ case (top
panel, right figure), where the damping time for the dust phase is $\sim$ one
period. The differential velocity between the two phases becomes more and more
efficiently damped as the drag coefficient increases (right panel, from top to
bottom), making the dust phase stick to the gas.
### 4.4 dustyshock: shock tube in a dust-gas mixture
Propagation of a shock in a two-fluid dust and gas mixture (the dustyshock
problem hereafter) has been studied both analytically (see e.g. Rudinger 1964)
and numerically (see e.g. Miura & Glass 1982; Saito et al. 2003), using grid
based methods. The dustyshock occurs in two stages: a transient stage (for
which no analytic solution is known and therefore studied numerically)
followed by a stationary stage which consists of the solution for a pure gas
solution propagating at a modified $\gamma$ and the modified sound speed (Eq.
99, see also Miura & Glass 1982). In an astrophysical context, simulations of
a dusty shock were used by Paardekooper & Mellema (2006) to test their
Godunov-type scheme using a Roe Solver developed to simulate astrophysical
dust and gas mixtures.
The hypothesis for the dust phase in these seminal studies are essentially the
same as the ones used in this paper. However, unnecessary additional
complications arise from their choice of the Stokes drag regime (a function of
the local Reynolds number for the particles), the addition of a heat transfer
term (depending on the dust conductivity and the Nusselt number of the system)
and a temperature-dependent gas viscosity. For the purposes of benchmarking of
our numerical scheme, we instead simulate a simplified problem: using a linear
drag regime with constant drag term $K$, no heat transfer between the phases
and no viscosity other than the standard shock-capturing terms used in SPH.
While the evolution during the transient stage may be different from those
considered in previous studies, the solution during the stationary stage
remains unchanged.
#### 4.4.1 dustyshock: setup
We setup the dustyshock problem as a two fluid version of the standard Sod
(1978) problem. Equal mass particles are placed in the 1D domain
$x\in[-0.5,0.5]$, where for $x<0$ we use
$\rho_{\mathrm{g}}=\rho_{\mathrm{d}}=1$, $v_{\rm{g}}=v_{\rm{d}}=0$ and
$P_{\rm{g}}=1$, while for $x>0$ $\rho_{\mathrm{g}}=\rho_{\mathrm{d}}=0.125$,
$v_{\rm{g}}=v_{\rm{d}}=0$ and $P_{\rm{g}}=0.1$. We use an ideal gas equation
of state $P=(\gamma-1)\rho u$ with $\gamma=5/3$. The density jump means that
for SPH the resolution is 8 times higher to the left of the shock than to the
right. We adopt the same initial resolution in both the gas and the dust. This
differs slightly from the setup used by Miura & Glass (1982) and Saito et al.
(2003) where the dust is only placed in the right half of the box. Standard
artificial viscosity and conductivity terms are employed for shock-capturing
in SPH as described in Price (2011) with constant coefficients $\alpha_{\rm
SPH}=1$, $\beta_{\rm SPH}=2$ and $\alpha_{\rm u}=1$.
#### 4.4.2 dustyshock: transient evolution
Figure 11: Results of the dustyshock problem with a moderate drag coefficient
$K=1$ and a dust-to-gas ratio of unity. This shows the dustyshock solution
during the transient stage where the analytic solution is not known (the
solution for the later stationary stage is shown by the dotted red line for
comparison). Results are similar to those obtained in previous studies (Miura
& Glass, 1982; Saito et al., 2003). Top panels show velocity and density in
both gas (solid points) and dust (open circles), while bottom panels show
thermal energy and pressure in the gas. Initial particle spacing to the left
of the shock in both fluids is $\Delta x=0.001$ while to the right it is
$\Delta x=0.008$, giving $569$ equal mass particles in each phase.
Figure 12: Results of the dustyshock problem with a high drag coefficient
($K=1000$) and a dust-to-gas ratio of unity, thus being in the stationary
phase where the analytic solution is known (solid/red lines). At low
resolution (left figure, same resolution as Fig. 11) the results are incorrect
due to the failure to satisfy the resolution criterion at high drag (Eq. 101).
With this criterion satisfied (right panel, using $2\times 11255$ particles)
the numerical solution faithfully reproduces the analytic result. Thus, as in
the dustywave test, extremely high resolution is required to obtain the
correct solution at high drag.
Fig. 11 shows the results of a simulation using $2\times 569$ particles (i.e.
a particle spacing of $\Delta x=0.001$ for $x<0$) with a moderate drag
coefficient ($K=1$) and a dust-to-gas ratio of unity, showing velocity and
density for both the gas and dust (top panels) and the thermal energy and
pressure in the gas (lower panels). With this choice of drag coefficient the
system remains in the transient regime at the time shown (since $t<t_{\rm
s}$). It should be noted that while there is no known analytic solution for
this stage of the problem, the shock profile we obtain is similar to those
found previously (see e.g. Miura & Glass 1982; Saito et al. 2003). Initially
as the shock propagates in the mixture, the dust (initially at rest)
dissipates the momentum and kinetic energy from the gas, lowering the
propagation velocity compared to the ideal (gas only) case (dotted red line).
The dust density ramps up roughly linearly behind the shock, reaching a
density near the contact discontinuity roughly twice the unshocked dust
density. Saito et al. (2003); Miura & Glass (1982) and Paardekooper & Mellema
(2006) also found a similar behaviour, also with a factor of 2 increase in the
dust density behind the shock. The gas-dust interaction to the left of the
contact discontinuity and in the rarefaction wave has not been previously
studied since the above authors place the dust only downwards of the shock
front. We find that the dust density decreases to near zero upstream of the
contact discontinuity, increasing sharply at the head of the rarefaction wave,
transitioning smoothly through the rarefaction wave to match the undisturbed
value. We have checked that increasing the resolution further does not change
the solution significantly for this choice of drag parameters.
#### 4.4.3 dustyshock: stationary regime
Fig. 12 shows the results of simulations with a high drag coefficient
($K=1000$) and a dust-to-gas ratio of unity. In this case, since $t>t_{\rm
s}$, the mixture quickly reaches the stationary regime. We are thus able to
compare the SPH results to the analytic solution given by the solid red line
(this corresponds to the standard hydrodynamic shock solution with modified
sound speeds given by Eq. 99). The left figure shows the results at low
resolution ($2\times 569$ particles, as in Fig. 11, while the right figure
shows the results at $20\times$ higher resolution ($2\times 11255$ particles).
As in the dustywave test, we find that at high drag an extremely high
resolution is required to obtain the correct solution, consistent with our
resolution criterion derived above (Eq. 101). If this criterion is not
satisfied the numerical shock solution is strongly inaccurate (left panel).
Figure 13: Cross-section slice showing density in the midplane in the 3D
dustysedov problem, for both the gas (left panel) and the dust (right panel)
at $t=0.1$. A dust-to-gas ratio of $0.01$ and a drag coefficient of $K=1$ have
been used with $100^{3}$ SPH particles in each phase. Note the slight
difference in the blast radius between the dust and the gas, consistent with
the response time ($t_{\rm s}$) of the dust to the gas drag.
Figure 14: Results of the 3D dustysedov test, showing the density in the gas
(left figure) and dust (right figure) from a Sedov blast wave propagating in
an astrophysical ($1\%$ dust-to-gas ratio) mixture of gas and dust with a
constant drag coefficient $K=1$. The low dust-to-gas ratio means that the gas
is only weakly affected by the drag from the dust, and is thus close to the
self-similar Sedov solution (dotted/red line). The dust density is affected by
the propagation of the blast, resulting in an overdensity that closely mirrors
the gas overdensity. Results are shown using $50^{3}$ (top panels) and
$100^{3}$ particles (bottom panels).
### 4.5 dustysedov: Sedov blast waves in a dust-gas mixture
The dustysedov test concerns the propagation of a Sedov blast wave in a dust-
gas mixture. Although the self-similar Sedov solution is known for the
propagation of a blast wave in a gas phase, the solution for a two fluid dust-
gas mixture is unknown (though at high drag as previously it may be expected
that the solution should revert to the gas-only solution using the modified
sound speed). We do not attempt to simulate this problem at high drag as it
would involve a prohibitive computational expense, instead adopting an
“astrophysical” dust-to-gas ratio of $1\%$ and a moderate drag coefficient
such that the presence of dust represents only a small perturbation to the gas
evolution. Results of purely hydrodynamic SPH solutions for this test can be
found e.g. in Rosswog & Price (2007) and Springel & Hernquist (2002).
#### 4.5.1 dustysedov: Setup
We setup the dustysedov problem in a 3D periodic box (the boundary conditions
are irrelevant for the times shown) at two different resolutions, filling the
box $x,y,z\in[-0.5,0.5]$ by $50^{3}$ and $100^{3}$ SPH particles for both the
gas and the dust. Gas particles are set up on a regular cubic lattice, with
the dust particles also on a cubic lattice but shifted by half of the lattice
step in each direction. We use $\alpha_{\rm SPH}=1$ and $\beta_{\rm SPH}=2$ in
the artificial viscosity terms, and $\alpha_{\rm u}=1$ in the artificial
conductivity term. An ideal gas equation of state $P=(\gamma-1)\rho u$ is
adopted with $\gamma=5/3$.
In the self-similar Sedov solution, the thermal energy of the gas is initially
concentrated at $r=0$. In the SPH simulation we distribute the internal energy
of the gas over the particles located inside a radius $r<r_{\rm{b}}$ where
$r_{\rm b}$ is set to 2h (i.e., the radius of the smoothing kernel). In code
units the total blast energy is $E=1$, with $\hat{\rho}_{\mathrm{g}}=1$ and
$\hat{\rho}_{\mathrm{d}}=0.01$. For $r>r_{\rm{b}}$, the gas sound speed is set
to be $2\times 10^{-5}$ in code units. The dust-to-gas ratio is set to $0.01$
to be consistent with the value measured for the interstellar medium. The drag
coefficient is set to $K=1$. Translated to physical units, assuming a box size
of $1$ pc, an ambient sound speed of $2\times 10^{4}$ cm/s and a gas density
of $\rho_{0}=6\times 10^{-23}$ g/cm3 the energy of the blast is $2\times
10^{51}$ erg and time is measured in units of $100$ years, roughly
corresponding to a supernova blast wave propagating into the interstellar
medium. Obviously in a real supernova the temperature inside the blast would
be much higher than the sublimation temperature of the dust, meaning that it
would be quickly evaporated, so the dustysedov test is mainly useful as a
benchmarking problem.
#### 4.5.2 dustysedov: Results
Fig. 14 shows the densities of both the gas (left figure) and the dust phase
(right figure) at $t=0.1$ using $50^{3}$ (top) and $100^{3}$ (bottom)
particles for both fluids. Fig. 13 shows a cross section of the density in the
midplane in the high resolution ($100^{3}$) simulation, showing the gas (left)
and dust (right). As the gas and dust densities are $1$ and $0.01$,
respectively and the drag coefficient is $K=1$ in code units, the stopping
time is $t_{\rm{s}}=0.01$, which represents $10\%$ of the time required for
the blast to fill the box. The response of the dust to the forcing by the gas
drag is therefore of order $10\%$. Consequently, an overdensity in the dust
phase forms due to the passage of the overdensity in the gas. The
overdensities in the gas and the dust phases are dephased slightly (seen by
comparing the position of the peak densities in the gas and dust in Fig. 14),
consistent with the finite time ($t_{\rm{s}}$) required for the dust to
respond to the gas forcing.
Figure 15: Rendering of the gas density of a typical T-Tauri Star
protoplanetary disc at two different resolutions: $10^{5}$ (left) and $10^{6}$
(right) gas particles. Increasing the resolution smoothens the gas phase. The
initial dust to gas ratio is $0.01$ so that the dust only slightly affects the
gas. Figure 16: Rendering of the dust surface density in a typical T-Tauri
Star protoplanetary disc with four different configurations: $2\times 10^{5}$
particles using a mixed smoothing length (top left) and $2\times 10^{5}$
particles (top right), $10^{5}$ dust and $10^{5}$ dust particles (bottom left)
and $2\times 10^{6}$ particles (bottom right) using the gas smoothing length.
Using the mixed smoothing length results in artificial structures in the dust
density. Smoother dust profiles are achieved by i) using the gas smoothing
length and ii) increasing the gas resolution. More accurate results are
obtained by increasing the number of both the gas and the dust particles.
### 4.6 dustydisc: Settling and migration in an accretion disc
Our final test, dustydisc, concerns the evolution of the dust and gas mixture
in a protoplanetary disc, where vertical settling and radial drift of the dust
particles (see e.g. Chiang & Youdin 2010b) are known to be crucial processes
in the early stages of planet formation. SPH is well suited to this problem
since i) free boundaries are trivial to implement and ii) the exact
conservation of angular momentum by both the gas and dust parts of the
algorithm means that the problem can be simulated for many dynamical times. We
used the standard linear Epstein regime given by Eq. 87. The drag force is
integrated explicitly. The key features — namely both vertical settling and
radial migration — are expected to occur. We focus here on the vertical
settling of the grains since the migration is extensively discussed in Ayliffe
et al. (2011). For this specific test, the ’artificial viscosity for a disc’
described in Lodato & Price (2010) and implemented in phantom is used.
#### 4.6.1 dustydisc: Setup
We setup $10^{5}$ gas particles and $10^{5}$ dust particles in a
$0.01M_{\odot}$ gas disc (with $0.0001M_{\odot}$ of dust) surrounding a
$1M_{\odot}$ star. The disc extends from 10 to 400 AU. Both gas and dust
particles are placed using a Monte-Carlo setup such that the surface density
profiles of both phases are $\Sigma\left(r\right)\propto r^{-1}$. The radial
profile of the gas temperature is taken to be $T\left(r\right)\propto
r^{-0.6}$ with a flaring $H/r=0.05$ at 100 AU. One code unit of time
corresponds to $10^{3}$ yrs. A uniform grain size of 1 cm is used.
#### 4.6.2 dustydisc: Resolution issue
Fig. 15 compares the evolution of the gas phase, varying the number of gas
particles from $10^{5}$ (left panel) to $10^{6}$ (right panel). The number of
dust particles does not affect the density profile of the gas given the small
initial dust to gas ratio. As expected, a smoother gas profile is achieved by
using a higher resolution.
Fig. 16 compares the evolution of the dust phase varying i) the smoothing
length used to compute the drag term and ii) the number of particles in each
phase. When the drag term is computed with a mixed smoothing length
($h=[h_{\rm g}+h_{\rm d}]/2$, top left panel), artificial structures develop
in the dust phase due to over-concentration of dust particles below the
resolution of the gas. These numerical artefacts are removed using instead the
gas smoothing length (top right panel). Indeed, the gas smoothing length is
larger than the dust smoothing length since the dust grains concentrate when
they reach the disc mid plane (see discussion in Sec. 2.4.2). Smoother dust
density profiles are achieved increasing the gas resolution ($10^{6}$ gas
particles keeping $10^{5}$ dust particles, bottom left panel). Increasing the
number of gas particles thus reduces the numerical noise in the dust phase.
Finally, the smoothest dust density profile is naturally obtained when the
resolution in the two phases is the highest ($2\times 10^{6}$ particles,
bottom right panel).
#### 4.6.3 dustydisc: Vertical settling of the particles
Figure 17: Vertical settling of a centimetre dust grain initially located at
$r_{0}=100$ AU and $z_{0}=2$ AU (solid/black). SPH results are compared to the
estimate given by the damped harmonic oscillator approximation (pointed/red).
The agreement between the numerical and the analytic solutions indicates that
the vertical settling of the dust grain is correctly reproduced by the SPH
algorithm. The analytic estimate neglects the radial drift of the grain and
the vertical stratification of the disc.
Fig. 17 compares the vertical settling of a dust particle obtained with the
SPH simulation and the analytic estimation given by the evolution of a damped
harmonic oscillator (see e.g. Garaud & Lin 2004). The particle is initially
located at $r=100$ AU and $z=2$ AU, and the evolution is computed for $50$
code units. The agreement between the numerical and the analytic solutions
indicates that the vertical settling of the dust grain is accurately
reproduced by the SPH algorithm. It should be noted that the analytic
estimation neglects the radial drift of the grain and the vertical
stratification of the disc. More precisely, this model assumes an expansion to
order zero in $(\frac{\partial P_{\rm g}}{\partial
r}/\hat{\rho}_{\mathrm{g}})/(\mathcal{G}M/r^{2})$ (meaning that the radial and
the vertical motion of the dust particles are decoupled) and to first order in
$z_{0}/H$ (the vertical stratification is neglected). The SPH results are thus
expected to slightly differ from the analytic approximation in both amplitude
and phase.
## 5 Discussion
In astrophysics, gas and dust mixtures have been predominantly studied with
grid-based codes. The gas phase is computed as usual whereas the dust is
treated by using superparticles (e.g. Youdin & Johansen 2007, Bai & Stone
2010). Computing the drag is usually divided into three steps: i)
interpolation of the gas velocities at the particle positions, ii) calculation
of the drag force on the particles and iii) attribution of the back-reaction
from the particles onto the nearby cells.
Our algorithm has two key advantages compared to the current grid-based
algorithms. First, the procedure used for interpolating the gas velocities at
the dust position conserves angular momentum exactly, avoiding artificial
local torques. Moreover, the interpolation in current grid-based codes is
performed with a standard bell-shaped kernel, regardless of the nature of the
error in the drag terms. We expect that these aspects of the drag computation
in grid-based codes may be improved by generalising the techniques involved in
our SPH algorithm. Second, the equations of motion of the mixture (considering
the drag terms only) in SPH are invariant when permuting the dust and the gas
indices, i.e. $\rm{g}\leftrightarrow\rm{d}$. This symmetry is broken in grid-
based schemes where the two phases are treated with two different methods
(super-particles superimposed to a grid). The SPH algorithm provides
rigorously identical results when interchanging the dust and the gas
properties (this has been verified on the dustybox problem which involves only
the drag, see above).
## 6 Conclusions
We have developed a new general SPH formalism for two-fluid dust and gas
mixtures, with the aim of simulating the dynamics of dusty gas systems in a
range of astrophysical contexts. In doing so we have generalised the standard
methods developed over 15 years ago by Monaghan & Kocharyan (1995) and
Monaghan (1997) for treating dusty gas in SPH. In Sec. 1, we highlighted seven
key issues. In this, Paper I, we have addressed five of these issues as
follows: 1) we have introduced a simple way to compute SPH densities on two
fluids with variable smoothing lengths; 2) the conservative part of the SPH
equations have been derived from a Lagrangian; 3) we have demonstrated how the
use of “double-hump” shaped kernels significantly improve the accuracy of the
SPH interpolation of drag terms; 4) we find a necessary criterion $h\lesssim
c_{\rm s}t_{\rm s}$ in order to correctly resolve differential motion between
gas and dust that becomes critical at high drag; we also find it important to
ensure $h_{\rm gas}\lesssim h_{\rm dust}$ to avoid artificial over-
concentration of dust particles, implying a higher resolution should be
employed in the gas phase relative to the dust.
Finally, to address issue 7), we have presented a comprehensive suite of
simple test problems that can be used to benchmark astrophysical dusty gas
codes. These consist of the dustybox, dustywave, dustyshock, dustysedov and
dustydisc problems. The first three of these have known (or partially known in
the case of dustyshock) analytic solutions and can be easily setup in any code
with standard boundary conditions. We have used these tests to explore the
issues raised above and have demonstrated that with the appropriate resolution
criteria satisfied, our formalism is robust and provides accurate results.
The two remaining issues — namely implicit timestepping and treatments of
astrophysical drag regimes —- are addressed in a companion paper (Paper II).
While this paper concentrates on two-fluid gas and dust mixtures, the
algorithm is general and can be applied easily to the treatment of other
multi-fluid systems in SPH (e.g. ambipolar diffusion).
## Acknowledgments
We thank Ben Ayliffe, Matthew Bate, Joe Monaghan and Laure Fouchet for useful
discussions and comments. Figures have been produced using splash (Price,
2007) with the new giza backend by DJP and James Wetter. We are grateful to
the Australian Research Council for funding via Discovery project grant
DP1094585.
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|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Guillaume Laibe (Monash), Daniel J. Price (Monash)",
"submitter": "Guillaume Laibe",
"url": "https://arxiv.org/abs/1111.3090"
}
|
1111.3356
|
# A note about the relation between fixed point theory on cone metric spaces
and fixed point theory on metric spaces
Ion Marian Olaru
###### Abstract
Let $Y$ be a locally convex Hausdorff space, $K\subset E$ a cone and
$\leq_{K}$ the partial order defined by $K$. Let $(X,p)$ be a $TVS-$ cone
metric space, $\varphi:K\rightarrow K$ a vectorial comparison function and
$f:X\rightarrow X$ such that
$p(f(x),f(y))\leq_{K}\varphi(p(x,y)),$
for all $x,y\in X$. We shall show that there exists a scalar comparison
function $\psi:\mathbb{R}_{+}\rightarrow\mathbb{R}_{+}$ and a metric
$d_{p}$(in usual sense) on $X$ such that
$d_{p}(f(x),f(y))\leq\psi(d_{p}(x,y)),$
for all $x,y\in X$. Our results extend the results of Du (2010) [Wei-Shih Du,
A note on cone metric fixed point theory and its equivalence, Nonlinear Anal.
72 (2010), 2259-2261].
2010 Mathematical Subject Classification: 47H10, 54H25
Keywords: K-metric spaces, cone metric space, TVS- cone metric spaces,
comparison function
## 1 Introduction and preliminaries
Fixed point theory in K-metric and K-normated spaces was developed by A.I.
Perov and his consortiums ([7], [8], [9]). The main idea consists to use an
ordered Banach space instead of the set of real numbers, as the codomain for a
metric. For more details on fixed point theory in K-metric and K-normed
spaces, we refer the reader to [15]. Without mentioning these previous works,
Huang and Zhang [6] reintroduced such spaces under the name of cone metric
spaces but went further, defining convergent and Cauchy sequences in the terms
of interior points of the underlying cone. They also proved some fixed point
theorems in such spaces in the same work. After that, fixed point results in
cone metric spaces have been studied by many other authors. References
[1],[3], [10], [11], [12], [13] are some works in this line of research.
However, very recently Wei-Shih Du in [5] used the scalarization function and
investigated the equivalence of vectorial versions of fixed point theorems in
cone metric spaces and scalar versions of fixed point theorems in metric
spaces. He showed that many of the fixed point results in ordered K-metric
spaces for maps satisfying contractive conditions of a linear type in K-metric
spaces can be considered as the corollaries of corresponding theorems in
metric spaces.
Let $E$ be a topological vector space (for short $t.v.s$) with its zero vector
$\theta_{E}$.
###### Definition 1.1.
( [5], [6]) A subset K of E is called a cone if:
* (i)
K is closed, nonempty and $K\neq\\{\theta_{E}\\}$;
* (ii)
$a,b\in\mathbb{R}$, $a,b\geq 0$ and $x,y\in K$ imply $ax+by\in K$;
* (iii)
$K\cap-K=\\{\theta_{E}\\}.$
For a given cone $K\subset E$, we can define a partial ordering $\leq_{K}$
with respect to $K$ by
(1.1) $x\leq_{K}y\ if\ and\ only\ if\ y-x\in K.$
We shall write $x<_{K}y$ to indicate that $x\leq_{K}y$ but $x\neq y$, while
$x\ll y$ will stand for $y-x\in intK$ (interior of $K$).
In the following, unless otherwise specified, we always suppose that $Y$ is a
locally convex Hausdorff with its zero vector $\theta$, $K$ a cone in $Y$ with
$intK\neq\emptyset$ , $e\in intK$ and $\leq_{K}$ a partial ordering with
respect to $K$.
###### Definition 1.2.
( [5]) Let $X$ be a nonempty set. Suppose that a mapping $d:X\times
X\rightarrow Y$ satisfies:
* (i)
$\theta\leq_{K}d(x,y)$ for all $x,y\in X$ and $d(x,y)=\theta$ if and only if
$x=y$;
* (ii)
$d(x,y)=d(y,x)$, for all $x,y\in X$ ;
* (iii)
$d(x,y)\leq_{K}d(x,z)+d(z,y)$ for all $x,y,z\in X$.
Then $d$ is called a TVS-cone metric on $X$ and $(X,d)$ is called a TVS-cone
metric space.
The nonlinear scalarization function $\xi_{e}:Y\rightarrow\mathbb{R}$ is
defined as follows
$\xi_{e}(y)=\inf\\{r\in\mathbb{R}\mid y\in r\cdot e-K\\}.$
###### Lemma 1.1.
( [4]) For each $r\in\mathbb{R}$ and $y\in Y$, the following statements are
satisfied:
* (i)
$\xi_{e}(y)\leq r$ if and only if $y\in r\cdot e-K$;
* (ii)
$\xi_{e}(y)>r$ if and only if $y\notin r\cdot e-K$;
* (iii)
$\xi_{e}(y)\geq r$ if and only if $y\notin r\cdot e-intK$;
* (iv)
$\xi_{e}(y)<r$ if and only if $y\in r\cdot e-intK$;
* (vi)
$\xi_{e}(\cdot)$ is positively homogeneous and continuous on Y;
* (vii)
if $y_{1}\in y_{2}+K$ then $\xi_{e}(y_{2})\leq\xi_{e}(y_{1})$;
* (viii)
$\xi_{e}(y_{1}+y_{2})\leq\xi_{e}(y_{1})+\xi_{e}(y_{2})$, for all
$y_{1},y_{2}\in Y$.
###### Theorem 1.1.
( [5]) Let $(X,p)$ be a $TVS-$cone metric space. Then
$d_{p}:X\times X\rightarrow[0,\infty)$
defined by $d_{p}=\xi_{e}\circ d$ is a metric.
## 2 Main results
###### Definition 2.1.
Let $K\subset Y$ be a cone. A function $\varphi:K\rightarrow K$ is called a
vectorial comparison function if
* (i)
$k_{1}\leq_{P}k_{2}$ implies $\varphi(k_{1})\leq_{P}\varphi(k_{2})$;
* (ii)
$\varphi(0)=0$ and $0<_{P}\varphi(k)<_{P}k$ for $k\in K-\\{0\\}$;
* (iii)
$k\in intK$ implies $k-\varphi(k)\in intK$;
* (iv)
if $t_{0}\geq 0$ then $\lim\limits_{t\rightarrow t_{0}^{+}}\varphi(t\cdot
e)=\varphi(t_{0}\cdot e)$.
###### Example 1.
* (i)
if K is an arbitrary cone in a Banach space E and $\lambda\in(0,1)$, then
$\varphi:K\rightarrow K$, defined by $\varphi(k)=\lambda k$ is a vectorial
comparison function;
* (ii)
Let $E=\mathbb{R}^{2}$, $K=\\{(x,y)\mid x,y\geq 0\\}$ and let
$\varphi_{1},\varphi_{2}:[0,\infty)\rightarrow[0,\infty)$ be such that
* (a)
$\varphi_{1},\varphi_{2}$ are increasing functions;
* (b)
if $t>0$ then $\varphi_{i}(t)<t$ for $i=\overline{1,2}$;
* (c)
$\varphi_{1},\varphi_{2}$ are right continuous.
Then $\varphi:K\rightarrow K$, defined by
$\varphi(x,y)=(\varphi_{1}(x),\varphi_{2}(y))$ is a vectorial comparison
function;
###### Definition 2.2.
( [14]) A function $\varphi:\mathbb{R}_{+}\rightarrow\mathbb{R}_{+}$ is called
a scalar comparison function if
* (i)
$t_{1}\leq t_{2}$ implies $\varphi(t_{1})\leq\varphi(t_{2})$;
* (ii)
$\varphi^{n}(t)\stackrel{{\scriptstyle n\rightarrow\infty}}{{\rightarrow}}0$
for all $t>0$
The following lemma will be useful in the sequel
###### Lemma 2.1.
( [14]) If $\varphi:\mathbb{R}_{+}\rightarrow\mathbb{R}_{+}$ is increasing and
right upper semicontinuous then the following assertions are equivalent:
* (a)
$\varphi^{n}(t)\stackrel{{\scriptstyle n\rightarrow\infty}}{{\rightarrow}}0$
for all $t>0$;
* (b)
$\varphi(t)<t$ for all $t>0$.
###### Lemma 2.2.
We consider $M:\mathbb{R}\rightarrow Y,M(r)=r\cdot e$. Then we have
* (i)
$M(0)=\theta$;
* (ii)
if $r_{1}\leq r_{2}$ then $M(r_{1})\leq_{K}M(r_{2})$;
* (iii)
$y\leq_{K}M\circ\xi_{e}(y)$ for all $y\in Y$;
* (iv)
$\xi\circ M(r)\leq r$ for all $r\in\mathbb{R}$;
* (v)
if $y_{1}\ll y_{2}$ then $\xi_{e}(y_{1})<\xi_{e}(y_{2})$.
Proof:
$(i)$ It is obvious;
$(ii)$ Let be $r_{1}\leq r_{2}$. Then $(r_{2}-r_{1})\cdot e\in K$. Thus
$M(r_{1})\leq_{K}M(r_{2})$;
$(iii)$ Since $\xi_{e}(y)=\inf\\{r\in\mathbb{R}\mid y\leq_{K}r\cdot e\\}$ it
follows that $y\leq_{K}\xi_{e}(y)\cdot e=M\circ\xi_{e}(y)$ for all $y\in Y$;
$(iv)$ Let be $r\in\mathbb{R}$. Since $\\{r^{\prime}\in\mathbb{R}\mid r\cdot
e\leq_{K}r^{\prime}\cdot e\\}\supseteq\\{r^{\prime}\in\mathbb{R}\mid r\leq
r^{\prime}\\}$ we get
$\xi_{e}(M(r))=\xi_{e}(r\cdot e)=\inf\\{r^{\prime}\in\mathbb{R}\mid r\cdot
e\leq_{K}r^{\prime}\cdot e\\}\leq\inf\\{r^{\prime}\in\mathbb{R}\mid r\leq
r^{\prime}\\}=r.$
$(v)$ Let be $y_{1}\ll y_{2}$. We remark that $y_{1}\ll
y_{2}\leq_{K}\xi_{e}(y_{2})\cdot e$. Then, via Remark 1.3 of Radenović and
Kadelburg [11], it follows that $y_{1}\ll\xi_{e}(y_{2})\cdot e$. Hence
$y_{1}\in\xi_{e}(y_{2})\cdot e-intK$. By using Lemma 1.1 (iv) we get
$\xi_{e}(y_{1})<\xi_{e}(y_{2})$.
###### Theorem 2.1.
Let $(X,p)$ be a TVS-cone metric and $\varphi:K\rightarrow K$ be a vectorial
comparison function such that
$p(f(x),f(y))\leq_{K}\varphi(p(x,y)),$
for all $x,y\in X$. Then there exists a scalar comparison function
$\psi:\mathbb{R}_{+}\rightarrow\mathbb{R}_{+}$ such that
$d_{p}(f(x),f(y))\leq\psi(d_{p}(x,y)),$
for all $x,y\in X$.
Proof: Let be $t\in\mathbb{R}_{+}$. Then $\theta\leq_{K}M(t)$. It follows that
$M(t)\in K$ for all $t\in\mathbb{R}_{+}$.
We define
$\psi:\mathbb{R}_{+}\rightarrow\mathbb{R}_{+},$
$\psi(t)=\xi_{e}\circ\varphi\circ M(t)$
First, we note that for all $t\in\mathbb{R}_{+}$ we have
$0\leq\xi_{e}\circ\varphi\circ M(t)\leq\xi_{e}\circ M(t)\leq t.$
Now, we remark that for each $x,y\in X$ we have
$d_{p}(f(x),f(y))\leq\xi_{e}\circ\varphi(p(x,y))\leq\xi_{e}\circ\varphi(M(\xi_{e}(p(x,y))))=\psi(d_{p}(x,y)).$
We claim that $\psi$ is a scalar comparison function. Since $\xi_{e}$,
$\varphi$ and $M$ are increasing functions, it follows that $\psi$ is
increasing function. In order to prove that
$\psi^{n}(t)\stackrel{{\scriptstyle n\rightarrow\infty}}{{\rightarrow}}0$ for
all $t>0$, we shall use Lemma 2.1. Next we show that $\psi(t)<t$ for all
$t>0$.
Let be $t_{0}>0$. Then $t_{0}\cdot e\in intK$. Therefore $\varphi(t_{0}\cdot
e)\ll t_{0}\cdot e$. It follows that
$\psi(t_{0})=\xi_{e}\circ\varphi(t_{0}\cdot e)<\xi_{e}\circ M(t_{0})\leq
t_{0}.$
Since $\lim\limits_{t\rightarrow t_{0}^{+}}\psi(t)=\lim\limits_{t\rightarrow
t_{0}^{+}}\xi_{e}\circ\varphi(t\cdot e)=\xi_{e}(\lim\limits_{t\rightarrow
t_{0}^{+}}\varphi(t\cdot e))=\xi_{e}\circ\varphi(t_{0}\cdot e)=\psi(t_{0})$ it
follows that $\psi$ is right upper semicontinuous. Hence
$\psi^{n}(t)\stackrel{{\scriptstyle n\rightarrow\infty}}{{\rightarrow}}0$.
###### Corollary 2.1.
Let $(X,p)$ be a complete TVS cone metric space and $\varphi:K\rightarrow K$ a
vectorial comparison function such that
$p(fx,fy)\leq_{K}\varphi(p(x,y)),$
for all $x,y\in X$. Then, f has a unique fixed point $x_{0}$.
Proof: We apply Theorem 2.1 and Theorem 1 pp 459 of Boyd and Wong ([2]).
###### Remark 2.1.
For $\varphi(k)=\lambda\cdot k$, $\lambda\in[0,1)$ we obtain, via Lemma 2.2
$(iv)$ and Corollary 2.1, the results of W.S. Du [5].
###### Remark 2.2.
Let $(X,p)$ a cone metric space. For $\varphi(k)=\lambda\cdot k$,
$\lambda\in[0,1)$ we obtain, via Remark 2.1, the results of L.G. Huang and
Zhang Xian [6].
Let $(X,d)$ be a TVS cone-metric space and let $\varphi:K\rightarrow K$ be a
vectorial comparison function. For a pair $(f,g)$ of self-mappings on $X$
consider the following conditions:
* (C)
for arbitrary $x,y\in X$ there exists $u\in\\{d(gx,gy),d(gx,fx),d(gy,fy)\\}$
such that $d(fx,fy)\leq_{P}\varphi(u)$.
* $(C_{1})$
for arbitrary $x,y\in X$ there exists
$w\in\\{d_{p}(gx,gy),d_{p}(gx,fx),d_{p}(gy,fy)\\}$ such that
$d_{p}(fx,fy)\leq\psi(u)$.
###### Remark 2.3.
The condition (C) imply the condition $(C_{1})$.
Indeed since the condition $(C)$ hold, it follows that at least one of the
following three cases holds:
* Case 1:
$u=d(gx,gy)$. Then
$\xi_{e}(p(fx,fy))\leq\xi_{e}\circ\varphi(u)\leq\xi_{e}\circ\varphi\circ
M(\xi_{e}(u))=\psi(d_{p}(gx,gy))$
* Case 2:
$u=d(gx,fx)$. Then
$\xi_{e}(p(fx,fy))\leq\xi_{e}\circ\varphi(u)\leq\xi_{e}\circ\varphi\circ
M(\xi_{e}(u))=\psi(d_{p}(gx,fx))$
* Case 3:
$u=d(gy,fy)$. Then
$\xi_{e}(p(fx,fy))\leq\xi_{e}\circ\varphi(u)\leq\xi_{e}\circ\varphi\circ
M(\xi_{e}(u))=\psi(d_{p}(gy,fy))$
## References
* [1] M. Abbas, G. Jungck, Common fixed point results for noncommuting mappings without continuity in cone metric spaces, J. Math.Anal.Appl. 341 (2008)
* [2] D.W. Boyd, J.S.W. Wong, On nonlinear contraction, Proceed. A.M.S., 20(1969).
* [3] Arandjelović I., Kadelburg Z, Radenović S., Boyd-Wong-type common fixed point results in cone metric spaces, Applied Mathematics and Computation 217(2011), 7167-7171.
* [4] G.Y. Chen, X.X. Huang, X.Q. Yang, Vector Optimization, Springer-Verlag, Berlin, Heidelberg, Germany, 2005.
* [5] Wei-Shih Du, A note on cone metric fixed point theory and its equivalence, Nonlinear Anal. 72 (2010), 2259-2261.
* [6] L.G. Huang, X. Zhang, Cone metric spaces and fixed Point theorems of contractive mappings , J. Math. Anal. Appl., 332(2007), 1468-1476.
* [7] B.V. Kvedaras, A.V. Kibenko and A.I. Perov, On some boundary value problems, Litov. matem. sbornik, 5 (1) (1965), 69-84.
* [8] A.I. Perov, The Cauchy problem for systems of ordinary diferential equations, Approximate methods of solving diferential equations, Kiev, Naukova Dumka, 1964, 115- 134 [Russian].
* [9] A.I. Perov and A.V. Kibenko, An approach to studying boundary value problems, Izvestija AN SSSR, Seria Math. 30 (2) (1966), 249-264 [Russian].
* [10] G. Jungck, S. Radenović, S. Radojević and V. Rakočević, Common fixed point theorems for weakly compatible pairs on cone metric spaces, Fixed Point Theory and Applications, 2009, Article ID 643640, 12 pages, doi: 10.1155/2009/643640.
* [11] Stojan Radenović, Zoran Kadelburg, Quasi-Contractions on symetric and cone symetric spaces, Banach J. Math. Anal.(2011), No.1 pp 38-50.
* [12] Raja P., Vaezpour S.M, Some extensions of Banach’s contraction principle in complete cone metric spaces, Fixed Point Theory and Applications Vol. 2008, article ID 768294, 11 pages, doi:101155/2008/768924.
* [13] Rezapour Sh., Hamblbarani R., Some notes on the paper ”Cone metric spaces and fixed point thorems of contractive mappings”, J. Math., Anal. Appl., 345(2008), 719-724.
* [14] I.A. Rus, Generalized contractions, Seminar on Fixed Point Theory, Preprint No. 3, 1983, pp 1-130.
* [15] P.P. Zabrejko, K-metric and K-normed linear spaces: survey, Collect. Math. 48 (4-6) (1997), 825-859.
Departament of Mathematics,
Faculty of Sciences,
University ”Lucian Blaga” of Sibiu,
Dr. Ion Ratiu 5-7, Sibiu, 550012, Romania
E-mail: marian.olaru@ulbsibiu.ro
|
arxiv-papers
| 2011-11-14T20:55:36 |
2024-09-04T02:49:24.329657
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Ion Olaru",
"submitter": "Ion Olaru",
"url": "https://arxiv.org/abs/1111.3356"
}
|
1111.3425
|
On leave of absence from Center for Nuclear Physics, ] Institute of Physics,
Hanoi, Vietnam
# Pairing reentrance in hot rotating nuclei
N. Quang Hung1 [ hung.nguyen@ttu.edu.vn N. Dinh Dang2,3 dang@riken.jp 1)
School of Engineering, TanTao University, TanTao University Avenue, TanDuc
Ecity, Duc Hoa, Long An Province, Vietnam
2) Theoretical Nuclear Physics Laboratory, RIKEN Nishina Center for
Accelerator-Based Science, 2-1 Hirosawa, Wako City, 351-0198 Saitama, Japan
3) Institute for Nuclear Science and Technique, Hanoi, Vietnam
###### Abstract
The pairing gaps, heat capacities and level densities are calculated within
the BCS-based quasiparticle approach including the effect of thermal
fluctuations on the pairing field within the pairing model plus noncollective
rotation along the $z$ axis for 60Ni and 72Ge nuclei. The analysis of the
numerical results obtained shows that, in addition to the pairing gap, the
heat capacity can also serve as a good observable to detect the appearance of
the pairing reentrance in hot rotating nuclei, whereas such signature in the
level density is rather weak.
Suggested keywords
###### pacs:
21.60.-n, 21.60.Jz, 24.60.-k, 24.10.Pa
## I Introduction
In the collective rotation of a deformed nucleus, the rotation axis is
perpendicular to the symmetry axis, and the Coriolis force, which breaks the
Cooper pairs, increases with the total angular momentum so that at a certain
critical angular momentum all Cooper pairs are broken. The nucleus undergoes
then a phase transition from the superfluid phase to the normal one (the SN
phase transition). This is called the Mottelson-Valatin effect Mottelson . A
similar effect is expected in spherical nuclei, although the rotation is no
longer collective. The total angular momentum is made up by those of the
nucleons from the broken pairs, which occupy the single-particle levels around
the Fermi surface and block them against the scattered pairs. The pairing
correlations decrease until a sufficiently large total angular momentum
$M_{c}$, where the pairing gap $\Delta$ completely vanishes. At finite
temperature ($T\neq$ 0), the increase of $T$ relaxes the tight packing of
quasiparticles around the Fermi surface, which is caused by a large angular
momentum $M\geq M_{c}$, and spreads them farther away from the Fermi level.
This makes some levels become partially unoccupied, therefore, available for
scattered pairs. As the result, when $T$ increases up to some critical value
$T_{1}$, the pairing correlations are energetically favored, and the pairing
gap reappears. As $T$ goes higher, the increase of a large number of
quasiparticles eventually breaks down the pairing gap at $T_{2}$ $(>T_{1})$.
This phenomenon, predicted by Kamuri Kammuri and Moretto Moretto , is called
thermally assisted pairing or anomalous pairing, and later as pairing
reentrance by Balian, Flocard and V$\acute{\rm e}$n$\acute{\rm e}$roni Balian
.
However, it has been shown already in the 1960s that the sharp SN-phase
transition at $M=M_{c}$ in the Mottelson-Valatin effect is an artifact of the
BCS method. As a matter of fact, a proper particle-number projection before
variation has removed the discontinuity in the pairing gap as a decreasing
function of the angular momentum Mang . Similarly, by taking the effect of
thermal fluctuations in the pairing field into account, the SN phase
transition predicted by the BCS theory is smoothed out. The gap $\Delta(T)$ of
a non-rotating nucleus does not collapse at $T_{c}\simeq 0.568\Delta(T=0)$,
but monotonically decreases with increasing $T$, remaining finite even at
$T\gg T_{c}$ Moretto2 ; MBCS . This result is reconfirmed by shell-model
calculations of pairing energy as a function of excitation energy Zele , and
by embedding the exact eigenvalues of the pairing problem into the canonical
ensemble Exact . By considering an exactly solvable cranked deformed shell
model Hamiltonian it has also been shown that the pairing gap, quenched at
$T=$ 0 and high rotational frequency, reappears at $T_{1}$ ($\ll T_{c}$) at
$M\geq M_{c}$ Frauendorf . However, different from the prediction by the BCS
theory, the pairing gap does not vanish at $T>T_{1}$.
The behavior of hot rotating nuclei can be put in correspondence with
superconductors in the presence of an external magnetic field, where the
magnetic field plays the role as that of the nuclear rotation. The reentrance
of superconducting correlations, which is also known as the unconventional
superconductivity, has been the subject of recent theoretical and experimental
studies in condensed matter. The Grenoble High Magnetic Field Laboratory has
recently discovered that URhGe becomes superconducting at low temperature in
the presence of a strong magnetic field (between 8 and 13 T), well above the
field value of 2 T, at which superconductivity is first destroyed Sci . The
reentrance of superconductivity under the magnetic field is interpreted as to
be caused by spin reorientation reent , which bears some similarity with the
reappearance of the scattered pairs in rotating hot nuclei discussed above.
Recently, we have developed an approach based on the finite-temperature BCS
(FTBCS) that includes the effects due to quasiparticle-number fluctuations in
the pairing field and the $z$ projection of angular momentum at $T\neq$ 0,
which we call as the FTBCS1 (with ”1” denoting the effect due to
quasiparticle-number fluctuations) SCQRPA . This approach reproduces well the
effect of smoothing out the SN phase transition at $T\neq$ 0 as well as the
pairing reentrance in hot (noncollectively) rotating nuclei. In the latter
case, for $M\geq M_{c}$ the pairing gap also reappears at $T_{1}$ and remains
finite at $T>T_{1}$. It has also been pointed out that the microscopic
mechanism of the nonvanishing gap at high $T$ is the quasiparticle-number
fluctuations, which are ignored in the conventional BCS theory. A refined
version of the FTBCS1 also includes the contribution of coupling to pair
vibrations within the self-consistent quasiparticle random-phase approximation
SCQRPA .
That the pairing reentrance is not an artifact of the mean field analysis, but
a robust physical effect, has been obtained by exact diagonalization of the 2D
attractive Hubbard model, where a nonmonotonic filed dependence of the pair
susceptibility in the presence of the external magnetic field was found for
various cluster sizes both in the weak and strong coupling limit Gorczyca .
Nonetheless, the experimental extraction of the pairing gap in hot nuclei is
not simple because one has to properly exclude the admixture with the
contribution of uncorrelated single-particle configurations from the odd-even
mass difference Ensemble . Therefore the detection of the pairing reentrance
effect by using the experimentally extracted pairing gaps seems to be elusive,
especially when the formally derived pairing gap has a value smaller than the
average spacing between the single-particle levels.
Meanwhile, the heat capacity has been extracted from the experimental level
densities EXP . The existence of a bump or an $S$ shape on the curve of the
heat capacity at $T\sim T_{c}$ allows one to discuss about the smoothing of
the SN phase transition in finite nuclei. In a recent calculation of the heat
capacity in 72Ge within the Shell Model Monte Carlo (SMMC) approach, by
reconfirming the pairing reentrance effect, the authors of Ref. Dean claimed
that they found a local dip in the heat capacity at rotation frequency of 0.5
MeV at $T\sim$ 0.45 MeV, and a corresponding local maximum on the temperature
dependence of the logarithm of level density. They associated such
irregularities in the heat capacity and level density as the signatures of the
pairing reentrance. There are, however, two concerns regarding these results.
The first one is that, as well-known, at such low temperature the SMMC
approach produces quite large error bars. As a consequence, instead of
approaching zero as it should be to fulfill the third law of thermodynamics,
the SMMC heat capacity at $T<$ 0.5 MeV jumps to 25 $\sim$ 30 (Fig. 4 of Dean
), which makes the statement on the signature of the pairing reentrance
ambiguous. The second one is that the SMMC in Ref. Dean used the same Fock-
space single-particle energies of the shells ($0f1p-0g1d2s$) for both neutron
and proton spectra. Hence, the difference between neutron and proton spectra
came solely from the difference in the valence particle numbers outside the
closed-shell core of 40Ca, which ignored altogether the Coulomb barrier in the
proton spectrum. Our FTBCS1 theory is free from such deficiency at low $T$,
and it works well with the schematic as well as single-particle spectra, which
are obtained from the realistic Woods-Saxon potential. Therefore, in the
present paper we will calculate the heat capacity as well as the level density
within the FTBCS1 theory to see if these quantities can be used to identify
the pairing reentrance phenomenon in realistic nuclei at finite temperature
and angular momentum.
The paper is organized as follows. The formalism for the calculations of
thermodynamic quantities such as pairing gap, heat capacity, and level density
in hot non-collectively rotating nuclei within the FTBCS and FTBCS1 theories
is presented in Sec. II. The results of numerical calculations are analyzed in
Sec. III. The paper is summarized in the last section, where conclusions are
drawn.
## II Formalism
We consider the pairing Hamiltonian describing a spherical system rotating
about the symmetry $z$ axis Moretto :
$H=H_{P}-\lambda\hat{N}-\gamma\hat{M}~{},$ (1)
where $\lambda$ and $\gamma$ are the chemical potential and rotation
frequency, respectively. $H_{P}$ is the standard pairing Hamiltonian of a
system, which consists $N$ particles interacting via a monopole pairing force
with the constant parameter $G$ (the BCS pairing Hamiltonian), namely
$H_{P}=\sum_{k}\epsilon_{k}(a_{+k}^{\dagger}a_{+k}+a_{-k}^{\dagger}a_{-k})-G\sum_{kk^{\prime}}{a_{k}^{\dagger}a_{-k}^{\dagger}a_{-k^{\prime}}a_{k^{\prime}}}~{},$
(2)
with $a_{\pm k}^{\dagger}(a_{\pm k})$ being the creation (annihilation)
operators of a particle (neutron or proton) with angular momentum $k$,
projection $\pm m_{k}$, and energy $\epsilon_{k}$. The particle-number
operator $\hat{N}$ and total angular momentum $\hat{M}$, which coincides with
its $z$ projection, are given as
$\hat{N}=\sum_{k}(a_{+k}^{\dagger}a_{+k}+a_{-k}^{\dagger}a_{-k})~{},\hskip
14.22636pt\hat{M}=\sum_{k}m_{k}(a_{+k}^{\dagger}a_{+k}-a_{-k}^{\dagger}a_{-k})~{}.$
(3)
After the Bogoliubov transformation from the particle operators,
$a_{k}^{\dagger}$ and $a_{k}$, to the quasiparticle ones,
$\alpha_{k}^{\dagger}$ and $\alpha_{k}$,
$a_{k}^{\dagger}=u_{k}\alpha_{k}^{\dagger}+v_{k}\alpha_{-k}~{},\hskip
14.22636pta_{-k}=u_{k}\alpha_{-k}-v_{k}\alpha_{k}^{\dagger}~{},$ (4)
the Hamiltonian (1) is transformed into the quasiparticle one $\cal H$, whose
explicit form can be found, e.g., in Refs. SCQRPA ; SCQRPA1 .
As has been discussed in Refs. Kammuri ; Moretto ; SCQRPA , for a spherically
symmetric system, the laboratory-frame $z$ axis, which is taken as the axis of
quantization, can always be made coincide with the body-fixed one, which is
aligned with the direction of the total angular momentum within the quantum
mechanical uncertainty. Therefore the total angular momentum is completely
determined by its $z$-projection $M$ alone. For systems of an axially
symmetric oblate shape rotating about the symmetry axis, which in this case is
the principal body-fixed one, this noncollective motion is known as “single-
particle” rotation. The pairing reentrance effect was originally obtained
within the BCS theory in Refs. Kammuri ; Moretto by considering such systems
described by Hamiltonian (1). Its physical interpretation based on the thermal
effect, which relaxes the tight packing of quasiparticles around the Fermi
surface due to a large angle momentum $M\geq M_{c}$, and spreads them farther
away from the Fermi level, fits well in the framework of this “single-
particle” rotation. However, as has been pointed out in Ref. Moretto , for
non-spherical nuclei, and specifically in the case of axially symmetric ones,
the spin and angular momentum projections on the symmetry axis are not good
quantum numbers. In this case the formalism used here is not completed because
it does not include the angular momentum’s component perpendicular to the
symmetry axis. Cranking model might serve as a better solution of the problem
in this situation. This remains to be investigated because the results
obtained within the Lipkin model with $J_{x}$ cranking did not reveal any
pairing reentrance so far Civi . On the other hand, in the region of high
level densities (at high excitation energies and/or high $T$) the values of
angular momentum projection on the symmetry axis will be mixed among the
levels, which worsen the axial symmetry. The melting of shell structure will
also eventually drive nuclei to their average spherical shape.
### II.1 FTBCS1 equations at finite angular momentum
The FTBCS1 includes a set of FTBCS-based equations, corrected by the effects
of quasiparticle-number fluctuations, for the level-dependent pairing gap
$\Delta_{k}$, average particle number $N$, and average angular momentum $M$.
The derivation of the FTBCS1 equations was reported in detail in Ref. SCQRPA ,
so we do present here only the final equations. The FTBCS1 equation for the
pairing gap is written as a sum of two parts, the level-independent part
$\Delta$ and the level-dependent part $\delta\Delta_{k}$, namely
$\Delta_{k}=\Delta+\delta\Delta_{k}~{},$ (5)
where
$\Delta=G\sum_{k^{\prime}}{u_{k^{\prime}}v_{k^{\prime}}(1-n_{k^{\prime}}^{+}-n_{k^{\prime}}^{-})}~{},\hskip
14.22636pt\delta\Delta_{k}=G\frac{\delta{\cal
N}_{k}^{2}}{1-n_{k}^{+}-n_{k}^{-}}u_{k}v_{k}~{},$ (6)
where
$\displaystyle u_{k}^{2}$ $\displaystyle=$
$\displaystyle\frac{1}{2}\left(1+\frac{\epsilon_{k}-Gv_{k}^{2}-\lambda}{E_{k}}\right)~{},\hskip
14.22636ptv_{k}^{2}=\frac{1}{2}\left(1-\frac{\epsilon_{k}-Gv_{k}^{2}-\lambda}{E_{k}}\right)~{},$
$\displaystyle E_{k}$ $\displaystyle=$
$\displaystyle\sqrt{(\epsilon_{k}-Gv_{k}^{2}-\lambda)^{2}+\Delta_{k}^{2}}~{},\hskip
14.22636ptn_{k}^{\pm}=\frac{1}{1+e^{\beta(E_{k}\mp\gamma m_{k})}}~{},\hskip
14.22636pt\beta=1/T~{}.$ (7)
with the quasiparticle-number fluctuations $\delta{\cal N}_{k}^{2}$ at nonzero
angular momentum
$\delta{\cal N}_{k}^{2}=(\delta{\cal N}_{k}^{+})^{2}+(\delta{\cal
N}_{k}^{-})^{2}=n_{k}^{+}(1-n_{k}^{+})+n_{k}^{-}(1-n_{k}^{-})~{}.$ (8)
The corrections due to coupling to pair vibration beyond the quasiparticle
mean field at finite temperature and angular momentum are significant only in
light nuclei like oxygen or neon isotopes, whereas they are negligible for
medium and heavy nuclei (See Figs. 6 - 8 of Ref. SCQRPA ). As the present
paper considers two medium nuclei, 60Ni and 72Ge, these corrections on the
FTBCS1 equations are neglected in the numerical calculations.
The equations for the particle number and total angular momentum are the same
as Eqs. (25) and (26) of Ref. SCQRPA , namely
$N=2\sum_{k}\left[v_{k}^{2}(1-n_{k}^{+}-n_{k}^{-})+\frac{1}{2}(n_{k}^{\dagger}+n_{k}^{-})\right],\hskip
14.22636ptM=\sum_{k}m_{k}(n_{k}^{+}-n_{k}^{-})~{}.$ (9)
The system of coupled equations (5) – (9) are called the FTBCS1 equations at
finite angular momentum. Once the FTBCS1 equations are solved, the total
energy ${\cal E}$, heat capacity $C$ and entropy $S$ of the system are
calculated
$\displaystyle{\cal E}$ $\displaystyle=$ $\displaystyle\langle{\cal
H}\rangle~{},\hskip 14.22636ptC=\frac{\partial{\cal E}}{\partial T}~{},$
$\displaystyle S$ $\displaystyle=$ $\displaystyle-\sum_{k}[n_{k}^{+}{\rm
ln}n_{k}^{+}+(1-n_{k}^{+}){\rm ln}(1-n_{k}^{+})+n_{k}^{-}{\rm
ln}n_{k}^{-}+(1-n_{k}^{-}){\rm ln}(1-n_{k}^{-})]~{}.$ (10)
### II.2 Level density
Within the conventional FTBCS, the level density is calculated as the invert
Laplace transformation of the grand partition function Moretto
$\rho(E,N,M)=\frac{1}{(2\pi i)^{3}}\oint d\beta\oint d\alpha\oint d\mu
e^{S}~{},$ (11)
where $\alpha=\beta\lambda$, $\mu=\beta\gamma$, and $S$ is entropy of the
system. The saddle-point approximation gives a good evaluation of the integral
(11). As the result, the total level density of a system with $N$ neutrons and
$Z$ protons is given as
$\rho(E,N,M)=\frac{e^{S}}{(2\pi)^{2}\sqrt{D}}~{},$ (12)
where $S=S_{N}+S_{Z}$ and
$D=\left|\begin{array}[]{cccc}\frac{\partial^{2}\Omega}{\partial\alpha_{N}^{2}}&\frac{\partial^{2}\Omega}{\partial\alpha_{N}\partial\alpha_{Z}}&\frac{\partial^{2}\Omega}{\partial\alpha_{N}\partial\mu}&\frac{\partial^{2}\Omega}{\partial\alpha_{N}\partial\beta}\\\
\frac{\partial^{2}\Omega}{\partial\alpha_{Z}\partial\alpha_{N}}&\frac{\partial^{2}\Omega}{\partial\alpha_{Z}^{2}}&\frac{\partial^{2}\Omega}{\partial\alpha_{Z}\partial\mu}&\frac{\partial^{2}\Omega}{\partial\alpha_{Z}\partial\beta}\\\
\frac{\partial^{2}\Omega}{\partial\mu\partial\alpha_{N}}&\frac{\partial^{2}\Omega}{\partial\mu\partial\alpha_{Z}}&\frac{\partial^{2}\Omega}{\partial\mu^{2}}&\frac{\partial^{2}\Omega}{\partial\mu\partial\beta}\\\
\frac{\partial^{2}\Omega}{\partial\beta\partial\alpha_{N}}&\frac{\partial^{2}\Omega}{\partial\beta\partial\alpha_{Z}}&\frac{\partial^{2}\Omega}{\partial\beta\partial\mu}&\frac{\partial^{2}\Omega}{\partial\beta^{2}}\end{array}\right|~{}.$
(13)
The logarithm of the grand-partition function of the systems is given as
$\Omega=\Omega_{N}+\Omega_{Z}=S+\alpha_{N}N+\alpha_{Z}Z+\mu M-\beta{\cal
E}~{}.$ (14)
The derivation of $\Omega$ with respect to $\alpha$ and $\mu$ can be seen
explicitly in Eqs. (25)-(35) of Ref. Moretto . Within the FTBCS1, the grand-
partition function has the same form as that given by Eq. (14) of the FTBCS.
Therefore, the first and second derivatives of the FTBCS1 grand-partition
function are the same as those of the FTBCS ones. The only difference comes
from the first derivatives of the pairing gap with respect to $\alpha$, $\mu$,
and $\beta$ because of the quasiparticle-number fluctuations in the FTBCS1 gap
equation (6). In this case, instead of the simple Eqs. [(33)-(35)] of Ref.
Moretto , the derivatives become rather complicate expressions, which are
obtained by taking the first derivatives of the left and right-hand sides of
Eq. (6) with respect to $\alpha$, $\mu$, and $\beta$. Amongst the three
derivatives of the FTBCS1 gap, ${\partial\Delta_{k}}/{\partial\beta}$ can be
obtained by using its definition, namely
$\frac{\partial\Delta_{k}}{\partial\beta}=-T^{2}\frac{\partial\Delta_{k}}{\partial
T}=-T^{2}\frac{\Delta_{k}(T+\delta T)-\Delta_{k}(T)}{\delta T}~{},$ (15)
and can be easily calculated numerically by choosing an appropriate value of
$\delta T$ as the input parameter. The other two derivatives,
${\partial\Delta_{k}}/{\partial\alpha}$ and
${\partial\Delta_{k}}/{\partial\mu}$, must be calculated from their explicit
analytic expressions because $\alpha$ and $\beta$ are two Lagrange
multipliers, which are obtained by solving the FTBCS1 equations. The final
equations for ${\partial\Delta_{k}}/{\partial\alpha}$ and
${\partial\Delta_{k}}/{\partial\mu}$ are derived as
$\sum_{k}\left(A_{k}\beta\frac{\partial\Delta_{k}}{\partial\alpha}+B_{k}\right)+\left(C_{k}^{+}+C_{k}^{-}-\frac{2}{G}\right)\beta\frac{\partial\Delta_{k}}{\partial\alpha}+(D_{k}^{+}+D_{k}^{-})=0~{},$
(16)
$\sum_{k}\left(A_{k}^{\prime}\beta\frac{\partial\Delta_{k}}{\partial\mu}+B_{k}^{\prime}\right)+\left(C_{k}^{\prime}+D_{k}^{\prime}-\frac{2}{G}\right)\beta\frac{\partial\Delta_{k}}{\partial\mu}+E_{k}^{\prime}=0~{},$
(17)
where
$A_{k}=\frac{1}{E_{k}^{3}}\left\\{(\lambda-\epsilon_{k})^{2}(1-n_{k}^{+}-n_{k}^{-})+\beta\Delta_{k}^{2}E_{k}^{2}[n_{k}^{+}+n_{k}^{-}-2(n_{k}^{+})^{2}-2(n_{k}^{-})^{2}]\right\\}~{},$
(18)
$B_{k}=-\frac{\Delta_{k}(\lambda-\epsilon_{k})}{E_{k}^{3}}\left\\{1-n_{k}^{+}-n_{k}^{-}-\beta
E_{k}^{2}(n_{k}^{+}+n_{k}^{-}-2(n_{k}^{+})^{2}-2(n_{k}^{-})^{2}\right\\}~{},$
(19)
$C_{k}^{+}=\frac{n_{k}^{+}(1-n_{k}^{+})}{(1-n_{k}^{+}-n_{k}^{-})^{2}E_{k}^{3}}\left\\{(\lambda-\epsilon_{k})^{2}(1-n_{k}^{+}-n_{k}^{-})-\beta\Delta_{k}^{2}E_{k}[(n_{k}^{+}-1)^{2}+n_{k}^{-}(2n_{k}^{+}-n_{k}^{-})]\right\\}~{},$
(20)
$C_{k}^{-}=\frac{n_{k}^{-}(1-n_{k}^{-})}{(1-n_{k}^{+}-n_{k}^{-})^{2}E_{k}^{3}}\left\\{(\lambda-\epsilon_{k})^{2}(1-n_{k}^{+}-n_{k}^{-})-\beta\Delta_{k}^{2}E_{k}[(n_{k}^{-}-1)^{2}+n_{k}^{+}(2n_{k}^{-}-n_{k}^{+})]\right\\}~{},$
(21)
$D_{k}^{+}=-\frac{n_{k}^{+}(1-n_{k}^{+})\Delta_{k}(\lambda-\epsilon_{k})}{(1-n_{k}^{+}-n_{k}^{-})^{2}E_{k}^{3}}\\{1-n_{k}^{+}-n_{k}^{-}+\beta
E_{k}[(n_{k}^{+}-1)^{2}+n_{k}^{-}(2n_{k}^{+}-n_{k}^{-})]\\}~{},$ (22)
$D_{k}^{-}=-\frac{n_{k}^{-}(1-n_{k}^{-})\Delta_{k}(\lambda-\epsilon_{k})}{(1-n_{k}^{+}-n_{k}^{-})^{2}E_{k}^{3}}\\{1-n_{k}^{+}-n_{k}^{-}+\beta
E_{k}[(n_{k}^{-}-1)^{2}+n_{k}^{+}(2n_{k}^{-}-n_{k}^{+})]\\}~{},$ (23)
$A_{k}^{\prime}=\frac{1}{E_{k}^{3}}\\{(\lambda-\epsilon_{k})^{2}(1-n_{k}^{+}-n_{k}^{-})+\beta\Delta_{k}^{2}E_{k}[n_{k}^{+}(1-n_{k}^{+})+n_{k}^{-}(1-n_{k}^{-})]\\}~{},$
(24) $B_{k}^{\prime}=\frac{\beta
m_{k}\Delta_{k}}{E_{k}}[n_{k}^{+}(1-n_{k}^{+})+n_{k}^{-}(1-n_{k}^{-})]~{},$
(25)
$C_{k}^{\prime}=\frac{(\lambda-\epsilon_{k})^{2}}{(1-n_{k}^{+}-n_{k}^{-})E_{k}^{3}}[n_{k}^{+}(1-n_{k}^{+})+n_{k}^{-}(1-n_{k}^{-})]~{},$
(26)
$D_{k}^{\prime}=-\frac{\beta\Delta_{k}^{2}}{(1-n_{k}^{+}-n_{k}^{-})^{2}E_{k}^{2}}\\{n_{k}^{+}(1-n_{k}^{+})[(n_{k}^{+}-1)^{2}+n_{k}^{-}(2n_{k}^{+}-n_{k}^{-})]$
$+n_{k}^{-}(1-n_{k}^{-})[(n_{k}^{-}-1)^{2}+n_{k}^{+}(2n_{k}^{-}-n_{k}^{+})]\\}~{},$
(27) $E_{k}^{\prime}=-\frac{\beta
m_{k}\Delta_{k}}{E_{k}}[n_{k}^{+}(1-n_{k}^{+})+n_{k}^{-}(1-n_{k}^{-})]~{}.$
(28)
By solving first the FTBCS1 equations, then Eqs. (16) and (17), one obtains
${\partial\Delta_{k}}/{\partial\alpha}$ and
${\partial\Delta_{k}}/{\partial\alpha}$ as functions of $T$ at a given value
of the total angular momentum $M$.
## III Analysis of numerical results
The numerical calculations are carried out for two realistic 60Ni and 72Ge
nuclei. The latter is considered in order to have a comparison with the
results obtained within the SMMC approach in Ref. Dean . The single-particle
spectra for these two nuclei are obtained within the axially deformed Woods-
Saxon potential WS , whose parameters are chosen to be the same as those given
in Ref. leveldens . All the bound (negative energy) single-particle states are
used in the calculations. The quadrupole deformation parameters $\beta_{2}$
are equal to 0 and -0.224 for 60Ni and 72Ge, respectively. The pairing
interaction parameters are adjusted so that the pairing gaps at $T=0$ fit the
experimental values obtained from the odd-even mass differences. These values
are $G_{N}=$ 0.347 MeV, which gives $\Delta_{N}=$ 1.7 MeV for neutrons in 60Ni
and $G_{N}=$ 0.291 MeV, which gives $\Delta_{N}=$ 1.7 MeV for neutrons in
72Ge. For protons in 72Ge, the value $G_{Z}=$ 0.34 MeV is chosen to give
$\Delta_{Z}=$ 1.5 MeV, whereas there is no pairing gap for the closed-shell
protons ($Z=$ 28) in 60Ni. Because the pairing gap (5) is level-dependent, the
level-weighted gap $\bar{\Delta}$ is considered, which is defined as
$\bar{\Delta}=\sum_{k}\Delta_{k}/\Omega$ with the total number $\Omega$ of
levels in the deformed basis [In the case of spherical basis it becomes
$\bar{\Delta}=\sum_{j}(2j+1)\Delta_{j}/\sum_{j}(2j+1)$ SCQRPA ].
Figure 1: (Color online) Level-weighted neutron pairing gap $\bar{\Delta}$
[(a), (d)], heat capacity $C$ [(b), (e)], and heat capacity divided by
temperature $C/T$ [(c), (f)] for 60Ni obtained at different values of angular
momentum $M$ as functions of $T$. Panels (a) - (c) show the FTBCS results,
whereas the predictions by the FTBCS1 are displayed in (d) - (f).
Shown in Fig. 1 are the neutron level-weighted pairing gap $\bar{\Delta}$, the
heat capacity $C$, and the ratio $C/T$ obtained as functions of $T$ at several
values of the total angular momentum $M$. The left column represents the
predictions by the standard FTBCS, whereas the results obtained within the
FTBCS1 are displayed in the right column. Both approaches show the pairing
reentrance in the gaps at $M=$ 4 and 6 $\hbar$, namely the gap increases with
$T$ up to $T\simeq$ 0.3 MeV, then decreases as $T$ increases further. Because
of the quasiparticle-number fluctuations, the FTBCS1 gap does not collapse at
$T_{c}$ as the FTBCS one, but decreases monotonically at high $T$. At $M=$ 14
$\hbar$, while the FTBCS gap completely vanishes at all $T$, the FTBCS1 gap
shows a spectacular reentrance effect, namely it increases from the zero value
at $T=$ 0 up to around 0.3 MeV at $T\simeq$ 0.8 MeV, and then slowly decreases
as $T$ further increases.
The heat capacities obtained within the FTBCS and FTBCS1 look alike, except
for the region around $T_{c}$, where the quasiparticle-number fluctuations
smooth out the sharp SN phase transition so that the sharp local maximum is
depleted to a broad bump. In the region, where the pairing reentrance takes
place, namely at $T\simeq$ 0.3 MeV and $M=$ 4 or 6 $\hbar$, a weak local
minimum is seen on the curve representing the temperature dependence of the
heat capacity similarly to the feature reported in Ref. Dean . This local
minimum is magnified by using the ratio $C/T$ so that the latter might be
useful in experiments as a quantity to identify the pairing reentrance.
However, when the gap is too small as in the pairing reentrance at $M=$ 14
$\hbar$, the heat capacity $C$ ($C/T$) obtained within FTBCS1 is almost
identical to that predicted by the FTBCS, where the gap is zero.
Figure 2: (Color online) Level-weighted pairing gaps for neutrons [(a), (e)],
protons [(b), (f)], heat capacity $C$ [(c), (g)], and heat capacity divided by
temperature $C/T$ [(d), (h)] for 72Ge obtained at different values of angular
momentum $M$ as functions of $T$. Panels (a) - (d) show the FTBCS results,
whereas the predictions by the FTBCS1 are displayed in (e) - (h).
For 72Ge, both neutron and proton gaps exist, which cause two peaks in the
temperature dependence of the heat capacity obtained within the FTBCS, as
shown in Fig. 2 (c). The overall features of $C$ and $C/T$ for 72Ge are
similar to those obtained for 60Ni. As compared with the results of Ref. Dean
, where the same single-particle energies in the ($0f1p-0g1d2s$) shells were
used for both neutrons and protons, and where the pairing reentrance was
predicted for neutrons, no pairing reentrance effect for neutrons is seen in
the results of our calculations. On the other hand, the pairing reentrance
takes place for protons at $M\geq$ 6 $\hbar$, as shown in Fig. 2 (b) and 2
(f).
Figure 3: (Color online) Level densities as functions of $T$ at several values
of total angular momentum $M$ obtained for 60Ni [(a), (b)] and 72Ge [(c), (d)]
within the FTBCS (left panels) and FTBCS1 (right panels).
Another experimentally measurable quantity, which may help to identify the
pairing reentrance effect, is the level density. In fact, the authors of Ref.
Dean claimed that the pairing reentrance causes an irregularity in a shape of
a small local maximum at low $T$ on the curve, which describes the temperature
dependence of the level density. The level densities obtained at several
values of the total angular momentum $M$ for 60Ni and 72Ge are displayed in
Fig. 3 as functions of $T$. These results show a trend of transition of the
level density from a convex function of $T$ to a concave function after the
pairing reentrance occurs. This is particularly clear for 60Ni by comparing
the FTBCS1 predictions for the level density at $M<$ 14 $\hbar$, which are
convex functions of $T$, with that obtained at $M=$ 14 $\hbar$, which is a
concave function of $T$. For 72Ge this trend is less obvious because of the
existence of proton and neutron pairing gaps with different values of $T_{c}$
within the FTBCS. However, contrary to the result shown in the inset of Fig. 4
in Ref. Dean , no pronounced local maximum that might correspond to the
pairing reentrance is seen here in the temperature dependence of the level
density for 72Ge. Since the results for the pairing gap, heat capacity, and
level density strongly depend on the selected single-particle energies, the
irregularity seen in the temperature dependence of the level density at
$\omega=$ 0.5 MeV in the inset of Fig. 4 of Ref. Dean might well be an
artifact caused by using the same single-particle energies for both neutron
and proton spectra.
Figure 4: (Color online) Level averaged neutron and proton gaps
$\bar{\Delta}_{N,Z}$, heat capacity $C$, $C/T$, excitation energy $E^{*}$ and
level density for 72Ge as functions of $T$ at several values of $M$ [shown in
(a)] obtained within the FTBCS1 in the test calculations by using the shells
$(0f1p-0g1d2s)$ atop the 40Ca core with the values of Woods-Saxon neutron
single-particle energies adopted for both neutron and proton spectra.
To show that it is indeed the case, we carried out the test calculations by
using only the $(0f1p-0g1d2s)$ shells on top of the 40Ca core with the values
of Woods-Saxon neutron single-particle energies adopted for both neutron and
proton spectra of 72Ge, as in Ref. Dean . The difference now solely comes from
that between the numbers of valence neutrons and protons (20 valence neutrons
and 12 valence protons). The results of these test calculations within the
FTBCS1 are shown in Fig. 4. They show the pairing reentrance in the neutron
pairing gap instead of the proton one, at low $M$. This is in qualitative
agreement with the pairing reentrance predicted for neutrons in Fig. 2 of Ref.
Dean , where no proton pairing reentrance is seen up to $\omega=$ 0.5 MeV. In
our calculations, however, the proton pairing reentrance takes place at rather
high $M=$ 22 $\hbar$. In either case, the pairing reentrance is so strong that
causes the excitation energy $E^{*}$ to decrease slightly with increasing $T$
at low $T$. This violates the second law of thermodynamics. As a consequence,
at $M\geq$ 6 $\hbar$, the heat capacity becomes negative at $T<$ 0.4 MeV. The
results for the level density remain essentially the same as compared to those
previously obtained by using all proton and neutron single-particle levels,
but with only bound-state single particle energies [Fig. 3 (d)], and no
irregularities such as a pronounced local maximum are found. In our opinion,
the pick on the dotted curve in the inset of Fig. 4 in Ref. Dean emerges
because of the two lower values of ${\rm ln}\rho$ at $T\simeq$ 0.42 and 0.45
MeV. These two lower values are the results obtained by calculating the level
density in the canonical ensemble $\rho(E)=\beta e^{S}/\sqrt{2\pi C}$, making
use of the two large values of $C$ equal to around 16 and 28 (with large error
bars). However, these large values of the heat capacity at low $T$ are the
artifacts of the SMMC calculations because the heat capacity must be zero (or
very small) at $T=$ 0 (or very low $T$) to avoid an infinite (or very large)
entropy, which would violate the third law of thermodynamics. Therefore, we
conclude that the neutron pairing reentrance effect, reported in Ref. Dean ,
is caused by the use of the same single-particle spectrum for both protons and
neutrons, whereas the irregularity seen on the curve of ${\rm ln}\rho$ in the
inset of Fig. 4 of Ref. Dean is caused by unphysically large values of the
heat capacity at low $T$ in the SMMC technique.
The present calculations within the FTBCS and FTBCS1 do not take into account
the effects of residual interactions beyond the monopole pairing one. It is
well known that these effects are responsible for strong collective motion in
finite nuclei, which leads to the increase of nuclear level density. In
spherical nuclei the collective enhancement of level density is caused by
vibrational excitations, whereas in deformed nuclei it comes from the
collective rotation. The contribution of collective motion to the increase of
nuclear level density has been studied in detail by Ignatyuk and collaborators
starting from the early 1970s Ignatyuk ; Junghans . Because the Hamiltonian
used in the SMMC calculations of Ref. Dean included the quadrupole-quadrupole
interaction, let us estimate the effect of the collective quadrupole vibration
on the increase of level density. By using the adiabatic approximation (3) of
Ref. Junghans for the enhancement coefficient $K_{vib}$ due to quadrupole
vibration, and the experimental energies $E(2^{+}_{1})$ of the lowest
quadrupole excitation in Ref. 1st2 , we found $K_{vib}\simeq$ 1.06 and 1.94
for 60Ni at $T=$ 0.3 MeV and 72Ge at $T=$ 0.4 MeV, respectively. These values
of T are those, at which the pairing reentrance starts to show up in these
nuclei. With the deformation parameter $\beta_{2}=$ -0.224 adopted in the
present calculations for 72Ge and by using Eq. (9) of Ref. Zagreb , we found
the enhancement coefficient $K_{coll}(\beta_{2})\simeq$ 1.96 for 72Ge at $T=$
0.4 MeV, whereas for the spherical nucleus, 60Ni, $K_{coll}=K_{vib}$ = 1.06 at
$T=$ 0.3 MeV. Therefore, for both nuclei, 60Ni and 72Ge, one can expect that
the collective quadrupole enhancement of level density is not dramatic at the
value of temperature, where the pairing reentrance is supposed to take place.
The contribution of collective motion generated by higher multipolarities to
the increase of level density is expected to be much smaller. In Ref. combi
the quasiparticle Tamm-Dancoff Approximation, which includes the isoscalar
quadrupole-quadrupole interaction and $J_{x}$ cranking, was used to calculate
the level density within the microcanonical ensemble. The authors of Ref.
combi found 4 $\leq K_{rot}\leq$ 6 and 1.002 $\leq K_{vib}\leq$ 1.012 at
excitation energy 3 $\leq E\leq$ 8 MeV (i.e. at around 0.38 $\leq T\leq$ 0.63
MeV) for 162Dy. They also found a monotonic decrease of the average pairing
gap with increasing the excitation energy up to $E=$ 8.5 MeV ($T\simeq$ 0.65
MeV), i.e. much higher than $T_{c}\simeq$ 0.34 MeV. These results are in good
qualitative agreement with our estimations.
Finally, it is worth noticing that the inclusion of the approximate particle-
number projection within the Lipkin-Nogami method does not significantly alter
the behavior of the paring reentrance obtained within the FTBCS1 theory (See
Fig. 6 of Ref. SCQRPA ). This does not diminish the value of an approach based
on exact particle-number and angular momentum projections. In Ref. Horoi the
exact solution of the nuclear shell model is used to study the SN phase
transition including residual interactions other than the pairing one. The
results of Ref. Horoi , which fully respect the particle number and angular
momentum conservations, confirm the presence of a long tail of pair
correlations far beyond the BCS phase transition region in agreement with the
prediction by the FTBCS1. The approach of Ref. Horoi does not use any
external heat bath, which determines the temperature of thermal equilibrium.
Therefore the nuclear temperature can only be extracted from the level density
by using the Clausius definition of thermodynamic entropy. This task is not
easy because of the discrete and finite nuclear spectra (See, e.g. Ref.
Ensemble and references therein). Nonetheless, instead of using temperature,
it would be interesting to see if the pairing reentrance takes place in the
pair correlator as a function of excitation energy at various values of
angular momentum within the method of Ref. Horoi .
## IV Conclusions
The present paper studies the temperature dependences of the heat capacity and
level density in hot medium-mass nuclei, which undergo a noncollective
rotation about the symmetry axis. The numerical calculations, carried out by
using the realistic Woods-Saxon single-particle energies for 60Ni and 72Ge
within the FTBCS and FTBCS1 theories, have shown the pairing reentrance in the
pairing gap at finite angular momentum $M$ and temperature $T$. Instead of
decreasing with increasing $T$, the gap first increases with $T$ then
decreases at higher $T$. It is demonstrated that the heat capacity $C$, or
rather $C/T$, and level density $\rho$ can be used to experimentally identify
the pairing reentrance effect. The pairing reentrance, when it occurs, leads
to a clear depletion in the temperature dependence of the heat capacity,
whereas the level density weakly changes from a convex function of $T$ to a
concave one.
Regarding the appearance of the local minimum in the heat capacity because of
the pairing reentrance, the results of the present paper agree with that of
the SMMC calculations in Ref. Dean . However, the present results show no
pronounced local maximum in the temperature dependence of the level density.
The pairing reentrance is seen in the proton pairing gap of 72Ge at low $M$,
whereas Ref. Dean reported this effect in the neutron pairing energy. The
test calculations by using the same single-particle configuration as that used
in Ref. Dean , but obtained within the Woods-Saxon potential, reveals that the
neutron pairing reentrance in 72Ge is an artifact, which is caused by the use
of the same single-particle spectrum for both protons and neutrons, whereas
the irregularity on the curve for the logarithm of level density, reported in
Ref. Dean , is caused by unphysically large values of the heat capacity at low
$T$ in the SMMC approach.
###### Acknowledgements.
The numerical calculations were carried out using the FORTRAN IMSL Library by
Visual Numerics on the RIKEN Integrated Cluster of Clusters (RICC) system. NQH
acknowledges the support by the National Foundation for Science and Technology
Development(NAFOSTED) of Vietnam through Grant No. 103.04-2010.02. He also
thanks the Theoretical Nuclear Physics Laboratory of RIKEN Nishina Center for
its hospitality during his visit in RIKEN.
## References
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|
arxiv-papers
| 2011-11-15T03:57:41 |
2024-09-04T02:49:24.336753
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "N. Quang Hung and N. Dinh Dang",
"submitter": "Nguyen Quang Hung",
"url": "https://arxiv.org/abs/1111.3425"
}
|
1111.3511
|
# Polygons of the Lorentzian plane and spherical simplexes
François Fillastre
University of Cergy-Pontoise
UMR CNRS 8088
Departement of Mathematics
F-95000 Cergy-Pontoise
FRANCE
francois.fillastre@u-cergy.fr
((v3) )
## 1 Introduction
It is a common occurence that sets of geometric objects themselves carry some
kind of geometric structure. A classical example for this is the set of all
conformal structures on a given compact surface. Riemann discovered that this
set, the “space” of conformal structures, can be described by a finite number
of parameters called moduli. The corresponding parameter or moduli space
turned out to be a very interesting geometric object in itself whose study is
the subject of Teichmüller theory.
On a more basic level, one can consider spaces consisting of objects of
elementary geometry like (shapes of) polyhedra in Euclidean space. Thurston
[Thu98] found that in this case, the corresponding moduli space carries the
structure of a complex hyperbolic manifold, and he established a link with
sets of triangulations of the 2-sphere.
Bavard and Ghys [BG92] considered sets of polygons in the Euclidean plane. Fix
a compact convex polygon $P$ with $n\geq 3$ edges and let ${\cal P}(P)$ be the
space of convex polygons with $n$ edges parallel to those of $P$. The elements
of ${\cal P}(P)$ are then determined by the distances of the lines containing
the edges from the origin, which gives $n$ parameters. Following [Thu98],
Bavard and Ghys proved that on the space of parameters, the area of the
polygons in ${\cal P}(P)$ is a quadratic form, and they computed its
signature. The kernel of the corresponding bilinear form has dimension 2 (due
to the fact that area is invariant under translations), and there is only one
positive direction. Hence, up to the kernel, one gets a Lorentzian signature.
As a consequence, the set of elements of ${\cal P}(P)$ with area equal to one,
considered up to translations, can be identified with a subset of the
hyperbolic space $\mathbb{H}^{n-3}$. This subset turns out to be a hyperbolic
convex polyhedron of a special kind: it is a simplex with the property that
each hyperplane containing a facet meets orthogonally all but two hyperplanes
containing the other facets. Such simplices are called hyperbolic
orthoschemes. The dihedral angles of the orthoscheme can be computed from the
angles of $P$, and [BG92] contains a list of convex polygons $P$ such that the
orthoscheme obtained from $P$ is of Coxeter type, i.e. has acute angles of the
form $\pi/k$, $k\in\mathbb{N}$. This list was previously known [IH85, IH90],
but it appeared it was incomplete [Fil11].
In this paper we consider a class of non-compact plane polygons whose moduli
space is a spherical orthoscheme. These polygons, the $t$-convex polygons
introduced in Section 3, are best described not in terms of the Euclidean
geometry on $\mathbb{R}^{2}$, but as subsets of the Lorentz plane. Instead of
the area we will consider a suitably defined coarea that turns out to be a
positive definite quadratic form on the parameter space, an $n$-dimensional
vector space. Restricting to coarea one we obtain a subset of the unit sphere
in that parameter space, and this subset is shown to be a spherical
orthoscheme. Moreover, any spherical orthoschem can be obtained in this way.
It is amusing that in [BG92] Euclidean polygons led to Lorentz metrics and
hyperbolic orthoschemes, while in the present paper Lorentzian polygons give
rise to Euclidean metrics and spherical orthoschemes. The author does not know
if there is a way to obtain Euclidean orthoschemes from spaces of plane convex
polygons.
## 2 Background on the Lorentz plane
Recall that the Lorentz plane is $\mathbb{R}^{2}$ equipped with the Lorentz
inner product, that is the bilinear form
$\langle\binom{x_{1}}{x_{2}},\binom{y_{1}}{y_{2}}\rangle_{1}=x_{1}y_{1}-x_{2}y_{2}.$
A non-zero vector $v$ can be _space-like_ ($\langle v,v\rangle_{1}>0$), _time-
like_ ($\langle v,v\rangle_{1}<0$) or _light-like_ ($\langle
v,v\rangle_{1}=0$). The set of time-like vectors has two connected components,
and we denote the upper one, the set of _future_ time-like vectors, by
${\mathcal{F}}:=\\{x\in{\mathbb{R}}^{2}|\langle x,x\rangle_{1}<0,x_{2}>0\\}.$
The set of unit future time-like vectors is
${\mathbb{H}}:=\\{x\in{\mathbb{R}}^{2}|\langle x,x\rangle_{1}=-1,x_{2}>0\\},$
which will be the analog of the circle in the Euclidean plane, see Figure 1.
In higher dimension, the generalization of ${\mathbb{H}}$ together with its
induced metric is a model of the hyperbolic space, in the same way that the
unit sphere for the Euclidean metric with its induced metric is a model of the
round sphere. In particular, if the angle between two unit vectors in the
Euclidean plane is seen as the distance between the two corresponding points
on the circle, the _(Lorentzian) angle_ between two future time-like vectors
$x$ and $y$ is the unique $\varphi>0$ such that
$\cosh\varphi=-\frac{\langle x,y\rangle_{1}}{\sqrt{\langle
x,x\rangle_{1}\langle y,y\rangle_{1}}}$ (1)
(see [Rat06, (3.1.7)] for the existence of $\varphi$). The angle $\varphi$ is
the distance on ${\mathbb{H}}$ (for the induced metric) between
$x/\sqrt{-\langle x,x\rangle_{1}}$ and $y/\sqrt{-\langle y,y\rangle_{1}}$.
Figure 1: The cone ${\mathcal{F}}$ of future time-like vectors and the curve
${\mathbb{H}}$ of unit future time-like vectors.
${\mathcal{F}}$ and ${\mathbb{H}}$ are globally invariant under the action of
the linear isometries of the Lorentzian plane, called _hyperbolic
translations_ :
$H_{t}:=\left(\begin{array}[]{cc}\cosh t&\sinh t\\\ \sinh t&\cosh
t\end{array}\right),t\in{\mathbb{R}}.$ (2)
In all the paper we fix a positive $t$. We denote by $<H_{t}>$ the free group
spanned by $H_{t}$.
## 3 $t$-convex polygons
Let $a\in{\mathcal{F}}$. We will denote by
$a^{\bot}:=\\{x\in{\mathbb{R}}^{2}|\langle x,a\rangle_{1}=\langle
a,a\rangle_{1}\\}$
the line that passes through $a$ and is parallel to the $1$-dimensional
subspace orthogonal to $a$ under $\langle\cdot,\cdot\rangle_{1}$.
###### Definition 3.1.
Let $(\eta_{1},\ldots,\eta_{n})$, $n\geq 1$, be pairwise distinct unit future
time-like vectors in the Lorentzian plane (i.e. $\eta_{i}\in{\mathbb{H}}$),
and let $h_{1},\ldots,h_{n}$ be positive numbers. A _$t$ -convex polygon_ $P$
is the intersection of the half-planes bounded by the lines
$(H_{t}^{k}(h_{i}\eta_{i}))^{\bot},\forall k\in\mathbb{Z},\forall
i=1,\ldots,n.$
The half-planes are chosen such that the vectors $\eta_{i}$ are inward
pointing. The positive numbers $h_{i}$ are the _support numbers_ of $P$.
A $t$-convex polygon is called _elementary_ if it is defined by a single
future time-like vector $\eta$ and a positive number $h$. Note that for each
$k$, $(H_{t}^{k}(h\eta))^{\bot}$ is tangent to $h{\mathbb{H}}$ (the upper
hyperbola with radius $h$). Hence a $t$-convex polygon is the intersection of
a finite number of elementary $t$-convex polygons.
###### Example 3.2.
Let $t_{0}=\sinh^{-1}(1)$, so
$H_{t_{0}}:=\left(\begin{array}[]{cc}\sqrt{2}&1\\\
1&\sqrt{2}\end{array}\right).$
Let us denote by $P_{1}$ the elementary $t_{0}$-convex polygon defined by the
vector $\eta=\binom{0}{1}$ and the number $h=1$, see Figure 2a. The elementary
$t_{0}$-convex polygon $P_{2}$ of Figure 2b is obtained from $p_{1}$ by a
slightly change of $\eta$ and $h$. Their intersection forms the $t_{0}$-convex
polygon of Figure 2c.
(a) A part of the $t_{0}$-convex polygon $P_{1}$. For the Lorentzian metric,
all the edges have equal length and all the angles between edges are equal.
(b) A part of the $t_{0}$-convex polygon $P_{2}$. For the Lorentzian metric,
all the edges have equal length and all the angles between edges are equal.
(c) A part of the $t_{0}$-convex polygon obtained as the intersection of
$P_{1}$ and $P_{2}$.
Figure 2: To Example 3.2.
###### Lemma 3.3.
A $t$-convex polygon $P$ is a proper convex subset of ${\mathbb{R}}^{2}$
contained in ${\mathcal{F}}$, bounded by a polygonal line with a countable
number of sides, and globally invariant under the action of $<H_{t}>$.
###### Proof.
The group invariance is clear from the definition. $P$ is the intersection of
a finite number of elementary $t$-convex polygons, so we only have to check
the other properties in the elementary case. Actually the only non-immediate
one is that an elementary $t$-convex polygon is contained in ${\mathcal{F}}$.
Let us consider an elementary $t$-convex polygon made from a single future
time-like vector $\eta$ and a number $h$. Without loss of generality, consider
that $h=1$. Let $u=H^{k}_{t}(\eta)$ and $v=H^{k^{\prime}}_{t}(\eta)$ and let
$x$ be the intersection between $u^{\bot}$ and $v^{\bot}$. As $\langle
x,u\rangle_{1}=\langle x,v\rangle_{1}=-1$, $x$ is orthogonal to $u-v$, which
is a space-like vector (compute its norm with the help of (1)). Hence $x$ is
time-like, and as $u^{\bot}$ and $v^{\bot}$ never meet the past cone, $x$ is
future. It is easy to deduce that the $t$-convex polygon is contained in
${\mathcal{F}}$. ∎
Note that as a convex surface, a $t$-convex polygon can also be a
$t^{\prime}$-convex polygon (for example it is also invariant under the action
of any subgroup of $<H_{t}>$), but we will only consider the action of a given
$<H_{t}>$.
Given a $t$-convex polygon $P$, we will require that the set of elementary
$t$-convex polygons such that their intersection gives $P$ is minimal, i.e
each $\eta_{i}$ is the inward unit normal of a genuine edge $e_{i}$ of $P$.
The edge at the left (resp. right) of $e_{i}$ is denoted by $e_{i-1}$ (resp.
$e_{i+1}$). Let $p_{i}$ be the foot of the perpendicular from the origin to
the line containing $e_{i}$ (in particular, $p_{i}=h_{i}\eta_{i}$). Let
$p_{ii+1}$ be the vertex between $e_{i}$ and $e_{i+1}$. We denote by
$h_{ii+1}$ (resp. $h_{ii-1}$) the signed distance from $p_{i}$ to $p_{ii+1}$
(resp. from $p_{i}$ to $p_{i-1i}$): it is non negative if $p_{i}$ is on the
same side of $e_{i+1}$ (resp. $e_{i-1}$) as $P$. The angle between $\eta_{i}$
and $\eta_{i+1}$ is denoted by $\varphi_{i}$. See Figure 3.
Figure 3: Notations for a $t$-convex polygon.
###### Lemma 3.4.
With the notations introduced above,
$h_{ii+1}=\frac{h_{i}\cosh\varphi_{i}-h_{i+1}}{\sinh\varphi_{i}},h_{ii-1}=\frac{h_{i}\cosh\varphi_{i-1}-h_{i-1}}{\sinh\varphi_{i-1}}.$
(3)
###### Proof.
By definition, $h_{ii+1}$ is non negative when $\langle
p_{i}-p_{i+1},\eta_{i+1}\rangle_{1}\leq 0$, i.e.
$-(h_{i+1}-h_{i}\cosh\varphi_{i})\geq 0.$
Hence
$h_{ii+1}=-\frac{h_{i+1}-h_{i}\cosh\varphi_{i}}{|h_{i+1}-h_{i}\cosh\varphi_{i}|}\sqrt{\langle
p_{ii+1}-p_{i},p_{ii+1}-p_{i}\rangle_{1}}.$
Up to an orientation and time orientation preserving linear isometry, one can
take $\eta_{i}=\left(0\atop 1\right)$. In particular $p_{i}=\left(0\atop
h_{i}\right)$ and $(p_{ii+1})_{2}=h_{i}$ hence
$\langle p_{ii+1}-p_{i},p_{ii+1}-p_{i}\rangle_{1}=(p_{ii+1})_{1}^{2}.$
We also have $\eta_{i+1}=\left(\sinh\varphi_{i}\atop\cosh\varphi_{i}\right)$,
and as $\langle p_{ii+1},\eta_{i+1}\rangle_{1}=-h_{i+1}$ we get
$(p_{ii+1})_{1}=\frac{-h_{i+1}+h_{i}\cosh\varphi_{i}}{\sinh\varphi_{i}}.$
The proof for $h_{ii-1}$ is similar, considering
$\eta_{i-1}=\left(-\sinh\varphi_{i}\atop\cosh\varphi_{i}\right)$. ∎
## 4 The cone of support vectors
Let $P$ be a $t$-convex polygon. Choose an edge and denote its inward unit
normal by $\eta_{1}$. We denote the inward unit normal of the edge on the
right by $\eta_{2}$, and so on until $\eta_{n+1}=H_{t}(\eta_{1})$. The edges
with normals $\eta_{1},\ldots,\eta_{n}$ are the _fundamental edges_ of $P$.
Note that with this labeling, if $\varphi_{i}$ is the angle between $\eta_{i}$
and $\eta_{i+1}$, we have
$\varphi_{1}+\varphi_{2}+\cdots+\varphi_{n}=t.$ (4)
The number $h_{i}(P)$ is the support number of the edge with normal
$\eta_{i}$, and $h(P)=(h_{1}(P),\ldots,h_{n}(P))$ is the _support vector_ of
$P$. So $P$ is identified with a vector of ${\mathbb{R}}^{n}$, in such a way
that $\eta_{1},\ldots,\eta_{n}$ are in bijection with the standard basis of
$\mathbb{R}^{n}$. Of course $P$ is uniquely determined by its support vector.
###### Definition 4.1.
Choose $\eta\in{\mathbb{H}}$ and let
$\varphi_{1},\varphi_{2},\cdots,\varphi_{n}$ be positive numbers satisfying
(4). The _cone of support vectors_
$\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ is the
set of support vectors of $t$-convex polygons with inward unit normals
$\eta_{1}=\eta$, $\eta_{i+1}=H_{\varphi_{i}}(\eta_{i})$.
A priori the definition of $\overline{\mathcal{P}}$ depends not only on the
angles $\varphi_{i}$ but also on the choice of $\eta$. Actually choosing
another starting $\eta^{\prime}\in{\mathbb{H}}$, the hyperbolic translation
from $\eta$ to $\eta^{\prime}$ gives a linear isomorphism between the two
resulting sets of support vectors. Hence
$\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ could be
defined as the set of $t$-convex polygons with ordered angles
$(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ up to hyperbolic translations.
Note also that if $s$ is a cyclic permutation, then
$\overline{\mathcal{P}}(\varphi_{s(1)},\ldots,\varphi_{s(n)})$ is the same as
$\overline{\mathcal{P}}(\varphi_{1},\ldots,\varphi_{n})$.
It is possible to prove that
$\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ is a
convex polyhedral cone with non-empty interior in ${\mathbb{R}}^{n}$, but this
will be easier after a suitable metrization of ${\mathbb{R}}^{n}$, that is the
subject of the next section.
## 5 Coarea
###### Definition 5.1.
Let $P\in\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$.
The _coarea_ of $P$ is
$\operatorname{coarea}(P)=\frac{1}{2}\sum_{i=1}^{n}h_{i}(P)\ell_{i}(P)$
where the sum is on the fundamental edges, and
$\ell_{i}(P)=h_{ii-1}(P)+h_{ii+1}(P)$ is the length of the $i$th fundamental
edge (hence positive).
Geometrically $\operatorname{coarea}(P)$ is the area (in the sense of the
Lebesgue measure) of a fundamental domain for the action of $H_{t}$ on the
complement of $P$ in ${\mathcal{F}}$. The main point is that hyperbolic
translations (2) have determinant $1$, so they preserve the area, which is
then independent of the choice of the fundamental domain, see Figure 4.
Moreover the area of a triangle with a space-like edge $e$ of length $l$ and
$0$ as a vertex has area $\frac{1}{2}lh$, if $h$ is the Lorentzian distance
between $0$ and the line containing $e$. (To see this, perform a hyperbolic
translation such that $e$ is horizontal and compute the area.) Note that the
coarea depends not only on the polygonal line $P$ but also on the group
$<H_{t}>$, so it would be more precise to speak about “$t$-coarea”, but as the
group is fixed from the beginning, no confusion is possible.
Figure 4: The two shaded regions have the same area. This area is the coarea
of the polygon.
For a given cone of support vectors, the coarea can be formally extended to
${\mathbb{R}}^{n}$ with the help of (3): for $h\in{\mathbb{R}}^{n}$,
$\operatorname{coarea}(h)=\frac{1}{2}\sum_{i=1}^{n}h_{i}\ell_{i}(h)$
with
$\ell_{i}(h):=\frac{h_{i}\cosh\varphi_{i-1}-h_{i-1}}{\sinh\varphi_{i-1}}+h_{i}\frac{h_{i}\cosh\varphi_{i}-h_{i+1}}{\sinh\varphi_{i}}.$
(5)
If $n=1$, there is only one angle between the unit inward normal $\eta$ and
its image under $H_{t}$, which is equal to $t$, and
$\operatorname{coarea}(h)=h^{2}\frac{\cosh t-1}{\sinh t}.$
If $n\geq 2$, we introduce the _mixed-coarea_
$\operatorname{coarea}(h,k)=\frac{1}{2}\sum_{i=1}^{n}h_{i}\frac{k_{i}\cosh\varphi_{i-1}-k_{i-1}}{\sinh\varphi_{i-1}}+h_{i}\frac{k_{i}\cosh\varphi_{i}-k_{i+1}}{\sinh\varphi_{i}},$
which is the polarization of the $\operatorname{coarea}$. Actually, it is
clearly a bilinear form, and
$\operatorname{coarea}(\eta_{k},\eta_{j})=\left\\{\begin{array}[]{ccc}0&\mbox{if
}&2\leq|j-k|\leq n+1\\\
\displaystyle{-\frac{1}{2}\frac{1}{\sinh\varphi_{k-1}}}&\mbox{if}&j=k-1\\\
\displaystyle{-\frac{1}{2}\frac{1}{\sinh\varphi_{k}}}&\mbox{if}&j=k+1\\\
\displaystyle{\frac{1}{2}\left(\frac{\cosh\varphi_{k-1}}{\sinh\varphi_{k-1}}+\frac{\cosh\varphi_{k}}{\sinh\varphi_{k}}\right)}&\mbox{if}&j=k\end{array}\right.$
(6)
so $\operatorname{coarea}$ is symmetric. We also obtain the following key
result.
###### Proposition 5.2.
The symmetric bilinear form $\operatorname{coarea}$ is positive definite.
###### Proof.
As $\cosh\varphi_{k}>1$, the matrix
$(\operatorname{coarea}(u_{k},u_{j}))_{kj}$ is strictly diagonally dominant,
and symmetric with positive diagonal entries, hence positive definite, see for
example [Var00, 1.22]. ∎
The Cauchy–Schwarz inequality applied to support vectors of $t$-convex
polygons gives the following _reversed Minkowski inequality_ :
###### Corollary 5.3.
Let $P,Q$ be $t$-convex polygons with parallel edges. Then
$\operatorname{coarea}(P,Q)^{2}\leq\operatorname{coarea}(P)\operatorname{coarea}(Q),$
with equality if and only if $P$ and $Q$ are homothetic:
$\exists\lambda>0,\forall i,h_{i}(P)=\lambda h_{i}(Q)$.
## 6 Spherical orthoschemes
$\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ is
clearly a cone in ${\mathbb{R}}^{n}$. Moreover it is the set of vectors of
positive edge lengths, for the edge lengths defined by (5). From the
definition of the coaera, for $h\in{\mathbb{R}}^{n}$,
$2\operatorname{coarea}(\eta_{i},h)=\ell_{i}(h)$, so $\eta_{i}$ is an inward
normal vector to the facet of $\overline{\mathcal{P}}$ defined by
$\ell_{i}=0$. So $\overline{\mathcal{P}}$ is polyhedral, and it is convex
because the $\eta_{i}$ form a basis of ${\mathbb{R}}^{n}$. Let us denote by
$\mathcal{P}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ the intersection of
$\overline{\mathcal{P}}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ with the
unit sphere of $({\mathbb{R}}^{n},\operatorname{coarea})$ (i.e. the set of
support vectors of $t$-convex polygons with coarea one). It follows that
$\mathcal{P}$ is a spherical simplex. If $n=1$, $\mathcal{P}$ is a point on a
line, so from now on assume that $n>1$.
When $n=2$, $\mathcal{P}$ is an arc on the unit circle with length $\theta$
satisfying
$\cos\theta=\frac{\sinh\varphi_{2}}{\sinh(\varphi_{1}+\varphi_{2})}.$
When $n=3$, $\mathcal{P}$ is a spherical triangle with acute inner angles,
whose cosines are given by:
$-\frac{\operatorname{coarea}(\eta_{k},\eta_{k+1})}{\sqrt{\operatorname{coarea}(\eta_{k},\eta_{k})}\sqrt{\operatorname{coarea}(\eta_{k+1},\eta_{k+1})}}=\sqrt{\frac{\sinh\varphi_{k-1}\sinh\varphi_{k+1}}{\sinh(\varphi_{k-1}+\varphi_{k})\sinh(\varphi_{k}+\varphi_{k+1})}}.$
(7)
When $n\geq 3$, from (6) we see that each facet has an acute interior dihedral
angle with exactly two other facets, and is orthogonal to the other facets.
Such spherical simplexes are called _acute spherical orthoschemes_. See
[Deb90, 5] for the history and main properties of these very particular
simplexes. Note that there are no spherical Coxeter orthoschemes, because the
Coxeter diagram of a spherical orthoscheme must be a cycle, and there is no
cycle in the list of Coxeter diagrams of spherical Coxeter simplexes. The list
can be found for example in [Rat06].
Let us denote by $U_{k}$ the line through $p_{k}$ (so the angle between
$U_{k}$ and $U_{k+1}$ is $\varphi_{k}$), and by $\lambda$ the cross ratio
$[U_{k-1},U_{k},U_{k+1},U_{k+2}]$, namely if $u_{k-1},u_{k},u_{k+1},u_{k+2}$
are the intersections of the lines $U_{i}$ with any line not passing through
zero and endowed with coordinates then (see [Ber94])
$\lambda=[U_{k-1},U_{k},U_{k+1},U_{k+2}]=\frac{u_{k+1}-u_{k-1}}{u_{k+1}-u_{k}}\frac{u_{k+2}-u_{k}}{u_{k+2}-u_{k-1}}.$
We have the formula (see [PY12])
$\frac{\sinh\varphi_{k-1}\sinh\varphi_{k+1}}{\sinh(\varphi_{k-1}+\varphi_{k})\sinh(\varphi_{k}+\varphi_{k+1})}=\frac{\lambda-1}{\lambda}=[U_{k-1},U_{k+2},U_{k},U_{k+1}].$
From a given $n$-dimensional acute spherical orthoscheme $O$ we can find
angles (positive real numbers) $(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$
such that $\mathcal{P}(\varphi_{1},\varphi_{2},\ldots,\varphi_{n})$ is
isometric to $O$. Let $0<A<1$ be the square of the cosine of an acute dihedral
angle of $O$. We have first to find ordered time-like lines
$U_{1},U_{2},U_{3},U_{4}$ such that $[U_{1},U_{2},U_{3},U_{4}]=\frac{1}{1-A}$,
i.e. we have to prove that the cross-ratio of the lines can reach any value
$>1$. Choose arbitrary distinct ordered time-like $U_{1},U_{2},U_{4}$. If
$U_{3}=U_{4}$ then $[U_{1},U_{2},U_{3},U_{4}]=1$, and if $U_{3}=U_{2}$ then
$[U_{1},U_{2},U_{3},U_{4}]=+\infty$, so by continuity any given value $>1$ can
be reached for a suitable $U_{3}$ between $U_{2}$ and $U_{4}$.
$U_{1},U_{2},U_{3},U_{4}$ give angles $\varphi_{1},\varphi_{2},\varphi_{3}$.
Now the other $\varphi_{k}$ are easily obtained as follows. Given the next
dihedral angle of $O$ (they can be ordered by ordering the unit normals to
$O$, see [Deb90]), the square of its cosine should be equal to
$\frac{\sinh\varphi_{2}\sinh\varphi_{4}}{\sinh(\varphi_{2}+\varphi_{3})\sinh(\varphi_{3}+\varphi_{4})}$
and $\varphi_{2},\varphi_{3}$ are known, so we get $\varphi_{4}$. And so on.
## 7 Spherical cone-manifolds
Let $n>2$ and consider the orthoscheme
$\mathcal{P}=\mathcal{P}(\varphi_{1},\ldots,\varphi_{n})$. A facet of
$\mathcal{P}$ is isometric to the space of $t$-convex polygons with
$\eta_{1},\ldots,\hat{\eta_{i}},\ldots,\eta_{n}$ ($\hat{\eta_{i}}$ means that
$\eta_{i}$ is deleted from the list) as normals to the fundamental edges. The
angles between the normals are
$\varphi_{1},\ldots,\varphi_{i-2},\varphi_{i-1}+\varphi_{i},\varphi_{i+1},\ldots,\varphi_{n}$.
This orthoscheme is also isometric to a facet of the orthoscheme
$\mathcal{P}^{\prime}$ obtained by permuting $\varphi_{i-1}$ and $\varphi_{i}$
in the list of angles. Hence we can glue $\mathcal{P}$ and
$\mathcal{P}^{\prime}$ isometrically along this common facet. We denote by
${\mathcal{C}}(\varphi_{1},\ldots,\varphi_{n})$ the $(n-1)$-dimensional
spherical cone-manifold obtained by gluing in this way all the $(n-1)!$
orthoschemes obtained by permutations of the list
$\varphi_{1},\ldots,\varphi_{n}$, up to cyclic permutations.
When $n=3$, ${\mathcal{C}}(\varphi_{1},\varphi_{2},\varphi_{3})$ is isometric
to a spherical cone-metric on the sphere with three conical singularities,
with cone-angles $<\pi$, obtained by gluing two isometric spherical triangles
along corresponding edges.
Let $n\geq 4$. Around the codimension 2 face of ${\mathcal{C}}$ isometric to
$N:=\mathcal{C}(\varphi_{1},\ldots,\varphi_{k}+\varphi_{k+1},\ldots,\varphi_{j}+\varphi_{j+1},\ldots,\varphi_{n+3})$
are glued four orthoschemes, corresponding to the four ways of ordering
$(\varphi_{k},\varphi_{k+1})$ and $(\varphi_{j},\varphi_{j+1}).$ As the
dihedral angle of each orthoscheme at such codimension $2$ face is $\pi/2$,
the total angle around $N$ in ${\mathcal{C}}$ is $2\pi$. Hence metrically $N$
is actually not a singular set. Around the codimension 2 face of
${\mathcal{C}}$ isometric to
$S:={\mathcal{C}}(\varphi_{1},\ldots,\varphi_{k}+\varphi_{k+1}+\varphi_{k+2},\ldots,\varphi_{n+3})$
are glued six orthoschemes corresponding to the six ways of ordering
$(\varphi_{k},\varphi_{k+1},\varphi_{k+2})$. Let $\Theta$ be the cone-angle
around $S$. It is the sum of the dihedral angles of the six orthoschemes glued
around it. As formula (7) is symmetric for two variables, $\Theta$ is two
times the sum of three different dihedral angles. A direct computation gives
($k=1$ in the formula)
$\cos(\Theta/2)=\textstyle\frac{\sinh\varphi_{1}\sinh\varphi_{2}\sinh\varphi_{3}-\sinh(\varphi_{1}+\varphi_{2}+\varphi_{3})(\sinh\varphi_{1}\sinh\varphi_{2}+\sinh\varphi_{2}\sinh\varphi_{3}+\sinh\varphi_{3}\sinh\varphi_{1})}{\sinh(\varphi_{1}+\varphi_{2})\sinh(\varphi_{2}+\varphi_{3})\sinh(\varphi_{3}+\varphi_{1})}.$
During the computation we used that
$\sinh(a+b)\sinh(b+c)-\sinh a\sinh c=\sinh b\sinh(a+b+c)$
which can be checked with $\frac{1}{2}\left(\cosh(x+y)-\cosh(x-y)\right)=\sinh
x\sinh y$. The analogous formula in the Euclidean convex polygons case was
obtained in [KNY99].
For example when $\varphi_{i}=\varphi\;\forall i$, we have
$\cos(\Theta/2)=-\frac{2\cosh(\varphi)^{2}+\sinh(\varphi)^{2}}{2\cosh(\varphi)^{3}}.$
The function on the right-hand side is a bijection from the positive numbers
to $]-1,0[$, hence all the $\Theta\in]2\pi,3\pi[$ (the dihedral angle
$\theta\in]\pi/3,\pi/2[$) are uniquely reached. In particular ${\mathcal{C}}$
is not an orbifold.
The cone-manifold ${\mathcal{C}}$ comes with an isometric involution which
consists of reversing the order of the angles
$(\varphi_{1},\ldots,\varphi_{n})$.
## 8 Higher dimensional generalization
The generalization of $t$-convex polygons to higher dimensional Minkowski
spaces is as follows. Let us consider the $d$-dimensional hyperbolic space
${\mathbb{H}}^{d}$ as a pseudo-sphere in the $d+1$-dimensional Minkowski space
$M^{d+1}$, and let $\Gamma$ be a discrete group of linear isometry of
$M^{d+1}$ such that ${\mathbb{H}}^{d}/\Gamma$ is a compact hyperbolic
manifold. A $\Gamma$-convex polyhedron is, given
$\eta_{1},\ldots,\eta_{n}\in{\mathbb{H}}^{d}$ and positive numbers
$h_{1},\ldots,h_{n}$, the intersection of the future sides of the space-like
hyperplanes $(\gamma(h_{i}\eta_{i}))^{\bot}$ $\forall
i,\forall\gamma\in\Gamma$. The mixed-coarea is generalized as a “mixed
covolume”. For details and computation of the signature, see [Fil13]. Actually
for a given set of $\eta_{i}$, many combinatorial types may appear, and one
has to restrict to type cones (cones of polyhedra with parallel facets and
same combinatorics). It should be interesting to investigate the kind of
spherical polytopes that appear.
Another related question is to look at the quadratic form given by the face
area of the polyhedra (in a fundamental domain) and its relations with the
moduli spaces of flat metric with conical singularities of negative curvature
on compact surfaces of genus $>1$ (the quotient of the boundary of a
$\Gamma$-convex polyhedron is isometric to such a metric).
The analogous questions in the convex polytopes case are the subject of
[FI13]. The moduli space of flat metrics on the sphere was studied in [Thu98].
## Acknowledgement
The author thanks anonymous referee and Haruko Nishi who helped to imporve the
redaction of the present text. Up to trivial changes, the introduction was
written by an anonymous referee. The polygons introduced in the present paper
are very particular cases of objects studied in [Fil13] and [FV13].
Work supported by the ANR GR Analysis-Geometry.
## References
* [Ber94] M. Berger. Geometry. I. Universitext. Springer-Verlag, Berlin, 1994. Translated from the 1977 French original by M. Cole and S. Levy, Corrected reprint of the 1987 translation.
* [BG92] C. Bavard and É. Ghys. Polygones du plan et polyèdres hyperboliques. Geom. Dedicata, 43(2):207–224, 1992.
* [Deb90] H. E. Debrunner. Dissecting orthoschemes into orthoschemes. Geom. Dedicata, 33(2):123–152, 1990.
* [FI13] F. Fillastre and I. Izmestiev. Shapes of polyhedra, mixed volumes, and hyperbolic geometry. In preparation, 2013.
* [Fil11] F. Fillastre. From spaces of polygons to spaces of polyhedra following Bavard, Ghys and Thurston. Enseign. Math. (2), 57(1-2):23–56, 2011.
* [Fil13] F. Fillastre. Fuchsian convex bodies: basics of Brunn–Minkowski theory. To appear _Geometric and Functional Analysis_ , 2013.
* [FV13] F. Fillastre and G. Veronelli. Lorentzian area measures and the Christoffel problem. 2013\.
* [IH85] H.-C. Im Hof. A class of hyperbolic Coxeter groups. Exposition. Math., 3(2):179–186, 1985.
* [IH90] H.-C. Im Hof. Napier cycles and hyperbolic Coxeter groups. Bull. Soc. Math. Belg. Sér. A, 42(3):523–545, 1990. Algebra, groups and geometry.
* [KNY99] S. Kojima, H. Nishi, and Y. Yamashita. Configuration spaces of points on the circle and hyperbolic Dehn fillings. Topology, 38(3):497–516, 1999.
* [PY12] A. Papadopoulos and S. Yamada. A Remark on the Projective Geometry of Constant Curvature Spaces. 2012\.
* [Rat06] J. Ratcliffe. Foundations of hyperbolic manifolds, volume 149 of Graduate Texts in Mathematics. Springer, New York, second edition, 2006.
* [Thu98] W. P. Thurston. Shapes of polyhedra and triangulations of the sphere. In The Epstein birthday schrift, volume 1 of Geom. Topol. Monogr., pages 511–549 (electronic). Geom. Topol. Publ., Coventry, 1998. Circulated as a preprint since 1987.
* [Var00] R. Varga. Matrix iterative analysis, volume 27 of Springer Series in Computational Mathematics. Springer-Verlag, Berlin, expanded edition, 2000.
|
arxiv-papers
| 2011-11-15T12:01:06 |
2024-09-04T02:49:24.346896
|
{
"license": "Public Domain",
"authors": "Fran\\c{c}ois Fillastre",
"submitter": "Fran\\c{c}ois Fillastre",
"url": "https://arxiv.org/abs/1111.3511"
}
|
1111.3569
|
# Dynamics of Hot Accretion Flow with Thermal Conduction
Kazem Faghei
School of Physics, Damghan University, Damghan, Iran
E-mail:kfaghei@du.ac.ir
###### Abstract
The purpose of this paper is to explore the dynamical behaviour of hot
accretion flow with thermal conduction. The importance of thermal conduction
on hot accretion flow is confirmed by observations of the hot gas that
surrounds Sgr A∗ and a few other nearby galactic nuclei. In this research, the
effect of thermal conduction is studied by a saturated form of it, as is
appropriate for weakly collisional systems. The angular momentum transport is
assumed to be a result of viscous turbulence and the $\alpha$-prescription is
used for the kinematic coefficient of viscosity. The equations of accretion
flow are solved in a simplified one-dimensional model that neglects the
latitudinal dependence of the flow. To solve the integrated equations that
govern the dynamical behaviour of the accretion flow, we have used an unsteady
self-similar solution. The solution provides some insights into the dynamics
of quasi-spherical accretion flow and avoids from limits of the steady self-
similar solution. In comparison to accretion flows without thermal conduction,
the disc generally becomes cooler and denser. These properties are
qualitatively consistent with performed simulations in hot accretion flows.
Moreover, the angular velocity increases with the magnitude of conduction,
while the radial infall velocity decreases. The mass accretion rate onto the
central object is reduced in the presence of thermal conduction. We found that
the viscosity and thermal conduction have the opposite effects on the physical
variables. Furthermore, the flow represents a transonic point that moves
inward with the magnitude of conduction or viscosity.
###### keywords:
accretion, accretion discs, conduction – hydrodynamics.
††pubyear: 2011
## 1 Introduction
Accretion is one of the most important physical processes in astrophysics. It
is widely accepted that the accreting matter toward the central object is a
source of power active galactic nuclei (AGN), and galactic X-ray sources (see
for a review e.g. Frank et al. 2002). This idea is also well-applicable to
interpret many observations of astrophysical phenomena, such as, prototype
stellar objects (Chang & Choi 2002), symbiotic stars (Lee & Park 1999), gamma-
ray bursts (Brown et al. 2000). There are some types of accretion flows that a
significant fraction of the generated heat by dissipation processes retains in
the fluid rather than being radiated away, have been the subject of
considerable attention in recent years (Ogilvie 1999; Chang 2005, Akizuki &
Fukue 2006, Khesali & Faghei 2009, Faghei 2011). These advection dominated
accretion flows (ADAF) place an intermediate position between the spherically
symmetric accretion flow of non-rotating fluid (Bondi 1952) and the cool thin
disc of classical accretion disc theory (e. g. Pringle 1981). These types of
accretion flows have been widely applied to explain observations of the
galactic black hole candidate (e. g. Narayan et al. 1996, Hameury et al.
1997), the spectral transition of Cyg x-1 (Esin 1996) and multi-wavelength
spectral properties of Sgr A∗ (Narayan & Yi 1995; Manmoto et al. 2000; Narayan
et al. 1997).
The X-ray observations of black holes imply that they are capable of accreting
gas under a variety of flow configurations. In particular, observational
evidences confirm existence of hot accretion flow, contrasted with the
classical cold and thin accretion disc scenario (Shakura & Sunyaev 1973). Hot
accretion flow can be found in the population of supermassive black holes in
galactic nuclei and during quiescent of accretion onto stellar-mass black
holes in X-ray transients (e.g., Narayan et al. 1998a; Lasota et al. 1996; Di
Matteo et al. 2000; Esin et al. 1997, 2001; Menou et al. 1999; see Narayan et
al. 1998b; Melia & Falcke 2001; Narayan 2002; Narayan & Quataert 2005 for
reviews). Chandra observations provide tight constraints on the density and
temperature of gas at or near the Bondi capture radius in Sgr A∗ and several
other nearby galactic nuclei. Tanaka & Menou (2006) used these constraints
(Loewenstein et al. 2001; Baganoff et al. 2003; Di Matteo et al. 2003; Ho et
al. 2003) to calculate the mean free path for the observed gas. They suggested
that accretion in these systems will be proceeded under the weakly collisional
condition. Furthermore, they suggested that thermal conduction can be as a
possible mechanism by which the sufficient extra heating is provided in hot
advection dominated accretion flows.
Generally, semi-analytical studies of hot accretion flows with thermal
conduction have been related to steady state models (e. g. Tanaka & Menou
2006; Johnson & Quataert 2007; Shadmehri 2008; Abbassi et al. 2008, 2010;
Ghanbari et al. 2009), and dynamics of such systems have been studied in
simulation models (e. g. Sharma et al. 2008; Wu et al. 2010). For example,
Tanaka & Menou (2006) have carried out a related analysis and found the
accretion flow can spontaneously produce thermal outflows driven in part by
conduction. Their analysis is two-dimensional but self-similar in radius.
Their assumption of self-similarity enforces a density profile that varies as
$r^{-3/2}$, whereas simulations of ADAFs consistently find density profiles
shallower than this (e.g., Stone et al. 1999; Igumenshchev & Abramowicz 1999;
Stone & Pringle 2001; Hawley & Balbus 2002; Igumenshchev et al. 2003). Johnson
& Quataert (2007) studied the effects of electron thermal conduction on the
properties of hot accretion flows, under the assumption of spherical symmetry.
Since, electron heat conduction is important for low accretion rate systems,
thus their model is applicable for Sgr A∗ in the Galactic centre. They show
that heat conduction leads to supervirial temperatures, implying that
conduction significantly modifies the structure of the accretion flow. Their
model similar to Tanaka & Menou (2006) was the steady state, but they solved
their equations numerically.
As mentioned, semi-analytical studies of hot accretion flows with thermal
conduction have been in a steady state. Thus, it will be interesting to study
dynamics of such systems. Ogilvie (1999) by the unsteady self-similar method
studied time-dependence of quasi-spherical accretion flow without thermal
conduction. The solutions of Ogilvie (1999) provided some insight into the
dynamics of quasi-spherical accretion and avoided many of the limits of the
steady self-similar solution. In this research, we want to explore how thermal
conduction can affect the dynamics of a rotating and accreting viscous gas. We
answer this question by solving Ogilvie (1999) model that is affected by
thermal conduction. This paper is organized as follows. In Section 2, we
define the general problem of constructing a model for hot quasi-spherical
accretion flow. In Section 3, we use the unsteady self-similar method to solve
the integrated equations that govern the dynamical behaviour of the accreting
gas, and numerical study of the model is brought in this section, too. We will
present a summary of the model in Section 4.
## 2 Basic Equations
We start with the approach adopted by Ogilvie (1999), who studied quasi-
spherical accretion flows without thermal conduction. Thus, we derive the
basic equations that describe the physics of accretion flow with thermal
conduction. We use the spherical coordinates $(r,\theta,\phi)$ centred on the
accreting object and make the following standard assumptions:
1. 1.
The gravitational force on a fluid element is characterized by the Newtonian
potential of a point mass, $\Psi=-GM_{*}/r$, with $G$ representing the
gravitational constant and $M_{*}$ standing for the mass of the central star.
2. 2.
The written equations in spherical coordinates are considered in the
equatorial plane $\theta=\pi/2$ and terms with any $\theta$ and $\phi$
dependence are neglected, hence all quantities will be expressed in terms of
spherical radius $r$ and time $t$.
3. 3.
For simplicity, self-gravity and general relativistic effects have been
neglected.
Under these assumptions, the dynamics of accretion flow describes by the
following equations:
the continuity equation
$\frac{\partial\rho}{\partial t}+\frac{1}{r^{2}}\frac{\partial}{\partial
r}(r^{2}\rho v_{r})=0,$ (1)
the radial force equation
$\frac{\partial v_{r}}{\partial t}+v_{r}\frac{\partial v_{r}}{\partial
r}=r(\Omega^{2}-\Omega_{K}^{2})-\frac{1}{\rho}\frac{\partial p}{\partial r},$
(2)
the azimuthal force equation
$\rho\left[\frac{\partial}{\partial
t}(r^{2}\Omega)+v_{r}\frac{\partial}{\partial
r}(r^{2}\Omega)\right]=\frac{1}{r^{2}}\frac{\partial}{\partial r}\left[\nu\rho
r^{4}\frac{\partial\Omega}{\partial r}\right],$ (3)
the energy equation
$\displaystyle\frac{1}{\gamma-1}\left[\frac{\partial p}{\partial
t}+v_{r}\frac{\partial p}{\partial
r}\right]+\frac{\gamma}{\gamma-1}\frac{p}{r^{2}}\frac{\partial}{\partial
r}\left(r^{2}v_{r}\right)=$ $\displaystyle Q_{vis}-Q_{rad}+Q_{cond}.$ (4)
Here $\rho$ the density, $v_{r}$ the radial velocity, $\Omega$ the angular
velocity, $\Omega_{K}[=(GM_{*}/r^{3})^{1/2}]$ Keplerian angular velocity, $p$
the gas pressure, $\gamma$ is the adiabatic index, $\nu$ the kinematic
viscosity coefficient and it is given as in Narayan & Yi (1995a) by an
$\alpha$-model
$\nu=\alpha\frac{p_{gas}}{\rho\Omega_{K}}.$ (5)
The parameter of $\alpha$ is assumed to be a constant less than unity. The
terms on the right-hand side of the energy equation, $Q_{vis}$ is the heating
rate of the gas by the viscous dissipation, $Q_{rad}$ represents the energy
loss through radiative cooling, and $Q_{cond}$ is the transported energy by
thermal conduction. For the right-hand side of the energy equation, we can
write
$Q_{adv}=Q_{vis}-Q_{rad}+Q_{cond}$ (6)
where $Q_{adv}$ is the advective transport of energy. We employ the advection
factor, $f=1-Q_{rad}/Q_{vis}$, that describes the fraction of the dissipation
energy which is stored in the accretion flow and advected into the central
object rather than being radiated away. The advection factor of $f$ in general
depends on the details of the heating and radiative cooling mechanism and will
vary with position (e.g. Watari 2006, 2007; Sinha et al. 2009). However, we
assume a constant $f$ for simplicity. Clearly, the case $f=1$ corresponds to
the extreme limit of no radiative cooling and in the limit of efficient
radiative cooling, we have $f=0$.
In a collisional plasma, mean free path for electron energy exchange,
$\lambda$, is shorter than temperature scale height, $L_{T}=T/|\nabla T|$, and
thus the heat flux due to thermal conduction can be written as
$F_{cond}=-\kappa\nabla T,$ (7)
where $\kappa$ is the thermal conductivity coefficient. Thermal conductivity
in a dense, fully ionized gas is given by the Spitzer (1962) formula,
$\kappa=\frac{1.84\times 10^{-5}T_{e}^{5/2}}{\ln\Lambda},$ (8)
where $T_{e}$ is the electron temperature ($T_{e}=T$ for a one-temperature
plasma) and $\ln\Lambda$ is Coulomb logarithm that for $T>4.2\times 10^{5}K$
is
$\ln\Lambda=29.7+\ln n^{-1/2}(T_{e}/10^{6}K).$ (9)
The heat is conducted by the electron, and equation (8) includes the effect of
the self-consistent electric required to maintain the electric current at
zero; this reduces $\kappa$ by a factor of about $0.4$ from the value it would
otherwise have (Cowie & McKee 1977, hereafter CM77).
As noted in the introduction, the inner regions of hot accretion flows are, in
many cases, collisionless with electron mean free path due to Coulomb
collision larger than the radius (e. g. Tanaka & Menou 2006). When the mean
free path of an electron becomes comparable to or larger than the temperature
gradient scale $\lambda\gtrsim T/|\nabla T|$, equation (7) for the heat flux
is no longer valid; CM77 described this effect as saturation. The maximum heat
flux in a plasma can be expressed as $(3/2)n_{e}kT_{e}v_{char}$, where
$v_{char}$ is a characteristic velocity which one might expect to be the order
of the electron thermal velocity (Parker 1963). Assuming Maxwellian
distribution for heat source, the characteristic velocity can be written as
(Williams 1971; CM77)
$V_{char}=(\frac{8}{9\,\pi})^{1/2}\,(\frac{k\,T}{m_{e}})^{1/2}.$ (10)
Similar to CM77, we assume that the heat flux is reduced by the same factor of
$0.4$ in the saturated case as in the classical (collisional) case so that the
saturated heat flux is
$F_{sat}=0.4\,n_{e}kT_{e}\,\sqrt{\frac{2kT_{e}}{\pi m_{e}}}.$ (11)
CM77 showed that the saturated heat flux is significantly less than
conjectured by Parker (1963). Thus, in order to explicitly allow for
uncertainty in the estimate of $F_{sat}$, they introduced a factor of
$\phi_{s}$, which was less than unity and rewrote equation (11) as
$F_{sat}=5\phi_{s}\rho c_{s}^{3}=5\phi_{s}p\sqrt{\frac{p}{\rho}},$ (12)
where $c_{s}$ is sound speed, which is defined as $c_{s}^{2}=p/\rho$. The
factor of $\phi_{s}$ is called as saturation constant (CM77). Now, the viscous
heating rate and the energy transport by thermal conduction are expressed as
$Q_{vis}=\nu\rho r^{2}\left(\frac{\partial\Omega}{\partial r}\right)^{2}$ (13)
$Q_{cond}=-\frac{1}{r^{2}}\frac{\partial}{\partial
r}\left(r^{2}F_{sat}\right)$ (14)
By using equations (13) and (14) for the advective transport of energy, we can
write
$Q_{adv}=f\nu\rho r^{2}\left(\frac{\partial\Omega}{\partial
r}\right)^{2}-\frac{1}{r^{2}}\frac{\partial}{\partial
r}\left(r^{2}F_{sat}\right)$ (15)
The mass accretion rate in a qausi-spherical accretion flow can be written as
$\dot{M}(r,t)=-4\pi r^{2}\rho v_{r}.$ (16)
We will use this quantity in the next section, and will investigate effects of
saturation constant and viscous parameter on it.
## 3 Self-Similar Solutions
### 3.1 analysis
Tanaka & Menou (2006) solved essentially equations (1)-(4) for the case of a
steady, radially self-similar flow. Here, we will try to find unsteady self-
similar solutions for these equations. Thus, we introduce a similarity
variable $\xi$ and assume that each physical quantity is given by the
following form:
$\xi=r(GM_{*}t^{2})^{-1/3},$ (17)
$\rho(r,t)=R(\xi)(\dot{M}_{0}/GM_{*})t^{-1},$ (18)
$p(r,t)=\Pi(\xi)(\dot{M}_{0}/(GM_{*})^{1/3})t^{-5/3},$ (19)
$v_{r}(r,t)=V(\xi)(GM_{*})^{1/3}t^{-1/3},$ (20)
$\Omega(r,t)=\omega(\xi)t^{-1},$ (21) $\dot{M}(r,t)=\dot{M}_{0}\dot{m}(\xi),$
(22)
where $\dot{M}_{0}$ is a constant and its value can be obtained by typical
values of the system. In addition, we assumed that $\dot{M}(r,t)$ under
similarity transformations is a function of $\xi$ only (Khesali & Faghei 2008,
2009). Substitution of above transformations into the basic equations (1)-(4),
yields dimensionless equations below,
$\left(V-\frac{2\xi}{3}\right)\frac{dR}{d\xi}-R=-\frac{R}{\xi^{2}}\frac{d}{d\xi}\left(\xi^{2}V\right),$
(23)
$\displaystyle\left(V-\frac{2\xi}{3}\right)\frac{dV}{d\xi}-\frac{V}{3}=\xi(\omega^{2}-\xi^{-3})-\frac{1}{R}\frac{d\Pi}{d\xi},$
(24) $\displaystyle
R\left[\left(V-\frac{2\xi}{3}\right)\frac{d}{d\xi}\left(\xi^{2}\omega\right)+\frac{1}{3}\left(\xi^{2}\omega\right)\right]~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}$
$\displaystyle=\frac{\alpha}{\xi^{2}}\frac{d}{d\xi}\left[\Pi\xi^{11/2}\frac{d\omega}{d\xi}\right],$
(25)
$\displaystyle\frac{1}{\gamma-1}\left[\left(V-\frac{2\xi}{3}\right)\frac{d\Pi}{d\xi}-\frac{5}{3}\Pi\right]+\frac{\gamma}{\gamma-1}\frac{\Pi}{\xi^{2}}\frac{d}{d\xi}\left(\xi^{2}V\right)$
$\displaystyle=\alpha
f\Pi\xi^{7/2}\left(\frac{d\omega}{d\xi}\right)^{2}-\frac{5\phi_{s}}{\xi^{2}}\frac{d}{d\xi}\left(\xi^{2}\Pi\sqrt{\frac{\Pi}{R}}\right).$
(26)
These equations provide a fourth-order system of non-linear ordinary
differential equations that must be solved numerically.
### 3.2 Inner limit
An appropriate asymptotic solution as $\xi\rightarrow 0$ is the form as
$R(\xi)\sim\xi^{-3/2}(R_{0}+R_{1}\xi+\cdot\cdot\cdot),$ (27)
$\Pi(\xi)\sim\xi^{-5/2}(\Pi_{0}+\Pi_{1}\xi+\cdot\cdot\cdot),$ (28)
$V(\xi)\sim\xi^{-1/2}(V_{0}+V_{1}\xi+\cdot\cdot\cdot),$ (29)
$\omega(\xi)\sim\xi^{-3/2}(\omega_{0}+\omega_{1}\xi+\cdot\cdot\cdot),$ (30)
in which undetermined coefficients of $R_{0}$, $R_{1}$, $\Pi_{0}$, and etc
must be specified. By substituting above relations in equations (23)-(26) and
choosing the significant sentences, we can write
$V_{0}^{2}+\frac{5\Pi_{0}}{R_{0}}-2+2\,\omega_{0}^{2}\approx 0,$ (31)
$2\,R_{0}\,V_{0}+3\,\alpha\,\Pi_{0}\approx 0,$ (32)
$6\,V_{0}\,(\gamma-\frac{5}{3})+(\gamma-1)\left[20\phi_{s}\sqrt{\frac{\Pi_{0}}{R_{0}}}-9\alpha
f\omega_{0}^{2}\right]\approx 0.$ (33)
Also, the dimensionless mass accretion rate, $\dot{m}(\xi)$, under the above
asymptotic solution becomes
$\dot{m}_{in}\approx-4\pi R_{0}V_{0},$ (34)
where $\dot{m}_{in}$ is the value of $\dot{m}$ at $\xi_{in}$, where $\xi_{in}$
is a point near to the centre. After algebraic manipulations for equations
(31)-(34), we obtain an algebraic equation for $R_{0}$:
$\displaystyle R_{0}^{2}-\frac{10}{27}\frac{\phi_{s}}{\alpha
f}\sqrt{\frac{6\dot{m}_{in}}{\alpha\pi}}\,R_{0}^{3/2}-\frac{5}{12}\frac{\dot{m}_{in}}{\alpha\pi}\times$
$\displaystyle\left(1-\frac{2}{5f}\frac{\gamma-5/3}{\gamma-1}\right)R_{0}-\frac{1}{32}\left(\frac{\dot{m}_{in}}{\pi}\right)^{2}\approx
0,$ (35)
and the rest of the physical variables are
$\Pi_{0}\approx\frac{\dot{m}_{in}}{6\pi\alpha},$ (36)
$V_{{0}}\approx-\frac{\dot{m}_{in}}{4\pi R_{0}},$ (37) $\omega_{0}^{2}\approx
1-\frac{5}{12}\frac{\dot{m}_{in}}{\pi\alpha
R_{0}}\left(1+\frac{3}{40}\frac{\alpha\dot{m}_{in}}{\pi R_{0}}\right).$ (38)
Without thermal conduction, $\phi_{s}=0$, equation (35) can be solved
analytically. Since, we want to consider systems with non-zero saturation
constant, $\phi_{s}\neq 0$, we will solve this equation numerically.
Figure 1: Time-dependent self-similar solution for $\gamma=1.3$, $\alpha=0.1$,
$f=1.0$, and $\dot{m}_{in}=0.001$. The solid, the dashed, and the short-dashed
lines represent $\phi_{s}=0$, $0.05$, and $0.1$, respectively.
Figure 2: Same as Figure 1, but $\phi_{s}=0.03$. The solid, dahsed, and short-
dashed lines represent $\alpha=0.05$, $0.1$, and $0.2$, respectively.
### 3.3 Numerical solution
If the value of $\xi_{in}$ is guessed, i.e. we take a point very near to the
centre, the equations (23)-(26) by Runge-Kutta-Fehlberg fourth-fifth order
method can be integrated from this point outward by the above expansion
[(27)-(30)]. Examples of such solutions are presented in Figs 1-5.
#### 3.3.1 The influences of saturation constant and viscous parameter on
physical quantities
The delineated quantity of $\Pi/R$ in Figs 1 and 2 is the sound speed square
in self-similar flow, which is rescaled in the course of time and represents
the flow temperature. The profiles of $\Pi/R$ in Fig. 1 show the flow
temperature decreased by adding saturation constant, $\phi_{s}$. Because, the
generated heat by viscous dissipation can be transfered by thermal conduction.
Furthermore, this temperature decrease is qualitatively consistent with
simulation results of Sharma et al. (2008) and Wu et al. (2010). We know the
viscous dissipation of the flow increases by adding $\alpha$ parameter. Thus,
the temperature increased by viscous parameter confirmed by the temperature
profiles in Fig 2. The density profiles show the gas density increased by
adding the $\phi_{s}$ parameter. It can be due to temperature fall of fluid.
Increasing density by adding saturation constant is another consistency of our
results with simulations of Wu et al. (2010). Also, the density profiles show
that it decreases by adding the viscous parameter. This also could be result
of the temperature rising. The viscous turbulence in this paper is
proportional to the gas temperature ($\nu\propto c_{s}^{2}\propto T$). Thus,
increase or decrease of temperature will affect dynamics of the accreting gas.
As we said the temperature decreased by saturation constant implies that the
viscous turbulence decreases, too. The decreasing of viscous turbulence
reduces the effect of negative viscous torque in angular momentum equation.
Thus, we expect the flow rotates faster by adding saturation constant that
confirmed by the angular velocity profiles in Fig. 1. Since, the efficiency of
the angular momentum transport decreases by adding the saturation constant, we
expect the decrease of radial infall velocity that the radial velocity
profiles in Fig. 1 confirm it. The efficiency of angular momentum transport
increases by adding the viscous parameter of $\alpha$. Thus, we expect the
flow rotates slower and accretes faster by adding the viscous parameter. The
profiles of radial and angular velocities in Fig. 2 confirm them.
Figure 3: Time-dependent self-similar solution of mass accretion rate. The
input parameters in left panel are same as Figure 1, but the solid, the
dashed, and the short-dashed lines represent $\phi_{s}=0.01$, $0.02$, and
$0.03$, respectively. The input parameters in right panel are same as Figure
2. but $\phi_{s}=0.01$.
#### 3.3.2 Mass accretion rate
The behaviour of mass accretion rate as a function of similarity variable
$\xi$ for several values of the viscous parameter and saturation constant are
plotted in Fig. 3. In the present model, the mass accretion rate is reduced by
radius. While, the mass accretion rate in steady hot accretion flows is a
constant (Tanaka & Menou 2006). There are some researches in steady hot
accretion flows that have studied power-law function of mass accretion rate
(Shadmehri 2008; Abbassi et al. 2008). However, the mass accretion rate in
their models is not dependent on important parameters such as saturation
constant and viscous parameter. The profiles of mass accretion rate in Fig. 3
show that it is reduced by adding the saturation constant. This property is
qualitatively consistent with numerical results of Johnson & Quataert (2007).
Also, the profiles of mass accretion rate imply that it increases by adding
the viscous parameter of $\alpha$. This property is qualitatively consistent
with previous works in accretion flows (e. g. Park 2009).
Figure 4: Time-dependent self-similar solution of Mach number. The input
parameters in left panel are same as Figure 1, but $\alpha=0.5$ and the solid,
the dashed, and the short-dashed lines represent $\phi_{s}=0.0$, $0.2$, and
$0.3$, respectively. The input parameters in right panel are same as Figure 1,
but $\phi_{s}=0.05$ and the solid, the dashed, and the short-dashed lines
represent $\alpha=0.2$, $0.3$, $0.4$, respectively.
#### 3.3.3 Mach number
Here, it will be interesting to investigate the existence of the transonic
point in hot accretion flow. The transonic point occurs in place that the
amount of Mach number becomes equal to unity. The Mach number referring to the
reference frame is defined as (Gaffet & Fukue 1983; Fukue 1984)
$\mu\equiv\frac{v_{r}-v_{F}}{c_{s}}=\frac{V-n\xi}{S}$ (39)
where
$v_{F}=\frac{dr}{dt}=n\frac{r}{t}$ (40)
is the velocity of the reference frame which is moving outward as time goes
by, which the sound speed can be subsequently expressed as
$c_{s}^{2}\equiv\frac{p}{\rho}=S^{2}(GM_{*}/t)^{2/3}$ (41)
and, $S=\left(\Pi/R\right)^{1/2}$ the sound speed in self-similar flow is
rescaled in the course of time. The Mach number introduced so far, represents
the _instantaneous_ and _local_ Mach number of the unsteady self-similar flow.
In steady self-similar solution (e. g. Tanaka & Menou 2006), the Mach number
does not vary by radii and is a constant. While, the Mach number in unsteady
self-similar varies by radii (see Fig. 4). As seen in Fig. 4, there is a
transonic point ($|\mu|=1$). The dependency of transonic point to saturation
constant shows that this point moves inward by adding the parameter of
$\phi_{s}$. Because, thermal conduction transfers the heat to larger radii, so
the sound speed/temperature decreases by adding the saturation constant. In
other words, whatever the radii smaller, the radial velocity relative to the
sound speed larger. Also, the Mach number profiles show the transonic point
decreasing by adding the viscous parameter which can be due to increase of
radial velocity along with adding the viscous parameter.
Figure 5: Same as Figure 1, but $\phi_{s}=0.1$. The solid, the dashed, and the
short-dashed lines represent $\gamma=1.3$, $1.4$, and $1.5$, respectively.
#### 3.3.4 Comparison of steady and unsteady self-similar solutions
As a comparison between steady and unsteady self-similar solution, the
physical quantities in unsteady self-similar solutions are divided into their
radial dependence in steady self-similar solution then plotted in terms of
$\xi$ in Fig. 5. Also, the effect of adiabatic index on hot accretion flow is
investigated in Fig. 5. The delineated quantities ($R\,\xi^{3/2}$,
$V\,\xi^{1/2}$, $\cdot$ $\cdot$ $\cdot$) in Fig 5 are constant in steady self-
similar solutions of hot accretion flows (Tanaka & Menou 2006; Shadmehri 2008;
Abbassi et al. 2008; Ghanbari et al. 2009). While they vary with position in
this research. Fig. 5 represents the density profile varies shallower than
$\xi^{-3/2}$, that this property is qualitatively consistent with simulation
results (e.g., Stone et al. 1999; Igumenshchev & Abramowicz 1999; Stone &
Pringle 2001; Hawley & Balbus 2002; Igumenshchev et al. 2003). The radial
dependency of the temperature and radial velocity in unsteady self-similar
solution show that they vary deeper than the steady self-similar solution.
Also, the study of the angular velocity show that it varies shallower than
$\xi^{-3/2}$. Thus, the physical quantities in the present model avoid the
limits of the steady self-similar solution. The physical quantities profiles
in Fig. 5 show that their radial dependency in the unsteady self-similar
limited to the steady self-similar solution by adding adiabatic index
$\gamma$.
## 4 Summary and Discussion
In hot accretion flows, the collision timescale between ions and electrons is
longer than the inflow timescale. Thus, the inflow plasma is collisionless,
and the transfer of energy by thermal conduction can be dynamically important.
The low collisional rate of the gas is confirmed by direct observation,
particularly in the case of the Galactic centre (Quataert 2004; Tanaka & Menou
2006) and in the intracluster medium of galaxy clusters (Sarazin 1986).
Here, we have investigated how thermal conduction affects dynamics of hot
quasi-spherical accretion flows. We adopted the presented solutions by Ogilvie
(1999) and Tanaka & Menou (2006). Thus, we assumed that angular momentum
transport is due to viscous turbulence and the $\alpha$-prescription is used
for the kinematic coefficient of viscosity. We also assumed the flow does not
have a good cooling efficiency and so a fraction of energy accretes along with
matter on to the central object. The effect of thermal conduction is studied
by a saturation form of it introduced by Cowie & McKee (1977). To solve the
equations that govern the dynamical behaviour of hot accretion flow, we have
used unsteady self-similar solution.
The effect of saturation constant and the viscous parameter on the present
model is investigated. The solutions show that with the increase of
conductivity, the equatorial density becomes denser and the temperature
becomes lower. These results are qualitatively consistent with simulation
results of Wu et al. (2010). Furthermore, the solutions show that by adding
the saturation constant, the angular velocity becomes larger and the radial
velocity decreases. The mass accretion rate is reduced by adding the
saturation constant that is qualitatively consistent with the result of
Johnson & Quataert (2007). The solutions imply that the viscous parameter has
opposite effects in comparison to saturation constant on physical quantities
of the system. Also, the study of physical quantities of the present model in
comparison to steady self-similar solution show that our results deviate from
steady self-similar solution and do not have its limits.
Here, we studied dynamical behaviour of hot accretion flow in one-dimensional
approach ignored by latitudinal dependence of physical quantities. Although,
some authors have shown that latitudinal dependence of physical quantities is
important for structure and dynamics of hot accretion flow (Tanaka & Menou
2006; Ghanbari et al. 2009; Wu et al. 2010). Thus, latitudinal behaviour of
the present model can be investigated in other studies.
## Acknowledgments
I would like to thank the anonymous referee for very useful comments that
helped me to improve the initial version of the paper.
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arxiv-papers
| 2011-11-15T16:27:39 |
2024-09-04T02:49:24.355053
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kazem Faghei",
"submitter": "Kazem Faghei",
"url": "https://arxiv.org/abs/1111.3569"
}
|
1111.3618
|
# Linking to Data - Effect on Citation Rates in Astronomy
Edwin A. Henneken1, Alberto Accomazzi1
###### Abstract
Is there a difference in citation rates between articles that were published
with links to data and articles that were not? Besides being interesting from
a purely academic point of view, this question is also highly relevant for the
process of furthering science. Data sharing not only helps the process of
verification of claims, but also the discovery of new findings in archival
data. However, linking to data still is a far cry away from being a
“practice”, especially where it comes to authors providing these links during
the writing and submission process. You need to have both a willingness and a
publication mechanism in order to create such a practice. Showing that
articles with links to data get higher citation rates might increase the
willingness of scientists to take the extra steps of linking data sources to
their publications. In this presentation we will show this is indeed the case:
articles with links to data result in higher citation rates than articles
without such links.
The ADS is funded by NASA Grant NNX09AB39G.
1Smithsonian Astrophysical Observatory, 60 Garden Street, Cambridge, MA 02138
## 1 Introduction
Furthering science depends to a large degree on knowledge and information
transfer. Therefore it critically relies on discoverability. This applies to
findings in publications and to the underlying data that led to these
findings. Therefore, significant amounts of energy (and funds) should be
invested in improving discoverability, of both publications and data. Major
progress has been made on the level of publications by improved visibility and
more sophisticated techniques for information discovery. The adoption of
faceted filtering, recommender systems and semantic interlinking of resources
are good examples of this (Accomazzi & Dave (2011), Henneken et al. (2011)).
It is time that exposure of data becomes common practice. A publication based
on a data set is just one expression of the potential of that data set. It
totally depends on the background and the interests of the researchers which
representation of that potential will be selected. However, there are many
other representations. The scientific community would also benefit greatly
from the ability to combine a data set with other available data sets. Also,
having data available publicly would greatly facilitate the verification of
claims (Fischer & Zigmond (2010)). The special session “The Literature-Data
Connection: Meaning, Infrastructure and Impact” at the 218th Meeting of the
American Astronomical Society (Boston, May 2011) was dedicated to this
discussion. As part of the discussion of how to create a practice of linking
data to publications, the question was raised whether such publications would
see a citation advantage. That would be like getting a tax benefit for “being
green”. Everybody agrees that “being green” is a sensible thing to do, but
having some kind of incentive definitely helps as additional motivation.
Motivation is an essential ingredient for creating a practice. Since citations
are a measure used for scientific impact, it is logical to ask whether
investing energy into making data available publicly results in a citation
advantage.
In this presentation we address the question whether there is a citation
advantage. We explore the question using the holdings and citation data of the
SAO/NASA Astrophysics Data System (ADS).
## 2 Results
With every record in the ADS holdings a number of possible attributes
(“links”) can be associated, giving access to information related to that
record. The attribute used for this analysis is the “D” link, associated with
access to on-line data. Currently these links point to data hosted at data
centers (like CDS, HEASARC and MAST). The following set of records was chosen
for this study: articles published in _The Astrophysical Journal_ (including
_Letters_ and _Supplement_), _The Astronomical Journal_ , _The Monthly Notices
of the R.A.S._ and _Astronomy & Astrophysics_ including _Supplement_), during
the period 1995 through 2000. Comparing publications with a “D” link to those
without such a link would, to a large degree, be comparing apples with
oranges, because of the range in subject matter. In order for the comparison
to make sense, the subject matter of the publications needs to be restricted.
We decided to use keywords as filter. We determined the set of 50 most
frequently used keywords in articles with data links. The articles to be used
for the analysis were obtained by requiring that they have at least 3 keywords
in common with that set of 50 keywords. This resulted in a set of 3814
articles with data links and 7218 articles without data links. The box diagram
in figure 1 characterizes the distribution of citations in the sets with and
without data links, for respectively 2 and 4 years after publication.
Figure 1.: Distribution of citations of articles published in _The
Astrophysical Journal_ (including _Letters_ and _Supplement_), _The
Astronomical Journal_ , _The Monthly Notices of the R.A.S._ and _Astronomy &
Astrophysics_ including _Supplement_), during the period 1995 through 2000.
The extent of the box corresponds with the interquartile range of the
citations and whiskers extend to 1.5 times the interquartile range. The
horizontal lines within the boxes correspond with the medians. From left to
right, the boxes correspond respectively with the citation distributions for
the article set with and without data links 2 years after publication, and 4
years after publication. The medians are respectively at 10, 8, 17 and 13
citations.
For this analysis, a random selection of 3814 articles was extracted from the
set of 7218 articles (without links to data). For both sets the citation
accumulation was determined for each article. From now on, we will refer to
the set with data links as $\emph{D}_{d}$ and the one without data links as
$\emph{D}_{n}$. These citation distributions were used to calculate the mean
citation accumulation for each set, normalized by the total number of
citations in the entire set of publications. The results are shown in figure
2.
Figure 2.: The normalized number of citations for data sets $\emph{D}_{d}$ and
$\emph{D}_{n}$. The citations have been normalized by the total number of
citations.
Figure 2 indicates that publications with a data link have a larger citation
rate than publications that do not. To get get an indication of how much more
citations a publication with a data link accumulates, on average, figure 3
shows the cumulative citation distribution, normalized by the total number of
citations for articles without data links, 10 years after publication.
Figure 3.: The cumulative citation distributions for data sets $\emph{D}_{d}$
and $\emph{D}_{n}$. The citation counts have been normalized by the total
number of citations for articles without data links, 10 years after
publication.
Figure 3 indicates that for this data set, articles with data links on average
acquired 20% more citations (compared to articles without these links) over a
period of 10 years. The fact that this increase is statistically significant
follows from a regression analysis performed on the entire data set. This
confirmed the increase of 20% in citation count (at a 95% confidence level).
## 3 Discussion
Our study seems to indicate that publications with links to on-line data seem
to have a higher citation rate than publications that do not. Could this
effect be attributed to another systematic effect? For example, studies have
shown that e-printing results in higher citation rates (see for example
Henneken et al. (2006)). However, both sets used to construct figures 2 and 3
turn out to be homogeneous in other publication attributes. For example, in
each set about 20% of the publications have e-prints associated with them. So,
the increased citation rates associated with e-printing contribute similarly
in both sets. Also, both sets are homogenous in links to object information
(NED and SIMBAD links). Lastly, could data centers, in attributing data links
to articles, have cherry-picked important (i.e. more citable) data sets? Both
sets of publications turn out to be homogenous in citation distributions as
well. This leads us to believe that the effect observed is real.
In a study of medical literature on cancer microarray clinical trials, Piwowar
et al. (2007) found that “publicly available data was significantly associated
with a 69% increase in citations”. Even though citation rates are different
for different disciplines, the qualitative observation still holds. Studies
and discussions in other disciplines show that data sharing is viewed as
important and highly relevant for the integrity and furthering of science, and
that the hurdles encountered have much in common between various disciplines
(Bruna (2010), Delamothe (1996), Kansa et al. (2010), Pisani et al. (2010),
South & Duke (2010), Vickers (2011), Vandewalle et al. (2009)).
## References
* Accomazzi & Dave (2011) Accomazzi, A., & Dave, R. 2011, in Astronomical Data Analysis Software and Systems XX, edited by I. N. Evans, A. Accomazzi, D. J. Mink, & A. H. Rots, vol. 442 of Astronomical Society of the Pacific Conference Series, 415
* Bruna (2010) Bruna, E. M. 2010, Biotropica, 42, 399
* Delamothe (1996) Delamothe, T. 1996, British Medical Journal, 312, 1241
* Fischer & Zigmond (2010) Fischer, B. A., & Zigmond, M. J. 2010, Science and Engineering Ethics, 16, 783
* Henneken et al. (2011) Henneken, E. A., Kurtz, M. J., Accomazzi, A., Grant, C., Thompson, D., Bohlen, E., Milia, G. D., Luker, J., & Murray, S. S. 2011, in Future Professional Communication in Astronomy II, Edited by Alberto Accomazzi, Astrophysics and Space Science Proceedings, Volume 1. ISBN 978-1-4419-8368-8. Springer Science+Business Media, LLC, 2011, p. 125, edited by A. Accomazzi, 125. arXiv:1005.2308
* Henneken et al. (2006) Henneken, E. A., Kurtz, M. J., Eichhorn, G., Accomazzi, A., Grant, C., Thompson, D., & Murray, S. S. 2006, Journal of Electronic Publishing, 9, 2. arXiv:cs/0604061
* Kansa et al. (2010) Kansa, E. C., Kansa, S. W., Burton, M. M., & Stankowski, C. 2010, Archaeologies, 6, 301
* Pisani et al. (2010) Pisani, E., Whitworth, J., Zaba, B., & Abou-Zahr, C. 2010, The Lancet, 375, 703
* Piwowar et al. (2007) Piwowar, H. A., Day, R. S., & Fridsma, D. B. 2007, PLoS One, 2
* R Development Core Team (2011) R Development Core Team 2011, R: A Language and Environment for Statistical Computing, R Foundation for Statistical Computing, Vienna, Austria. ISBN 3-900051-07-0, URL http://www.R-project.org
* South & Duke (2010) South, D. B., & Duke, C. S. 2010, Journal of Forestry, 108, 370
* Vandewalle et al. (2009) Vandewalle, P., Kovacevic, J., & Vetterli, M. 2009, IEEE Signal Processing Magazine, 26, 37
* Vickers (2011) Vickers, A. J. 2011, British Medical Journal, 342
|
arxiv-papers
| 2011-11-15T19:40:32 |
2024-09-04T02:49:24.363394
|
{
"license": "Public Domain",
"authors": "Edwin A. Henneken, Alberto Accomazzi",
"submitter": "Edwin Henneken",
"url": "https://arxiv.org/abs/1111.3618"
}
|
1111.3647
|
# Flavor SU(4) breaking between effective couplings
Bruno El-Bennich Universidade Cruzeiro do Sul, Rua Galvão Bueno, 868,
01506-000 São Paulo, SP, Brazil Instituto de Física Teórica, Universidade
Estadual Paulista, Rua Dr. Bento Teobaldo Ferraz, 271, 01140-070 São Paulo,
SP, Brazil Gastão Krein Instituto de Física Teórica, Universidade Estadual
Paulista, Rua Dr. Bento Teobaldo Ferraz, 271, 01140-070 São Paulo, SP, Brazil
Lei Chang Physics Division, Argonne National Laboratory, Argonne, Illinois
60439, USA Craig D. Roberts Physics Division, Argonne National Laboratory,
Argonne, Illinois 60439, USA Institut für Kernphysik, Forschungszentrum
Jülich, D-52425 Jülich, Germany Department of Physics, Illinois Institute of
Technology, Chicago, Illinois 60616-3793, USA David J. Wilson Physics
Division, Argonne National Laboratory, Argonne, Illinois 60439, USA
(4 November 2011)
###### Abstract
Using a framework in which all elements are constrained by Dyson-Schwinger
equation studies in QCD, and therefore incorporates a consistent, direct and
simultaneous description of light- and heavy-quarks and the states they
constitute, we analyze the accuracy of $SU(4)$-flavor symmetry relations
between $\pi\rho\pi$, $K\rho K$ and $D\rho D$ couplings. Such relations are
widely used in phenomenological analyses of the interactions between matter
and charmed mesons. We find that whilst $SU(3)$-flavor symmetry is accurate to
20%, $SU(4)$ relations underestimate the $D\rho D$ coupling by a factor of
five.
###### pacs:
14.40.Lb, 11.15.Tk, 12.39.Ki, 24.85.+p
I. Introduction. Hadrons in-medium are the focus of intense theoretical and
experimental activity. The chief motivation in heavy-ion collisions is a
better understanding of QCD’s deconfined phase, viz. the putative quark-gluon
plasma, its chiral restoration phase transition and associated order
parameters. Whilst an enhancement of charm and strangeness in the quark-gluon
phase is predicted to lead to the copious production of $D_{(s)}$ mesons
Cacciari:2005rkKuznetsova:2006bh at the large hadron collider, $J/\psi$
suppression has long been suggested as an unambiguous signature for quark-
gluon plasma formation Matsui:1986dk . Notwithstanding ongoing debates about
charmonia production mechanisms and a wide range of suppression effects, much
effort is sensibly dedicated to understanding the complicated final-state
interactions which occur after hadronization of the plasma; see, e.g., Ref.
Bracco:2011pg .
Charmed-meson interactions with nuclear matter will also be studied at the
future Facility for Antiproton and Ion Research (FAIR) and possibly at
Jefferson Laboratory (JLab). Low-momentum charmonia, such as $J/\psi$ and
$\psi$, and $D^{(*)}$ mesons can be produced by annihilation of antiprotons on
nuclei (FAIR) or by scattering electrons from nuclei (JLab). Since charmonia
do not share valence quarks in common with the surrounding nuclear medium,
proposed interaction mechanisms include: QCD van der Waals forces, arising
from the exchange of two or more gluons between color-singlet states
Peskin:1979vaBrodsky:1989jd ; and intermediate charmed hadron states
Brodsky:1997ghKo:2000jx , such that $\bar{D}^{(\ast)}D^{(\ast)}$ hadronic
vacuum polarization components of the $J/\psi$ interact with the medium via
meson exchanges Krein:2010vp .
A kindred approach is applied to low-energy interactions of open-charm mesons
with nuclei, which may create a path to the production of charmed nuclear
bound states ($D$-mesic nuclei) Tsushima:1998ru ; Haidenbauer:2007jq ;
Haidenbauer:2010ch ; Yamaguchi:2011xb . These studies rely on model
Lagrangians, within which effective interactions are expressed through
couplings between $D^{(\ast)}$\- and light-pseudoscalar- and vector-mesons.
The models are typically an $SU(4)$ extension of light-flavor chirally-
symmetric Lagrangians. Most recently, exotic states formed by heavy mesons and
a nucleon were investigated, based upon heavy-meson chiral perturbation theory
Yamaguchi:2011xb . In that study a universal coupling, $g_{\pi}$, between a
heavy quark and a light pseudoscalar or vector meson was inferred from the
strong decay $D^{*}\to D\pi$, cf. Ref. ElBennich:2010ha .
In the context of chiral Lagrangians, it is natural to question the
reliability of couplings based on $SU(4)$ symmetry. Flavor breaking effects
are already known to occur in the strange sector and should only be expected
to increase when including charm quarks. The order of magnitude of this larger
symmetry breaking is signalled by the compilation of charmed couplings in Ref.
Bracco:2011pg , where $SU(4)$ relations are shown to be violated at various
degrees (ranging from 7% to 70%) in couplings between two heavy mesons and one
light meson. No states containing a $s$-quark were considered.
Herein, we study a different quantitative measure, based upon ratios between
the $D\rho D$, $K\rho K$ and $\pi\rho\pi$ couplings; namely, a difference
between the same coupling involving either a $c$-, $s$\- or light-quark. We
are motivated by the notion that the $K\rho K$ and $D\rho D$ systems are
dynamically equivalent in the sense that the heavier quark acts as a spectator
and contributes predominantly to the static properties of the mesons, whereas
the exchange dynamics is mediated by the light quarks. In practice, the
symmetry idea is expressed by implementing $g_{D\rho D}\simeq g_{K\rho K}$ in
the meson-exchange models Haidenbauer:2007jq ; Haidenbauer:2010ch . The
$\pi\rho\pi$ coupling provides a well-constrained benchmark.
II. DSE Framework. Our primary object of interest is a phenomenological
coupling that relates the transition amplitude of an initial pseudoscalar
$H=Qf$-meson, $Q=c,s$ and $f=u,d$, to an identical meson via emission of an
off-shell $\rho$. The matrix element for this transition is
$\langle H(p_{2})|\,\rho(P,\lambda)\,|H(p_{1})\rangle=g_{H\rho H}\
\bm{\epsilon}_{\lambda}\\!\cdot P\,,$ (1)
an expression which defines the dimensionless coupling of the two pseudoscalar
mesons to a vector meson with momentum $P=p_{2}-p_{1}$ and polarization state
$\lambda$. The decay $\rho\to\pi\pi$ is also described by such a matrix
element. However, there is no associated physical process when
$m_{\rho}^{2}<4m_{H}^{2}$ and $p_{1}^{2}=p_{2}^{2}=-m_{H}^{2}$. (N.B. A
Euclidean metric is used:
$\\{\gamma_{\mu},\gamma_{\nu}\\}=2\,\delta_{\gamma\nu};\,\gamma_{\mu}^{\dagger}=\gamma_{\mu};\;a\\!\cdot\\!b=\sum^{4}_{i=1}a_{i}b_{i}$;
and
$\mathrm{tr}[\gamma_{5}\gamma_{\mu}\gamma_{\nu}\gamma_{\rho}\gamma\sigma]=-4\epsilon_{\mu\nu\rho\sigma},\,\epsilon_{1234}=1$.
For a space-like vector $P_{\mu},P^{2}>0$.) Nevertheless, a coupling of this
sort is employed in defining $\rho$-meson-mediated exchange-interactions
between a nucleon and pseudoscalar strange- or charm-mesons. In such
applications: the off-shell $\rho$-meson’s momentum is necessarily spacelike;
and a coupling and form factor may be defined once one settles on a definition
of the off-shell $\rho$-meson.
Symmetry-preserving models built upon predictions of QCD’s Dyson-Schwinger
equations (DSEs) provide a sound framework within which to examine heavy-meson
observables ElBennich:2010ha ; Ivanov:1997yg ; Ivanov:1998ms ; Ivanov:2007cw ;
ElBennich:2009vx . Such studies describe quark propagation via fully dressed
Schwinger functions, which has a material impact on light-quark
characteristics Chang:2011vu .
At leading-order in a systematic, symmetry-preserving truncation scheme
Bender:1996bb , one may express Eq. (1) as
$\displaystyle g_{H\\!\rho H}\ \bm{\epsilon}^{\lambda}\\!\cdot P$
$\displaystyle=$ $\displaystyle\mathrm{tr}_{\mathrm{CD}}\\!\
\\!\int\\!\frac{d^{4}k}{(2\pi)^{4}}\,\Gamma_{H}(k;k_{1})S_{Q}(k_{Q})$ (2)
$\displaystyle\times\
\bar{\Gamma}_{H}(k;-k_{2})S_{f}(k_{f}^{\prime})\,\bm{\epsilon}^{\lambda*}\\!\cdot\bar{\Gamma}_{\rho}(k;-P)S_{f}(k_{f})\;,$
where $S$ represent dressed-quark propagators for the indicated flavor and
$\Gamma_{H}$ are meson Bethe-Salpeter amplitudes (BSAs), with $H=\pi,K,D$. In
Eq. (2): the trace is over color and spinor indices;
$k_{Q}=k+w_{1}p_{1},k_{f}^{\prime}=k+w_{1}p_{1}-p_{2}$, $k_{f}=k-w_{2}p_{1}$,
where the relative- momentum partitioning parameters satisfy $w_{1}+w_{2}=1$;
and $\bm{\epsilon}^{\lambda}_{\mu}$ is the vector-meson polarization four-
vector. This approximation has been employed successfully; see, for instance,
applications in Refs. Ivanov:2007cw ; Chang:2011vu ; Roberts:1994hh ;
Tandy:1997qf ; Jarecke:2002xd ; Maris:2003vk ; Roberts:2007jh .
We simultaneously calculate the $D$-, $K$\- and $\rho$-meson leptonic decay
constants via Ivanov:1998ms :
$\displaystyle P_{\mu}f_{H}$ $\displaystyle=$
$\displaystyle\mathrm{tr}_{\mathrm{CD}}\\!\
\int\\!\frac{d^{4}k}{(2\pi)^{4}}\,\gamma_{5}\gamma_{\mu}\,\chi_{H}(k;P)\,,$
(3) $\displaystyle M_{\rho}f_{\rho}$ $\displaystyle=$
$\displaystyle\frac{1}{3}\mathrm{tr}_{\mathrm{CD}}\\!\
\int\\!\frac{d^{4}k}{(2\pi)^{4}}\,\gamma_{\mu}\,\chi_{\mu}^{\rho}(k;P)\,,$ (4)
where $\chi(k;P)=S_{f_{1}}(k+w_{1}P)\Gamma(k;P)S_{f_{2}}(k-w_{2}P)$. The BSAs
are canonically normalized; viz., for pseudoscalars
$\displaystyle 2\,P_{\mu}$ $\displaystyle=$
$\displaystyle\left[\frac{\partial}{\partial
K_{\mu}}\Pi(P,K)\right]_{K=P}^{P^{2}=-m^{2}_{0^{-}}}\ ,$ (5)
$\displaystyle\Pi(P,K)$ $\displaystyle=$
$\displaystyle\mathrm{tr}_{\mathrm{CD}}\\!\
\int\\!\frac{d^{4}k}{(2\pi)^{4}}\,\bar{\Gamma}_{0^{-}}(k;-P)S_{f_{1}}(k+w_{1}K)$
(6) $\displaystyle\times\ \Gamma_{0^{-}}(k;P)S_{f_{2}}(k-w_{2}K)\,,$
with an analogous expression for the $\rho$ Ivanov:1998ms .
The solution of QCD’s gap equation is the dressed-quark propagator, which has
the general form
$S(p)=-i\gamma\cdot p\,\sigma_{V}(p^{2})+\sigma_{S}(p^{2})=1/[i\gamma\cdot
p\,A(p^{2})+B(p^{2})]\,.$ (7)
For light-quarks, it is a longstanding DSE prediction that both the wave-
function renormalization, $Z(p^{2})=1/A(p^{2})$, and dressed-quark mass-
function, $M(p^{2})=B(p^{2})/A(p^{2})=\sigma_{S}(p^{2})/\sigma_{V}(p^{2})$,
receive strong momentum-dependent modifications at infrared momenta:
$Z(p^{2})$ is suppressed and $M(p^{2})$ enhanced. These features are
characteristic of dynamical chiral symmetry breaking (DCSB) and, plausibly, of
confinement. (N.B. Eqs. (8), (9) represent the quark propagator $S(p)$ as an
entire function, which entails the absence of a Lehmann representation and is
a sufficient condition for confinement Krein:1990sf ; Roberts:2007ji .) The
significance of this infrared dressing has long been emphasized Roberts:1994hh
; e.g., it is intimately connected with the appearance of Goldstone modes
Chang:2011vu . The predicted behavior of $Z(p^{2})$, $M(p^{2})$ has been
confirmed in numerical simulations of lattice-regularized QCD Roberts:2007ji ;
Bowman:2005vxBhagwat:2006tu .
Whilst numerical solutions of the quark DSE are readily obtained, the utility
of an algebraic form for $S(p)$, when calculations require the evaluation of
numerous integrals, is self-evident. An efficacious parametrization,
exhibiting the aforementioned features and used extensively Ivanov:1998ms ;
Ivanov:2007cw ; Roberts:1994hh ; Cloet:2008re , is expressed via
$\displaystyle\bar{\sigma}_{S}(x)$ $\displaystyle=$ $\displaystyle
2\,\bar{m}\,{\cal F}(2(x+\bar{m}^{2}))$ (8) $\displaystyle+{\cal
F}(b_{1}x)\,{\cal F}(b_{3}x)\,\left[b_{0}+b_{2}{\cal F}(\epsilon x)\right]\,,$
$\displaystyle\bar{\sigma}_{V}(x)$ $\displaystyle=$
$\displaystyle\frac{1}{x+\bar{m}^{2}}\,\left[1-{\cal
F}(2(x+\bar{m}^{2}))\right]\,,$ (9)
with $x=p^{2}/\lambda^{2}$, $\bar{m}$ = $m/\lambda$, ${\cal
F}(x)=[1-\exp(-x)]/x$, $\bar{\sigma}_{S}(x)=\lambda\,\sigma_{S}(p^{2})$ and
$\bar{\sigma}_{V}(x)=\lambda^{2}\,\sigma_{V}(p^{2})$. The parameter values
were fixed Ivanov:1998ms by requiring a least-squares fit to a wide range of
light- and heavy-meson observables, and take the values:
$\begin{array}[]{llcccc}f&\bar{m}_{f}&b_{0}^{f}&b_{1}^{f}&b_{2}^{f}&b_{3}^{f}\\\
\hline\cr u=d&0.00948&0.131&2.94&0.733&0.185\\\
s&0.210&0.105&3.18&0.858&0.185\end{array}\,.$ (10)
At a scale $\lambda=0.566\,$GeV, the current-quark masses take the values
$m_{u}=5.4\,$MeV and $m_{s}=119\,$MeV, and one obtains the following Euclidean
constituent-quark masses Maris:1997tm : $\hat{M}_{u}^{E}=0.36\,$GeV and
$\hat{M}_{s}^{E}=0.49\,$GeV. (N.B. $\epsilon=10^{-4}$ in Eq. (8) acts only to
decouple the large- and intermediate-$p^{2}$ domains Roberts:1994hh .)
We note that studies which do not or cannot implement light-quark dressing in
this QCD-consistent manner invariably encounter problems arising from the need
to employ large constituent-quark masses and the associated poles in the
light-quark propagators ElBennich:2008xyElBennich:2008qa . This typically
translates into considerable model sensitivity for computed observables
ElBennich:2009vx .
Whereas the impact of DCSB on light-quark propagators is significant, the
effect diminishes with increasing current-quark mass (see, e.g., Fig. 1 in
Ref. Ivanov:1998ms ). This can be explicated by considering the dimensionless
and renormalization-group-invariant ratio
$\varsigma_{f}:=\sigma_{f}/M^{E}_{f}$, where $\sigma_{f}$ is a constituent-
quark $\sigma$-term: $\varsigma_{f}$ measures the effect of explicit chiral
symmetry breaking on the dressed-quark mass-function compared with the sum of
the effects of explicit and dynamical chiral symmetry breaking. Calculation
reveals Roberts:2007jh : $\varsigma_{u}=0.02$, $\varsigma_{s}=0.23$,
$\varsigma_{c}=0.65$, $\varsigma_{b}=0.8$. Plainly, $\varsigma_{f}$ vanishes
in the chiral limit and remains small for light quarks, since the magnitude of
their constituent mass owes primarily to DCSB. On the other hand, for heavy
quarks, $\varsigma_{f}\to 1$ because explicit chiral symmetry breaking is the
dominant source of their mass. Notwithstanding this, confinement remains
important for the heavy-quarks. These considerations are balanced in the
following simple form for the $c$-quark propagator:
$S_{c}(k)=\frac{-i\gamma\cdot k+\hat{M}_{c}}{\hat{M}_{c}^{2}}{\cal
F}(k^{2}/\hat{M}_{c}^{2})\,,$ (11)
which implements confinement but produces a momentum-independent c-quark mass-
function; namely, $\sigma_{V}^{c}(k^{2})/\sigma_{S}^{c}(k^{2})=\hat{M}_{c}$.
We use $\hat{M}_{c}=1.32\,{\rm GeV}$ Ivanov:1998ms .
A meson is described by the amplitude obtained from a homogeneous Bethe-
Salpeter equation. In solving that equation the simultaneous solution of the
gap equation is required. Since we have already chosen to simplify the
calculations by parametrizing $S(p)$, we follow Refs. ElBennich:2010ha ;
Ivanov:1998ms ; Ivanov:2007cw ; ElBennich:2009vx and also employ that
expedient with $\Gamma_{H(\rho)}$.
In this connection, the quark-level Goldberger-Treiman relations derived in
Ref. Maris:1997hd motivate and support the following parametrization of the
$\pi$ and $K$ BSAs:
$\Gamma_{\pi,K}(k;P)=i\gamma_{5}\,\frac{\surd
2}{f_{\pi,K}}\,B_{\pi,K}(k^{2})\,,\\\ $ (12)
where $B_{\pi,K}:=\left.B_{u}\right|_{m_{u}\to 0}^{b_{0}^{u}\to
b_{0}^{\pi,K}}$ and are obtained from Eqs. (7) – (9) through the replacements
$b_{0}^{u}\rightarrow b_{0}^{\pi}=0.204$, $b_{0}^{u}\rightarrow
b_{0}^{K}=0.319$, which yield computed values $f_{\pi}=146\,$MeV,
$f_{K}=178\,$MeV Ivanov:1998ms . Equation (12) expresses the fact that the
dominant invariant function in a pseudoscalar meson’s BSA is closely related
to the scalar piece of the dressed-quark self energy owing to the axial-vector
Ward-Takahashi identity and DCSB.
Regarding the $\rho$ meson, DSE studies Jarecke:2002xd ; Pichowsky:1999mu
indicate that, in applications such as ours, one may effectively use
$\Gamma^{\mu}_{\rho}(k;P)=\left(\gamma^{\mu}-P^{\mu}\,\frac{\gamma\cdot
P}{P^{2}}\right)\frac{\exp(-k^{2}/\omega_{\rho}^{2})}{\mathcal{N}_{\rho}}\,,$
(13)
namely, a function whose support is greatest in the infrared. Similarly, for
the $D$ meson we choose:
$\Gamma_{D}(k;P)=i\gamma_{5}\,\frac{\exp(-k^{2}/\omega_{D}^{2})}{\mathcal{N}_{D}}\;.$
(14)
The normalizations, $\mathcal{N}_{\rho}$, $\mathcal{N}_{D}$, are obtained from
Eqs. (5), (6) and simultaneous calculation of the weak decay constant in Eqs.
(3), (4). In the expression for the coupling, Eq. (1), as well as in Eqs.
(3)–(5), we follow the momentum-partitioning prescription of Ref.
ElBennich:2010ha , which leads to $w_{1}^{c}=0.79$; viz., most but not all the
heavy-light-meson’s momentum is carried by the $c$-quark. We note that
Poincaré covariance is a hallmark of the direct application of DSEs to the
calculation of hadron properties. In such an approach, no physical observable
can depend on the choice of momentum partitioning. However, that feature is
compromised if, as herein, one does not retain the complete structure of
hadron bound-state amplitudes Maris:1997tm . Any sensitivity to the
partitioning is an artifact arising from our simplifications Ivanov:2007cw ;
ElBennich:2010ha .
III. Results. The $D$-meson’s width parameter is determined via analysis of
relevant leptonic and strong decays: $\omega_{D}=1.63\pm 0.10\,$GeV for
$m_{D}=1.865\,$GeV yields $f_{D}=206\pm 9\,$MeV Eisenstein:2008sq and
$g_{D^{\ast}D\pi}=18.7^{+2.5}_{-1.4}$ _cf_. $17.9\pm 1.9$ Anastassov:2001cw .
For the $\rho$, we use $\omega_{\rho}=0.56\pm 0.01\,$GeV and
$w_{2}^{\rho}=0.38$, both determined Ivanov:2007cw via a least-squares fit to
an array of light-light- and heavy-light-meson observables with
$m_{\rho}=0.77\,$GeV. Using Eqs. (3), (5) and (6), one therewith obtains
$f_{\rho}=209\,$MeV, _cf_. experiment $216\,$MeV, which follows from the
$e^{+}e^{-}$ decay width Nakamura:2010zzi .
With the width parameters fixed, we computed the $D\rho D$, $K\rho K$ and
$\pi\rho\pi$ couplings in impulse approximation, following Eq. (2). Our
results are depicted in Fig. 1. Notably, we compute the amplitude directly: at
all values of $P^{2}$ and current-quark mass. We do not need to resort to
extrapolations, neither from spacelike$\,\to\,$timelike momenta nor in
current-quark mass, expedients which are necessary in some other approaches
Bracco:2011pg ; Becirevic:2009xp .
Figure 1: _Upper panel_ – Dimensionless couplings: $g_{D\rho D}$ (solid
curve); $g_{K\rho K}$ (dashed curve); and $g_{\pi\rho\pi}$ (dotted curve) –
all computed as a function of the $\rho$-meson’s off-shell four-momentum-
squared, with the pseudoscalar mesons on-shell. Recall that with our Euclidean
metric, $P^{2}>0$ is spacelike. _Lower panel_ – Ratios of couplings: $g_{K\rho
K}/g_{D\rho D}$ (solid curve); and $g_{K\rho K}/g_{\pi\rho\pi}$ (dashed
curve). In the case of exact $SU(4)$ symmetry, these ratios take the values,
respectively, $1$ (dot-dashed line) and $(1/2)$ (dotted line). The vertical
dotted line marks the $\rho$-meson’s on-shell point in both panels. (N.B. In
GeV: $m_{D}=1.865$, $m_{\rho}=0.77$, $m_{K}=0.494$, $m_{\pi}=0.138$.)
The behavior of $g_{\pi\rho\pi}(P^{2})$ provides a context for our results.
Experimentally Nakamura:2010zzi , $g_{\pi\rho\pi}(-m_{\rho}^{2})=6.0$; and the
best numerically-intensive DSE computation available produces Jarecke:2002xd
$g_{\pi\rho\pi}(-m_{\rho}^{2})=5.2$. Our algebraically-simplified framework
produces $g_{\pi\rho\pi}(-m_{\rho}^{2})=4.8$, just 8% smaller than the latter,
and a $P^{2}$-dependence for the coupling which closely resembles that in Ref.
Mitchell:1996dn ; e.g., both are smooth, monotonically decreasing functions
and our value of
$g_{\pi\rho\pi}(-m_{\rho}^{2})/g_{\pi\rho\pi}(m_{\rho}^{2})=0.14$ is just 10%
smaller. On the domain $P^{2}\in[-m_{\rho}^{2},m_{\rho}^{2}]$
$g_{\pi\rho\pi}(s=P^{2})=\frac{1.84-1.45s}{1+0.75s+0.085s^{2}}$ (15)
provides an accurate interpolation of our result. If one insists on a monopole
parametrization at spacelike-$P^{2}$, then a monopole mass of
$\Lambda_{\pi\rho\pi}=0.61\,$GeV provides a fit with relative-error-standard-
deviation$\,=5$%.
In the case of exact $SU(3)$ symmetry, one would have $g_{K\rho
K}=g_{\pi\rho\pi}/2$. It is clear from the figure that the assumption provides
a fair approximation to our result on a domain which one can reasonably
consider as relevant to meson-exchange model phenomenology; viz., on
$P^{2}\in[-m_{\rho}^{2},m_{\rho}^{2}]$ the error ranges from
$(-10)\,$–$40\,$%. On this domain an accurate interpolation is provided by
$g_{K\rho K}(s)=\frac{0.94-0.62s}{1+0.55s-0.16s^{2}}.$ (16)
If one insists on a monopole parametrization at spacelike-$P^{2}$, then a
monopole mass of $\Lambda_{K\rho K}=0.77\,$GeV provides a fit with relative-
error-standard-deviation$\,=4$%.
With $SU(4)$ symmetry, the picture is different. We have a numerical result
that is reliably interpolated via
$g_{D\rho D}(s)=\frac{5.05-4.26s}{1+0.36s-0.060s^{2}}.$ (17)
A monopole parametrization at spacelike-$P^{2}$, with mass-scale
$\Lambda_{D\rho D}=0.69\,$GeV, provides a fit with relative-error-standard-
deviation$\,=5$%. Our computed value $g_{D\rho D}(0)=5.05$ is 75% larger than
an estimate obtained using QCD sum rules ($3.0\pm 0.02$ Bracco:2011pg ) and
100% larger than a vector-meson-dominance estimate ($2.52$ Lin:1999ad ).
Moreover, if $SU(4)$ symmetry were exact, then $g_{D\rho D}=g_{K\rho
K}=g_{\pi\rho\pi}/2$, but it is plain from Eq. (16) that $g_{K\rho
K}(0)=0.92$, a result which exposes a symmetry violation of $440$% at
$P^{2}=0$. Furthermore, on the entire domain
$P^{2}\in[-m_{\rho}^{2},m_{\rho}^{2}]$, the symmetry-based expectation
$g_{D\rho D}=g_{K\rho K}$ is always violated, at a level of between
$360\,$–$\,440$%. The second identity, $g_{D\rho D}=g_{\pi\rho\pi}/2$, is
violated at the level of $320\,$–$\,540$%. (N.B. In connection with heavy-
quark symmetry, corrections of this order have also been encountered $c\to d$
transitions Ivanov:1998ms .)
These conclusions are dramatic, so it is important to explain why we judge
them to be robust. The computations of $g_{\pi\rho\pi}$ and $g_{K\rho K}$ are
considered reliable because we can smoothly take the limit
$s$-quark$\,\to\,u$-quark and thereby recover a unique function that agrees
with earlier computations by other groups.
This leaves the possibility of uncertainties connected with $S_{c}(k)$, Eq.
(11); $\Gamma_{D}(k;P)$, Eq. (14); and the momentum partitioning parameter,
$w_{1}^{c}$. To explore sensitivity to the $c$-quark propagator we used an
even simpler, non-confining constituent-like form; viz.,
$S_{C}(k)=1/(i\gamma\cdot k+\hat{M}_{c})$. The effect at spacelike-$P^{2}$ is
modest. However, the impact is large at timelike-$P^{2}$ because thereupon the
$\rho$-meson momentum-squared begins to explore a neighborhood of the spurious
pole in $S_{C}(k)$. Thus, the simpler propagator serves to _increase_ the
violation of $SU(4)$ symmetry. Regarding $\Gamma_{D}(k;P)$, uncertainty is
implicit in the value of $\omega_{D}=1.63\pm 0.10\,$GeV, constrained by the
weak decay constant $f_{D^{+}}=206\pm 9\,$MeV Eisenstein:2008sq . However,
variations of even 20% in $\omega_{D}$ have no material impact on our results.
Connected with that, a 20% change in $w_{1}^{c}$ produces only a 4% variation
in $\omega_{D}$ via the fit to $f_{D^{+}}$, hence any possibility of an effect
from $w_{1}^{c}$ can be discounted owing to the previous consideration.
IV. Discussion. Predictions for bound-states and resonances derived from
meson-exchange models are sensitive to the values of couplings in their
Lagrangians. In these non-relativistic models the couplings are commonly fixed
to reproduce some known experimental data, e.g. the scattering length of a
physical system. The most prominent such coupling, namely $g_{\pi\\!N}$, has
long been used in nucleon-nucleon potentials and serves to define the strength
of the pion’s coupling to a nucleon. It also determines the scale of the long-
range force in the nucleon-nucleon interaction and associated scattering cross
sections. Analogously, the strength of the couplings $D^{(\ast)}\\!D\pi$,
$D^{(\ast)}D^{(\ast)}\rho$ between $D$ mesons and a light pion or $\rho$-meson
plays a crucial role in the formation of charmed-nuclei. However, whereas
$g_{\pi\\!N}$ can be extracted from $\pi N$-scattering data Ericson:2000md ,
no such information is available for charmed-meson interactions with nucleons.
In our approach, which is based on an internally consistent use of impulse
approximation and unifies the description of light- and heavy-mesons, we
compute these couplings from the transition amplitude between two $D$ mesons
and an off-shell light meson. We find that $SU(4)$ symmetry is a very poor
guide to the couplings. On the other hand, in relation to such models it
provides a constructive suggestion that one might reasonably employ
$F^{\rm ME}_{D\rho D}(|\vec{q}|^{2})=g^{\rm ME}_{D\rho D}\frac{\Lambda_{D\rho
D}^{{\rm ME}\,2}}{\Lambda_{D\rho D}^{{\rm ME}\,2}+|\vec{q}|^{2}},$ (18)
with $g^{\rm ME}_{D\rho D}\approx 5$, $\Lambda^{\rm ME}_{D\rho D}\approx
0.7\,$GeV, to describe $D\,D$ scattering via $\rho(\vec{q})$-meson exchange.
This might be compared with the parametrization Haidenbauer:2007jq :
$F^{H}_{D\rho D}(|\vec{q}|^{2})=g^{H}_{D\rho D}\frac{\Lambda^{H\,2}_{D\rho
D}}{\Lambda_{D\rho D}^{H\,2}+|\vec{q}|^{2}},$ (19)
$\Lambda^{H}_{D\rho D}=1.4\,$GeV, $g^{H}_{D\rho D}\approx 2$, based on the
notion of $SU(4)$ symmetry, which our analysis has discredited. The coupling
in Eq. (19) is smaller than that in Eq. (18) but the evolution is harder.
These effects cancel to some degree, but here the magnitudes are such that our
result, Eq. (18), provides an integrated interaction
$V_{0}=\int d^{3}\vec{q}\;F^{H}_{D\rho
D}(|\vec{q}|^{2})^{2}\frac{1}{|\vec{q}|^{2}+m_{\rho}^{2}}$ (20)
that is roughly 40% greater. (N.B. If $g^{H}_{D\rho D}\to
2.6\approx(1/2)g^{\rm ME}_{D\rho D}$, then $V_{0}^{H}\approx V_{0}^{ME}$.) By
the same measure, our $D\rho D$ interaction is 20% stronger than that in Ref.
Yamaguchi:2011xb , which uses $\Lambda^{Y}_{D\rho D}=1.14\,$GeV, $g_{V}=5.8$
and hence
$g_{D\rho D}^{Y}=0.9g_{V}[1-m_{\rho}^{2}/\Lambda^{Y\,2}_{D\rho D}]=2.85\,.$
(21)
Whilst our results argue against hard form factors, the interaction
enhancement they produce is abundantly clear. Notably, a large value for the
interaction strength entails an inflated cross-section in $DN$ scattering. In
particular, in the meson-exchange model of Ref. Haidenbauer:2007jq (single-
meson exchange version), the $I=1$ $\bar{D}N$ cross-section is inflated by a
factor of $\sim 5$, when using the our result, Eq. (18), for $\omega$ and
$\rho$, instead of Eq. (19). Hence, implementation of our results could have
material consequences on, e.g., the possibility for formation of charmed-
resonances or -bound-states in nuclei.
###### Acknowledgements.
We acknowledge useful input from A. Hosaka and S. M. Schmidt. This work was
supported by: Conselho Nacional de Desenvolvimento Científico e Tecnológico,
grant no. 305894/2009-9, Fundação de Amparo à Pesquisa do Estado de São Paulo,
grant nos. 2009/50180-0, 2009/51296-1 and 2010/05772-3; United States
Department of Energy, Office of Nuclear Physics, contract no. DE-
AC02-06CH11357; and Forschungszentrum Jülich GmbH.
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|
arxiv-papers
| 2011-11-15T21:01:08 |
2024-09-04T02:49:24.369314
|
{
"license": "Public Domain",
"authors": "Bruno El-Bennich, Gast\\~ao Krein, Lei Chang, Craig D. Roberts and\n David J. Wilson",
"submitter": "Bruno El-Bennich",
"url": "https://arxiv.org/abs/1111.3647"
}
|
1111.3812
|
# On exterior moduli of quadrilaterals and special functions
Matti Vuorinen and Xiaohui Zhang Department of Mathematics and Statistics,
University of Turku, 20014 Turku, Finland vuorinen@utu.fi, xiazha@utu.fi
###### Abstract.
In this paper two identities involving a function defined by the complete
elliptic integrals of the first and second kinds are proved. Some functional
inequalities and elementary estimates for this function are also derived from
the properties of monotonicity and convexity of this function. As
applications, some functional inequalities and the growth of the exterior
modulus of a rectangle are studied.
††footnotetext: File: vz20130411.tex, printed: 2024-8-27, 18.38
Keywords. Exterior modulus, complete elliptic integral, inequality
2010 Mathematics Subject Classification. 33E05, 31A15
## 1\. Introduction
###### 1.1.
Exterior modulus of a quadrilateral. For $h>0$ consider the rectangle $D$ with
vertices $1+ih$, $ih$, $0$, $1$ in the upper half plane
$\mathbb{H}^{2}=\\{x+iy:y>0\\}$ and a bounded harmonic function
$u:\mathbb{C}\setminus D\to\mathbb{R}$ satisfying the Dirichlet-Neumann
boundary value problem $u(z)=0$ for $z\in[0,1]$, $u(z)=1$ for $z\in[ih,1+ih]$,
$\frac{\partial u}{\partial n}(z)=0$ for $z\in[1,1+ih]\cup[0,ih]$ where $n$ is
the direction of the exterior normal to $\partial D$. The number
$\mathcal{M}(1+ih,ih,0,1)=\int_{\mathbb{C}\setminus D}|\nabla u|^{2}dm$
is called the exterior modulus of the rectangle $D(1+ih,ih,0,1)$.
This quantity also has an interpretation as the modulus of the family of all
curves, joining the segments $[1+ih,ih]$ and $[0,1]$ in the complement of the
rectangle $D$, which also is equal to $\mathcal{M}(1+ih,ih,0,1)$ (cf. [A]).
For a polygonal quadrilateral $D(a,b,0,1)$ with vertices
$a,b\in\mathbb{H}^{2}$ and base $[0,1]$, the exterior modulus
$\mathcal{M}(a,b,0,1)$ can be defined in the same way.
As far as we know there is no analytic formula for $\mathcal{M}(a,b,0,1)$.
Numerical methods for the computation of $\mathcal{M}(a,b,0,1)$ were recently
studied by H. Hakula, A. Rasila, and M. Vuorinen in [HRV2] which motivate the
present study. They used numerical methods such as hp-FEM and the Schwarz-
Christoffel mapping. Similar problems for the interior modulus have been
studied in [HQR, HRV1]. The literature and software dealing with numerical
conformal mapping problems are very wide, see e.g. [DT, PS].
Here we study the above problem for the case of a rectangle. In this case an
explicit formula involving complete elliptic integrals was given by P. Duren
and J. Pfaltzgraff [DP], and our goal is to analytically study the dependence
of the formula on $h$.
###### 1.2.
Complete elliptic integrals. Let $\mathcal{K}(r)$ and $\mathcal{E}(r)$ stand
for the complete elliptic integrals of the first and second kind, respectively
(see (2.1)). Let $r^{\prime}=\sqrt{1-r^{2}}$ for $r\in(0,1)$. We often denote
$\mathcal{K}^{\prime}(r)=\mathcal{K}(r^{\prime}),\quad\mathcal{E}^{\prime}(r)=\mathcal{E}(r^{\prime})$.
Define the function $\psi$ as follows
(1.3)
$\psi(r)=\frac{2(\mathcal{E}(r)-(1-r)\mathcal{K}(r))}{\mathcal{E}^{\prime}(r)-r\mathcal{K}^{\prime}(r)},\quad
r\in(0,1).$
The function $\psi:(0,1)\to(0,\infty)$ is a homeomorphism, see Theorem 3.1 or
[DP]. In particular, $\psi^{-1}:(0,\infty)\to(0,1)$ is well-defined.
###### 1.4.
Duren-Pfaltzgraff formula for a rectangle. In [DP], P. Duren and J.
Pfaltzgraff studied the modulus $\mathcal{M}(\Gamma)$ of the family of curves
$\Gamma$ joining the opposite sides of length $b$ of the rectangle with sides
$a$ and $b$, in the exterior of the rectangle, and gave the formula [DP,
Theorem 5]
(1.5)
$\mathcal{M}(\Gamma)=\dfrac{\mathcal{K}^{\prime}(r)}{2\mathcal{K}(r)},\quad\mbox{where}\quad
r=\psi^{-1}(a/b).$
The exterior modulus $\mathcal{M}(\Gamma)$ is a conformal invariant of a
quadrilateral. In [ADV], the authors gave a sharp comparison between the
function $\psi$ and Robin modulus of a given rectangle. Their result can be
rewritten as the following inequality
(1.6) $\dfrac{\pi r}{(1-r)^{2}}<\psi(r)<\dfrac{16r}{\pi(1-r)^{2}},\quad
r\in(0,1).$
In this paper two identities involving the function $\psi$ are proved, and
some functional inequalities and elementary estimates for the function $\psi$
are also derived from the monotonicity and convexity of the combinations of
the function $\psi$ and some elementary functions. As applications, we will
study the growth of the exterior modulus with respect to the length of one
side of the rectangle. The main results are listed as follows.
###### Theorem 1.7.
For $r\in(0,1)$, the function $\psi$ satisfies the identities
$\psi(r^{2})\psi\left(\left(\dfrac{1-r}{1+r}\right)^{2}\right)=1,\quad\psi\left(\dfrac{1-r}{1+r}\right)\psi\left(\dfrac{1-r^{\prime}}{1+r^{\prime}}\right)=1.$
###### Theorem 1.8.
The function $f(r)=(1-\sqrt{r})^{2}\psi(r)/r$ is strictly decreasing from
$(0,1)$ onto $(4/\pi,\pi)$. In particular, for all $r\in(0,1)$
$\dfrac{4r}{\pi(1-\sqrt{r})^{2}}<\psi(r)<\dfrac{\pi r}{(1-\sqrt{r})^{2}}.$
###### Theorem 1.9.
The function $f(x)=\psi(1/\operatorname{ch}(x))$ is decreasing and convex from
$(0,\infty)$ onto $(0,\infty)$. In particular, for $r,s\in(0,1)$,
(1.10)
$2\psi\left(\frac{\sqrt{2rs}}{\sqrt{1+rs+r^{\prime}s^{\prime}}}\right)\leq\psi(r)+\psi(s)$
with equality in the above inequality if and only if $r=s$.
###### Theorem 1.11.
For $x,y\in(0,1)$,
$\psi\left(H_{p}(x,y)\right)\leq
H_{p}\left(\psi(x),\psi(y)\right){\quad}\mbox{if}{\quad}p\geq 0,$
and
$\psi\left(H_{p}(x,y)\right)\geq
H_{p}\left(\psi(x),\psi(y)\right){\quad}\mbox{if}{\quad}p\leq-1.$
The equality holds in each case if and only if $x=y$. Here $H_{p}$ is the
power mean defined as
$H_{p}(x,y)=\left\\{\begin{array}[]{ll}\left(\dfrac{x^{p}+y^{p}}{2}\right)^{1/p},&p\neq
0\\\ \sqrt{xy},&p=0.\end{array}\right.$
## 2\. Preliminaries
For $0<r<1$, the functions
(2.1)
$\mathcal{K}(r)=\int_{0}^{\pi/2}\dfrac{dt}{\sqrt{1-r^{2}\sin^{2}t}},\quad\mathcal{E}(r)=\int_{0}^{\pi/2}\sqrt{1-r^{2}\sin^{2}t}\,dt$
with limiting values $\mathcal{K}(0)=\pi/2=\mathcal{E}(0)$,
$\mathcal{K}(1-)=\infty$ and $\mathcal{E}(1)=1$ are known as Legendre’s
complete elliptic integrals of the first and second kind, respectively. These
two functions are connected by Legendre’s relation [BF, 110.10]
(2.2)
$\mathcal{E}\mathcal{K}^{\prime}+\mathcal{E}^{\prime}\mathcal{K}-\mathcal{K}\mathcal{K}^{\prime}=\dfrac{\pi}{2}.$
Some derivative formulas involving these elliptic integrals are as follows
[Bo, p.21]:
(2.3)
$\left\\{\begin{array}[]{ll}\dfrac{d\mathcal{K}}{dr}=\dfrac{\mathcal{E}-r^{\prime
2}\mathcal{K}}{rr^{\prime
2}},&\dfrac{d\mathcal{E}}{dr}=\dfrac{\mathcal{E}-\mathcal{K}}{r},\vspace{1mm}\\\
\dfrac{d}{dr}\left(\mathcal{E}-r^{\prime
2}\mathcal{K}\right)=r\mathcal{K},&\dfrac{d}{dr}\left(\mathcal{K}-\mathcal{E}\right)=\dfrac{r\mathcal{E}}{r^{\prime
2}}.\\\ \end{array}\right.$
The functions $\mathcal{K}$ and $\mathcal{E}$ satisfy the following identities
due to Landen [BF, 163.01, 164.02]
(2.4) $\mathcal{K}\left(\dfrac{2\sqrt{r}}{1+r}\right)=(1+r)\mathcal{K}(r),$
(2.5)
$\mathcal{K}\left(\dfrac{1-r}{1+r}\right)=\dfrac{1}{2}(1+r)\mathcal{K}^{\prime}(r),$
(2.6)
$\mathcal{E}\left(\dfrac{2\sqrt{r}}{1+r}\right)=\dfrac{2\mathcal{E}(r)-r^{\prime
2}\mathcal{K}(r)}{1+r},$ (2.7)
$\mathcal{E}\left(\dfrac{1-r}{1+r}\right)=\dfrac{\mathcal{E}^{\prime}(r)+r\mathcal{K}^{\prime}(r)}{1+r}.$
Using Landen’s transformation formulas, we have the following identities.
###### Lemma 2.8.
For $r\in(0,1)$, let $t=(1-r)/(1+r)$. Then
(2.9) $\mathcal{K}(t^{2})=\dfrac{(1+r)^{2}}{4}\mathcal{K}^{\prime}(r^{2}),$
(2.10) $\mathcal{K}^{\prime}(t^{2})={(1+r)^{2}}\mathcal{K}(r^{2}),$ (2.11)
$\mathcal{E}(t^{2})=\dfrac{\mathcal{E}^{\prime}(r^{2})+(r+r^{2}+r^{3})\mathcal{K}^{\prime}(r^{2})}{(1+r)^{2}},$
(2.12)
$\mathcal{E}^{\prime}(t^{2})=\dfrac{4\mathcal{E}(r^{2})-(3-2r^{2}-r^{4})\mathcal{K}(r^{2})}{(1+r)^{2}}.$
###### Proof.
By Landen’s transformations (2.4) and (2.5), we have
$\dfrac{2(1+r^{2})}{(1+r)^{2}}\mathcal{K}(t^{2})=\mathcal{K}\left(\dfrac{1-r^{2}}{1+r^{2}}\right)=\dfrac{1}{2}(1+r^{2})\mathcal{K}^{\prime}(r^{2}).$
This implies (2.9).
For (2.10),
$\mathcal{K}^{\prime}(t^{2})=\dfrac{(1+r)^{2}}{1+r^{2}}\mathcal{K}\left(\dfrac{2r}{1+r^{2}}\right)=(1+r)^{2}\mathcal{K}(r^{2})$
where the first equality is Landen’s transformation (2.5) with the parameter
$t^{2}$ and the second equality follows from (2.4) with the parameter $r^{2}$.
Using Landen’s transformation (2.6) with the change of parameter $r\mapsto
t^{2}$ and the formula (2.9), we get
(2.13)
$\mathcal{E}\left(\dfrac{1-r^{2}}{1+r^{2}}\right)=\dfrac{(1+r)^{2}\mathcal{E}(t^{2})-r(1+r^{2})\mathcal{K}^{\prime}(r^{2})}{1+r^{2}}.$
On the other hand, by (2.7)
(2.14)
$\mathcal{E}\left(\dfrac{1-r^{2}}{1+r^{2}}\right)=\dfrac{\mathcal{E}^{\prime}(r^{2})+r^{2}\mathcal{K}^{\prime}(r^{2})}{1+r^{2}}.$
Hence (2.11) follows from (2.13) and (2.14).
For (2.12), by the change of parameter $r\mapsto t^{2}$ in Landen’s
transformation (2.7) and the formula (2.10), we have
(2.15)
$\mathcal{E}\left(\dfrac{2r}{1+r^{2}}\right)=\dfrac{(1+r)^{2}\mathcal{E}^{\prime}(t^{2})+(1-r^{2})^{2}\mathcal{K}(r^{2})}{2(1+r^{2})}.$
On the other hand, by (2.6)
(2.16)
$\mathcal{E}\left(\dfrac{2r}{1+r^{2}}\right)=\dfrac{2\mathcal{E}(r^{2})-(1-r^{4})\mathcal{K}(r^{2})}{1+r^{2}}.$
Hence (2.12) follows from (2.15) and (2.16). ∎
The next lemma is a monotone form of l’Hôpital’s rule and will be useful in
deriving monotonicity properties and obtaining inequalities [AVV1, Theorem
1.25].
###### Lemma 2.17 (Monotone form of l’Hôpital’s Rule).
Let $-\infty<a<b<\infty$, and let $f,g:[a,b]\to\mathbb{R}$ be continuous on
$[a,b]$, differentiable on $(a,b)$. Let $g^{\prime}(x)\neq 0$ on $(a,b)$.
Then, if $f^{\prime}(x)/g^{\prime}(x)$ is increasing (decreasing) on $(a,b)$,
so are
$\dfrac{f(x)-f(a)}{g(x)-g(a)}\qquad\mbox{and}\qquad\dfrac{f(x)-f(b)}{g(x)-g(b)}.$
If $f^{\prime}(x)/g^{\prime}(x)$ is strictly monotone, then the monotonicity
on the conclusion is also strict.
The following Lemma 2.18 is from [AVV1, Theorem 3.21 (1),(7)].
###### Lemma 2.18.
_(1)_ $r^{-2}(\mathcal{E}-r^{\prime 2}\mathcal{K})$ is strictly increasing and
convex from $(0,1)$ onto $(\pi/4,1)$.
_(2)_ For each $c\in[1/2,\infty)$, $r^{\prime c}\mathcal{K}$ is decreasing
from $[0,1)$ onto $(0,\pi/2]$.
###### Lemma 2.19.
_(1)_ $f_{1}(r)=\mathcal{E}-(1-r)\mathcal{K}$ is strictly increasing and
concave from $(0,1)$ onto $(0,1)$.
_(2)_ $f_{2}(r)=(\mathcal{E}-(1-r)\mathcal{K})/r$ is strictly decreasing from
$(0,1)$ onto $(1,\pi/2)$.
_(3)_ $f_{3}(r)=\mathcal{E}^{\prime}-r\mathcal{K}^{\prime}$ is strictly
decreasing and convex from $(0,1)$ onto $(0,1)$.
_(4)_ $f_{4}(r)=(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})/(1-r)$ is
strictly decreasing from $(0,1)$ onto $(0,1)$.
_(5)_ $f_{5}(r)=(\mathcal{E}-r^{\prime}\mathcal{K})/(1-\sqrt{r^{\prime}})^{2}$
is strictly decreasing from $(0,1)$ onto $(1,\pi/2)$.
_(6)_ $f_{6}(r)=(3-r)\mathcal{E}^{\prime}-(1+r)\mathcal{K}^{\prime}$ is
increasing form $(0,1)$ onto $(-\infty,0)$.
_(7)_ $f_{7}(r)=(1+r)(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})/(1-r)$ is
strictly deceasing from $(0,1)$ onto $(0,1)$.
_(8)_ $f_{8}(r)=(\mathcal{E}-(1-r)\mathcal{K})/(\sqrt{r}(1-r)\mathcal{K})$ is
strictly increasing from $(0,1)$ onto $(0,\infty)$. $f_{8}(0.479047\cdots)=1.$
###### Proof.
(1) By differentiation and the derivative formulas (2.3),
$f_{1}^{\prime}(r)=\dfrac{\mathcal{E}}{1+r}$
which is positive and decreasing. Then the properties of monotonicity and
concavity of $f_{1}$ follow. The limiting value $f_{1}(0)=0$ is clear and
$f_{1}(1-)=\mathcal{E}(1)-\lim\limits_{r\to 1-}(r^{\prime
2}\mathcal{K}/(1+r))=1$ by Lemma 2.18(2).
(2) Since $f_{1}$ is concave and $f_{2}(r)=f_{1}(r)/r$, $f_{2}$ is decreasing
by the monotone form of l’Hôpital’s rule. By l’Hôpital’s rule
$f_{2}(0)=f^{\prime}_{1}(0)=\pi/2$, and $f_{2}(1)=1$ is clear.
(3) By (2.3), we have
$f_{3}^{\prime}(r)=-\dfrac{\mathcal{K}^{\prime}-\mathcal{E}^{\prime}}{1+r},$
which is negative and increasing in $(0,1)$. Then $f_{3}$ is decreasing and
convex in $(0,1)$. The limiting value $f_{3}(0)=1$ follows from Lemma 2.18(2),
and $f_{3}(1)=0$ is clear.
(4) Let $h(r)=1-r$. Since $f_{3}$ is convex, $f_{3}^{\prime}(r)/h^{\prime}(r)$
is decreasing. Thus $f_{4}(r)=f_{3}(r)/h(r)$ is also decreasing by the
monotone form of l’Hôpital’s rule. By l’Hôpital’s rule
$f_{4}(1)=-f^{\prime}_{3}(1)=0$, and $f_{4}(0)=f_{3}(0)=1$.
(5) For the proof we first make the change of variable $r=2\sqrt{x}/(1+x)$.
The Landen transformations (2.4) and (2.6) lead to
$h(x)=f_{5}\left(\dfrac{2\sqrt{x}}{1+x}\right)=\dfrac{\mathcal{E}(x)-x^{\prime
2}\mathcal{K}(x)}{1-x^{\prime}}=\dfrac{h_{1}(x)}{h_{2}(x)},$
where $h_{1}(x)=\mathcal{E}(x)-x^{\prime 2}\mathcal{K}(x)$ and
$h_{2}(x)=1-x^{\prime}$ with $h_{1}(0)=0=h_{2}(0)$. Then by (2.3) we have
$\dfrac{h_{1}^{\prime}(x)}{h_{2}^{\prime}(x)}=\dfrac{x\mathcal{K}(x)}{x/x^{\prime}}=x^{\prime}\mathcal{K}(x),$
which is strictly decreasing by Lemma 2.18(2). This implies that $h$ is
decreasing by the monotone form of l’Hôpital’s rule, and hence $f_{5}$ is also
decreasing in $(0,1)$.
(6) By differentiation, we have
$f_{6}^{\prime}(r)=\dfrac{(1-r)(2r\mathcal{E}^{\prime}+\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})}{r(1+r)}>0,$
and hence $f_{6}$ is increasing. The limiting values are clear.
(7) By simple computation, $f_{7}^{\prime}(r)=f_{6}(r)/(1-r)^{2}<0$ and hence
$f_{7}$ is decreasing. The limiting values follow from part (4).
(8) Differentiation and simplification give that
$f_{8}^{\prime}(r)=\dfrac{((1+r)\mathcal{K}-\mathcal{E})(2\mathcal{E}-r^{\prime
2}\mathcal{K})}{2r^{3/2}(1+r)(1-r)^{2}\mathcal{K}^{2}}>0,$
and hence $f_{8}$ is strictly increasing. The limit
$f_{8}(0+)=\lim_{r\to
0+}\dfrac{\mathcal{E}-(1-r)\mathcal{K}}{\sqrt{r}\mathcal{K}}=\lim_{r\to
0+}\dfrac{\mathcal{E}-(1-r)\mathcal{K}}{r}\dfrac{\sqrt{r}}{\mathcal{K}}=0$
follows from the part (2). The limit $f_{8}(1-)=\infty$ is clear. ∎
Let
(2.20) $\mu(r)=\frac{\pi}{2}\frac{\mathcal{K}^{\prime}(r)}{\mathcal{K}(r)}$
be the modulus of Grötzsch’s ring $\mathbb{B}^{2}\setminus[0,r]$ (see
[LV],[AVV1]).
###### Lemma 2.21.
The function $f(r)=\mu(r)\psi(r)$ is strictly increasing from $(0,1)$ onto
$(0,\infty)$.
###### Proof.
Since the function $f$ can be rewritten as
$f(r)=\pi\sqrt{r}\mathcal{K}^{\prime}\dfrac{\mathcal{E}-(1-r)\mathcal{K}}{\sqrt{r}(1-r)\mathcal{K}}\dfrac{1-r}{\mathcal{E}^{\prime}-r\mathcal{K}^{\prime}},$
the conclusion follows from Lemma 2.18(2), Lemma 2.19(4) and (8). ∎
## 3\. Proofs of Main Results
In this section we will prove two identities involving the function $\psi$,
and some functional inequalities and elementary estimates for the function
$\psi$ are also derived from the monotonicity and convexity of the
combinations of the function $\psi$ and some elementary functions.
###### Theorem 3.1.
The function $\psi(r)$ is strictly increasing and convex from $(0,1)$ onto
$(0,\infty)$, and the function $\psi(r)/r$ is strictly increasing from $(0,1)$
onto $(0,\infty)$.
###### Proof.
By differentiation, and using (2.3) and Legendre’s identity (2.2), we have
(3.2)
$\dfrac{d\psi}{dr}=\dfrac{2(1-r)}{1+r}\dfrac{\mathcal{E}^{\prime}\mathcal{K}+\mathcal{E}\mathcal{K}^{\prime}-\mathcal{K}\mathcal{K}^{\prime}}{(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})^{2}}=\dfrac{\pi}{1-r^{2}}\left(\frac{1-r}{\mathcal{E}^{\prime}-r\mathcal{K}^{\prime}}\right)^{2},$
which is positive and strictly increasing by Lemma 2.19(4). Hence $\psi(r)$ is
strictly increasing and convex, and consequently $\psi(r)/r$ is strictly
increasing by the monotone form of l’Hôpital’s rule. ∎
###### Proof of Theorem 1.7.
By simple calculations, the first identity follows from the definition of
$\psi$ and Lemma 2.8. The second identity follows from the first one with the
change of parameter $r\mapsto\sqrt{(1-r)/(1+r)}$. ∎
###### Corollary 3.3.
$\psi(3-2\sqrt{2})=1.$
###### Proof.
Let $r=1/\sqrt{2}$. Then
$(1-r)/(1+r)=3-2\sqrt{2}=(1-r^{\prime})/(1+r^{\prime})$, and the second
identity in the Theorem 1.7 implies $\psi(3-2\sqrt{2})=1.$ ∎
###### Remark 3.4.
Let $\Delta$ be the family of curves lying outside the rectangle $R$ and
joining the opposite sides of length $a$. Then a basic fact is
$\mathcal{M}(\Gamma)=1/\mathcal{M}(\Delta).$
By (1.5) and (2.20), we have
$\mathcal{M}(\Gamma)=\mu(r)/\pi,\quad\mbox{and}\quad\mathcal{M}(\Delta)=\mu(s)/\pi$
with $r=\psi^{-1}(a/b)$ and $s=\psi^{-1}(b/a).$ By the identity [AVV1,
Exercises 5.68(2)]
$\mu(r^{2})\mu\left(\left(\dfrac{1-r}{1+r}\right)^{2}\right)=\pi^{2},$
it is easy to see that $\mathcal{M}(\Gamma)=1/\mathcal{M}(\Delta)$ is
equivalent to $s=((1-\sqrt{r})/(1+\sqrt{r}))^{2}.$ Since
$s=\psi^{-1}(b/a)=\psi^{-1}(1/\psi(r))$, we have
$s=((1-\sqrt{r})/(1+\sqrt{r}))^{2}$ which is equivalent to
$\dfrac{1}{\psi(r)}=\psi\left(\left(\dfrac{1-\sqrt{r}}{1+\sqrt{r}}\right)^{2}\right).$
###### Proof of Theorem 1.8.
The theorem follows from
$f(r)=\dfrac{(1-\sqrt{r})^{2}\psi(r)}{r}=2\dfrac{\mathcal{E}-(1-r)\mathcal{K}}{r}\dfrac{(1-\sqrt{r})^{2}}{\mathcal{E}^{\prime}-r\mathcal{K}^{\prime}},$
since $(\mathcal{E}-(1-r)\mathcal{K})/r$ and
$(1-\sqrt{r})^{2}/(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})$ are both
decreasing by Lemma 2.19(2), (5), respectively. The limiting values are clear
by Lemma 2.19(2), (5). ∎
###### Corollary 3.5.
The function $g(r)=(1-\sqrt{r})\operatorname{arth}(1-\sqrt{r})\psi(r)/r$ is
strictly decreasing from $(0,1)$ onto $(4/\pi,\infty)$.
###### Proof.
This follows from Theorem 1.8, since
$g(r)=f(r)\operatorname{arth}(1-\sqrt{r})/(1-\sqrt{r})$ where $f(r)$ is as in
Theorem 1.8 and $\operatorname{arth}(1-\sqrt{r})/(1-\sqrt{r})$ is strictly
decreasing from $(0,1)$ onto $(1,\infty)$. ∎
Since the bounds for $\psi$ in (1.6) and the Theorem 1.8 are not comparable in
the whole interval $(0,1)$, we could combine them to get the following
inequalities:
###### Corollary 3.6.
For $0<r<1$,
$\max\left\\{\dfrac{\pi
r}{(1-r)^{2}},\dfrac{4r}{\pi(1-\sqrt{r})^{2}}\right\\}<\psi(r)<\min\left\\{\dfrac{16r}{\pi(1-r)^{2}},\dfrac{\pi
r}{(1-\sqrt{r})^{2}}\right\\}.$
###### Proof of Theorem 1.9.
Let $r=1/\operatorname{ch}(x)$ and $s=1/\operatorname{ch}(y)$. Then
$dr/dx=-\operatorname{sh}(x)/\operatorname{ch}^{2}(x)=-rr^{\prime}$ and
$f^{\prime}(x)=-\pi
rr^{\prime}\dfrac{1-r}{1+r}\dfrac{1}{(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})^{2}}=-\pi
g(r),$
where
$g(r)=rr^{\prime}(1-r)/((1+r)(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})^{2})$.
By the change of variable $r=(1-t)/(1+t)$ and using Landen’s transformations
(2.5) and (2.7), we have
$g\left(\dfrac{1-t}{1+t}\right)=\dfrac{1-t}{2}\frac{t^{3/2}}{(\mathcal{E}-t^{\prime
2}\mathcal{K})^{2}}$
which is decreasing in $t$ by Lemma 2.18(1), and consequently, $f^{\prime}(x)$
is increasing in $x$. Therefore, $f$ is decreasing and convex on $(0,\infty)$.
In particular, we have $f((x+y)/2)\leq(f(x)+f(y))/2$, with equality if and
only if $x=y$. Now
$\operatorname{ch}^{2}\left(\dfrac{x+y}{2}\right)=\dfrac{1+rs+r^{\prime}s^{\prime}}{2rs}.$
Hence
$f\left(\frac{x+y}{2}\right)\leq\frac{f(x)+f(y)}{2}$
gives
$\psi(r)+\psi(s)\geq
2\psi\left(\frac{\sqrt{2rs}}{\sqrt{1+rs+r^{\prime}s^{\prime}}}\right),$
with equality if and only if $r=s$. ∎
###### Remark 3.7.
It is clear that $f(x)$ is decreasing and $2f(x+y)\leq f(x)+f(y)$. Since
$\operatorname{ch}(x+y)=\dfrac{1+r^{\prime}s^{\prime}}{rs},$
we have
$2\psi\left(\dfrac{rs}{1+r^{\prime}s^{\prime}}\right)\leq\psi(r)+\psi(s)$
which is weaker than the inequality (1.10).
A function $f:I\to J$ is called $H_{p,q}-$convex (concave) if it satisfies
$f(H_{p}(x,y))\leq(\geq)H_{q}(f(x),f(y))$
for all $x,y\in I$ and strictly $H_{p,q}-$convex (concave) if the inequality
is strict, except for $x=y$. Recently, many authors investigated the
$H_{p,q}-$convexity (concavity) of special functions, see [AVV2, BalPV, Ba,
BaPV, CWZQ, WZJ]. The following theorems give some functional inequalities by
studying the generalized convexity (concavity) of the function $\psi$.
###### Theorem 3.8.
The function $f(x)=\log(1/\psi(e^{-x}))$ is strictly increasing and concave
from $(0,\infty)$ onto $(-\infty,\infty)$. In particular, for $r,s\in(0,1)$,
$\psi(\sqrt{rs})\leq\sqrt{\psi(r)\psi(s)}$
with equality if and only if $r=s$.
###### Proof.
Let $r=e^{-x}$ and $s=e^{-y}$. Then $dr/dx=-r$ and
$f^{\prime}(x)=\dfrac{r\psi^{\prime}(r)}{\psi(r)}=\frac{\pi}{2}\dfrac{r}{\mathcal{E}-(1-r)\mathcal{K}}\dfrac{1-r}{(1+r)(\mathcal{E}^{\prime}-r\mathcal{K}^{\prime})}$
which is positive and increasing in $r$ by Lemma 2.19(2) and (7), hence
decreasing in $x$. Therefore, $f$ is strictly increasing and concave on
$(0,\infty)$. In particular, we have $f((x+y)/2)\geq(f(x)+f(y))/2$, with
equality if and only if $x=y$. This gives
$\psi(\sqrt{rs})\leq\sqrt{\psi(r)\psi(s)}$
with equality if and only if $r=s$. ∎
###### Proof of Theorem 1.11.
For $p=0$, the inequality is from Theorem 3.8. Now we assume that $p\neq 0$.
Let $0<x<y<1$ and $t=((x^{p}+y^{p})/2)^{1/p}>x$. Define
$f(x)=\psi(t)^{p}-\dfrac{\psi(x)^{p}+\psi(y)^{p}}{2}.$
By differentiation, we have $dt/dx=\frac{1}{2}(x/t)^{p-1}$ and
$\displaystyle f^{\prime}(x)$ $\displaystyle=$
$\displaystyle\dfrac{1}{2}p\psi(t)^{p-1}\psi^{\prime}(t)\left(\dfrac{x}{t}\right)^{p-1}-\dfrac{1}{2}p\psi(x)^{p-1}\psi^{\prime}(x)$
$\displaystyle=$
$\displaystyle\dfrac{p}{2}x^{p-1}\left(\left(\dfrac{\psi(t)}{t}\right)^{p-1}\psi^{\prime}(t)-\left(\dfrac{\psi(x)}{x}\right)^{p-1}\psi^{\prime}(x)\right).$
We first consider the case of $p>0$. Previous calculation gives
$\displaystyle f^{\prime}(x)$ $\displaystyle=$
$\displaystyle\dfrac{p}{2}x^{p-1}\left(\left(\dfrac{\psi(t)}{t}\right)^{-1}\psi^{\prime}(t)\left(\dfrac{\psi(t)}{t}\right)^{p}-\left(\dfrac{\psi(x)}{x}\right)^{-1}\psi^{\prime}(x)\left(\dfrac{\psi(x)}{x}\right)^{p}\right)$
$\displaystyle=$
$\displaystyle\frac{{\pi}p\,x^{p-1}}{4}\left(\dfrac{t}{\mathcal{E}(t)-(1-t)\mathcal{K}(t)}\dfrac{1-t}{(1+t)(\mathcal{E}^{\prime}(t)-t\mathcal{K}^{\prime}(t))}\left(\dfrac{\psi(t)}{t}\right)^{p}\right.$
$\displaystyle\qquad\left.-\dfrac{x}{\mathcal{E}(x)-(1-x)\mathcal{K}(x)}\dfrac{1-x}{(1+x)(\mathcal{E}^{\prime}(x)-x\mathcal{K}^{\prime}(x))}\left(\dfrac{\psi(x)}{x}\right)^{p}\right)$
which is positive by Lemma 2.19(2),(7) and Theorem 3.1 since $t>x$ and $p>0$.
Hence $f$ is strictly increasing and $f(x)<f(y)=0$. This implies that
$\psi\left(\left(\frac{x^{p}+y^{p}}{2}\right)^{1/p}\right)\leq\left(\frac{\psi(x)^{p}+\psi(y)^{p}}{2}\right)^{1/p}.$
For the case of $p\leq-1$, by previous calculation we have
$f^{\prime}(x)=\dfrac{p}{2}x^{p-1}\left(\left(\dfrac{\psi(t)}{t}\right)^{-2}\psi^{\prime}(t)\left(\dfrac{\psi(t)}{t}\right)^{p+1}-\left(\dfrac{\psi(x)}{x}\right)^{-2}\psi^{\prime}(x)\left(\dfrac{\psi(x)}{x}\right)^{p+1}\right).$
Since $(\psi(x)/x)^{p+1}$ is decreasing, we only need to prove
$(\psi(x)/x)^{-2}\psi^{\prime}(x)$ is strictly decreasing in $(0,1)$. In fact,
with the change of variable $x\mapsto(1-t)/(1+t)$,
$\left(\dfrac{\psi(x)}{x}\right)^{-2}\psi^{\prime}(x)=\dfrac{\pi}{4}\left(\dfrac{x(1-x)}{x^{\prime}(\mathcal{E}-(1-x)\mathcal{K})}\right)^{2}=\dfrac{\pi}{4}\left(\dfrac{t^{\prime
2}}{\mathcal{E}^{\prime}-t^{2}\mathcal{K}^{\prime}}\dfrac{\sqrt{t}}{1+t}\right)^{2}$
which is a product of two positive and strictly increasing functions of $t$ by
Lemma 2.18(1). Hence $f^{\prime}(x)>0$ and $f$ is strictly increasing in
$(0,1)$. Now we have $f(x)<f(y)=0$, and consequently
$\psi\left(\left(\frac{x^{p}+y^{p}}{2}\right)^{1/p}\right)\geq\left(\frac{\psi(x)^{p}+\psi(y)^{p}}{2}\right)^{1/p}$
since $p$ is negative.
The equality case is obvious. This completes the proof. ∎
## 4\. Applications
In this section we always denote $R=[0,1]\times[0,b]$. Let $\Gamma_{b}$ and
$\Delta_{b}$ be the families of curves joining the opposite sides of length
$b$ of the rectangle $R$, in the exterior and interior of the rectangle,
respectively. It is well-known that $\mathcal{M}(\Delta_{b})=b$. By the
formula of Duren and Pfaltzgraff (1.5), we have
$\mathcal{M}(\Gamma_{b})=\dfrac{1}{\pi}\mu(\psi^{-1}(1/b)).$
Setting $r=\sqrt{2}-1$ in (2.14), we get
$\mathcal{K}^{\prime}(3-2\sqrt{2})=2\mathcal{K}(3-2\sqrt{2})$. By Corollary
3.3,
$\mathcal{M}(\Gamma_{1})=\dfrac{1}{\pi}\mu(\psi^{-1}(1))=1=\mathcal{M}(\Delta_{1}).$
Now we will study the behavior of the modulus $\mathcal{M}(\Gamma_{b})$ with
respect to the sides of length $b$. The following Theorem 4.2 shows
(4.1)
$\left\\{\begin{array}[]{ll}\mathcal{M}(\Gamma_{b})>\mathcal{M}(\Delta_{b}),&\mbox{for}\quad
0<b<1,\\\ \mathcal{M}(\Gamma_{b})<\mathcal{M}(\Delta_{b}),&\mbox{for}\quad
b>1.\end{array}\right.$
###### Theorem 4.2.
There exists a number $r_{0}=8.24639\ldots$ such that the function
$f(r)=\dfrac{1}{\pi}\mu(\psi^{-1}(r))-\frac{1}{r}$
is strictly increasing in $(0,r_{0})$ and decreasing in $(r_{0},\infty)$, with
the limiting value $f(\infty)=0$. In particular,
$\dfrac{1}{\pi}\mu(\psi^{-1}(r))<\frac{1}{r},\quad\mbox{for}\quad 0<r<1,$
and
$\dfrac{1}{\pi}\mu(\psi^{-1}(r))>\frac{1}{r},\quad\mbox{for}\quad r>1.$
###### Proof.
Let $s=\psi^{-1}(r)$. Then $r=\psi(s)$ and, by the derivative formula (3.2),
$\dfrac{ds}{dr}=(\dfrac{dr}{ds})^{-1}=\dfrac{s^{\prime
2}}{\pi}\left(\dfrac{\mathcal{E}^{\prime}(s)-s\mathcal{K}^{\prime}(s)}{1-s}\right)^{2}.$
By differentiation and
$\dfrac{d\mu(s)}{ds}=\dfrac{-\pi^{2}}{4ss^{\prime 2}\mathcal{K}(s)^{2}},$
we have
$\displaystyle f^{\prime}(r)$ $\displaystyle=$
$\displaystyle\dfrac{1}{\pi}\dfrac{d\mu}{ds}\dfrac{ds}{dr}+\dfrac{1}{r^{2}}$
$\displaystyle=$
$\displaystyle\dfrac{1}{r^{2}}-\dfrac{1}{4s\mathcal{K}(s)^{2}}\left(\dfrac{\mathcal{E}^{\prime}(s)-s\mathcal{K}^{\prime}(s)}{1-s}\right)^{2}$
$\displaystyle=$
$\displaystyle\dfrac{1}{\psi(s)^{2}}\left(1-\left(\dfrac{\mathcal{E}(s)-(1-s)\mathcal{K}(s)}{\sqrt{s}(1-s)\mathcal{K}(s)}\right)^{2}\right).$
which is positive in $(0,r_{0})$ and negative in $(r_{0},\infty)$ with
$r_{0}=\psi(0.479047\ldots)=8.24639\ldots$ by Lemma 2.19(8). Hence $f$ is
strictly increasing in $(0,r_{0})$ and decreasing in $(r_{0},\infty)$. Since
$f(1)=0$ and $f(\infty)=0$, we have $f(r)<0$ for $r\in(0,1)$ and $f(r)>0$ for
$r\in(1,\infty)$. ∎
The next theorem shows that the modulus $\mathcal{M}(\Gamma_{b})$ has a
logarithmic growth with respect to the length of side $b$.
###### Theorem 4.3.
For $b\in(0,\infty)$,
(4.4) $L(b)<\mathcal{M}(\Gamma_{b})<U(b),$
where
(4.5) $\displaystyle L(b)$ $\displaystyle:=$
$\displaystyle\dfrac{2}{\pi}\left(1-\left(1+\sqrt{4b/\pi}\right)^{-4}\right)^{1/4}\,\log\left(2\left(1+\sqrt{{4b}/{\pi}}\right)\right)$
$\displaystyle>$
$\displaystyle\dfrac{2}{\pi}\left(1-\left(1+\sqrt{4b/\pi}\right)^{-1}\right)\,\log\left(2\left(1+\sqrt{{4b}/{\pi}}\right)\right),$
and
(4.6) $\displaystyle U(b)$ $\displaystyle:=$
$\displaystyle\dfrac{1}{\pi}\log\left(2\left(1+\sqrt{\pi
b}\right)^{2}\left(1+\sqrt{1-\left(1+\sqrt{\pi b}\right)^{-4}}\right)\right)$
$\displaystyle<$ $\displaystyle\dfrac{2}{\pi}\log\left(2\left(1+\sqrt{\pi
b}\right)\right).$
###### Proof.
By Theorem 1.8 we have
$\left(\dfrac{\sqrt{r}}{\sqrt{\pi}+\sqrt{r}}\right)^{2}<s=\psi^{-1}(r)<\left(\dfrac{\sqrt{r}}{\sqrt{4/\pi}+\sqrt{r}}\right)^{2},\quad
r\in(0,\infty).$
By [AVV1, Theorem 5.13(4),(5)],
$\sqrt{s^{\prime}}\log{\dfrac{4}{s}}<\mu(s)<\log{\dfrac{2(1+s^{\prime})}{s}},\quad
s\in(0,1).$
Combining the above inequalities and replacing $r$ with $1/b$, we get the
inequalities (4.4). The inequality (4.5) follows from the inequality
$1-a^{x}>(1-a)^{x}$ for $a\in(0,1)$ and $x\in(1,\infty)$. The inequality (4.6)
is obvious. ∎
###### Theorem 4.7.
For $a,b\in(0,\infty)$,
1. (1)
$\mathcal{M}(\Gamma_{2ab/(a+b)})\leq\sqrt{\mathcal{M}(\Gamma_{a})\mathcal{M}(\Gamma_{b})}\leq\dfrac{\mathcal{M}(\Gamma_{a})+\mathcal{M}(\Gamma_{b})}{2}\leq\mathcal{M}(\Gamma_{(a+b)/2})$;
2. (2)
$\left\\{\begin{array}[]{ll}\mathcal{M}(\Gamma_{H_{p}(a,b)})\leq
H_{p}(\mathcal{M}(\Gamma_{a}),\mathcal{M}(\Gamma_{b})),&p\leq-1,\vspace{1mm}\\\
\mathcal{M}(\Gamma_{H_{p}(a,b)})\geq
H_{p}(\mathcal{M}(\Gamma_{a}),\mathcal{M}(\Gamma_{b})),&p\geq 1.\\\
\end{array}\right.$
Equality holds in each case if and only if $a=b$.
###### Proof.
In part (1), the second inequality is clear. For the third inequality, let
$s=\psi^{-1}(1/a)$, $t=\psi^{-1}(1/b)$. Then
$\displaystyle\dfrac{\mathcal{M}(\Gamma_{a})+\mathcal{M}(\Gamma_{b})}{2}$
$\displaystyle=$ $\displaystyle\dfrac{1}{\pi}\dfrac{\mu(s)+\mu(t)}{2}$
$\displaystyle\leq$
$\displaystyle\dfrac{1}{\pi}\mu(\sqrt{st})\leq\dfrac{1}{\pi}\mu(H_{-1}(s,t))$
$\displaystyle\leq$
$\displaystyle\dfrac{1}{\pi}\mu(\psi^{-1}(H_{-1}(\psi(s),\psi(t))))$
$\displaystyle=$ $\displaystyle\dfrac{1}{\pi}\mu(\psi^{-1}(H_{-1}(1/a,1/b)))$
$\displaystyle=$ $\displaystyle\mathcal{M}(\Gamma_{(a+b)/2}),$
where the first inequality follows from [AVV1, Theorem 5.12(1)] (also see
[WZJ, Theorem]) and the third inequality follows from Theorem 1.11. Let
$m(a)=\mu(\psi^{-1}(a))/\pi$ and $u=\psi^{-1}(a)$. By logarithmic
differentiation, we have
$\dfrac{d}{da}\log
m(a)=-\dfrac{1}{2u\mathcal{K}^{\prime}(u)\mathcal{K}(u)}\left(\dfrac{\mathcal{E}^{\prime}(u)-u\mathcal{K}^{\prime}(u)}{1-u}\right)^{2},$
which is strictly increasing in $u$ by Lemma 2.18(2) and Lemma 2.19(4), and
hence strictly increasing in $a$. Since $m(a)$ is logarithmic convex, we have
$m\left(\dfrac{a+b}{2}\right)\leq\sqrt{m(a)m(b)},$
which implies the first inequality in part (1) by replacing $a,b$ with
$1/a,1/b$, respectively.
For the part (2), let $M(x):=\mathcal{M}(\Gamma_{x})$. Let $0<x<y<1$ and
$t=((x^{p}+y^{p})/2)^{1/p}>x$. Define
$f(x)=M(t)^{p}-\dfrac{M(x)^{p}+M(y)^{p}}{2}.$
By differentiation, we have $dt/dx=\frac{1}{2}(x/t)^{p-1}$ and
(4.8)
$f^{\prime}(x)=\dfrac{p}{2}x^{p-1}\left(\left(\dfrac{M(t)}{t}\right)^{p-1}M^{\prime}(t)-\left(\dfrac{M(x)}{x}\right)^{p-1}M^{\prime}(x)\right).$
Let $M(x)=m(a)$, then $a=1/x=\psi(u)$. Now we have
$\displaystyle\left(\dfrac{M(x)}{x}\right)^{p-1}M^{\prime}(x)$
$\displaystyle=$ $\displaystyle(m(a)a)^{p-1}m^{\prime}(a)(-a^{2})$
$\displaystyle=$
$\displaystyle\left(\dfrac{\mu(u)\psi(u)}{\pi}\right)^{p-1}\psi(u)^{2}\dfrac{1}{4u\mathcal{K}(u)^{2}}\left(\dfrac{\mathcal{E}^{\prime}(u)-u\mathcal{K}^{\prime}(u)}{1-u}\right)^{2}$
$\displaystyle=$
$\displaystyle\left(\dfrac{\mu(u)\psi(u)}{\pi}\right)^{p-1}\left(\dfrac{\mathcal{E}(u)-(1-u)\mathcal{K}(u)}{\sqrt{u}(1-u)\mathcal{K}(u)}\right)^{2},$
which is strictly increasing in $u$ by Lemmas 2.21 and 2.19(8), and hence
strictly decreasing in $x$ for each $p\geq 1$. This implies that
$f^{\prime}(x)<0$ if $p\geq 1$.
For the case of $p\leq-1$, we have
$\displaystyle\left(\dfrac{M(x)}{x}\right)^{p-1}M^{\prime}(x)$
$\displaystyle=$
$\displaystyle(m(a)a)^{p-1}m^{\prime}(a)(-a^{2})=-(m(a)a)^{p+1}m(a)^{-2}m^{\prime}(a)$
$\displaystyle=$
$\displaystyle\left(\dfrac{\mu(u)\psi(u)}{\pi}\right)^{p+1}\dfrac{1}{u\mathcal{K}^{\prime}(u)^{2}}\left(\dfrac{\mathcal{E}^{\prime}(u)-u\mathcal{K}^{\prime}(u)}{1-u}\right)^{2},$
which is strictly decreasing in $u$ by Lemmas 2.21, 2.18(2) and 2.19(4), and
hence strictly increasing in $x$ for each $p\leq-1$. Since $p$ is negative,
this still implies that $f^{\prime}(x)<0$.
It is easy to see that $f^{\prime}(x)<0$ implies the inequalities in the part
(2). ∎
###### Open problem 4.9.
What is the exact domain of $p$ for which the function $\psi$ is
$H_{p,p}$-convex (concave)? More generally, find the exact $(p,q)$ domain for
which the function $\psi$ is $H_{p,q}$-convex (concave). The same questions
can be asked for the modulus $\mathcal{M}(\Gamma_{b})$.
### Acknowledgments
The research of Matti Vuorinen was supported by the Academy of Finland,
Project 2600066611. Xiaohui Zhang is indebted to the CIMO (Grant TM-09-6629)
and the Finnish National Graduate School of Mathematics and its Applications
for financial support. Both authors wish to thank Árpád Baricz and the referee
for their helpful comments on the manuscript.
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|
arxiv-papers
| 2011-11-16T14:26:14 |
2024-09-04T02:49:24.379551
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Matti Vuorinen and Xiaohui Zhang",
"submitter": "Xiaohui Zhang",
"url": "https://arxiv.org/abs/1111.3812"
}
|
1111.3824
|
Rev. 10/III/12 JM
Higher-order Erdős–Szekeres theorems
###### Abstract
Let $P=(p_{1},p_{2},\ldots,p_{N})$ be a sequence of points in the plane, where
$p_{i}=(x_{i},y_{i})$ and $x_{1}<x_{2}<\cdots<x_{N}$. A famous 1935
Erdős–Szekeres theorem asserts that every such $P$ contains a monotone
subsequence $S$ of $\lceil\sqrt{N}\,\rceil$ points. Another, equally famous
theorem from the same paper implies that every such $P$ contains a convex or
concave subsequence of $\Omega(\log N)$ points.
Monotonicity is a property determined by pairs of points, and convexity
concerns triples of points. We propose a generalization making both of these
theorems members of an infinite family of Ramsey-type results. First we define
a $(k+1)$-tuple $K\subseteq P$ to be _positive_ if it lies on the graph of a
function whose $k$th derivative is everywhere nonnegative, and similarly for a
_negative_ $(k+1)$-tuple. Then we say that $S\subseteq P$ is _$k$ th-order
monotone_ if its $(k+1)$-tuples are all positive or all negative.
We investigate quantitative bound for the corresponding Ramsey-type result
(i.e., how large $k$th-order monotone subsequence can be guaranteed in every
$N$-point $P$). We obtain an $\Omega(\log^{(k-1)}N)$ lower bound
($(k-1)$-times iterated logarithm). This is based on a quantitative Ramsey-
type theorem for _transitive colorings_ of the complete $(k+1)$-uniform
hypergraph (these were recently considered by Pach, Fox, Sudakov, and Suk).
For $k=3$, we construct a geometric example providing an $O(\log\log N)$ upper
bound, tight up to a multiplicative constant. As a consequence, we obtain
similar upper bounds for a Ramsey-type theorem for _order-type homogeneous_
subsets in ${\mathbb{R}}^{3}$, as well as for a Ramsey-type theorem for
hyperplanes in ${\mathbb{R}}^{4}$ recently used by Dujmović and Langerman.
## 1 Introduction
In this paper we mainly consider sets $P=\\{p_{1},p_{2},\ldots,p_{N}\\}$ of
points in the plane, where $p_{i}=(x_{i},y_{i})$. We always assume that no two
of the $x$-coordinates coincide, and unless stated otherwise, we also assume
that the $p_{i}$ are numbered so that $x_{1}<x_{2}<\cdots<x_{N}$ (the same
also applies to subsets of $P$, which we will enumerate in the order of
increasing $x$-coordinates).
Two theorems of Erdős and Szekeres. Among simple results in combinatorics,
only few can compete with the following one in beauty and usefulness:
###### Theorem 1.1 (Erdős–Szekeres on monotone subsequences [ES35])
For every positive integer $n$, among every $N=(n-1)^{2}+1$ points
$p_{1},\ldots,p_{N}\in{\mathbb{R}}^{2}$ as above, one can always choose a
_monotone subset_ of at least $n$ points, i.e., indices
$i_{1}<i_{2}<\cdots<i_{n}$ such that either $y_{i_{1}}\leq
y_{i_{2}}\leq\cdots\leq y_{i_{n}}$ or $y_{i_{1}}\geq y_{i_{2}}\geq\cdots\geq
y_{i_{n}}$.
See, for example, Steele [Ste95] for a collection of six nice proofs and some
applications. For many purposes, it is more natural to view the above theorem
as a purely combinatorial result about permutations, but here we prefer the
geometric formulation (which is also similar to the one in the original
Erdős–Szekeres paper).
Another result of the same paper of Erdős and Szekeres is the following well-
known gem in discrete geometry:111Somewhat unfortunately, the name
Erdős–Szekeres theorem refers to Theorem 1.1 in some sources and to Theorem
1.2 or similar statements in other sources.
###### Theorem 1.2 (Erdős–Szekeres on convex/concave configurations [ES35])
For every positive integer $n$, among every $N={2n-4\choose n-2}+1\approx
4^{n}/\sqrt{n}$ points $p_{1},\ldots,p_{N}\in{\mathbb{R}}^{2}$ as above, one
can always choose a _convex configuration_ or a _concave configuration_ of $n$
points, i.e., indices $i_{1}<i_{2}<\cdots<i_{n}$ such that the slopes of the
segments $p_{i_{j}}p_{i_{j+1}}$, $j=1,2,\ldots,n-1$, are either monotone
nondecreasing or monotone nonincreasing.
See, e.g., [MS00, Mat02] for proofs and surveys of developments around this
result.
$k$-general position. To simplify our forthcoming discussion, at some places
it will be convenient to assume that the considered point sets are in a
“sufficiently general” position. Namely, we define a set $P$ to be in _$k$
-general position_ if no $k+1$ points of $P$ lie on the graph of a polynomial
of degree at most $k-1$. In particular, $1$-general position requires that no
two $y$-coordinates coincide, and $2$-general position means the usual general
position, i.e., no three points collinear.
$k$th-order monotone subsets. Here we propose a view of Theorems 1.1 and 1.2
as the first two members in an infinite sequence of Ramsey-type results about
planar point sets.222There is also a (trivial) 0th member, namely, the
statement that in every $P$, at least half of the points either have all
$y$-coordinates nonnegative or have or all $y$-coordinates nonpositive.
In Theorem 1.1, monotonicity of a subset is a property of _pairs_ of points of
the subset, and actually, it suffices to look at pairs of consecutive points.
Similarly, convexity or concavity of a configuration in Theorem 1.2 is a
property of triples, and again it is enough to look at consecutive triples.
In the former case, we are considering the slope of the segment determined by
a pair of points, which can be thought of as the first derivative. In the
latter case, a triple is convex iff its points lie on the graph of a smooth
convex function, i.e., one with nonnegative second derivative everywhere.
With this point of view, it is natural to define a $(k+1)$-tuple $K\subseteq
P$ to be _positive_ if it lies on the graph of a function whose $k$-th
derivative (exists and) is everywhere nonnegative, and similarly for a
_negative_ $(k+1)$-tuple (in Section 2, we will provide several other,
equivalent characterizations of these properties). Then we say that an
arbitrary subset $S\subseteq P$ is _$k$ th-order monotone_ if its
$(k+1)$-tuples are all positive or all negative.
First-order monotonicity is obviously equivalent to monotonicity as in Theorem
1.1, and second-order monotonicity is equivalent to convexity/concavity as in
Theorem 1.2. We will also see (Lemma 2.5) that, to certify $k$th-order
monotonicity, it is enough to consider all $(k+1)$-tuples of _consecutive_
points.
Let us remark that every $(k+1)$-tuple $K$ is positive or negative, and
moreover, if $K$ is in $k$-general position, it cannot be both positive and
negative (Corollary 2.3). We will write $\mathop{\rm sgn}\nolimits(K)=+1$ if
$K$ is positive and $\mathop{\rm sgn}\nolimits(K)=-1$ if $K$ is negative.
Ramsey’s theorem, quantitative bounds, and transitive colorings. Using the
just mentioned facts, one can immediately derive a Ramsey-type theorem for
$k$th-order monotone subsets from Ramsey’s theorem.
###### Proposition 1.3
For every $k$ and $n$ there exists $N$ such that every $N$-point planar set in
$k$-general position contains an $n$-point $k$th-order monotone subset.
Proof. We recall Ramsey’s theorem (for two colors; see, e.g., Graham,
Rothschild, and Spencer [GRS90]): for every $\ell$ and $n$ there exists $N$
such that for every coloring of the set ${X\choose\ell}$ of all $\ell$-element
subsets of an $N$-element set $X$ there exists an $n$-element _homogeneous_
set $Y\subseteq X$, i.e., a subset in which all $\ell$-tuples have the same
color. The smallest $N$ for which the claim holds is usually denoted by
$R_{\ell}(n)$.
In our case, we set $X=P$ and color each $(k+1)$-tuple $K\subseteq P$ with the
color $\mathop{\rm sgn}\nolimits(K)\in\\{\pm 1\\}$. Then homogeneous subsets
are exactly $k$th-order monotone subsets. $\Box$
Let us denote by $\operatorname{ES}_{k}(n)$ the smallest value of $N$ for
which the claim in this proposition holds. We have
$\operatorname{ES}_{1}(n)\leq(n-1)^{2}+1$ and
$\operatorname{ES}_{2}(n)\leq{2n-4\choose n-2}+1$ according to Theorems 1.1
and 1.2, respectively; moreover, these inequalities actually hold with
equality [ES35]. Our main goal is to estimate the order of magnitude of
$\operatorname{ES}_{k}(n)$ for $k\geq 3$.
The above proof gives $\operatorname{ES}_{k}(n)\leq R_{k+1}(n)$. However, for
$k=1$, and most likely for all $k$, the order of magnitude of $R_{k+1}(n)$ is
much larger than that of $\operatorname{ES}_{k}(n)$. Indeed, considering $k$
fixed and $n$ large, the best known lower and upper bounds of $R_{k+1}(n)$ are
of the form333We employ the usual asymptotic notation for comparing functions:
$f(n)=O(g(n))$ means that $|f(n)|\leq C|g(n)|$ for some $C$ and all $n$, where
$C$ may depend on parameters declared as constants (in our case on $k$);
$f(n)=\Omega(g(n))$ is equivalent to $g(n)=O(f(n))$; and $f(n)=\Theta(g(n))$
means that both $f(n)=O(g(n))$ and $f(n)=\Omega(g(n))$.
$R_{2}(n)=2^{\Theta(n)}$ and, for $k\geq 2$,
$\mathop{\rm twr}\nolimits_{k}(\Omega(n^{2}))\leq R_{k+1}(n)\leq\mathop{\rm
twr}\nolimits_{k+1}(O(n)),$
where the tower function $\mathop{\rm twr}\nolimits_{k}(x)$ is defined by
$\mathop{\rm twr}\nolimits_{1}(x)=x$ and $\mathop{\rm
twr}\nolimits_{i+1}(x)=2^{\mathop{\rm twr}\nolimits_{i}(x)}$. It is widely
believed that the upper bound is essentially the truth. This belief is
supported by known bounds for more than two colors, where the lower bound for
$(k+1)$-tuples is also a tower of height $k+1$; see Conlon, Fox, and Sudakov
[CFS11] for a recent improvement and more detailed overview of the known
bounds.
The coloring of the $(k+1)$-tuples in the above proof of Proposition 1.3 is
not arbitrary. In particular, it has a property we call _transitivity_ (see
Lemma 2.5). Transitive colorings were introduced earlier in the recent
preprint Fox et al. [FPSS11, Section 6], under the same name.
To define a transitive coloring in general, we need to consider a hypergraph
whose vertex set is linearly ordered; w.l.o.g. we can identify it with the set
$[N]:=\\{1,2,\ldots,N\\}$. A coloring $c\colon{[N]\choose\ell}\to[m]$ is
_transitive_ if, for every $i_{1},\ldots,i_{\ell+1}\in[N]$,
$i_{1}<\cdots<i_{\ell+1}$, whenever the $\ell$-tuples
$\\{i_{1},\ldots,i_{\ell}\\}$ and $\\{i_{2},\ldots,i_{\ell+1}\\}$ have the
same color, then _all_ $\ell$-element subsets of
$\\{i_{1},\ldots,i_{\ell+1}\\}$ have the same color. Let $R^{\rm
trans}_{\ell}(n)$ denote the Ramsey number for transitive colorings, i.e., the
smallest $N$ such that any transitive coloring of the complete $\ell$-uniform
hypergraph on $[N]$ contains an $n$-element homogeneous subset. We have the
following bound.444By inspecting the proof of the next theorem, it is easy to
verify that the transitivity condition is not used in full strength—it
suffices to assume only that the subsets obtained by omitting one of $i_{2}$,
$i_{3}$ have the same color.
###### Theorem 1.4
For $k=1,2$, we have $R^{\rm trans}_{k+1}(n)=\operatorname{ES}_{k}(n)$, and
for every fixed $k\geq 3$,
$\operatorname{ES}_{k}(n)\leq R^{\rm trans}_{k+1}(n)\leq\mathop{\rm
twr}\nolimits_{k}(O(n)).$
We note that Fox et al. [FPSS11] proved the slightly weaker upper bound
$R^{\rm trans}_{k+1}(n)\leq\mathop{\rm twr}\nolimits_{k}(O(n\log n))$.
The proof of Theorem 1.4 is given in Section 3. The inequality
$\operatorname{ES}_{k}(n)\leq R^{\rm trans}_{k+1}(n)$ is clear (since every
$N$-point set in $k$-general position provides a transitive coloring of
$[N]\choose k+1$). The upper bounds for $R^{\rm trans}_{2}(n)$ and $R^{\rm
trans}_{3}(n)$ follow by translating the proofs of Theorem 1.1 and 1.2 to the
setting of transitive colorings almost word by word, and they are contained in
[FPSS11]. The upper bound on $R^{\rm trans}_{k+1}(n)$ is then obtained by
induction on $k$, with $k=3$ as the base case, following one of the usual
proofs of Ramsey’s theorem.
A set with no large third-order monotone subsets. For $k\leq 2$, the numbers
$\operatorname{ES}_{k}(n)$ (and thus $R^{\rm trans}_{k+1}(n)$) are known
exactly. Our perhaps most interesting result is an asymptotically matching
lower bound for $\operatorname{ES}_{3}(n)$.
###### Theorem 1.5
For all $n\geq 2$ we have $R^{\rm
trans}_{4}(2n+1)\geq\operatorname{ES}_{3}(2n+1)\geq 2^{2^{n-1}}+1$.
Consequently, $\operatorname{ES}_{3}(n)=2^{2^{\Theta(n)}}$.
The proof is given in Section 4. A Ramsey function with known doubly
exponential growth seems to be rare in geometric Ramsey-type problems (a
notable example is a result of Valtr [Val04]).
Order types. Here we change the setting from the plane to ${\mathbb{R}}^{d}$
and we consider an ordered sequence $P=(p_{1},p_{2},\ldots,p_{N})$ in
${\mathbb{R}}^{d}$. This time we do _not_ assume the first coordinates to be
increasing. For simplicity, we assume $P$ to be in general position, which now
means that no $d+1$ points of $P$ lie on a common hyperplane.
We recall that _order type_ of $P$ specifies the orientation of every
$(d+1)$-tuple of points of $P$, and it this way, it describes purely
combinatorially many of the geometric properties of $P$. More formally, the
order type of $P$ is the mapping $\chi\colon{[N]\choose d+1}\to\\{-1,+1\\}$,
where for a $(d+1)$-tuple $I=\\{i_{1},\ldots,i_{d+1}\\}$,
$i_{1}<i_{2}<\cdots<i_{d+1}$, $\chi(I):=\mathop{\rm sgn}\nolimits\det
M(p_{i_{1}},p_{i_{2}},\ldots,p_{i_{d+1}})$, where $M(q_{1},\ldots,q_{d+1})$ is
the $(d+1)\times(d+1)$ matrix whose $j$th column is $(1,q_{j})$, i.e., $1$
followed by the vector of the $d$ coordinates of $q_{j}$. See, e.g., Goodman
and Pollack [GP93] or [Mat02] for more background about order types.
From Ramsey’s theorem for $(d+1)$-tuples, we can immediately derive a Ramsey-
type result for order types: for every $d$ and $n$ there exists $N$ such that
every $N$-point sequence contains an $n$-point subsequence in which all the
$(d+1)$-tuples have the same orientation (we call such a subsequence _order-
type homogeneous_). Let us write $\mathop{\rm OT}\nolimits_{d}(n)$ for the
smallest such $N$.
In Section 5 we first observe that, by simple and probably well known
considerations, $\mathop{\rm OT}\nolimits_{1}(n)=(n-1)^{2}+1$ and $\mathop{\rm
OT}\nolimits_{2}(n)=2^{\Theta(n)}$. For $d\geq 3$, the best upper bound for
$\mathop{\rm OT}\nolimits_{d}(n)$ we are aware of is the one from the Ramsey
argument above, i.e., $\mathop{\rm OT}\nolimits_{d}(n)\leq
R_{d+1}(n)\leq\mathop{\rm twr}\nolimits_{d+1}(O(n))$. In particular, for
$\mathop{\rm OT}\nolimits_{3}(n)$ this upper bound is triply exponential; in
Section 5 we prove a doubly exponential lower bound.
###### Proposition 1.6
For all $d$ and $n$, $\mathop{\rm
OT}\nolimits_{d}(n)\geq\operatorname{ES}_{d}(n)$. In particular, $\mathop{\rm
OT}\nolimits_{3}(n)=2^{2^{\Omega(n)}}$.
A Ramsey-type result for hyperplanes. Let us consider a finite set $H$ of
hyperplanes in ${\mathbb{R}}^{d}$ in general position (every $d$ intersecting
at a single point). Let us say that $H$ is _one-sided_ if $V(H)$, the vertex
set of the arrangement of $H$, lies completely on one side of the coordinate
hyperplane $x_{d}=0$.
Let $\mathop{\rm OSH}\nolimits_{d}(n)$ be the smallest $N$ such that every set
$H$ of $N$ hyperplanes in ${\mathbb{R}}^{d}$ in general position contains a
one-sided subset of $n$ hyperplanes. Ramsey’s theorem for $d$-tuples
immediately gives $\mathop{\rm OSH}\nolimits_{d}(n)\leq R_{d}(n)$ (a $d$-tuple
gets color $+1$ if its intersection has a positive last coordinate, and color
$-1$ otherwise).
Matoušek and Welzl [MW92] observed that, actually, $\mathop{\rm
OSH}\nolimits_{2}(n)=\operatorname{ES}_{1}(n)=(n-1)^{2}+1$, and applied this
in a range-searching algorithm. Recently Dujmović and Langerman [DL11] used
the existence of $\mathop{\rm OSH}\nolimits_{d}(n)$ (essentially Lemma 9 in
the arXiv version of their paper) to prove several interesting results, such
as a ham-sandwich and centerpoint theorems for hyperplanes.
In Section 5 we show that lower bounds for $k$th-order monotone subsets in the
plane can be translated into lower bounds for $\mathop{\rm OSH}\nolimits_{d}$.
###### Proposition 1.7
We have $\mathop{\rm OSH}\nolimits_{d}(n)\geq\operatorname{ES}_{d-1}(n)$, and
in particular, $\mathop{\rm OSH}\nolimits_{3}(n)=2^{\Omega(n)}$ and555An
exponential lower bound for $\mathop{\rm OSH}\nolimits_{3}$ was known to the
authors of [MW92], and perhaps to others as well, but as far as we know, it
hasn’t appeared in print. $\mathop{\rm
OSH}\nolimits_{4}(n)=2^{2^{\Omega(n)}}$.
The lower bounds for $\mathop{\rm OSH}\nolimits_{d}(n)$ can also be translated
into lower bounds in the theorems of Dujmović and Langerman. For example, in
their ham-sandwich theorem, we have $d$ collections $H_{1},\ldots,H_{d}$ of
hyperplanes in ${\mathbb{R}}^{d}$, each of size $N$, and we want a hyperplane
$g$ such that in each $H_{i}$, we can find disjoint subsets $A_{i},B_{i}$ of
$n$ hyperplanes each such $V(A_{i})$ lies on one side of $g$ and $V(B_{i})$ on
the other side.
To derive a lower bound for the smallest necessary $N$, we fix $d$ affinely
independent points $p_{1},\ldots,p_{d}$ in the $x_{d}=0$ hyperplane, and a set
$H$ of $N$ hyperplanes in general position with no one-sided subset of size
$n$. We let $H_{i}$ be an affinely transformed copy of $H$ such that all of
$V(H_{i})$ lies very close to $p_{i}$. Then every potential ham-sandwich
hyperplane $g$ for these $H_{i}$ has to be almost parallel to the $x_{d}=0$
hyperplane, and thus there cannot be $A_{i},B_{i}$ of size $n$ for all $i$.
The work of Fox et al. While preparing a draft of the present paper, we
learned about a recent preprint of Fox, Pach, Sudakov, and Suk [FPSS11]. They
investigated various combinatorial and geometric problems inspired by Theorems
1.1 and 1.2, and as was mentioned above, among others, they introduced
transitive colorings,666With still another geometric source of such colorings
besides the Erdős–Szekeres theorems, namely, noncrossing convex bodies in the
plane but mainly they studied a related but different Ramsey-type quantity:
let $N_{\ell}(q,n)$ be the smallest integer $N$ such that, for every coloring
of ${[N]\choose\ell}$ with $q$ colors, there exists an $n$-element
$I=\\{i_{1},\ldots,i_{n}\\}\subseteq[N]$, $i_{1}<\cdots<i_{n}$, inducing a
_monochromatic monotone path_ , i.e., such that all the $\ell$-tuples of the
form $\\{i_{j},i_{j+1},\ldots,i_{j+\ell-1}\\}$, $j=1,2,\ldots,n-\ell+1$, have
the same color.
They note that $R^{\rm trans}_{\ell}(n)\leq N_{\ell}(2,n)$, and they obtained
the following bounds for $N_{\ell}(2,n)$:
$N_{2}(2,n)=\operatorname{ES}_{1}(n)$, $N_{3}(2,n)=\operatorname{ES}_{2}(n)$,
and for every fixed $k\geq 3$,
$\mathop{\rm twr}\nolimits_{k}(\Omega(n))\leq N_{k+1}(2,n)\leq\mathop{\rm
twr}\nolimits_{k}(O(n\log n)).$
As we mentioned after Theorem 1.4, this also yields an upper bound for $R^{\rm
trans}_{k+1}(n)$ only slightly weaker than the one in that theorem.
Open problems.
1. 1.
We have obtained reasonably tight bounds for $\operatorname{ES}_{3}(n)$, but
the gaps are much more significant for $\operatorname{ES}_{k}(n)$ with $k\geq
4$. According to the cases $k=1,2,3$, one may guess that
$\operatorname{ES}_{k}(n)$ is of order $\mathop{\rm
twr}\nolimits_{k}(\Theta(n))$, and thus that stronger lower bounds are needed,
but a possibility of a better upper bound shouldn’t also be overlooked. This
question looks both interesting and challenging.
2. 2.
A perhaps more manageable task might be a better lower bound for $R^{\rm
trans}_{k}(n)$, $k\geq 4$. A natural approach would be to imitate the
Stepping-Up Lemma used for lower bounds for the Ramsey numbers $R_{k}(n)$
(see, e.g., [CFS11]). But so far we have not succeeded in this, since even if
we start with a transitive coloring of $k$-tuples, we could not guarantee
transitivity for the coloring of $(k+1)$-tuples.
3. 3.
As for order-type homogeneous sequences, for $\mathop{\rm OT}\nolimits_{3}(n)$
we have the lower bound of $2^{2^{\Omega(n)}}$, but upper bound only
$\mathop{\rm twr}\nolimits_{4}(O(n))$ directly from Ramsey’s theorem. It seems
that the colorings given by the order type are not transitive in any
reasonable sense, and we have no good guess of which of the upper and lower
bounds should be closer to the truth. Similar comments apply to the problem
with one-sided subsets of planes in ${\mathbb{R}}^{3}$ (concerning
$\mathop{\rm OSH}\nolimits_{3}(n)$), and the higher-dimensional cases are even
more widely open.
4. 4.
Another interesting question is whether $n\log n$ can be replaced by $n$ in
the upper bound for the quantity $N_{\ell}(2,n)$ considered by Fox et al.
[FPSS11].
5. 5.
In our definition of $k$th-order positivity, every $(k+1)$-tuple of points
should lie on the graph of a function with a nonnegative $k$th derivative, and
different functions can be used for different $(k+1)$-tuples. In an earlier
version of this paper, we conjectured that, assuming $k$-general position, a
single function should suffice for all $(k+1)$-tuples; in other words, that
every $k$th-order monotone finite set finite set in $k$-general position lies
on a graph of a $k$-times differentiable function
$f\colon{\mathbb{R}}\to{\mathbb{R}}$ whose $k$th derivative is everywhere
nonnegative or everywhere nonpositive.
However, Rote [Rot12] disproved this for $k=3$ (while the cases $k=1,2$ do
hold, as is not hard to check). With his kind permission, we reproduce his
example at the end of Section 2.
Naturally, this opens up interesting new questions: How can one characterize
point sets lying on the graph of a function whose $k$th derivative is positive
everywhere? Is there a Ramsey-type theorem for such sets, and if yes, how
large is the corresponding Ramsey function?
## 2 On the definition of $k$th-order monotonicity
Here we provide several equivalent characterizations of $k$th-order
monotonicity of planar point sets and some of their properties. First we
recall several known results.
Divided differences and Newton’s interpolation. Let
$p_{1},p_{2},\ldots,p_{k+1}$ be points in the plane, $p_{i}=(x_{i},y_{i})$,
where the $x_{i}$ are all distinct (but not necessarily increasing). We recall
that the _$k$ th divided difference_
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},p_{2},\ldots,p_{k+1})$
is defined recursively as follows:
$\displaystyle\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{0}(p_{1})$
$\displaystyle:=$ $\displaystyle y_{1}$
$\displaystyle\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},p_{2},\ldots,p_{k+1})$
$\displaystyle:=$
$\displaystyle\frac{\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k-1}(p_{2},p_{3},\ldots,p_{k+1})-\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k-1}(p_{1},p_{2},\ldots,p_{k})}{x_{k+1}-x_{1}}.$
For example,
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{1}(p_{1},p_{2})$
equals the slope of the line $p_{1}p_{2}$. In general, the $k$th divided
difference is related to the $k$th derivative as follows (see, e.g., [Phi03,
Eq. 1.33]; note that the case $k=1$ is the Mean Value Theorem):
###### Lemma 2.1 (Cauchy)
Let the points $p_{1},\ldots,p_{k+1}$, $a:=x_{1}<x_{2}<\cdots<b:=x_{k+1}$, lie
on the graph of a function $f$ such that the $k$th derivative $f^{(k)}$ exists
everywhere on the interval $(a,b)$. Then there exists $\xi\in(a,b)$ such that
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})=\frac{f^{(k)}(\xi)}{k!}.$
We will also need the following result (see, e.g., [Phi03, Eq. 1.11–1.19]).
###### Lemma 2.2 (Newton’s interpolation)
Let $p_{1},\ldots,p_{k+1}\in{\mathbb{R}}^{2}$ be points with distinct
$x$-coordinates (here we need not assume that the $x$-coordinates are
increasing). Then the unique polynomial $f$ of degree at most $k$ whose graph
contains $p_{1},\ldots,p_{k+1}$ is given by
$f(x)=\sum_{i=1}^{k+1}\biggl{(}\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}(p_{1},\ldots,p_{i})\prod_{j=1}^{i-1}(x-x_{j})\biggr{)}$
In particular, the coefficient of $x^{k}$ is
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}(p_{1},\ldots,p_{k+1})$,
and it equals $f^{(k)}(x)/k!$ (which is a constant function).
We recall that a $(k+1)$-tuple $K=\\{p_{1},\ldots,p_{k+1}\\}$ was defined to
be positive if it is contained in the graph of a function having a nonnegative
$k$th derivative everywhere. We obtain the following equivalent
characterization:
###### Corollary 2.3
A $(k+1)$-tuple $K=\\{p_{1},\ldots,p_{k+1}\\}$ is positive iff
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})\geq
0$ (and similarly for a negative $(k+1)$-tuple). If $K$ is in $k$-general
position, we have $\mathop{\rm sgn}\nolimits K=\mathop{\rm
sgn}\nolimits\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})$.
Proof. If $K$ is contained in the graph of $f$ with $f^{(k)}\geq 0$
everywhere, then
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})\geq
0$ by Lemma 2.1.
Conversely, if
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})\geq
0$, then by Lemma 2.2, the unique polynomial of degree at most $k$ whose graph
contains $K$ is the required function with nonnegative $k$th derivative.
If, moreover, $K$ is in $k$-general position, then
$\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{k+1})\neq
0$, and so $K$ cannot be both $k$th-order positive and $k$th-order negative by
Lemma 2.1. $\Box$
We will also need the following criterion for the sign of a $(k+1)$-tuple.
###### Lemma 2.4
Let $K=\\{p_{1},p_{2},\ldots,p_{k+1}\\}$ be a $(k+1)$-tuple of points in
$k$-general position, $x_{1}<\cdots<x_{k+1}$, let $i\in[k+1]$, and let $f_{i}$
be the (unique) polynomial of degree at most $k-1$ whose graph passes through
the points of $K\setminus\\{p_{i}\\}$. Then $\mathop{\rm sgn}\nolimits
K=(-1)^{k-i}$ if $p_{i}$ lies below the graph of $f_{i}$, and $\mathop{\rm
sgn}\nolimits K=(-1)^{k+1-i}$ if $p_{i}$ lies above the graph.
Let $f$ be the polynomial of degree at most $k$ passing through all of $K$. We
use Newton’s interpolation (Lemma 2.2), but with the points reordered so that
$p_{i}$ comes last, and we get that
$f(x)=f_{i}(x)+\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{i-1},p_{i+1},\ldots,p_{k+1},p_{i})\prod_{j\in[k+1]\setminus\\{i\\}}(x-x_{j}).$
Using this with $x=x_{i}$, we get
$\displaystyle\mathop{\rm sgn}\nolimits(y_{i}-f_{i}(x_{i}))$ $\displaystyle=$
$\displaystyle\mathop{\rm sgn}\nolimits(f(x_{i})-f_{i}(x_{i}))$
$\displaystyle=$ $\displaystyle\mathop{\rm
sgn}\nolimits\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{i-1},p_{i+1},\ldots,p_{k+1},p_{i})\cdot\mathop{\rm
sgn}\nolimits\prod_{j\in[k+1]\setminus\\{i\\}}(x_{i}-x_{j}).$
Divided differences are invariant under permutations of the points (as can be
seen, e.g., from Lemma 2.2, since the interpolating polynomial does not depend
on the order of the points), and so $\mathop{\rm
sgn}\nolimits\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{k}(p_{1},\ldots,p_{i-1},p_{i+1},\ldots,p_{k+1},p_{i})=\mathop{\rm
sgn}\nolimits K$. Finally, the product
$\prod_{j\in[k+1]\setminus\\{i\\}}(x_{i}-x_{j})$ has $k+1-i$ negative factors,
thus its sign is $(-1)^{k+1-i}$, and the lemma follows. $\Box$
It remains to prove transitivity.
###### Lemma 2.5
Let $P=\\{p_{1},\ldots,p_{N}\\}$ be a point set in $k$-general position. Then
the $2$-coloring of $(k+1)$-tuples $K\in{P\choose k+1}$ by their sign is
transitive.
Proof. We consider a $(k+2)$-tuple $L=\\{p_{1},\ldots,p_{k+2}\\}$ with
$\mathop{\rm sgn}\nolimits\\{p_{1},\ldots,p_{k+1}\\}=\mathop{\rm
sgn}\nolimits\\{p_{2},\ldots,p_{k+2}\\}=+1$, and we fix
$i\in\\{2,\ldots,{k+1}\\}$. Let $f_{i,k+2}$ be the polynomial of degree at
most $k-1$ passing through $L\setminus\\{p_{i},p_{k+2}\\}$, and similarly for
$f_{1,k+2}$. Our goal is to show that $f_{i,k+2}(x_{k+2})<y_{k+2}$, since this
gives $\mathop{\rm sgn}\nolimits(L\setminus\\{p_{i}\\})=+1$ by Lemma 2.4.
Since $\mathop{\rm sgn}\nolimits(L\setminus\\{p_{1}\\})=+1$, we have
$f_{1,k+2}(x_{k+2})<y_{k+2}$ (Lemma 2.4 again), and so it suffices to prove
$f_{i,k+2}(x_{k+2})<f_{1,k+2}(x_{k+2})$.
Let us consider the polynomial $g:=f_{1,k+2}-f_{i,k+2}$; as explained above,
our goal is proving $\mathop{\rm sgn}\nolimits g(x_{k+2})=+1$. To this end, we
first determine $\mathop{\rm sgn}\nolimits g(x_{1})$: We have
$f_{i,k+2}(x_{1})=y_{1}$ and $\mathop{\rm
sgn}\nolimits(y_{1}-f_{1,k+2}(x_{1}))=(-1)^{k}$ (using $\mathop{\rm
sgn}\nolimits(L\setminus\\{p_{1}\\})=+1$ and Lemma 2.4). Hence $\mathop{\rm
sgn}\nolimits g(x_{1})=(-1)^{k-1}$.
Next, we observe that $g$ is a polynomial of degree at most $k-1$, and it
vanishes at $x_{2},\ldots,x_{i-1},x_{i+1},\ldots,x_{k+1}$. These are $k-1$
distinct values; thus, they include all roots of $g$, and each of them is a
simple root. Consequently, $g$ changes sign $(k-1)$-times between $x_{1}$ and
$x_{k+2}$. Hence, finally, $\mathop{\rm sgn}\nolimits
g(x_{k+2})=(-1)^{k-1}\mathop{\rm sgn}\nolimits g(x_{1})=+1$ as claimed. $\Box$
Rote’s example. Fig. 1 shows a 6-point set $P=\\{p_{1},\ldots,p_{6}\\}$ in
3-general position (no four points on a parabola). It is easy to check 3rd-
order positivity using Lemma 2.4: By transitivity, it suffices to look at
$4$-tuples of consecutive points. For $p_{1},\ldots,p_{4}$ we use the parabola
through $p_{1},p_{2},p_{3}$ (which actually degenerates to the $x$-axis); for
$p_{2},\ldots,p_{5}$ we use the dashed parabola through $p_{2},p_{3},p_{4}$
(which is very close to the $x$-axis in the relevant region); and for
$p_{3},\ldots,p_{6}$, the parabola through $p_{4},p_{5},p_{6}$ (drawn full).
Figure 1: Rote’s example: a 6-point 3rd-order positive set in 3-general
position that does not lie on the graph of any function with nonnegative 3rd
derivative.
It remains to check that $P$ does not lie on the graph of a function $f$ with
$f^{(3)}\geq 0$ everywhere. Assuming for contradiction that there is such an
$f$, we consider the point $q:=(x_{0},f(x_{0}))$, where $x_{0}$ is such that
the full parabola is below the $x$-axis at $x_{0}$. For the $4$-tuple
$\\{p_{1},p_{2},p_{3},q\\}$ to be positive, $q$ has to lie above the $x$-axis,
but the $4$-tuple $\\{q,p_{4},p_{5},p_{6}\\}$ is positive only if $q$ lies
below the parabola through $p_{4},p_{5},p_{6}$—a contradiction.
## 3 Upper bounds on the Ramsey numbers for transitive colorings
In this section we prove Theorem 1.4. As we mentioned in the remark following
that theorem, it suffices to establish the case $k\geq 3$.
Thus, we want to prove that $R^{\rm trans}_{k+1}(n)\leq\mathop{\rm
twr}\nolimits_{k}(C_{k}n)$ for all $n$ and for every $k\geq 3$, with suitable
constants $C_{k}$ depending on $k$. As the base of the induction we use
$R^{\rm trans}_{3}(n)\leq 4^{n}$, which, as was remarked earlier, follows by
imitating the proof of Theorem 1.2.
Thus, let $k\geq 3$ be fixed, let $n$ be given, and let us set $M:=R^{\rm
trans}_{k}(n)$. We will prove that
$R^{\rm trans}_{k+1}(n)\leq N:=2^{M^{k}}.$ (1)
Theorem 1.4 then follows from this recurrence and from the fact that
$2^{\mathop{\rm twr}\nolimits_{k-1}(n)^{k}}\leq\mathop{\rm
twr}\nolimits_{k}(kn)$ for $k\geq 3$, which is easy to check.
To prove (1), we follow an inductive proofs of Ramsey’s theorem going back to
Erdős and Rado [ER52]. Let $\chi\colon{[N]\choose k+1}\to\\{1,2\\}$ be an
arbitrary transitive 2-coloring. We set $A_{k-1}:=\\{1,2,\ldots,k-1\\}$ and
$X_{k-1}:=[N]\setminus A_{k-1}$. For $i=k,k+1,\ldots,M$ we will inductively
construct sets $A_{i},X_{i}\subseteq[N]$ such that
1. (i)
$A_{i}<X_{i}$ (i.e., all elements of $A_{i}$ precede all elements of $X_{i}$);
2. (ii)
$|A_{i}|=i$ and $|X_{i}|\geq|X_{i-1}|/2^{M^{k-1}}$; and
3. (iii)
the color of a $(k+1)$-tuple whose first $k$ elements all belong to $A_{i}$
does not depend on its last element; in other words, for $K\in{A_{i}\choose
k}$ and $x,y\in A_{i}\cup X_{i}$ with $K<\\{x,y\\}$, we have
$\chi(K\cup\\{x\\})=\chi(K\cup\\{y\\})$.
For the inductive step, suppose that $A_{i}$ and $X_{i}$ have already been
constructed. We let $x_{i}$ be the smallest element of $X_{i}$, we set
$A_{i+1}:=A_{i}\cup\\{x_{i}\\}$, and we write
$X^{\prime}_{i}:=X_{i}\setminus\\{x_{i}\\}$.
Let us call two elements $x,y\in X^{\prime}_{i}$ _equivalent_ if we have, for
every $K\in{A_{i-1}\choose k-1}$,
$\chi(K\cup\\{x_{i},x\\})=\chi(K\cup\\{x_{i},y\\})$. There are ${i\choose
k-1}$ possible choices of $K$, and hence there are at most $2^{i\choose
k-1}<2^{M^{k-1}}$ equivalence classes. We choose $X_{i+1}\subseteq
X^{\prime}_{i}$ as the largest equivalence class. Then (i), (iii) obviously
hold for $A_{i+1}$ and $X_{i+1}$, and we have
$|X_{i+1}|\geq(|X_{i}|-1)/(2^{M^{k-1}}-1)\geq|X_{i}|/2^{M^{k-1}}$ (since
$i\leq M$ and thus we have $|X_{i}|\geq
N/(2^{M^{k}-1})^{i-1}=2^{M^{k}-(i-1)M^{k-1}}\geq 2^{M^{k-1}}$). This finishes
the inductive construction of $A_{i}$ and $X_{i}$.
In this way, we construct the sets $A:=A_{M}$ and $X_{M}$ (note that
$|X_{M}|\geq 1$ by (ii)). Let $x$ be the first element of $X_{M}$, and let us
define a 2-coloring $\chi^{*}\colon{A\choose k}\to\\{1,2\\}$ of the $k$-tuples
of $A$ by $\chi^{*}(K):=\chi(K\cup\\{x\\})$.
We claim that, crucially, $\chi^{*}$ is transitive (which is not entirely
obvious). So we consider elements $a_{1}<a_{2}<\cdots<a_{k+1}$ of $A$, and we
suppose that
$\chi^{*}(\\{a_{1},\ldots,a_{k}\\})=\chi^{*}(\\{a_{2},\ldots,a_{k+1}\\})=:c$.
We want to show that
$\chi^{*}(\\{a_{1},\ldots,a_{k+1}\\}\setminus\\{a_{i}\\})=c$ for every
$i=2,3,\ldots,k$. We have
$c=\chi^{*}(\\{a_{1},\ldots,a_{k}\\})=\chi(\\{a_{1},\ldots,a_{k},x\\})=\chi(\\{a_{1},\ldots,a_{k+1}\\})$
(by definition and by the independence of $\chi$ of the last element), and
$c=\chi^{*}(\\{a_{2},\ldots,a_{k+1}\\})=\chi(\\{a_{2},\ldots,a_{k+1},x\\})$.
Next we use the transitivity of $\chi$ on the $(k+2)$-tuple
$(a_{1},\ldots,a_{k+1},x)$, obtaining
$\chi(\\{a_{1},\ldots,a_{k+1},x\\}\setminus\\{a_{i}\\})=c=\chi^{*}(\\{a_{1},\ldots,a_{k+1}\\}\setminus\\{a_{i}\\})$
as needed.
Now we can apply the inductive hypothesis to $A$, which yields an $n$-element
subset of $A$ homogeneous w.r.t. $\chi^{*}$, and this subset is homogeneous
w.r.t. $\chi$ as well, finishing the proof of Theorem 1.4. $\Box$
## 4 A lower bound for $\operatorname{ES}_{3}$
Here we prove Theorem 1.5, a lower bound for $\operatorname{ES}_{3}(2n+1)$. We
proceed by induction on $n$; the goal is to construct a set $P_{n}$ of
$N:=2^{2^{n-1}}$ points with no $(2n+1)$-point third-order monotone subset.
The induction starts for $n=2$ with an arbitrary $P_{2}$ of size
$2^{2^{1}}=4$.
In the inductive step, given $P_{n}$, we will construct $P_{n+1}$ so that
$|P_{n+1}|=|P_{n}|^{2}$; then the bound on the size of $P_{n}$ clearly holds.
We may assume that $P=P_{n}$ is in $3$-general position (this can always be
achieved by a small perturbation). By an affine transformation we also make
sure that $P\subset[1,2]\times[0,1]$; or actually,
$P\subset[1,1.9]\times[0,1]$ so that there is some room for perturbation.
Moreover, there is a small $\delta>0$ such that if $P^{\prime}$ is obtained
from $P$ by moving each point arbitrarily by at most $\delta$, then
$P^{\prime}$ is still in $3$-general position, the order of the points of
$P^{\prime}$ along the $x$-axis is the same as that for $P$, and the sign of
every $4$-tuple in $P^{\prime}$ is the same as the sign of the corresponding
$4$-tuple in $P$.
The construction. The construction of $P_{n+1}$ from $P=P_{n}$ as above
proceeds in the following steps.
1. 1.
We choose a sufficiently large number $A=A(P)$ (the requirements on it will be
specified later), and we set $\varepsilon:=1/A^{2}$.
2. 2.
For every point $p\in P$, let $Q_{p}$ be the image of $P$ under the affine map
that sends the square $[1,2]\times[0,1]$ to the axis-parallel rectangle of
width $\varepsilon$, height $\varepsilon^{2}$, and with the lower left corner
at $p$; see Fig. 2.
3. 3.
Let $\psi_{p}(x)=Ax^{2}+C_{p}$ be a quadratic function, where $A$ is as above
and $C_{p}$ is chosen so that $\psi_{p}(x(p))=0$ (where $x(p)$ is the
$x$-coordinate of $p$). Let $\breve{Q}_{p}$ be the set obtained by “adding
$\psi_{p}$ to $Q_{p}$”, i.e., by shifting each point $(x,y)\in Q_{p}$
vertically upwards by $\psi_{p}(x)$. We set $P_{n+1}:=\bigcup_{p\in
P}\breve{Q}_{p}$. We call the $\breve{Q}_{p}$ the _clusters_ of $P_{n+1}$.
Figure 2: A schematic illustration of the construction of $P_{n+1}$.
First we check that each cluster $\breve{Q}_{p}$ lies close to $p$.
###### Lemma 4.1
Each $\breve{Q}_{p}$ is contained in an $O(\sqrt{\varepsilon}\,)$-neighborhood
of $p$.
Proof. Writing $p=(x_{0},y_{0})$, the set $Q_{p}$ obviously lies in the
$2\varepsilon$-neighborhood of $p$, and the maximum amount by which a point of
$Q_{p}$ was translated upwards is at most
$\psi_{p}(x_{0}+\varepsilon)=A\left((x_{0}+\varepsilon)^{2}-x_{0}^{2}\right)=A(2x_{0}\varepsilon+\varepsilon^{2})=O(\sqrt{\varepsilon}\,).$
$\Box$
Here is a key property of the construction.
###### Lemma 4.2 (Slope lemma)
Let $\lambda$ be a parabola passing through three points of $P_{n+1}$ that
belong to three different clusters, or a line passing through two points of
different clusters. Let $\mu$ be a parabola passing through three points of a
single cluster $\breve{Q}_{p}$, or a line passing through two such points.
Then the maximum slope (first derivative) of $\lambda$ on the interval $[1,2]$
is smaller than the minimum slope of $\mu$ on $[1,2]$, provided that $A$ was
chosen sufficiently large.
Proof. Clearly, the maximum slope of any such $\lambda$ can be bounded from
above by some finite number depending only on $P$ but not on $A$. Thus, it
suffices to show that, with $A$ large, for every $\mu$ as in the lemma, the
minimum slope is bounded from below by $A$.
First let us assume that $\mu$ is a parabola passing through three points of
$\breve{Q}_{p}$, where $p=(x_{0},y_{0})$, let $\tilde{\mu}$ be the parabola
passing through the corresponding three points of $P$, and let the equation of
$\tilde{\mu}$ be $y=ax^{2}+bx+c$.
By the construction of $\breve{Q}_{p}$, the affine map transforming $P$ to
$Q_{p}$ sends a point with coordinates $(x,y)$ to the point
$(\varepsilon(x-1)+x_{0},\varepsilon^{2}y+y_{0})$. Calculation shows that the
image of $\tilde{\mu}$ under this affine map has the equation
$y=ax^{2}+(2a\varepsilon+b\varepsilon-2ax_{0})x+c^{\prime}$, where the value
of the absolute term $c^{\prime}$ need not be calculated since it doesn’t
matter. Hence the minimum slope of this curve on $[1,2]$ is bounded from below
by $-(8|a|+4|a|\varepsilon+2|b|\varepsilon+8|a|)$. Finally, $\mu$ is obtained
by adding $\psi_{p}(x)=Ax^{2}+C_{p}$ to this curve, and the minimum slope of
$\psi_{p}$ on $[1,2]$ is at least $2A$.
Next, let $\mu$ be a line passing through two points $q,r\in\breve{Q}_{p}$.
Let us choose another point $s\in\breve{Q}_{p}$ and consider the parabola
$\mu^{\prime}$ through $q,r,s$. By the Mean Value Theorem, the slope of $\mu$
equals the slope of $\mu^{\prime}$ at some point between $q$ and $r$, and the
latter is at least $A$ by the above. The lemma is proved. $\Box$
Let $K=\\{p_{1},p_{2},p_{3},p_{4}\\}\subseteq P_{n+1}$ be a 4-tuple,
$p_{i}=(x_{i},y_{i})$, $x_{1}<\cdots<x_{4}$. We assign a _type_ to $K$, which
is an ordered partition of $4$ given by the distribution of $K$ among the
clusters; for example, $K$ has type $1+1+2$ if the first point $p_{1}$ lies in
some $\breve{Q}_{p}$, $p_{2}$ lies in $\breve{Q}_{p^{\prime}}$ for
$p^{\prime}\neq p$, and $p_{3},p_{4}\in\breve{Q}_{p^{\prime\prime}}$,
$p^{\prime\prime}\neq p,p^{\prime}$.
The next lemma shows that the sign $K$ is determined by its type. We provide a
complete classification, although we will not use all of the types in the
subsequent proof.
###### Lemma 4.3
Let $K=\\{p_{1},p_{2},p_{3},p_{4}\\}\subseteq P_{n+1}$ be a $4$-tuple. If $K$
is of type $1+1+1+1$ or $4$, then the sign of $K$ is the same as that of the
corresponding $4$-tuple in $P$. Otherwise, the sign of $K$ is determined by
its type as follows:
* •
for types $3+1$ and $1+3$ it is $-1$;
* •
for types $1+1+2$ and $2+1+1$ it is $+1$;
* •
for type $1+2+1$ it is $-1$; and
* •
for type $2+2$ it is $+1$.
Proof. Since the transformation that converts $P$ into $\breve{Q}_{p}$
preserves the types of $4$-tuples, the statement for type 4 is clear. The
statement for type $1+1+1+1$ follows since, by Lemma 4.1, $K$ is obtained by a
sufficiently small perturbation of the corresponding $4$-tuple in $P$ (this
gives one of the lower bounds on $A$, since we need the bound in Lemma 4.1 to
be smaller than the $\delta$ considered at the beginning of our description of
the construction).
The statements for the remaining types are obtained by simple application of
the slope lemma (Lemma 4.2) together with Lemma 2.4. Namely, for type $3+1$,
we get that the parabola through $p_{1},p_{2},p_{3}$ lies above $p_{4}$ (by
comparing its slope to the slope of the line $p_{3}p_{4}$); see Fig. 3. For
type $1+3$ we similarly get that $p_{1}$ lies above the parabola through
$p_{2},p_{3},p_{4}$, and so the sign is $-1$ in both of these cases.
Figure 3: Determining the signs of $4$-tuples by type.
For type $1+1+2$, the segment $p_{3}p_{4}$ is steeper than the parabola
through $p_{1}p_{2}p_{3}$, and so the sign is $+1$. Similarly for type $2+1+1$
we get that $p_{1}$ lies below the parabola through $p_{2},p_{3},p_{4}$, which
again gives sign $+1$. For type $1+2+1$, $p_{3}$ lies above the parabola
through $p_{1},p_{2},p_{4}$, giving sign $-1$. Finally, for type $2+2$, the
segment $p_{1}p_{2}$ is steeper than $p_{2}p_{3}$, thus the parabola through
$p_{1},p_{2},p_{3}$ is concave, and hence its slope at $p_{3}$ and after it is
no larger than the slope of the segment $p_{2}p_{3}$. Thus, $p_{4}$ lies above
this parabola and the sign is $+1$ as claimed. $\Box$
Finishing the proof of Theorem 1.5. It remains to show that $P_{n+1}$
contains no $(2n+3)$-point third-order monotone subset.
For contradiction, suppose that $M\subseteq P_{n+1}$ is such a $(2n+3)$-point
subset. Let $2n+3=n_{1}+n_{2}+\cdots+n_{s}$ be the type of $M$ (i.e., $M$ has
$n_{i}\geq 1$ points in the $i$th leftmost cluster it intersects). By the
inductive assumption we have $s\leq 2n$ and $n_{i}\leq 2n$ for all $i$.
Let $n_{a}=\max_{i}n_{i}$ and $n_{b}=\max_{i\neq a}n_{i}$ be the two largest
among the $n_{i}$. For convenience, let us assume $a<b$; the case $a>b$ is
handled symmetrically. We distinguish three cases.
First, if $n_{a}\geq 3$ and $n_{b}\geq 2$, then we can select 4-tuples of
types $3+1$ and $2+2$ from the corresponding two clusters, which have
different signs, and so $M$ is not homogeneous.
Second, if $n_{a}\geq 3$ and $n_{b}=1$, then we have at least three $n_{i}$
equal to 1 (since $n_{a}\leq 2n$), and at least two of them lie on the same
side of the cluster corresponding to $n_{a}$, say to the right of it. Then we
can select 4-tuples of types $3+1$ and $2+1+1$, again of opposite signs.
Third, if $n_{a}=2$, then there are at least two other clusters of size 2.
From these three 2-element clusters, we can select 4-tuples of types 2+2 and
1+2+1, again of opposite signs.
This exhausts all possibilities ($n_{a}=1$ cannot happen, because $s\leq 2n$),
and Theorem 1.5 is proved. $\Box$
## 5 Order types and one-sided sets of hyperplanes
First we substantiate the two claims made above Proposition 1.6, concerning
$\mathop{\rm OT}\nolimits_{1}$ and $\mathop{\rm OT}\nolimits_{2}$. For $d=1$,
an order-type homogeneous sequence in ${\mathbb{R}}^{1}$ is just a monotone
sequence of real numbers, so $\mathop{\rm OT}\nolimits_{1}(n)=(n-1)^{2}+1$ by
Theorem 1.1.
In a similar spirit, it is easy to check that a planar order-type homogeneous
sequence corresponds to the vertices of a convex $n$-gon, enumerated in a
clockwise or counterclockwise order. Thus, $\mathop{\rm
OT}\nolimits_{2}(n)\geq\operatorname{ES}_{2}(\lceil n/2\rceil)=2^{\Omega(n)}$.
On the other hand, given any $N$-point sequence, we can first select a
subsequence of $\lceil\sqrt{N}\,\rceil$ points with increasing or decreasing
$x$-coordinates, and then we select a convex or concave configuration from it.
Thus, by Theorem 1.2, we have $\mathop{\rm OT}\nolimits_{2}(n)=2^{O(n)}$.
Proof of Proposition 1.6. For a point $p=(x,y)\in{\mathbb{R}}^{2}$, we define
the point $\tilde{p}:=(x,x^{2},\ldots,x^{d-1},y)\in{\mathbb{R}}^{d}$.
To prove that $\operatorname{ES}_{d}(n)\leq\mathop{\rm OT}\nolimits_{d}(n)$,
we consider a set $P=\\{p_{1},\ldots,p_{N}\\}\subset{\mathbb{R}}^{2}$ in
$d$-general position, $p_{i}=(x_{i},y_{i})$, where
$N=\operatorname{ES}_{d}(n)-1$ and $x_{1}<\cdots<x_{N}$, with no $d$th-order
monotone subset of $n$ points. It suffices to prove that the sequence
$\tilde{P}:=(\tilde{p}_{1},\tilde{p}_{2},\ldots,\tilde{p}_{N})$ has no
$n$-point order-type homogeneous subsequence. This follows from the next
lemma.
###### Lemma 5.1
For every $(d+1)$-tuple $(p_{1},\ldots,p_{d+1})$ of points in
${\mathbb{R}}^{2}$, $x_{1}<\cdots<x_{d+1}$, we have $\mathop{\rm
sgn}\nolimits(\\{p_{1},\ldots,p_{d+1}\\})=\mathop{\rm sgn}\nolimits\det
M(\tilde{p}_{1},\tilde{p}_{2},\ldots,\tilde{p}_{d+1})$, where
$M(q_{1},\ldots,q_{d+1})$ is the matrix from the definition of order type
above Proposition 1.6.
Proof. By Lemma 2.2 and Corollary 2.3, the sign of
$\\{p_{1},\ldots,p_{d+1}\\}$ equals the sign of the coefficient $a_{d}$ of the
unique polynomial $f(x)=\sum_{j=0}^{d}a_{j}x^{j}$ of degree at most $d$ whose
graph passes through the points $p_{1},\ldots,p_{d+1}$.
The vector $a=(a_{0},\ldots,a_{d})$ can be expressed as the solution of the
linear system $Va=y$, where $y=(y_{1},\ldots,y_{d+1})$ and $V$ is the
_Vandermonde matrix_ with $v_{ij}=x_{i}^{j-1}$, $i,j=1,2,\ldots,d+1$. By
Cramer’s rule, we obtain
$a_{d}=\frac{\det W}{\det V},$
where $W$ stands for the matrix $V$ with the last column replaced with the
vector $y$. As is well known, $\det V=\prod_{1\leq i<j\leq d+1}(x_{j}-x_{i})$,
and since $x_{1}<\cdots<x_{d+1}$, we have $\det V>0$. Thus, $\mathop{\rm
sgn}\nolimits a_{d}=\mathop{\rm sgn}\nolimits\det W$. Finally, we have
$W=\left(\begin{array}[]{cccccc}1&x_{1}&x_{1}^{2}&\ldots&x_{1}^{d-1}&y_{1}\\\
\vdots&\vdots&\vdots&\vdots&\vdots&\vdots\\\
1&x_{d+1}&x_{d+1}^{2}&\ldots&x_{d+1}^{d-1}&y_{d+1}\end{array}\right)=M(\tilde{p}_{1},\tilde{p}_{2},\ldots,\tilde{p}_{d+1})^{T}.$
The lemma follows, and Proposition 1.6 is proved. $\Box$
Proof of Proposition 1.7. The proof is very similar to the previous one. This
time we start with a set $P=\\{p_{1},\ldots,p_{N}\\}\subset{\mathbb{R}}^{2}$
in $(d-1)$-general position, $p_{i}=(x_{i},y_{i})$, where
$N=\operatorname{ES}_{d-1}(n)-1$ and $x_{1}<\cdots<x_{N}$, with no $(d-1)$th-
order monotone subset of $n$ points. We define a collection
$H=\\{h_{1},\ldots,h_{N}\\}$ of $N$ hyperplanes in ${\mathbb{R}}^{d}$, where
$h_{i}$ is given by
$h_{i}=\biggl{\\{}(\xi_{1},\ldots,\xi_{d})\in{\mathbb{R}}^{d}:\sum_{j=1}^{d}x_{i}^{j-1}\xi_{j}=y_{i}\biggr{\\}}.$
The intersection point $\xi=(\xi_{1},\ldots,\xi_{d})$ of, say,
$h_{1},\ldots,h_{d}$ is the solution of the linear system $V\xi=y$, where $V$
is the $d\times d$ Vandermonde matrix this time, $v_{ij}=x_{i}^{j-1}$.
Cramer’s rule then gives that the $d$th coordinate $\xi_{d}$, whose sign we
are interested in, equals $(\det W)/(\det V)$, where $W$ is obtained from $V$
by replacing the last column with $y$.
As we saw in the proof of Proposition 1.6, $(\det W)/(\det V)$ also expresses
the leading coefficient in the polynomial of degree $d-1$ passing through
$p_{1},\ldots,p_{d}$, and thus its sign equals $\mathop{\rm
sgn}\nolimits\mbox{$\Delta\\!\\!\\!\\!\\!\;\raisebox{2.15277pt}{\mbox{\boldmath$\scriptscriptstyle\mid$}}\,\,$}_{d-1}(p_{1},\ldots,p_{d})$.
It follows that one-sided subsets of $H$ precisely correspond to $(d-1)$st-
order monotone subsets in $P$, and the proposition is proved. $\Box$
## Acknowledgment
We would like to thank János Pach for kindly discussing some of the results of
Fox et al. [FPSS11] with us. We also thank Günter Rote for informing us about
about his refutation of our conjecture and for permission to present it in
this paper.
## References
* [CFS11] D. Conlon, J. Fox, and B. Sudakov. An improved bound for the stepping-up lemma. Discrete Applied Mathematics, 2011. In press.
* [DL11] V. Dujmović and S. Langerman. A center transversal theorem for hyperplanes and applications to graph drawing. In Proc. 27th ACM Symposium on Computational Geometry, pages 117–124, 2011. Full version arXiv:1012.0548.
* [ER52] P. Erdős and R. Rado. Combinatorial theorems on classifications of subsets of a given set. Proc. London Math. Soc., 3:417–439, 1952.
* [ES35] P. Erdős and G. Szekeres. A combinatorial problem in geometry. Compositio Math., 2:463–470, 1935.
* [FPSS11] J. Fox, J. Pach, B. Sudakov, and A. Suk. Erdős–Szekeres-type theorem for monotone paths and convex bodies. Arxiv preprint 1105.2097v1, 2011. _Proc. London Math. Soc._ , in press.
* [GP93] J. E. Goodman and R. Pollack. Allowable sequences and order types in discrete and computational geometry. In J. Pach, editor, New Trends in Discrete and Computational Geometry, volume 10 of Algorithms and Combinatorics, pages 103–134. Springer, Berlin etc., 1993.
* [GRS90] R. L. Graham, B. L. Rothschild, and J. Spencer. Ramsey Theory. J. Wiley & Sons, New York, 1990.
* [Mat02] J. Matoušek. Lectures on Discrete Geometry. Springer, New York, 2002.
* [MS00] W. Morris and V. Soltan. The Erdős–Szekeres problem on points in convex position—a survey. Bull. Amer. Math. Soc., New Ser., 37(4):437–458, 2000.
* [MW92] J. Matoušek and Emo Welzl. Good splitters for counting points in triangles. J. Algorithms, 13:307–319, 1992.
* [Phi03] George M. Phillips. Interpolation and approximation by polynomials. Springer, Berlin etc., 2003.
* [Rot12] G. Rote. Private communication, February 2012.
* [Ste95] M. J. Steele. Variations on the monotone subsequence theme of Erdős and Szekeres. In D. Aldous et al., editors, Discrete Probability and Algorithms, IMA Volumes in Mathematics and its Applications 72, pages 111–131. Springer, Berlin etc., 1995.
* [Val04] P. Valtr. Open caps and cups in planar point sets. Discr. Comput. Geom., 37:365–567, 2004.
|
arxiv-papers
| 2011-11-16T14:55:15 |
2024-09-04T02:49:24.387802
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Marek Elias, Jiri Matousek",
"submitter": "Marek Elias",
"url": "https://arxiv.org/abs/1111.3824"
}
|
1111.3828
|
Curves on Oeljeklaus-Toma Manifolds
Sima Verbitsky
Abstract
Oeljeklaus-Toma manifolds are complex non-Kähler manifolds constructed by
Oeljeklaus and Toma from certain number fields, and generalizing the Inoue
surfaces $S_{m}$. We prove that Oeljeklaus-Toma manifolds contain no compact
complex curves.
###### Contents
1. 1 Introduction
1. 1.1 Oeljeklaus-Toma manifolds
2. 2 Curves on the Oeljeklaus-Toma manifolds
1. 2.1 The exact semipositive (1,1)-form on the Oeljeklaus-Toma manifold
2. 2.2 The $(1,1)$-form $\omega$ and curves on the Oeljeklaus-Toma manifold
3. 3 Closing remarks
## 1 Introduction
Oeljeklaus-Toma manifolds (defined in [O–T]) are compact complex manifolds
that are a generalization of Inoue surfaces (defined in [I]). Let us describe
them in detail.
### 1.1 Oeljeklaus-Toma manifolds
Let $K$ be a number field (i.e. a finite extension of $\mathbb{Q}$), $s>0$ be
the number of its real embeddings and $2t>0$ be the number of its complex
embeddings. One can easily prove that for each $s$ and $t$ there exists a
field $K$ which has these numbers of real and complex embeddings (see e.g.
[O–T]).
Definition 1.1: The ring of algebraic integers $O_{K}$ is a subring of $K$
that consists of all roots of polynomials with integer coefficients which lie
in $K$. Unit group $O^{*}_{K}$ is the multiplicative subgroup of invertible
elements of $O_{K}$.
Let $m$ be $s+t$. Let $\sigma_{1},\ldots,\sigma_{s}$ be real embeddings of the
field $K$, $\sigma_{s+1},\ldots,\sigma_{s+2t}$ be complex embeddings such that
$\sigma_{s+i}$ and $\sigma_{s+t+i}$ are complex conjugate for each $i$ from
$1$ to $t$. Now we can define a map $l:O_{K}^{*}\rightarrow\mathbb{R}^{m}$
where
$l(u)=(\ln|\sigma_{1}(u)|,\ldots,\ln|\sigma_{s}(u)|,2\ln|\sigma_{s+1}(u)|,\ldots,2\ln|\sigma_{m}(u)|)$.
Denote $O_{K}^{*,+}=\\{a\in O_{K}^{*}:\sigma_{i}(a)>0,i=1,\ldots,s\\}$. Let us
consider following definitions:
Definition 1.2: A lattice $\Lambda$ in $\mathbb{R}^{n}$ is a discrete additive
subgroup such that $\Lambda\otimes\mathbb{R}=\mathbb{R}^{n}$.
Definition 1.3: [O–T] The group $U\subset O_{K}^{*,+}$ of rank $s$ is called
admissible for the field $K$ if the projection of $l(U)$ to the first $s$
components is a lattice in $\mathbb{R}^{s}$.
Consider a linear space $L=\\{x\in\mathbb{R}^{m}\mid\sum_{i=1}^{m}x_{i}=0\\}$.
The projection of $L\subset R^{m}$ to the first $s$ coordinates is surjective,
because $s<m$. Using the Dirichlet unit theorem (see e.g. [Mil09]) one can
prove that $l(O_{K}^{*,+})$ is a full lattice in $L$. Therefore there exists a
group $U$ that is admissible.
Let $\mathbb{H}=\\{z\in\mathbb{C}\mid\operatorname{im}z>0\\}$. Let $U\subset
O_{K}^{*,+}$ be a group which is admissible for $K$. The group $U$ acts on
$O_{K}$ multiplicatively. This defines a structure of semidirect product
$U^{\prime}:=U\ltimes O_{K}$. Define the action of $U^{\prime}$ on
$\mathbb{H}^{s}\times\mathbb{C}^{t}$ as follows. The element $u\in U$ acts on
$\mathbb{H}^{s}\times\mathbb{C}^{t}$ mapping $(z_{1},\ldots,z_{m})$ to
$(\sigma_{1}(u)z_{1},\ldots,\sigma_{m}(u)z_{m})$. Since $U$ lies in
$O_{K}^{*,+}$, the action $U$ on the first $s$ coordinates preserves
$\mathbb{H}$.
The additive group $O_{K}$ acts on $\mathbb{H}^{s}\times\mathbb{C}^{t}$ by
parallel translations: $a\in O_{K}$ is mapping $(z_{1},\ldots,z_{m})$ to
$(\sigma_{1}(a)+z_{1},\ldots,\sigma_{m}(a)+z_{m})$. Since the first $s$
embeddings are real, this action preserves $\mathbb{H}$ in the first $s$
coordinates.
One can see that $(u,a)\in U\ltimes O_{K}$ maps $(z_{1},\ldots,z_{m})$ to
$(\sigma_{1}(u)z_{1}+\sigma_{1}(a),\ldots,\sigma_{m}(u)z_{m}+\sigma_{m}(a))$.
One can easily show that this action is compatible with the group operation in
the semidirect product.
Definition 1.4: An Oeljeklaus-Toma manifold is the quotient of
$\mathbb{H}^{s}\times\mathbb{C}^{t}$ by the action of the group $U\ltimes
O_{K}$, which was defined above.
This quotient exists because $U\ltimes O_{K}$ acts properly discontiniously on
$\mathbb{H}^{s}\times\mathbb{C}^{t}$. Additionally
$\mathbb{H}^{s}\times\mathbb{C}^{t}/U\ltimes O_{K}$ is a compact complex
manifold. To prove it, let $U$ be admissible for $K$. The quotient
$\mathbb{H}^{s}\times\mathbb{C}^{t}/O_{K}$ is obviously diffeomorphic to the
trivial toric bundle $(\mathbb{R}_{>0})^{s}\times(S^{1})^{n}$. The group $U$
acts properly discontinuously on the base $(\mathbb{R}_{>0})^{s}$. Therefore
it acts properly discontinuously on
$\mathbb{H}^{s}\times\mathbb{C}^{t}/O_{K}$. Also, the groups $U$ and $O_{K}$
act holomorphically on $\mathbb{H}^{s}\times\mathbb{C}^{t}$. Therefore the
quotient has a holomorphic structure.
## 2 Curves on the Oeljeklaus-Toma manifolds
In this section we shall prove that there are no complex curves on the
Oeljeklaus-Toma manifolds, just as on Inoue surfaces of type $S_{M}$ (see
[I]).
### 2.1 The exact semipositive (1,1)-form on the Oeljeklaus-Toma manifold
The $(1,1)$-form we will be using was previously introduced in the paper
Subvarieties in Oeljeklaus-Toma manifolds by Ornea and Verbitsky. Authors use
this form to prove that Oeljeklaus-Toma manifolds with $t=1$ (that means that
the corresponding number field has only two complex embeddings) do not contain
any submanifolds. Our result works for all Oeljeklaus-Toma manifolds, but we
consider only curves instead of submanifolds of any dimension. Later on we
will explain how our method differs from the one in [O–V] and how this implies
the difference between our results.
Define the notion of $(1,1)$-form. Let $M$ be a smooth complex manifold,
$z_{1},\ldots,z_{n}$ — local complex coordinates in the open neighborhood of
the point $y\in M$.
Definition 2.1: A $(1,1)$-form on a complex manifold $M$ is a 2-form $\omega$,
such that $\omega(Iu,v)=-\omega(u,Iv)=\sqrt{-1}\omega(u,v)$ for each $u,v\in
T_{y}M$, where $I$ is the almost complex structure on $M$.
Definition 2.2: A $(1,1)$-form $\omega$ on a complex manifold $M$ is
semipositive if $\omega(u,Iu)\geqslant 0$ for each tangent vector $u\in
T_{y}M$.
As in [O–V], we consider a certain semipositive $(1,1)$-form on the
Oeljeklaus-Toma manifold $M=\mathbb{H}^{s}\times\mathbb{C}^{t}/(U\ltimes
O_{K})$. We introduce a $(1,1)$-form $\widetilde{\omega}$ on
$\widetilde{M}=\mathbb{H}^{s}\times\mathbb{C}^{t}$ which is preserved by the
action of the group $\Gamma=(U\ltimes O_{K})$ and since then it would be a
$(1,1)$-form on $M$.
Let $(z_{1},\ldots,z_{m})$ be complex coordinates on $\widetilde{M}$. Define
$\varphi(z)=\Pi_{i=1}^{s}\operatorname{im}(z_{i})^{-1}$. Since the first $s$
components of $\widetilde{M}$ correspond to upper half-planes
$\mathbb{H}\subset\mathbb{C}$, this function is positive on $\widetilde{M}$.
Let us now consider the form
$\widetilde{\omega}=\sqrt{-1}\partial\bar{\partial}\log\varphi$. Using
standard coordinates on $\widetilde{M}$ one can write this form as
$\widetilde{\omega}=\sqrt{-1}\sum_{i=1}^{s}\frac{dz_{i}\wedge
d\bar{z}_{i}}{4(\operatorname{im}z_{i})^{2}}$. Therefore $\widetilde{\omega}$
is a semipositive $(1,1)$-form on $\widetilde{M}$.
Let us show that this form is $\Gamma$-invariant.
The group $\Gamma$ is a semidirect product of the additive group $O_{K}$ and
the multiplicative group $U$.The additive group acts on the first $s$
components of $\widetilde{M}$ (which correspond to upper half-planes
$\mathbb{H}\subset\mathbb{C}$) by translations along the real line. Therefore
it does not change $\operatorname{im}z_{i}$ for $i=1\ldots s$. Hence the
function $\log\varphi$ is preserved by the action of the additive component.
The multiplicative component acts on the first $s$ coordinates of
$\widetilde{M}$ by multiplying them by a real number (since the first $s$
embeddings of the number field $K$ are real). Then every
$\operatorname{im}z_{i}$ is multiplied by a real number and so there is a real
number added to $\log(\operatorname{im}z_{i})$. Since
$\log\varphi(z)=-\sum_{i=1}^{s}\log(\operatorname{im}z_{i})$, there is a real
number added to $\log\varphi$. The operator $\bar{\partial}$ is zero on the
constants, so $\widetilde{\omega}=\sqrt{-1}\partial\bar{\partial}\log\varphi$
is preserved by action of the group $\Gamma$.
Since the $(1,1)$-form $\widetilde{\omega}$ is $\Gamma$-invariant it is the
pullback of $(1,1)$-form $\omega$ on the Oeljeklaus-Toma manifold
$M=\widetilde{M}/\Gamma$.
Let us now show that the form $\widetilde{\omega}$ is exact on
$\widetilde{M}$. For that we define the operator $d^{c}$.
Definition 2.3: Define the twisted differential $d^{c}=I^{-1}dI$ where $d$ is
a De Rham differential and $I$ is the almost complex structure.
Since $dd^{c}=2\sqrt{-1}\partial\bar{\partial}$ (see [G–H]), one can see that
$\widetilde{\omega}=\sqrt{-1}\partial\bar{\partial}\log\varphi=\frac{1}{2}dd^{c}\log\varphi$
and so $\widetilde{\omega}$ is exact as a form on $\widetilde{M}$. Also since
the operator $d^{c}$ vanishes on constants the form $d^{c}\log\varphi$ is
$\Gamma$-invariant, so $\omega$ is exact on $M$.
### 2.2 The $(1,1)$-form $\omega$ and curves on the Oeljeklaus-Toma manifold
Since the form $\omega$ on the manifold $M$ is semipositive, its integral on
any complex curve $C\subset M$ is nonnegative. The form $\omega$ is exact.
Hence Stokes’ theorem implies that its integral on any complex curve vanishes.
Therefore if $C\subset M$ is a closed complex curve, $\omega$ vanishes on it.
To find out on which curves $\omega$ vanishes, let us define the zero
foliation of the form $\omega$.
Definition 2.4: An involutive distribution (or foliation) on $M$ is a
subbundle $B\subset TM$ of the tangent bundle that is closed under the Lie
bracket: $[B,B]\subset B$.
Definition 2.5: A leaf of a foliation $B$ is a connected submanifold of $M$
such that its dimension is equal to $\dim B$ and that is tangent to $B$ at
every point.
Theorem 2.6: (Frobenius) Let $B\subset TM$ be an involutive distribution. Then
for each point of the manifold $M$, there is exactly one leaf of this
distribution that contains this point (see e.g. [Boo] Section IV. 8. Frobenius
Theorem).
Theorem 2.7: Let $N\subset M$ be a connected submanifold such that its tangent
space at every point lies in a foliation $F\subset TM$. Then $N$ lies in a
leaf of the foliation $F$ (see e.g. [Boo] Section IV. 8. Theorem 8.5).
Definition 2.8: The zero foliation of a semipositive $(1,1)$-form $\omega$ on
$M$ is the subundle of $TM$ that consists of tangent vectors $u\in T_{y}M$
such that $\omega(u,Iu)=0$, where $I$ is the almost complex structure on $M$.
Consider the zero foliation of $\widetilde{\omega}$ on $\widetilde{M}$.
The form $\widetilde{\omega}$ is strictly positive on each vector
$v=(z_{1},\ldots,z_{m})$ such that at least one of $z_{i}$ for $i=1,\ldots,s$
is nonzero. Such a vector cannot be tangent to a leaf of the zero foliation.
Therefore on each leaf of the zero foliation of the form $\widetilde{\omega}$
the first $s$ coordinates are constant.
Hence a leaf of the zero foliation of $\widetilde{\omega}$ on $\widetilde{M}$
is isomorphic to $\mathbb{C}^{t}$.
Let us now consider the zero foliation of $\omega$ on $M$.
We show that the non-trivial image of the action of $\Gamma$ on any leaf $L$
of the zero foliation of the form $\widetilde{\omega}$ does not intersect with
$L$.
One can see that $L$ is $(z_{1},\ldots,z_{s})\times\mathbb{C}^{t}$ for some
fixed $(z_{1},\ldots,z_{s})$. Therefore, for any $\gamma\in\Gamma$ such that
$L\cap\gamma(L)\neq\emptyset$, the first $s$ coordinates of the points in $L$
coincide with the first $s$ coordinates of the points in $\gamma(L)$. Then for
such $\gamma$ have the following system of equations:
$\sigma_{i}(u)z_{i}+\sigma_{i}(a)=z_{i},\quad i=1\ldots s,$
where $\gamma=(u,a)$.
These equations imply that $z_{i}=\frac{\sigma_{i}(a)}{1-\sigma_{i}(u)}$.
Therefore $z_{i}$ are real but $\mathbb{H}$ does not have real elements.
We showed that $L\cap\gamma(L)=\emptyset$ for every $\gamma\neq 1$ in
$\Gamma$.
Since $\omega$ vanishes on each compact curve $C\subset M$, each curve is
contained in some leaf of the zero foliation of $\omega$. Since
$\widetilde{\omega}$ is $\Gamma$-invariant, each leaf of the zero foliation of
$\omega$ on $M$ is isomorphic to a component of the leaf of the zero foliation
of $\widetilde{\omega}$ on $\widetilde{M}$. Therefore, it is isomorphic to
$\mathbb{C}^{t}$. And $\mathbb{C}^{t}$ does not contain any compact complex
submanifolds.
We proved the following theorem:
Theorem 2.9: There are no compact complex curves on the Oeljeklaus-Toma
manifolds.
## 3 Closing remarks
Let us now briefly explain the connection between our work and [O–V]. As in
[O–V] we use the zero foliation of a certain $(1,1)$-form. The leaves of this
zero foliation are $t$-dimensional complex manifolds. In [O–V] authors
consider $t=1$, and submanifolds of any dimension. We consider any $t\geq 1$,
and submanifolds of dimension 1. In both works one uses the semipositivity of
the $(1,1)$-form to prove that a submanifold and a leaf of the zero foliation,
which contains a point in this submanifold, could not intersect transversely.
In both cases, one of the manifolds is of dimension one. Non-transversality
implies that it is contained in the second one. Authors of [O–V] use the fact
that the leaves of the zero foliation are Zariski dense, and therefore could
not be contained in any submanifold. In our work we prove that each leaf is
isomorphic to $\mathbb{C}^{t}$ and therefore could not contain any
submanifolds.
The arguments used in our work and [O–V] could not be applied to the case of
higher dimensions.
It seems obvious that there should be Oeljeklaus-Toma manifols which contain
submanifolds — other Oeljeklaus-Toma manifolds, corresponding to smaller
number fields, but we do not have a formal proof.
In [O–T] it was proved that Oeljeklaus-Toma manifolds do not admit non-trivial
meromorphic functions. Therefore all the divisors on these manifolds are fixed
in their linear systems (which are 0-dimensional). It is conjectured that
Oeljeklaus-Toma manifolds admit no divisors.
## References
* [Boo] Boothby W.M. An Introduction to Differentiable Manifolds and Riemannian Geometry. Academic Press, San Diego, California, 2003.
* [G–H] Griffiths Ph., Harris J. Principles of Algebraic Geometry. Wiley-Interscience, 1994.
* [I] Inoue M. On surfaces of Class VII0, Invent. Math. 24 (1974), 269-310.
* [Mil08] Milne J.S. Fields and Galois Theory, September 2008.
This paper can be found on http://www.jmilne.org/math/CourseNotes/ft.html,
version 4.21
* [Mil09] Milne J.S. Algebraic Number Theory, April 2009.
This paper can be found on http://www.jmilne.org/math/CourseNotes/ant.html,
version 3.02
* [O–T] Oeljeklaus K., Toma M. Non-Kähler compact complex manifolds associated to number fields. Ann. Inst. Fourier 55 (2005), 1291-1300.
* [O–V] Ornea L., Verbitsky M. Subvarieties in Oeljeklaus-Toma manifolds.
* [P–V] Parton M., Vuletescu V. Examples of non-trivial rank in locally conformal Kähler geometry. Math. Z. (2010), DOI 10.1007/s00209-010-0791-5, arXiv:1001.4891.
* [R] Raghunathan M.S. Discrete subgroups of Lie groups. Springer 1972.
* [V] Voisin C. Hodge Theory and Complex Algebraic Geometry Volume 1. Cambridge University Press, 2002.
Sima Verbitsky
Moscow State University, Faculty of Mathematics and Mechanics
GSP-1 1, Leninskie gory, 119991 Moscow, Russia.
sverb57@gmail.com
|
arxiv-papers
| 2011-11-16T15:12:56 |
2024-09-04T02:49:24.396903
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Sima Verbitsky",
"submitter": "Sima Verbitsky",
"url": "https://arxiv.org/abs/1111.3828"
}
|
1111.3919
|
# Recipe recommendation using ingredient networks
Chun-Yuen Teng
Yu-Ru Lin
Lada A. Adamic
School of Information University of Michigan Ann Arbor, MI, USA
chunyuen@umich.edu IQSS, Harvard University CCS, Northeastern University
Boston, MA yuruliny@gmail.com School of Information University of Michigan
Ann Arbor, MI, USA ladamic@umich.edu
(30 July 1999)
###### Abstract
The recording and sharing of cooking recipes, a human activity dating back
thousands of years, naturally became an early and prominent social use of the
web. The resulting online recipe collections are repositories of ingredient
combinations and cooking methods whose large-scale and variety yield
interesting insights about both the fundamentals of cooking and user
preferences. At the level of an individual ingredient we measure whether it
tends to be essential or can be dropped or added, and whether its quantity can
be modified. We also construct two types of networks to capture the
relationships between ingredients. The complement network captures which
ingredients tend to co-occur frequently, and is composed of two large
communities: one savory, the other sweet. The substitute network, derived from
user-generated suggestions for modifications, can be decomposed into many
communities of functionally equivalent ingredients, and captures users’
preference for healthier variants of a recipe. Our experiments reveal that
recipe ratings can be well predicted with features derived from combinations
of ingredient networks and nutrition information.
###### category:
H.2.8 Database Management Database applications
###### keywords:
Data mining
###### keywords:
ingredient networks, recipe recommendation
††terms: Measurement; Experimentation
## 1 Introduction
The web enables individuals to collaboratively share knowledge and recipe
websites are one of the earliest examples of collaborative knowledge sharing
on the web. Allrecipes.com, the subject of our present study, was founded in
1997, years ahead of other collaborative websites such as the Wikipedia.
Recipe sites thrive because individuals are eager to share their recipes, from
family recipes that had been passed down for generations, to new concoctions
that they created that afternoon, having been motivated in part by the ability
to share the result online. Once shared, the recipes are implemented and
evaluated by other users, who supply ratings and comments.
The desire to look up recipes online may at first appear odd given that tombs
of printed recipes can be found in almost every kitchen. The Joy of Cooking
[12] alone contains 4,500 recipes spread over 1,000 pages. There is, however,
substantial additional value in online recipes, beyond their accessibility.
While the Joy of Cooking contains a single recipe for Swedish meatballs,
Allrecipes.com hosts “Swedish Meatballs I”, “II”, and “III”, submitted by
different users, along with 4 other variants, including “The Amazing Swedish
Meatball”. Each variant has been reviewed, from 329 reviews for “Swedish
Meatballs I" to 5 reviews for “Swedish Meatballs III". The reviews not only
provide a crowd-sourced ranking of the different recipes, but also many
suggestions on how to modify them, e.g. using ground turkey instead of beef,
skipping the “cream of wheat” because it is rarely on hand, etc.
The wealth of information captured by online collaborative recipe sharing
sites is revealing not only of the fundamentals of cooking, but also of user
preferences. The co-occurrence of ingredients in tens of thousands of recipes
provides information about which ingredients go well together, and when a
pairing is unusual. Users’ reviews provide clues as to the flexibility of a
recipe, and the ingredients within it. Can the amount of cinnamon be doubled?
Can the nutmeg be omitted? If one is lacking a certain ingredient, can a
substitute be found among supplies at hand without a trip to the grocery
store? Unlike cookbooks, which will contain vetted but perhaps not the best
variants for some individuals’ tastes, ratings assigned to user-submitted
recipes allow for the evaluation of what works and what does not.
In this paper, we seek to distill the collective knowledge and preference
about cooking through mining a popular recipe-sharing website. To extract such
information, we first parse the unstructured text of the recipes and the
accompanying user reviews. We construct two types of networks that reflect
different relationships between ingredients, in order to capture users’
knowledge about how to combine ingredients. The complement network captures
which ingredients tend to co-occur frequently, and is composed of two large
communities: one savory, the other sweet. The substitute network, derived from
user-generated suggestions for modifications, can be decomposed into many
communities of functionally equivalent ingredients, and captures users’
preference for healthier variants of a recipe. Our experiments reveal that
recipe ratings can be well predicted by features derived from combinations of
ingredient networks and nutrition information (with accuracy .792), while most
of the prediction power comes from the ingredient networks (84%).
The rest of the paper is organized as follows. Section 2 reviews the related
work. Section 3 describes the dataset. Section 4 discusses the extraction of
the ingredient and complement networks and their characteristics. Section 5
presents the extraction of recipe modification information, as well as the
construction and characteristics of the ingredient substitute network. Section
6 presents our experiments on recipe recommendation and Section 7 concludes.
## 2 Related work
Recipe recommendation has been the subject of much prior work. Typically the
goal has been to suggest recipes to users based on their past recipe ratings
[15][3] or browsing/cooking history [16]. The algorithms then find similar
recipes based on overlapping ingredients, either treating each ingredient
equally [4] or by identifying key ingredients [19]. Instead of modeling
recipes using ingredients, Wang et al. [17] represent the recipes as graphs
which are built on ingredients and cooking directions, and they demonstrate
that graph representations can be used to easily aggregate Chinese dishes by
the flow of cooking steps and the sequence of added ingredients. However,
their approach only models the occurrence of ingredients or cooking methods,
and doesn’t take into account the relationships between ingredients. In
contrast, in this paper we incorporate the likelihood of ingredients to co-
occur, as well as the potential of one ingredient to act as a substitute for
another.
Another branch of research has focused on recommending recipes based on
desired nutritional intake or promoting healthy food choices. Geleijnse et al.
[7] designed a prototype of a personalized recipe advice system, which
suggests recipes to users based on their past food selections and nutrition
intake. In addition to nutrition information, Kamieth et al. [9] built a
personalized recipe recommendation system based on availability of ingredients
and personal nutritional needs. Shidochi et al. [14] proposed an algorithm to
extract replaceable ingredients from recipes in order to satisfy users’
various demands, such as calorie constraints and food availability. Their
method identifies substitutable ingredients by matching the cooking actions
that correspond to ingredient names. However, their assumption that
substitutable ingredients are subject to the same processing methods is less
direct and specific than extracting substitutions directly from user-
contributed suggestions.
Ahn et al. [1] and Kinouchi et al [10] examined networks involving ingredients
derived from recipes, with the former modeling ingredients by their flavor
bonds, and the latter examining the relationship between ingredients and
recipes. In contrast, we derive direct ingredient-ingredient networks of both
compliments and substitutes. We also step beyond characterizing these networks
to demonstrating that they can be used to predict which recipes will be
successful.
## 3 Dataset
Allrecipes.com is one of the most popular recipe-sharing websites, where
novice and expert cooks alike can upload and rate cooking recipes. It hosts 16
customized international sites for users to share their recipes in their
native languages, of which we study only the main, English, version. Recipes
uploaded to the site contain specific instructions on how to prepare a dish:
the list of ingredients, preparation steps, preparation and cook time, the
number of servings produced, nutrition information, serving directions, and
photos of the prepared dish. The uploaded recipes are enriched with user
ratings and reviews, which comment on the quality of the recipe, and suggest
changes and improvements. In addition to rating and commenting on recipes,
users are able to save them as favorites or recommend them to others through a
forum.
We downloaded 46,337 recipes including all information listed from
allrecipes.com, including several classifications, such as a region (e.g. the
midwest region of US or Europe), the course or meal the dish is appropriate
for (e.g.: appetizers or breakfast), and any holidays the dish may be
associated with. In order to understand users’ recipe preferences, we crawled
1,976,920 reviews which include reviewers’ ratings, review text, and the
number of users who voted the review as useful.
### 3.1 Data preprocessing
The first step in processing the recipes is identifying the ingredients and
cooking methods from the freeform text of the recipe. Usually, although not
always, each ingredient is listed on a separate line. To extract the
ingredients, we tried two approaches. In the first, we found the maximal match
between a pre-curated list of ingredients and the text of the line. However,
this missed too many ingredients, while misidentifying others. In the second
approach, we used regular expression matching to remove non-ingredient terms
from the line and identified the remainder as the ingredient. We removed
quantifiers, such as e.g. “1 lb” or “2 cups”, words referring to consistency
or temperature, e.g. chopped or cold, along with a few other heuristics, such
as removing content in parentheses. For example “1 (28 ounce) can baked beans
(such as Bush’s Original®)" is identified as “baked beans". By limiting the
list of potential terms to remove from an ingredient entry, we erred on the
side of not conflating potentially identical or highly similar ingredients,
e.g. “cheddar cheese”, used in 2450 recipes, was considered different from
“sharp cheddar cheese”, occurring in 394 recipes.
We then generated an ingredient list sorted by frequency of ingredient
occurrence and selected the top 1000 common ingredient names as our finalized
ingredient list. Each of the top 1000 ingredients occurred in 23 or more
recipes, with plain salt making an appearance in 47.3% of recipes. These
ingredients also accounted for 94.9% of ingredient entries in the recipe
dataset. The remaining ingredients were missed either because of high
specificity (e.g. yolk-free egg noodle), referencing brand names (e.g.
Planters almonds), rarity (e.g. serviceberry), misspellings, or not being a
food (e.g. “nylon netting").
The remaining processing task was to identify cooking processes from the
directions. We first identified all heating methods using a listing in the
Wikipedia entry on cooking [18]. For example, baking, boiling, and steaming
are all ways of heating the food. We then identified mechanical ways of
processing the food such as chopping and grinding, and other chemical
techniques such as marinating and brining.
### 3.2 Regional preferences
Choosing one cooking method over another appears to be a question of regional
taste. 5.8% of recipes were classified into one of five US regions: Mountain,
Midwest, Northeast, South, and West Coast (including Alaska and Hawaii).
Figure 1 shows significantly ($\chi^{2}$ test p-value < 0.001) varying
preferences in the different US regions among 6 of the most popular cooking
methods. Boiling and simmering, both involving heating food in hot liquids,
are more common in the South and Midwest. Marinating and grilling are
relatively more popular in the West and Mountain regions, but in the West more
grilling recipes involve seafood (18/42 = 42%) relative to other regions
combined (7/106 = 6%). Frying is popular in the South and Northeast. Baking is
a universally popular and versatile technique, which is often used for both
sweet and savory dishes, and is slightly more popular in the Northeast and
Midwest. Examination of individual recipes reflecting these frequencies shows
that these differences in preference can be tied to differences in
demographics, immigrant culture and availability of local ingredients, e.g.
seafood.
Figure 1: The percentage of recipes by region that apply a specific heating
method.
## 4 Ingredient complement network
Can we learn how to combine ingredients from the data? Here we employ the
occurrences of ingredients across recipes to distill users’ knowledge about
combining ingredients.
We constructed an ingredient complement network based on pointwise mutual
information (PMI) defined on pairs of ingredients $(a,b)$:
$\mathrm{PMI(a,b)}=log\frac{p(a,b)}{p(a)p(b)},$
where
$p(a,b)=\frac{\mathrm{\\#\>of\>recipes\>containing}\>a\mathrm{\>and}\>b}{\mathrm{\\#\>of\>recipes}},$
$p(a)=\frac{\mathrm{\\#\>of\>recipes\>containing}\>a}{\mathrm{\\#\>of\>recipes}},$
$p(b)=\frac{\mathrm{\\#\>of\>recipes\>containing}\>b}{\mathrm{\\#\>of\>recipes}}.$
The PMI gives the probability that two ingredients occur together against the
probability that they occur separately. Complementary ingredients tend to
occur together far more often than would be expected by chance.
Figure 2 shows a visualization of ingredient complementarity. Two distinct
subcommunities of recipes are immediately apparent: one corresponding to
savory dishes, the other to sweet ones. Some central ingredients, e.g. egg and
salt, actually are pushed to the periphery of the network. They are so
ubiquitous, that although they have many edges, they are all weak, since they
don’t show particular complementarity with any single group of ingredients.
Figure 2: Ingredient complement network. Two ingredients share an edge if they
occur together more than would be expected by chance and if their pointwise
mutual information exceeds a threshold.
We further probed the structure of the complementarity network by applying a
network clustering algorithm [13]. The algorithm confirmed the existence of
two main clusters containing the vast majority of the ingredients. An
interesting satellite cluster is that of mixed drink ingredients, which is
evident as a constellation of small nodes located near the top of the sweet
cluster in Figure 2. The cluster includes the following ingredients: lime,
rum, ice, orange, pineapple juice, vodka, cranberry juice, lemonade, tequila,
etc.
For each recipe we recorded the minimum, average, and maximum pairwise
pointwise mutual information between ingredients. The intuition is that
complementary ingredients would yield higher ratings, while ingredients that
don’t go together would lower the average rating. We found that while the
average and minimum pointwise mutual information between ingredients is
uncorrelated with ratings, the maximum is very slightly positively correlated
with the average rating for the recipe ($\rho=0.09$, p-value < $10^{-10}$).
This suggests that having at least two complementary ingredients very slightly
boosts a recipe’s prospects, but having clashing or unrelated ingredients does
not seem to do harm.
## 5 Recipe modifications
Co-occurrence of ingredients aggregated over individual recipes reveals the
structure of cooking, but tells us little about how flexible the ingredient
proportions are, or whether some ingredients could easily be left out or
substituted. An experienced cook may know that apple sauce is a low-fat
alternative to oil, or may know that nutmeg is often optional, but a novice
cook may implement recipes literally, afraid that deviating from the
instructions may produce poor results. While a traditional hardcopy cookbook
would provide few such hints, they are plentiful in the reviews submitted by
users who implemented the recipes, e.g. “This is a great recipe, but using
fresh tomatoes only adds a few minutes to the prep time and makes it taste so
much better", or another comment about the same salsa recipe “This is by far
the best recipe we have ever come across. We did however change it just a
little bit by adding extra onion.”
As the examples illustrate, modifications are reported even when the user
likes the recipe. In fact, we found that 60.1% of recipe reviews contain words
signaling modification, such as “add", “omit", “instead”, “extra" and 14
others. Furthermore, it is the reviews that include changes that have a
statistically higher average rating (4.49 vs. 4.39, t-test p-value
$<10^{-10}$), and lower rating variance (0.82 vs. 1.05, Bartlett test p-value
$<10^{-10}$), as is evident in the distribution of ratings, shown in Fig. 3.
This suggests that flexibility in recipes is not necessarily a bad thing, and
that reviewers who don’t mention modifications are more likely to think of the
recipe as perfect, or to dislike it entirely.
In the following, we describe the recipe modifications extracted from user
reviews, including adjustment, deletion and addition. We then present how we
constructed an ingredient substitute network based on the extracted
information.
Figure 3: The likelihood that a review suggests a modification to the recipe
depends on the star rating the review is assigning to the recipe.
### 5.1 Adjustments
Some modifications involve increasing or decreasing the amount of an
ingredient in the recipe. In this and the following analyses, we split the
review on punctuation such as commas and periods. We used simple heuristics to
detect when a review suggested a modification: adding/using more/less of an
ingredient counted as an increase/decrease. Doubling or increasing counted as
an increase, while reducing, cutting, or decreasing counted as a decrease.
While it is likely that there are other expressions signaling the adjustment
of ingredient quantities, using this set of terms allowed us to compare the
relative rate of modification, as well as the frequency of increase vs.
decrease between ingredients. The ingredients themselves were extracted by
performing a maximal character match within a window following an adjustment
term.
Figure 4 shows the ratios of the number of reviews suggesting modifications,
either increases or decreases, to the number of recipes that contain the
ingredient. Two patterns are immediately apparent. Ingredients that may be
perceived as being unhealthy, such as fats and sugars, are, with the exception
of vegetable oil and margarine, more likely to be modified, and to be
decreased. On the other hand, flavor enhancers such as soy sauce, lemon juice,
cinnamon, Worcestershire sauce, and toppings such as cheeses, bacon and
mushrooms, are also likely to be modified; however, they tend to be added in
greater, rather than lesser quantities. Combined, the patterns suggest that
good-tasting but “unhealthy" ingredients can be reduced, if desired, while
spices, extracts, and toppings can be increased to taste.
Figure 4: Suggested modifications of quantity for the 50 most common
ingredients, derived from recipe reviews. The line denotes equal numbers of
suggested quantity increases and decreases.
### 5.2 Deletions and additions
Recipes are also frequently modified such that ingredients are omitted
entirely. We looked for words indicating that the reviewer did not have an
ingredient (and hence did not use it), e.g. “had no" and “didn’t have". We
further used “omit/left out/left off/bother with” as indication that the
reviewer had omitted the ingredients, potentially for other reasons. Because
reviewers often used simplified terms, e.g. “vanilla" instead of “vanilla
extract", we compared words in proximity to the action words by constructing
4-character-grams and calculating the cosine similarity between the n-grams in
the review and the list of ingredients for the recipe.
To identify additions, we simply looked for the word “add", but omitted
possible substitutions. For example, we would use “added cucumber", but not
“added cucumber instead of green pepper", the latter of which we analyze in
the following section. We then compared the addition to the list of
ingredients in the recipes, and considered the addition valid only if the
ingredient does not already belong in the recipe.
Table 1 shows the correlation between ingredient modifications. As might be
expected, the more frequently an ingredient occurs in a recipe, the more times
its quantity has the opportunity to be modified, as is evident in the strong
correlation between the the number of recipes the ingredient occurs in and
both increases and decreases recommended in reviews. However, the more common
an ingredient, the more stable it appears to be. Recipe frequency is
negatively correlated with deletions/recipe ($\rho=-0.22$), additions/recipe
($\rho=-0.25$), and increases/recipe ($\rho=-0.26$). For example, salt is so
essential, appearing in over 21,000 recipes, that we detected only 18 reviews
where it was explicitly dropped. In contrast, Worcheshire sauce, appearing in
1,542 recipes, is dropped explicitly in 148 reviews.
As might also be expected, additions are positively correlated with increases,
and deletions with decreases. However, additions and deletions are very weakly
negatively correlated, indicating that an ingredient that is added frequently
is not necessarily omitted more frequently as well.
Table 1: Correlations between ingredient modifications | addition | deletion | increase | decrease
---|---|---|---|---
# recipes | 0.41 | 0.22 | 0.61 | 0.68
addition | | -0.15 | 0.79 | 0.11
deletion | | | 0.09 | 0.58
increase | | | | 0.39
### 5.3 Ingredient substitute network
Replacement relationships show whether one ingredient is preferable to
another. The preference could be based on taste, availability, or price. Some
ingredient substitution tables can be found online111e.g.,
http://allrecipes.com/HowTo/common-ingredient-substitutions/detail.aspx, but
are neither extensive nor contain information about relative frequencies of
each substitution. Thus, we found an alternative source for extracting
replacement relationships – users’ comments, e.g. “I replaced the butter in
the frosting by sour cream, just to soothe my conscience about all the fatty
calories".
To extract such knowledge, we first parsed the reviews as follows: we
considered several phrases to signal replacement relationships: “replace $a$
with $b$”, “substitute $b$ for $a$”, “$b$ instead of $a$”, etc, and matched
$a$ and $b$ to our list of ingredients.
We constructed an ingredient substitute network to capture users’ knowledge
about ingredient replacement. This weighted, directed network consists of
ingredients as nodes. We thresholded and eliminated any suggested
substitutions that occurred fewer than 5 times. We then determined the weight
of each edge by $p(b|a)$, the proportion of substitutions of ingredient $a$
that suggest ingredient $b$. For example, 68% of substitutions for white sugar
were to splenda, an artificial sweetener, and hence the assigned weight for
the $sugar\rightarrow splenda$ edge is 0.68.
Figure 5: Ingredient substitute network. Nodes are sized according to the number of times they have been recommended as a substitute for another ingredient, and colored according to their indegree. Table 2: Clusters of ingredients that can be substituted for one another. A maximum of 5 additional ingredients for each cluster are listed, ordered by PageRank. main | other ingredients
---|---
chicken | turkey, beef, sausage, chicken breast, bacon
olive oil | butter, apple sauce, oil, banana, margarine
sweet | yam, potato, pumpkin, butternut squash,
potato | parsnip
baking | baking soda, cream of tartar
powder |
almond | pecan, walnut, cashew, peanut, sunflower s.
apple | peach, pineapple, pear, mango, pie filling
egg | egg white, egg substitute, egg yolk
tilapia | cod, catfish, flounder, halibut, orange roughy
spinach | mushroom, broccoli, kale, carrot, zucchini
italian | basil, cilantro, oregano, parsley, dill
seasoning |
cabbage | coleslaw mix, sauerkraut, bok choy
| napa cabbage
The resulting substitution network, shown in Figure 5, exhibits strong
clustering. We examined this structure by applying the map generator tool by
Rosvall et al. [13], which uses a random walk approach to identify clusters in
weighted, directed networks. The resulting clusters, and their relationships
to one another, are shown in Fig. 6. The derived clusters could be used when
following a relatively new recipe which may not receive many reviews, and
therefore many suggestions for ingredient substitutions. If one does not have
all ingredients at hand, one could examine the content of one’s fridge and
pantry and match it with other ingredients found in the same cluster as the
ingredient called for by the recipe. Table 2 lists the contents of a few such
sample ingredient clusters, and Fig. 7 shows two example clusters extracted
from the substitute network.
Figure 6: Ingredient substitution clusters. Nodes represent clusters and edges indicate the presence of recommended substitutions that span clusters. Each cluster represents a set of related ingredients which are frequently substituted for one another. |
---|---
(a) milk substitutes | (b) cinammon substitutes
Figure 7: Relationships between ingredients located within two of the clusters
from Fig. 6.
Finally, we examine whether the substitution network encodes preferences for
one ingredient over another, as evidenced by the relative ratings of similar
recipes, one which contains an original ingredient, and another which
implements a substitution. To test this hypothesis, we construct a “preference
network", where one ingredient is preferred to another in terms of received
ratings, and is constructed by creating an edge $(a,b)$ between a pair of
ingredients, where $a$ and $b$ are listed in two recipes $X$ and $Y$
respectively, if recipe ratings $R_{X}>R_{Y}$. For example, if recipe $X$
includes beef, ketchup and cheese, and recipe $Y$ contains beef and pickles,
then this recipe pair contributes to two edges: one from pickles to ketchup,
and the other from pickles to cheese. The aggregate edge weights are defined
based on PMI. Because PMI is a symmetric quantity
($\mathrm{PMI}(a;b)=\mathrm{PMI}(b;a)$), we introduce a directed PMI measure
to cope with the directionality of the preference network:
$\mathrm{PMI}(a\to b)=\mathrm{log}\frac{p(a\to b)}{p(a)p(b)},$
where
$p(a\to
b)=\frac{\mathrm{\\#\>of\>recipe\>pairs\>from}\>a\mathrm{\>to}\>b}{\mathrm{\\#\>of\>recipe\>pairs}},$
and $p(a)$, $p(b)$ are defined as in the previous section.
We find high correlation between this preference network and the substitution
network ($\rho=0.72,p<0.001$). This observation suggests that the substitute
network encodes users’ ingredient preference, which we use in the recipe
prediction task described in the next section.
## 6 Recipe recommendation
We use the above insights to uncover novel recommendation algorithms suitable
for recipe recommendations. We use ingredients and the relationships encoded
between them in ingredient networks as our main feature sets to predict recipe
ratings, and compare them against features encoding nutrition information, as
well as other baseline features such as cooking methods, and preparation and
cook time. Then we apply a discriminative machine learning method, stochastic
gradient boosting trees [6], to predict recipe ratings.
In the experiments, we seek to answer the following three questions. (1) Can
we predict users’ preference for a new recipe given the information present in
the recipe? (2) What are the key aspects that determine users’ preference? (3)
Does the structure of ingredient networks help in recipe recommendation, and
how?
### 6.1 Recipe Pair Prediction
The goal of our prediction task is: given a pair of similar recipes, determine
which one has higher average rating than the other. This task is designed
particularly to help users with a specific dish or meal in mind, and who are
trying to decide between several recipe options for that dish.
Recipe pair data. The data for this prediction task consists of pairs of
similar recipes. The reason for selecting similar recipes, with high
ingredient overlap, is that while apples may be quite comparable to oranges in
the context of recipes, especially if one is evaluating salads or desserts,
lasagna may not be comparable to a mixed drink. To derive pairs of related
recipes, we computed similarity with a cosine similarity between the
ingredient lists for the two recipes, weighted by the inverse document
frequency,
$log(\\#\>of\>recipes/\\#\>of\>recipes\>containing\>the\>ingredient)$. We
considered only those pairs of recipes whose cosine similarity exceeded 0.2.
The weighting is intended to identify higher similarity among recipes sharing
more distinguishing ingredients, such as Brussels sprouts, as opposed to
recipes sharing very common ones, such as butter.
A further challenge to obtaining reliable relative rankings of recipes is
variance introduced by having different users choose to rate different
recipes. In addition, some users might not have a sufficient number of reviews
under their belt to have calibrated their own rating scheme. To control for
variation introduced by users, we examined recipe pairs where the same users
are rating both recipes and are collectively expressing a preference for one
recipe over another. Specifically, we generated 62,031 recipe pairs ($a,b$)
where $rating_{i}(a)$ > $rating_{i}(b)$, for at least 10 users $i$, and over
50% of users who rated both recipe $a$ and recipe $b$. Furthermore, each user
$i$ should be an active enough reviewer to have rated at least 8 other
recipes.
Features. In the prediction dataset, each observation consists of a set of
predictor variables or features that represent information about two recipes,
and the response variable is a binary indicator of which gets the higher
rating on average. To study the key aspects of recipe information, we
constructed different set of features, including:
* •
Baseline: This includes cooking methods, such as chopping, marinating, or
grilling, and cooking effort descriptors, such as preparation time in minutes,
as well as the number of servings produced, etc. These features are considered
as primary information about a recipe and will be included in all other
feature sets described below.
* •
Full ingredients: We selected up to 1000 popular ingredients to build a “full
ingredient list”. In this feature set, each observed recipe pair contains a
vector with entries indicating whether an ingredient from the full list is
present in either recipe in the pair.
* •
Nutrition: This feature set does not include any ingredients but only
nutrition information such the total caloric content, as well as quantities of
fats, carbohydrates, etc.
* •
Ingredient networks: In this set, we replaced the full ingredient list by
structural information extracted from different ingredient networks, as
described in Sections 4 and 5.3. Co-occurrence is treated separately as a raw
count, and a complementarity, captured by the PMI.
* •
Combined set: Finally, a combined feature set is constructed to test the
performance of a combination of features, including baseline, nutrition and
ingredient networks.
To build the ingredient network feature set, we extracted the following two
types of structural information from the co-occurrence and substitution
networks, as well as the complement network derived from the co-occurrence
information:
Network positions are calculated to represent how a recipe’s ingredients
occupy positions within the networks. Such position measures are likely to
inform if a recipe contains any “popular” or “unusual” ingredients. To
calculate the position measures, we first calculated various network
centrality measures, including degree centrality, betweenness centrality,
etc., from the ingredient networks. A centrality measure can be represented as
a vector $\vec{g}$ where each entry indicates the centrality of an ingredient.
The network position of a recipe, with its full ingredient list represented as
a binary vector $\vec{f}$, can be summarized by $\vec{g}^{T}\cdot\vec{f}$,
i.e., an aggregated centrality measure based on the centrality of its
ingredients.
Network communities provide information about which ingredient is more likely
to co-occur with a group of other ingredients in the network. A recipe
consisting of ingredients that are frequently used with, complemented by or
substituted by certain groups may be predictive of the ratings the recipe will
receive. To obtain the network community information, we applied latent
semantic analysis (LSA) on recipes. We first factorized each ingredient
network, represented by matrix $W$, using singular value decomposition (SVD).
In the matrix $W$, each entry $W_{ij}$ indicates whether ingredient $i$ co-
occurrs, complements or substitues ingredient $j$.
Suppose $W_{k}=U_{k}\Sigma_{k}V_{k}^{T}$ is a rank-$k$ approximation of $W$,
we can then transform each recipe’s full ingredient list using the low-
dimensional representation, $\Sigma_{k}^{-1}V_{k}^{T}\vec{f}$, as community
information within a network. These low-dimensional vectors, together with the
vectors of network positions, constitute the ingredient network features.
Figure 8: Prediction performance. The nutrition information and ingredient
networks are more effective features than full ingredients. The ingredient
network features lead to impressive performance, close to the best
performance. Figure 9: Relative importance of features in the combined set.
The individual items from nutrition information are very indicative in
differentiating highly rated recipes, while most of the prediction power comes
from ingredient networks.
Learning method. We applied discriminative machine learning methods such as
support vector machines (SVM) [2] and stochastic gradient boosting trees [5]
to our prediction problem. Here we report and discuss the detailed results
based on the gradient boosting tree model. Like SVM, the gradient boosting
tree model seeks a parameterized classifier, but unlike SVM that considers all
the features at one time, the boosting tree model considers a set of features
at a time and iteratively combines them according to their empirical errors.
In practice, it not only has competitive performance comparable to SVM, but
can serve as a feature ranking procedure [11].
In this work, we fitted a stochastic gradient boosting tree model with 8
terminal nodes under an exponential loss function. The dataset is roughly
balanced in terms of which recipe is the higher-rated one within a pair. We
randomly divided the dataset into a training set (2/3) and a testing set
(1/3). The prediction performance is evaluated based on accuracy, and the
feature performance is evaluated in terms of relative importance [8]. For each
single decision tree, one of the input variables, $x^{j}$, is used to
partition the region associated with that node into two subregions in order to
fit to the response values. The squared relative importance of variable
$x^{j}$ is the sum of such squared improvements over all internal nodes for
which it was chosen as the splitting variable, as:
$imp(j)=\sum_{k}\hat{i}_{k}^{2}I(\mbox{splits on }x^{j})$
where $\hat{i}_{k}^{2}$ is the empirical improvement by the $k$-th node
splitting on $x^{j}$ at that point.
### 6.2 Results
Figure 10: Relative importance of features representing the network
structure. The substitution network has the strongest contribution ($39.8\%$)
to the total importance of network features, and it also has more influential
features in the top 100 list, which suggests that the substitution network is
complementary to other features. Figure 11: Relative importance of features
from nutrition information. The carbs item is the most influential feature in
predicting higher-rated recipes.
The overall prediction performance is shown in Fig. 8. Surprisingly, even with
a full list of ingredients, the prediction accuracy is only improved from .712
(baseline) to .746. In contrast, the nutrition information and ingredient
networks are more effective (with accuracy .753 and .786, respectively). Both
of them have much lower dimensions (from tens to several hundreds), compared
with the full ingredients that are represented by more than 2000 dimensions
(1000 ingredients per recipe in the pair). The ingredient network features
lead to impressive performance, close to the best performance given by the
combined set (.792), indicating the power of network structures in recipe
recommendation.
Figure 9 shows the influence of different features in the combined feature
set. Up to 100 features with the highest relative importance are shown. The
importance of a feature group is summarized by how much the total importance
is contributed by all features in the set. For example, the baseline
consisting of cooking effort and cooking methods contribute $8.9\%$ to the
overall performance. The individual items from nutrition information are very
indicative in differentiating highly-rated recipes, while most of the
prediction power comes from ingredient networks ($84\%$).
Figure 10 shows the top 100 features from the three networks. In terms of the
total importance of ingredient network features, the substitution network has
slightly stronger contribution ($39.8\%$) than the other two networks, and it
also has more influential features in the top 100 list. This suggests that the
structural information extracted from the substitution network is not only
important but also complementary to information from other aspects.
Looking into the nutrition information (Fig. 11), we found that carbohydrates
are the most influential feature in predicting higher-rated recipes. Since
carbohydrates comprise around $50\%$ or more of total calories, the high
importance of this feature interestingly suggests that a recipe’s rating can
be influenced by users’ concerns about nutrition and diet. Another interesting
observation is that, while individual nutrition items are powerful predictors,
a higher prediction accuracy can be reached by using ingredient networks
alone, as shown in Fig. 8. This implies the information about nutrition may
have been encoded in the ingredient network structure, e.g. substitutions of
less healthful ingredients with “healthier” alternatives.
Figure 12: Prediction performance over reduced dimensionality. The best
performance is given by reduced dimension $k=50$ when combining all three
networks. In addition, using the information about the complement network
alone is more effective in prediction than using other two networks. Figure
13: Influential substitution communities. The matrix shows the most
influential feature dimensions extracted from the substitution network. For
each dimension, the six representative ingredients with the highest intensity
values are shown, with colors indicating their intensity. These features
suggest that the communities of ingredient substitutes, such as the sweet and
oil in the first dimension, are particularly informative in prediction.
Constructing the ingredient network feature involves reducing high-dimensional
network information through SVD, as described in the previous section. The
dimensionality can be determined by cross-validation. As shown in Fig. 12,
features with a very large dimension tend to overfit the training data. Hence
we chose $k=50$ for the reduced dimension of all three networks. The figure
also shows that using the information about the complement network alone is
more effective in prediction than using either the co-occurrence and
substitute networks, even in the case of low dimensions. Consistently, as
shown in terms of relative importance (Fig. 10), the substitution network
alone is not the most effective, but it provides more complementary
information in the combined feature set.
In Figure 13 we show the most representative ingredients in the decomposed
matrix derived from the substitution network. We display the top five
influential dimensions, evaluated based on the relative importance, from the
SVD resultant matrix $V_{k}$, and in each of these dimensions we extracted six
representative ingredients based on their intensities in the dimension (the
squared entry values). These representative ingredients suggest that the
communities of ingredient substitutes, such as the sweet and oil substitutes
in the first dimension or the milk substitutes in the second dimesion (which
is similar to the cluster shown in Fig. 6), are particularly informative in
predicting recipe ratings.
To summarize our observations, we find we are able to effectively predict
users’ preference for a recipe, but the prediction is not through using a full
list of ingredients. Instead, by using the structural information extracted
from the relationships among ingredients, we can better uncover users’
preference about recipes.
## 7 Conclusion
Recipes are little more than instructions for combining and processing sets of
ingredients. Individual cookbooks, even the most expansive ones, contain
single recipes for each dish. The web, however, permits collaborative recipe
generation and modification, with tens of thousands of recipes contributed in
individual websites. We have shown how this data can be used to glean insights
about regional preferences and modifiability of individual ingredients, and
also how it can be used to construct two kinds of networks, one of ingredient
complements, the other of ingredient substitutes. These networks encode which
ingredients go well together, and which can be substituted to obtain superior
results, and permit one to predict, given a pair of related recipes, which one
will be more highly rated by users.
In future work, we plan to extend ingredient networks to incorporate the
cooking methods as well. It would also be of interest to generate region-
specific and diet-specific ratings, depending on the users’ background and
preferences. A whole host of user-interface features could be added for users
who are interacting with recipes, whether the recipe is newly submitted, and
hence unrated, or whether they are browsing a cookbook. In addition to
automatically predicting a rating for the recipe, one could flag ingredients
that can be omitted, ones whose quantity could be tweaked, as well as
suggested additions and substitutions.
## 8 Acknowledgments
This work was supported by MURI award FA9550-08-1-0265 from the Air Force
Office of Scientific Research. The methodology used in this paper was
developed with support from funding from the Army Research Office, Multi-
University Research Initiative on Measuring, Understanding, and Responding to
Covert Social Networks: Passive and Active Tomography. The authors gratefully
acknowledge D. Lazer for support.
## References
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* [2] Cortes, C., and Vapnik, V. Support-vector networks. Machine learning 20, 3 (1995), 273–297.
* [3] Forbes, P., and Zhu, M. Content-boosted matrix factorization for recommender systems: Experiments with recipe recommendation. Proceedings of Recommender Systems (2011).
* [4] Freyne, J., and Berkovsky, S. Intelligent food planning: personalized recipe recommendation. In IUI, ACM (2010), 321–324.
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* [10] Kinouchi, O., Diez-Garcia, R., Holanda, A., Zambianchi, P., and Roque, A. The non-equilibrium nature of culinary evolution. New Journal of Physics 10 (2008), 073020.
* [11] Lu, Y., Peng, F., Li, X., and Ahmed, N. Coupling feature selection and machine learning methods for navigational query identification. In CIKM, ACM (2006), 682–689.
* [12] Rombauer, I., Becker, M., Becker, E., and Maestro, L. Joy of cooking. Scribner Book Company, 1997.
* [13] Rosvall, M., and Bergstrom, C. Maps of random walks on complex networks reveal community structure. PNAS 105, 4 (2008), 1118.
* [14] Shidochi, Y., Takahashi, T., Ide, I., and Murase, H. Finding replaceable materials in cooking recipe texts considering characteristic cooking actions. In Proc. of the ACM multimedia 2009 workshop on Multimedia for cooking and eating activities, ACM (2009), 9–14.
* [15] Svensson, M., Höök, K., and Cöster, R. Designing and evaluating kalas: A social navigation system for food recipes. ACM Transactions on Computer-Human Interaction (TOCHI) 12, 3 (2005), 374–400.
* [16] Ueda, M., Takahata, M., and Nakajima, S. User’s food preference extraction for personalized cooking recipe recommendation. Proc. of the Second Workshop on Semantic Personalized Information Management: Retrieval and Recommendation (2011).
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* [18] Wikipedia. Outline of food preparation, 2011. [Online; accessed 22-Oct-2011].
* [19] Zhang, Q., Hu, R., Mac Namee, B., and Delany, S. Back to the future: Knowledge light case base cookery. In Proc. of The 9th European Conference on Case-Based Reasoning Workshop (2008), 15.
|
arxiv-papers
| 2011-11-16T19:32:57 |
2024-09-04T02:49:24.404707
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Chun-Yuen Teng, Yu-Ru Lin, and Lada A. Adamic",
"submitter": "Lada A. Adamic",
"url": "https://arxiv.org/abs/1111.3919"
}
|
1111.3925
|
# A Low-Delay Low-Complexity EKF Design for Joint Channel and CFO Estimation
in Multi-User Cognitive Communications
Pengkai Zhao Electrical Engineering, UCLA, CA, USA Cong Shen Qualcomm Inc.,
San Diego, CA, USA
###### Abstract
Parameter estimation in cognitive communications can be formulated as a multi-
user estimation problem, which is solvable under maximum likelihood solution
but involves high computational complexity. This paper presents a time-sharing
and interference mitigation based EKF (Extended Kalman Filter) design for
joint CFO (carrier frequency offset) and channel estimation at multiple
cognitive users. The key objective is to realize low implementation complexity
by decomposing high-dimensional parameters into multiple separate low-
dimensional estimation problems, which can be solved in a time-shared manner
via pipelining operation. We first present a basic EKF design that estimates
the parameters from one TX user to one RX antenna. Then such basic design is
time-shared and reused to estimate parameters from multiple TX users to
multiple RX antennas. Meanwhile, we use interference mitigation module to
cancel the co-channel interference at each RX sample. In addition, we further
propose adaptive noise variance tracking module to improve the estimation
performance. The proposed design enjoys low delay and low buffer size (because
of its online real-time processing), as well as low implementation complexity
(because of time-sharing and pipeling design). Its estimation performance is
verified to be close to Cramer-Rao bound.
## I Introduction
Cognitive communication system is widely accepted as a perspective way in
increasing the spectrum efficiency of wireless networks, where primary links
and secondary links can usually co-exist in the network, resulting in an
interference limited environment. Parameter estimation in cognitive
communications is a challenging problem because of (i) the existence of co-
channel interference, and (ii) the high-dimensional parameters from multiple
TX users to multiple RX antennas. In particular, note that different TX users
often have independent carrier frequency offset (CFO) values (including both
oscillator offsets and doppler offsets), which usually introduce serious
nonlinear components within the observed signal, complicating the estimation
problem. Meanwhile, channel responses from multiple TX users to multiple RX
antennas can result in a set of high-dimensional parameters, which are also
difficult to estimate. Finally, due to the existence of multi-user
interference, CFO and channel parameters usually have to be treated together
and be estimated in a joint way so as to approach the optimal performance,
which further increases the estimation complexity.
Without loss of generality, this paper assumes Orthogonal Frequency Division
Multiplexing (OFDM) system, which is an overwhelming choice for modern
wireless systems. The classical CFO and channel estimation method in a single-
user OFDM system is based on two repeated training symbols [1]. It has low
implementation complexity and near-optimal performance, but only applies to a
single-user scenario111It is also applicable for multi-user scenario with non-
overlapping training symbols, but this is not the case considered in this
paper.. In multi-user OFDMA systems with unique subcarrier set per TX user,
CFO and channel parameters can be recovered by exploiting distinct subcarrier
structures among TX users ([2, 3]). But this method requires separate
subcarrier allocation for different users. Consequently, in a general multi-
user cognitive system without specific subcarrier allocation per user and with
overlapped training symbols, ML and EM related methods seem to be the only
applicable choice, where all TX users’ parameters have to be formulated into a
maximum likelihood (ML) estimation problem [4], which is solvable under
Expectation Maximization (EM) method [5] in an iterative way. However, since
the entire OFDM block is stored offline and is iteratively processed multiple
times, these ML and EM approaches often require high computational complexity
and high processing delay.
Based on above considerations, this paper will focus on using Kalman filter
structure to estimate the CFO and channel parameters in multi-user cognitive
communications. Our major objective is to achieve low-complexity and low-delay
estimation performance in cognitive systems. In general, Kalman filter is a
good candidate for low delay and low complexity parameter estimation primarily
due to its real-timing processing property. It has been conventionally used
for CFO and channel estimation in multi-user OFDM systems, e.g., the FFT-Block
EKF design in [6], the parallel EKF design in [7], and the particle filter
design in [7]. However, these existing designs inherently suffer from multiple
issues related to complexity, delay and buffer size as follows:
1. 1.
Block EKF design in [6] operates on an FFT-block basis, which grows
increasingly complex as FFT size becomes large (e.g., 2048 FFT size). Also,
parameters estimated in this method are handled in a high dimension manner.
2. 2.
Parallel EKF design in [7] also operates on an FFT-block basis, marking it
complex under large FFT size. Parameters in this design are jointly estimated
by calculating the covariance information between different users, leading to
a high matrix dimension.
3. 3.
Beyond FFT size and parameter dimension issues similar to item 1 and 2,
particle filter design in [7] needs to repeat the Kalman operation at multiple
particle samples, yielding a multiplicative effect on complexity.
To summarize, the major challenge in implementing a low-complexity EKF design
lies in the factors of: (i) multiple TX users; (ii) multiple RX antennas;
(iii) high parameter dimension; and (iv) large FFT size (e.g., 2048 size).
With low-complexity and low-delay requirement in mind, this paper will present
a time-sharing and interference mitigation based Extended Kalman Filter (EKF)
design for multi-user cognitive communications, which can estimate the CFO and
channel parameters from multiple TX users to multiple RX antennas in a time-
sharing manner. Here low delay property is achieved by using Kalman filter
estimation at each RX sample in a real-time manner, and low complexity
property is achieved by reusing a single user EKF design in a time-
sharing222Time-sharing in this paper indicates that the same hardware module
can be reused by different processes at separate time slots. and pipeling way.
We first present a fundamental EKF design that estimates the CFO and channel
parameters from one TX user to one RX antenna. Then such basic EKF design is
reused in a time-shared way to estimate the parameters from multiple TX users
to multiple RX antennas. Meanwhile, at each RX sample, an interference
mitigation strategy is developed to estimate and remove the expected multi-
user interference. In addition, we provide an adaptive noise variance tracking
module to further enhance the estimation performance. Because of the usage of
EKF structure, our design is essentially different from the particle filters
in [7] and the EM method in [5]. Our design is also different from the
Parallel-EKF design in [7] and the FFT-Block EKF design in [6] at the
following perspectives: (i) our design runs at each time domain RX sample, not
at an FFT-block basis; (ii) our design treats each user separately, not
jointly; (iii) our design can be implemented in a time-sharing way, which is
less considered in [6] and [7]; (iv) system model in our design is different
from the ones in [6, 7] by integrating CFO parameter into channel response
(see Eqn. (5) in section II). Analysis and simulations results validate that
our proposed design can closely approach the Cramer-Rao bound, and has lower
computational complexity than the ones in [6, 7]. Finally, although cognitive
communication is a typical application scenario for our proposed design, it is
also applicable in many other multi-user systems that satisfy the conditions
presented in section II.
## II System Model
### II-A Problem Formulation
We consider a total of $Q$ TX users in the cognitive network. One of them is
the primary TX user (i.e., base station), and the rest are all secondary TX
users. Primary TX user’s transmission is based on a time division MAC
protocol, where time is divided into different time frames with equal
duration. Secondary users can maintain time synchronization with the primary
TX user by learning and synchronizing with its time frames. Each secondary RX
user is equipped with $N_{A}$ multiple antennas to decode the packets. Without
loss of generality, we assume that every TX user has only one spatial stream,
and there exists $Q\leq N_{A}$. Also, every TX user has a distinct training
symbol333This unique training symbol can be determined according to either the
unique user ID in the network, or the access order in the current time frame.
$s_{q}(n)$ with $1\leq q\leq Q$ and $0\leq n\leq N_{F}-1$. Here $N_{F}$ is the
FFT size of OFDM system.
Each TX user has an independent carrier frequency offset (CFO) that is caused
by both the oscillator offset and the doppler offset. Denote TX user $q$’s CFO
value as $\varepsilon_{q}$. For a given secondary RX user, the channel from TX
user $q$ to the $m$th RX antenna of this secondary user is denoted as
$h_{q,m}(p_{l}^{q,m})$, $\forall\ 1\leq l\leq L_{\rm max}$. Here $L_{\rm max}$
is the number of time domain paths in the channel response, and $p_{l}^{q,m}$
is an integer value representing the relative delay of the $l$th path. We
assume that all $p_{l}^{q,m}$ values have already been determined at an early
stage (e.g., using PN sequences at coarse synchronization). The received
signal at the $m$th RX antenna is derived as:
$\displaystyle y_{m}(n)$ $\displaystyle=$
$\displaystyle\sum_{q=1}^{Q}\exp\left(j\frac{2\pi\varepsilon_{q}n}{N_{F}}\right)\sum_{l=1}^{L_{\rm
max}}h_{q,m}(p_{l}^{q,m})s_{q}\left[(n-p_{l}^{q,m})_{N_{F}}\right]$ (1)
$\displaystyle+z_{m}(n),0\leq n\leq N_{F}-1$
where $(n-p_{l}^{q,m})_{N_{F}}=\left\\{(n-{p_{l}^{q,m}}){\ \rm mod\
}N_{F}\right\\}$ is circular shift, and $z_{m}(n)$ is the background noise at
the $m$th RX antenna.
The task in this work is to estimate CFO parameter $\varepsilon_{q}$ and
channel parameter $h_{q,m}(p_{l}^{q,m})$ for all users ($1\leq q\leq Q$) and
all antennas ($1\leq m\leq N_{A}$). Obviously, the optimal estimation is the
solution to this maximum likelihood (ML) problem:
$\displaystyle\min\sum_{m=1}^{N_{A}}{\big{|}}y_{m}(n)-$ (2)
$\displaystyle\sum_{q=1}^{Q}\exp\left(j\frac{2\pi{\widehat{\varepsilon}}_{q}n}{N_{F}}\right)\sum_{l=1}^{L_{\rm
max}}{\widehat{h}}_{q,m}({p_{l}^{q,m}})s_{q}\left[(n-p_{l}^{q,m})_{N_{F}}\right]{\big{|}}^{2}$
where ${\widehat{\varepsilon}}_{q}$ and ${\widehat{h}}_{q,m}({p_{l}^{q,m}})$
represent the estimated values. There are a total of $(L_{\rm max}N_{A}+1)Q$
parameters in Eqn. (2), which constitutes a high-dimensional parameter
estimation problem.
### II-B State-Space Formulation
ML solution can generally approach the optimal performance but it requires
huge computations, which are highly undesirable in most systems. Instead, this
paper proposes an EKF design for the estimation of the CFO and channel
parameters, which can sequentially update the estimation results at each RX
sample, resulting in low buffer size and low estimation delay. Initially, it
is straightforward to directly apply an EKF design at Eqn. (2) by building all
CFO and channel parameters into one state vector, whose dimension is as high
as $(L_{\rm max}N_{A}+1)Q$. This method will significantly increase the
complexity of the derived Kalman filter. With such complexity consideration in
mind, we first propose a low-dimensional EKF design that can estimate the
parameters from one TX user to one RX antenna, which has only $(L_{\rm
max}+1)$ parameters. Then we reuse this fundamental EKF design in a time-
shared manner to estimate the parameters from multiple TX users to multiple RX
antennas. In this way, high-dimensional parameters are estimated by
sequentially reusing a low-dimensional estimator, which reduces the complexity
of the proposed EKF design.
We first present an RX signal formulation from the perspective of TX user $q$
and the $m$th RX antenna as:
$\displaystyle y_{q,m}(n)$ $\displaystyle=$
$\displaystyle\exp\left(j\frac{2\pi\varepsilon_{q}n}{N_{F}}\right)\sum_{l=1}^{L_{\rm
max}}h_{q,m}(p_{l}^{q,m})s_{q}\left[(n-p_{l}^{q,m})_{N_{F}}\right]$ (3)
$\displaystyle+z_{q,m}(n),$
here $y_{q,m}(n)$ is extracted from $y_{m}(n)$ with the aid of interference
mitigation module, and $z_{q,m}(n)$ represents the residual noise at TX user
$q$ and the $m$th RX antenna, which includes both the residual co-channel
interference from other TX users and the background noise at the $m$th RX
antenna. Additionally, the initial value of $y_{q,m}(n)$ without any
interference mitigation is set to $y_{q,m}(n)=y_{m}(n)$. Details about the
interference mitigation module will be given in section III.
Now we define the associated state vector as:
$\displaystyle{\rm{\bf X}}_{q,m}(n)=\left[\varepsilon_{q},{\rm{\bf
H}}_{q,m}(n)\right]^{T},$ (4) $\displaystyle{\rm{\bf
H}}_{q,m}(n)=\exp\left(j\frac{2\pi
n\varepsilon_{q}}{N_{F}}\right)\left[h_{q,m}(p^{q,m}_{1}),...,h_{q,m}(p^{q,m}_{L_{\rm
max}})\right].$ (5)
The state-space model for TX user $q$ and the $m$th RX antenna can be derived
from (3) as:
$\displaystyle{\rm{\bf X}}_{q,m}(n)=f\left\\{{\rm{\bf
X}}_{q,m}(n-1)\right\\}={\rm{\bf D}}_{q,m}(\varepsilon_{q}){\rm{\bf
X}}_{q,m}(n-1),$ (6) $\displaystyle{\rm{\bf
D}}_{q,m}(\varepsilon_{q})=\left[\begin{array}[]{cc}1&{\rm{\bf 0}}_{1\times
L_{\rm max}}\\\ {\rm{\bf 0}}_{L_{\rm max}\times
1}&\exp\left(j2\pi\varepsilon_{q}/N_{F}\right){\rm{\bf I}}_{L_{\rm max}\times
L_{\rm max}}\end{array}\right]$ (9) $\displaystyle y_{q,m}(n)={\rm{\bf
G}}_{q,m}(n){\rm{\bf X}}_{q,m}(n)+z_{q,m}(n),$ (10) $\displaystyle{\rm{\bf
G}}_{q,m}(n)=\big{[}0,s_{q}\left[(n-p^{q,m}_{1})_{N_{F}}\right],s_{q}\left[(n-p^{q,m}_{2})_{N_{F}}\right],...,$
$\displaystyle s_{q}\left[(n-p^{q,m}_{L_{\rm max}})_{N_{F}}\right]\big{]}.$
(11)
Finally, it is worth mentioning that the derived state-space formulation is a
nonlinear model, since there is a nonlinear component
$\exp\left(j2\pi\varepsilon_{q}/N_{F}\right)$ in the matrix ${\rm{\bf
D}}_{q,m}(\varepsilon_{q})$.
## III Interference Mitigation based EKF Design
This section sequentially describes (i) the basic EKF design that aims at only
one TX user and one RX antenna, (ii) the interference mitigation module that
cancels co-channel interference at each RX sample, (iii) the proposed adaptive
noise variance tracking module, and (iv) the overall paradigm of the proposed
design.
### III-A Fundamental EKF Design
The key idea behind the EKF design is using Jacobian derivation to linearize
the nonlinear matrix ${\rm{\bf D}}_{q,m}(\varepsilon_{q})$ at local estimates:
$\displaystyle{\rm{\bf F}}_{q,m}(n-1)=\left.\frac{\partial f\left({\rm{\bf
X}}_{q,m}(n-1)\right)}{\partial{\rm{\bf
X}}_{q,m}(n-1)}\right|_{\widehat{{\rm{\bf X}}}_{q,m}(n-1|n-1)}$
$\displaystyle=\left[\begin{array}[]{cc}1&{\rm{\bf 0}}_{1\times L_{\rm
max}}\\\ \alpha(n-1){\widehat{\rm{\bf
H}}}_{q,m}^{T}(n-1|n-1)&\exp\left(\alpha(n-1)\right){\rm{\bf I}}_{L_{\rm
max}\times L_{\rm max}}\end{array}\right]$ (14) (15)
$\displaystyle\alpha(n-1)=j2\pi\widehat{\varepsilon}_{q,m}(n-1|n-1)/{N_{F}}$
(16)
where ${\widehat{{\rm{\bf X}}}_{q,m}(n-1|n-1)}$ represents the estimated state
vector after processing the $(n-1)$th RX sample. Based on (III-A), the
prediction steps in our fundamental EKF design are:
$\displaystyle{\widehat{{\rm{\bf X}}}_{q,m}(n|n-1)}=$ $\displaystyle{\rm{\bf
D}}_{q,m}(\widehat{\varepsilon}_{q,m}(n-1|n-1)){\widehat{{\rm{\bf
X}}}_{q,m}(n-1|n-1)},$ (17) $\displaystyle{\rm{\bf P}}_{q,m}(n|n-1)=$
$\displaystyle{\rm{\bf F}}_{q,m}(n-1){\rm{\bf P}}_{q,m}(n-1|n-1){\rm{\bf
F}}_{q,m}^{H}(n-1).$ (18)
And the updating steps are as follows:
$\displaystyle{\rm{\bf K}}_{q,m}(n)={\rm{\bf P}}_{q,m}(n|n-1){\rm{\bf
G}}_{q,m}^{H}(n)\times$ $\displaystyle\left[{\rm{\bf G}}_{q,m}(n){\rm{\bf
P}}_{q,m}(n|n-1){\rm{\bf G}}_{q,m}^{H}(n)+\sigma_{q,m}^{2}(n)\right]^{-1},$
(19) $\displaystyle{\rm{\bf P}}_{q,m}(n|n)=$ $\displaystyle\left[{\rm{\bf
I}}-{\rm{\bf K}}_{q,m}(n){\rm{\bf G}}_{q,m}(n)\right]{\rm{\bf
P}}_{q,m}(n|n-1),$ (20) $\displaystyle{\widehat{{\rm{\bf
X}}}_{q,m}(n|n)}={\widehat{{\rm{\bf X}}}_{q,m}(n|n-1)}+$
$\displaystyle{\rm{\bf K}}_{q,m}(n)\left[y_{q,m}(n)-{\rm{\bf
G}}_{q,m}(n){\widehat{{\rm{\bf X}}}_{q,m}(n|n-1)}\right].$ (21)
Here $\sigma_{q,m}^{2}(n)$ represents the variance of the observation noise
$z_{q,m}(n)$ in Eqn. (3). Also, CFO estimate
$\widehat{\varepsilon}_{q,m}(n|n)$ in the state vector ${\widehat{{\rm{\bf
X}}}_{q,m}(n|n)}$ should only use its real part as
$\widehat{\varepsilon}_{q,m}(n|n)={\rm
Real}\left[\widehat{\varepsilon}_{q,m}(n|n)\right]$. Finally, the EKF design
presented above is only used for the parameter estimation of one TX user and
one RX antenna. This basic EKF design is then iterated in a time-shared manner
to estimate the parameters of all TX users ($1\leq q\leq Q$) and all RX
antennas ($1\leq m\leq N_{A}$).
### III-B Interference Mitigation and Refined CFO Estimation
Before describing the interference mitigation module, we first look at the
refinement of the CFO estimates. Although TX user $q$ has only one CFO
parameter, our proposed EKF design can result in $N_{A}$ different estimates
that are derived from $N_{A}$ RX antennas, which are denoted as
$\widehat{\varepsilon}_{q,m}(n|n-1)$, $1\leq m\leq N_{A}$. As a result, we can
use these $N_{A}$ different estimates to get a refined result
$\widehat{\varepsilon}_{q}(n|n-1)$, which is calculated as:
$\displaystyle\widehat{\varepsilon}_{q}(n|n-1)=$
$\displaystyle\sum_{m=1}^{N_{A}}\frac{1/{\rm{\bf
P}}_{q,m}(n|n-1)_{(1,1)}}{\sum_{r=1}^{N_{A}}1/{\rm{\bf
P}}_{q,r}(n|n-1)_{(1,1)}}\widehat{\varepsilon}_{q,m}(n|n-1),$ (22)
where ${\rm{\bf P}}_{q,m}(n|n-1)_{(1,1)}$ denotes ${\rm{\bf
P}}_{q,m}(n|n-1)$’s element located at the $1^{\rm st}$ row and $1^{\rm st}$
column.
Recall that $y_{q,m}(n)$ involved in (3) and (III-A) is extracted from the
original RX signal $y_{m}(n)$ with the help of an interference mitigation
strategy. Having derived the refined CFO estimate
$\widehat{\varepsilon}_{q}(n|n-1)$, now the interference mitigation process
can be applied at $y_{q,m}(n)$ as follows:
$\displaystyle y_{q,m}(n)=y_{m}(n)-$ $\displaystyle\sum_{u=1,u\neq
q}^{Q}\exp\left(j\frac{2\pi\widehat{\varepsilon}_{u}(n|n-1)}{N_{F}}\right)\cdot{\rm{\bf
G}}_{u,m}(n){\widehat{{\rm{\bf X}}}_{u,m}(n|n-1)}.$
### III-C Adaptive Noise Variance Tracking
Since $y_{q,m}(n)$ is extracted from $y_{m}(n)$ via interference mitigation
module, the variance of $z_{q,m}(n)$, (i.e., $\sigma_{q,m}^{2}(n)$ used in
Eqn. (III-A)), is varying during the convergence process of the interference
mitigation module, which has to be adaptively tracked. Such variance tracking
is based on the following observation:
$\displaystyle\mathbb{E}\left|y_{q,m}(n)-{\rm{\bf
G}}_{q,m}(n){\widehat{\rm{\bf X}}}_{q,m}(n|n-1)\right|^{2}\approx$
$\displaystyle{\rm{\bf G}}_{q,m}(n){\rm{\bf P}}_{q,m}(n|n-1){\rm{\bf
G}}^{H}_{q,m}(n)+\sigma_{q,m}^{2}(n).$ (24)
Using (19), noise variance $\sigma_{q,m}^{2}(n)$ can be tracked as:
$\displaystyle\sigma_{q,m}^{2}(n)$ $\displaystyle=$
$\displaystyle\left[1-\frac{1-b}{1-b^{n+1}}\right]\cdot\sigma_{q,m}^{2}(n-1)+$
(25)
$\displaystyle\frac{1-b}{1-b^{n+1}}\cdot\left\\{\max\left[e_{q,m}(n),0\right]\right\\},$
$\displaystyle e_{q,m}(n)$ $\displaystyle=$
$\displaystyle\left|y_{q,m}(n)-{\rm{\bf G}}_{q,m}(n){\widehat{\rm{\bf
X}}}_{q,m}(n|n-1)\right|^{2}-$ (26) $\displaystyle{\rm{\bf
G}}_{q,m}(n){\rm{\bf P}}_{q,m}(n|n-1){\rm{\bf G}}^{H}_{q,m}(n),$
where $b=0.99$ is the decay factor used to exponentially weight the history
values.
### III-D Block Diagram
The complete functional diagram of our proposed design is shown in Fig. 1. In
this paradigm, received samples at all RX antennas are first processed in the
interference mitigation module. Then the resultant samples are sequentially
processed in the basic EKF module and noise variance tracking module. For the
ease of description, all components in Fig. 1 are depicted in a parallel
manner. However, in practice these design components can be implemented in a
time-shared manner, and only one single EKF module is physically required.
Figure 1: Block diagram of our proposed design.
## IV Simulations and Discussions
### IV-A Parameters Setup
The OFDM system built in the simulation has a bandwidth of 20MHz and a FFT
size $N_{F}=2048$. Each TX user’s CFO value is independently and randomly
generated within the range of [-2, 2]444Since integer frequency offsets can
generally be estimated during the coarse synchronization stage, CFO value at
fine synchronization stage is usually between -0.5 and 0.5. But here we use
range 2 to demonstrate our design’s estimation performance.. Wireless channels
are generated using the SUI-3 channel model [8], which has $L_{\rm max}=3$
non-zero paths at the time domain. SNR in this paper is defined as the ratio
of the signal power to the noise power at one RX antenna, i.e., ${\rm
SNR}=\sigma_{R}^{2}/\sigma_{Z}^{2}$ where $\sigma_{R}^{2}$ is the total
received signal power at one RX antenna that is coming from all TX users, and
$\sigma_{Z}^{2}$ is the power of the background noise. In the simulation, CFO
and channel parameters are estimated using one OFDM training symbol with
$N_{F}=2048$ samples. Estimation results are validated via the mean square
error (MSE) performance. Specifically, MSE for channel estimation is defined
as a normalized version:
$\displaystyle{\rm MSE}(h_{q,m})=\frac{\sum_{l=1}^{L_{\rm
max}}\left|{\widehat{h}}_{q,m}(p_{l}^{q,m})-h_{q,m}(p_{l}^{q,m})\right|^{2}}{\sum_{l=1}^{L_{\rm
max}}\left|h_{q,m}(p_{l}^{q,m})\right|^{2}}.$ (27)
Cramer-Rao bounds for the MSE results of CFO and channel estimation can be
derived according to [9] as:
$\displaystyle{\rm CRB}_{\rm CFO}({\rm SNR})=\frac{3Q}{2\pi^{2}\cdot
N_{F}\cdot{\rm SNR}\cdot N_{A}},$ (28) $\displaystyle{\rm CRB}_{\rm Chan}({\rm
SNR})=\left(\frac{L_{\rm max}}{N_{F}}+\frac{3}{2N_{F}}\right)\frac{Q}{{\rm
SNR}}.$ (29)
TABLE I: Complexity Comparison (Number of Complex Multiplications) Design Name | Number of Complex Multiplications
---|---
Proposed Design | $\approx L_{\rm max}^{3}+10L_{\rm max}^{2}+14L_{\rm max}+2$
Full-State EKF | $\mathcal{O}\left\\{Q^{3}(L_{\rm max}N_{A}+1)^{3}\right\\}$
FFT-Block EKF [6] | | $\mathcal{O}\\{N_{A}N_{F}(QN_{A})^{2}(L_{\rm max}+1)^{2}$
---
$+QN_{A}(N_{A}N_{F})^{2}(L_{\rm max}+1)^{2}$
$+(QN_{A})^{3}(L_{\rm max}+1)^{3}\\}$
Parallel EKF [7] | | $\mathcal{O}\\{Q^{2}N_{F}(L_{\rm max}+1)^{2}+QN_{F}^{2}(L_{\rm max}+1)^{2}$
---
$+QN_{F}(L_{\rm max}+1)^{3}\\}$
Particle Filter [7] | | $\mathcal{O}\\{N_{P}Q^{2}N_{F}(L_{\rm max}+1)^{2}+N_{P}QN_{F}^{2}(L_{\rm max}+1)^{2}$
---
$+N_{P}QN_{F}(L_{\rm max}+1)^{3}\\}$
$N_{P}$ is the number of particle samples.
EM method [5] | | $\mathcal{O}\\{N_{L}QN_{F}^{2}L_{\rm max}+N_{L}QN_{F}L_{\rm max}^{2}\\}$
---
$N_{L}$ is the number of iterations.
### IV-B Simulation Results
We consider a cognitive network with one primary link and three secondary
links (a total of 4 links). We investigate the CFO and channel estimations at
one secondary RX user with $N_{A}=4$ RX antennas. Without loss of generality,
we assume that this secondary RX user has the same received power from all TX
users. We plot the MSE results of CFO and channel estimation in Fig. 2 and
Fig. 3, respectively. The results show that our proposed design can closely
approach the Cramer-Rao bounds. In addition, we repeat our simulation by
disabling interference mitigation module, or noise variance tracking module.
The corresponding results (Fig. 2 and Fig. 3) indicate that without
interference mitigation, the estimation performance can be dramatically
degraded. And without noise variance tracking module, there could be an error
floor at high SNR values because of the inaccurate tracking of noise variance
information. Using the values in Fig. 2 and Fig. 3, it is feasible to further
investigate the BER/PER performance. But such discussions heavily depend on
the designed receiver structure, which is omitted here for page limitation.
### IV-C Delay, Buffer Size and Complexity Analysis
This subsection evaluates the issues of complexity, delay and buffer size in
the considered designs. We first look at the complexity issue. In particular,
we count the number of complex multiplications involved in our proposed EKF
design, which is listed in Table I. And for comparison, in that table, we also
list the complexity results of Full-State EKF, FFT-Block EKF [6], Parallel EKF
[7], Particle filter [7], and EM method [5]. Here Full-State EKF represents
the EKF that builds all $(L_{\rm max}N_{A}+1)Q$ states in (2) into one state
vector, yielding high state dimension. We can see that our proposed design
enjoys the lowest computation complexity, which is only at the order of
$L_{\rm max}^{3}$. But Full-State EKF’s complexity is around $Q^{3}N_{A}^{3}$
higher than our design. Moreover, FFT-Block EKF, Parallel EKF, Particle
Filter, and EM method’s complexities555Particle filters in [7] and EM designs
in [5] have even higher complexity because of either the number of particle
samples, or the number of iterations. all rely on FFT size $N_{F}$, which is
significantly large in our case ($N_{F}=2048$).
Now we further look at the delay and buffer size in the proposed design. Since
our EKF scheme updates the Kalman estimate at each RX sample (not at each FFT
block) in an online and real-time manner, it has low estimation delay and
requires low buffer size. However, Particle Filter [7], Parallel EKF [7], and
EM approach [5] all operate at an FFT-block basis with buffer size
$N_{F}=2048$ samples, resulting in both a large delay and a large buffer size.
Even worse, particle filter and EM method both need to process the FFT-block
multiple times (e.g., particle samples in particle filter, and iterations in
EM method), leading to additional estimation delay.
## V Conclusion
This paper has presented a low-delay and low-complexity EKF design that can
estimate the CFO and channel parameters in multi-user cognitive
communications. We first present a fundamental EKF design that works at one TX
user and one RX antenna. Then this basic EKF design is reused in a time-shared
way to estimate the parameters for multiple TX users at multiple RX antennas.
Besides, an interference mitigation strategy is proposed to estimate and
cancel the multi-user interference at each RX sample. Moreover, adaptive noise
variance tracking module is further employed to further enhance the estimation
performance. Compared with existing related designs (FFT-Block EKF [6],
Parallel EKF [7], Particle filter [7], and EM method [5]), our proposed design
enjoys low computation complexity (because of pipelining and time-sharing
design), low delay and low buffer size (due to its online and run-time
estimation). Besides, its estimation performance can closely approach the
Cramer-Rao bound.
## References
* [1] R. van Nee and R. Prasad, _OFDM for Wireless Multimedia Communications_. Artech House, Inc., 2000\.
* [2] Z. Cao, U. Tureli, and Y.-D. Yao, “Deterministic multiuser carrier-frequency offset estimation for interleaved OFDMA uplink,” _IEEE Transactions on Communications_ , vol. 52, no. 9, pp. 1585–1594, Sept. 2004.
* [3] Y. Ma and R. Tafazolli, “Channel estimation for OFDMA uplink: a hybrid of linear and BEM interpolation approach,” _IEEE Transactions on Signal Processing_ , vol. 55, no. 4, pp. 1568–1573, Apr. 2007.
* [4] M.-O. Pun, M. Morelli, and C.-C. J. Kuo, “Maximum-likelihood synchronization and channel estimation for OFDMA uplink transmissions,” _IEEE Transactions on Communications_ , vol. 54, no. 4, pp. 726–736, Apr. 2006.
* [5] M.-O. Pun, S.-H. Tsai, and C.-C. J. Kuo, “An EM-based joint maximum likelihood estimation of carrier frequency offset and channel for uplink OFDMA systems,” in _IEEE 60th VTC_ , vol. 1, 26-29 2004, pp. 598–602.
* [6] T. Roman, M. Enescu, and V. Koivunen, “Joint time-domain tracking of channel and frequency offsets for MIMO OFDM systems,” _Wireless Personal Communications_ , vol. 31, pp. 181–200, 2004.
* [7] K. J. Kim, M.-O. Pun, and R. A. Iltis, “Joint carrier frequency offset and channel estimation for uplink MIMO-OFDMA systems using Parallel Schmidt Rao-Blackwellized particle filters,” _Communications, IEEE Transactions on_ , vol. 58, no. 9, pp. 2697 –2708, sep. 2010.
* [8] _Channel models for fixed wireless applications_ , IEEE 802.16 Broadband Wireless Access Working Group IEEE 802.16a-03/01, 2003.
* [9] P. Stoica and O. Besson, “Training sequence design for frequency offset and frequency-selective channel estimation,” _IEEE Transactions on Communications_ , vol. 51, no. 11, pp. 1910–1917, Nov. 2003.
Figure 2: CFO estimation’s MSE results under different SNR values. Figure 3:
Channel estimation’s MSE results under different SNR values.
|
arxiv-papers
| 2011-11-16T19:54:22 |
2024-09-04T02:49:24.414652
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Pengkai Zhao, Cong Shen",
"submitter": "Pengkai Zhao",
"url": "https://arxiv.org/abs/1111.3925"
}
|
1111.4103
|
# Quantum correlations by four–wave mixing in an atomic vapor in a non-
amplifying regime: a quantum beam splitter for photons.
Quentin Glorieux quentin.glorieux@univ-paris-diderot.fr Luca Guidoni Samuel
Guibal Jean-Pierre Likforman Thomas Coudreau Univ Paris Diderot, Sorbonne
Paris Cité, Laboratoire Matériaux et Phénomènes Quantiques, UMR 7162 CNRS,
F-75205 Paris, France
###### Abstract
We study the generation of intensity quantum correlations using four–wave
mixing in a rubidium vapor. The absence of cavity in these experiments allows
to deal with several spatial modes simultaneously. In the standard,
amplifying, configuration, we measure relative intensity squeezing up to 9.2
dB below the standard quantum limit. We also theoretically identify and
experimentally demonstrate an original regime where, despite no overall
amplification, quantum correlations are generated. In this regime a four–wave
mixing set–up can therefore play the role of a photonic beam splitter with
non–classical properties, i.e. a device that splits a coherent state input
into two quantum correlated beams.
###### pacs:
42.50.Dv, 42.50.Lc, 42.50.Nn
Non–classical ”intense” beam have been widely studied in a large variety of
contexts including potential applications to quantum information protocols
Braunstein05 ; Sangouard10 , fundamental issues in quantum mechanics such as
entanglement and non–locality Reid09 , quantum imaging QuantumImaging or
enhancement of the sensitivity of gravitational wave interferometers Caves81 .
Quantum correlated beams are usually obtained through optical non–linear
effects described as $\chi^{(2)}$ or $\chi^{(3)}$ non linearities present in a
variety of media (see BachorBook for a review). In this paper, we study the
generation of quantum correlation by using four–wave mixing (4WM) in a hot
atomic vapor.
Based on $\chi^{(3)}$ non linearity, 4WM is well known to generate intense non
classical beams Yuen79 ; Slusher1985 ; Shelby86 ; Maeda87 . However, over the
last 20 years, attention has been focused mainly on $\chi^{(2)}$ media Wu87 ;
Heidmann87 ; Laurat05 ; Mehmet10 mainly because of their low losses
(availability of high quality optical crystals). On the contrary, in hot
vapors the presence of an atomic resonance enhances the non-linearity but also
usually increases the losses. Yet, recently, it was shown that non–degenerate
4WM in atomic vapors can produce very large amounts of quantum correlations
between intense beams McCormick2007 ; McCormick08 ; GlorieuxSPIE10 . Such a
set–up has a significant advantage over $\chi^{(2)}$ media in that it does not
require an optical cavity to enhance the nonlinearity and the related quantum
effects. This is particularly important in the case of quantum imaging where
spatially multimode quantum effects are involved QuantumImaging ; Boyer08 .
Furthermore, the generated beams directly match the atomic resonance frequency
of an atom–based quantum memory, a key requirement for quantum communications
Sangouard10 .
As noted above, the large nonlinear and quantum effects observed in 4WM
originate from the presence of an atomic resonance. This resonance also
induces incoherent effects, most notably absorption and spontaneous emission,
which, in general, decrease the degree of quantum correlations. These possible
drawbacks are often reduced by increasing the detuning from resonance,
resulting in an overall amplification of the probe and conjugate beams.
However, as we show, there exists a regime where quantum correlations can be
observed despite the fact that the probe beam is de-amplified by propagation
through the atomic vapor. In this regime, a 4WM setup then behaves as a
beam–splitter separating an incoming beam in two different beams, without
overall amplification.
However, when the input beam is in a coherent state, the two output states are
quantum correlated : we thus call this new device a quantum beam–splitter for
photons. On the two input ports of the device are respectively sent the vacuum
state and a coherent state and through the two output ports are emitted the
quantum correlated states. This denomination omits the role of the pump which
is crucial in this scheme as naturally no classical beam splitter can generate
non classical states starting from coherent states. The simplest way to model
theoretically such a device is to chain an ideal linear phase–insensitive
amplifier with a partially transmitting medium. Despite the introduction of
large losses, up to a level that cancels the gain, we show that quantum
correlated beams can be generated in such a configuration. We then introduce
the gemellity Treps05 , a criterion well adapted to describe experiments with
non balanced beams.
Figure 1: (Color Online) a) 4WM in hot atomic vapor schematic setup. b)
Relevant levels of the Rb D1 line described as a double-$\Lambda$ system.
$\Delta$ is the so called one–photon detuning and $\delta$ the two–photon
detuning.
We demonstrate, using a microscopic modelGlorieuxPRA10 , that 4WM in a hot
atomic vapor can efficiently implement a quantum beam–splitter and we show
that the limit for the maximum gemellity predicted in the linear amplifier
model can be theoretically exceeded in this new regime. Finally we test these
predictions experimentally.
## I 4WM in the amplifying regime
The experiment is based on McCormick2007 and is described in detail in
GlorieuxSPIE10 so that we only recall here its main features. A linearly
polarized intense pump beam, frequency locked near the 85Rb $D_{1}$ line, is
mixed with an orthogonally polarized weak probe beam inside an isotopically
pure cell of length $L$. The relevant levels are shown in Fig. 1, a. At the
output of the cell, due to 4WM, the probe beam is amplified and a conjugate
beam is generated (see Fig. 1, b). After filtering out the pump beam with a
polarizing beam splitter, intensity correlations between the probe and
conjugate beams are measured by a pair of high quantum-efficiency photodiodes
coupled to a spectrum analyzer.
A high gain can be observed for a relatively large set of experimental
parameters. The use of a heated cell yields a large number of atoms: for a
temperature $T$ ranging from 100∘C to 150∘C, the atomic density $\mathcal{N}$
calculated from the Clausius–Clapeyron formula Alcock84 varies from 6$\times
10^{12}~{}$cm-3 to $10^{14}~{}$cm-3. Thus, the equivalent optical depth,
$\mathcal{N}\sigma L$ varies between $5\times 10^{3}$ and $10^{5}$ where
$\sigma$ is the atomic cross section for the 5$S^{1/2}$ $\to$ 5$P^{1/2}$
transition in 85Rb. These atoms interact with beams close to resonance : the
single photon detuning $\Delta$ is typically 1 GHz (on the order of the
Doppler broadening) while the two–photon detuning $\delta$ is smaller than 10
MHz.
Within these domains of parameters, explored systematically GlorieuxSPIE10 ,
we have identified an optimal noise reduction regime. For $\Delta=+750$ MHz,
$\delta=+6$ MHz, $T=118^{\circ}$C, $P_{pump}=1200$ mW (corresponding to a Rabi
frequency $\Omega=1$ GHz), gain on the incoming probe beam up to 20 can be
observed. In these conditions, Fig. 2 shows the noise power of the intensity
difference of the probe and conjugate as a function of the analysis frequency
after correcting for the electronic noise:significant noise reduction is
observed in the range 500 kHz to 5 MHz and with a maximal noise reduction of
9.2 dB$\pm$0.5 dB below the SQL is observed between 1 and 2 MHz. This value is
slightly larger than the best results obtained to date with 4WM McCormick08
and very close to those obtained with OPOs Laurat05 . The matching of the
atomic resonance of Rb turns this setup into an ideal source of non-classical
light to interact with Rb vapor quantum memory PL1 ; PL2 .
Figure 2: (Color Online) Noise power of the intensity difference between the
probe and conjugate beams as a function of the frequency after correcting for
the electronic noise. A reduction of 9.2 dB$\pm$0.5 dB below the SQL is
reached at 1 MHz analysis frequency.
## II Quantum beam splitter regime
In the previously described regime,the larger was the gain, the larger were
the quantum correlations. However, we will show here that this not a necessary
condition and that one can, somewhat counter-intuitively, observe significant
quantum correlations in the absence of overall gain.
### II.1 Ideal linear amplifier model
In an ideal phase–insensitive amplifier, an input probe beam is amplified
while a conjugate beam is generated. At the output of the amplifier,
neglecting the contribution of the noise to the average number of photons, the
probe beam has an intensity $GI_{0}$ and the conjugate beam has an intensity
$(G-1)I_{0}$ where $G$ denotes the gain and $I_{0}$ the input probe beam
intensity. Taking into account the ideal character of the amplifier, no noise
is added and the intensity difference at the output has a noise ratio
$1/(2G-1)$ with respect to the input. For probe and conjugate at the input
respectively coherent and vacuum states, this noise ratio is equal to the
quantum correlations at the output of the amplifier. If we now extend this
model by including losses at the output of the medium on the probe and/or
conjugate beams, one would expect a reduction in these correlations as it is
well known that losses are detrimental to squeezing. Let us recall that this
is not always the case as the beams intensity are not balanced: a small amount
of extra losses on the probe beam will tend to make the two beams more
balanced and thus improve the noise reduction on the intensity difference as
noted e.g. in Jasperse11 .
Contrary to the case of Optical Parametric Oscillators above threshold
Laurat05 , 4WM naturally generates unbalanced beams. Unbalanced beams may
exhibit strong quantum correlations but the measurement of the noise on the
intensity difference is not an ideal criterion in this case. It is useful to
introduce the gemellity $\mathcal{G}$ Laurat03 ; Laurat04 ; Treps05 ;
Laurat05b defined by :
$\mathcal{G}=\frac{F_{a}+F_{b}}{2}-\sqrt{C_{ab}^{2}F_{a}F_{b}+\left(\frac{F_{a}-Fb}{2}\right)^{2}},$
(1)
where $F_{i}=\langle\hat{X}_{i}\hat{X}_{i}\rangle$ with $i$ used for $a$
(probe) and $b$ (conjugate),
$C_{ab}=\frac{\langle\hat{X}_{a}\hat{X}_{b}\rangle}{\sqrt{F_{a}F_{b}}}$ and
$\hat{X}_{i}$ is the amplitude quadrature of the related field as defined in
Treps05 . In case of balanced beams the gemellity is equal to the normalized
noise on the difference between the fluctuations of the two measurements :
$\mathcal{G}=\frac{\langle(\hat{X}_{i}-\hat{X}_{j})^{2}\rangle}{2}$, which is
the quantitative measure of the maximal “non–classicality” that can be
extracted from the correlated beams Treps05 . For balanced beams, such as the
ones produced in the limit of infinite gain, this value is equal to the
standard criterion, namely the intensity noise difference. In the conditions
above where the intensity difference noise is -9.2 dB, the noise on the
individual beams are $F_{a}=F_{b}$ = +12 dB at 1 MHz yielding a gemellity
$\mathcal{G}$ = -9.8 dB $\pm$ 0.5 dB. This value is comparable with record
values measured with an OPO above threshold Laurat05 and moreover a large
number of spatial modes (estimated to 100 in this particular configuration)
are squeezed simultaneously Boyer08 .
Using this criteria and introducing losses on the probe ($T_{a}$) and
conjugate ($T_{b}$) beams so that the overall transmission is equal to one
($T_{a}G+T_{b}(G-1)=1$), it is straightforward to show that there always
exists a region in the parameter space where a gemellity lower that one is
expected. To our knowledge, this phenomenon, albeit simple, has neither been
discussed nor observed. The larger quantum correlations reachable with no
overall amplification corresponds to the situation of a gain $G=1.23$, a
transmission of 0.62 on the probe beam and perfect transmission on the
conjugate beam. This configuration gives the limit for the gemellity reachable
by this simple model : $\mathcal{G}=-2.8$ dB.
### II.2 Microscopic model
To investigate further this effect, we have studied the 4WM process using a
microscopic model based on the cold–atom model described extensively in
GlorieuxPRA10 . This model assumes the simplified double–$\Lambda$ level
structure of Fig. 1, right. The Heisenberg–Langevin approach is used to obtain
the relevant classical quantities (probe gain $G_{a}$, conjugate gain $G_{b}$
defined with respect to $I_{0}$) as well as the quantum properties of the
output beams. In particular, it is possible to calculate noise spectra that
allow for quantifying quantum correlations both in terms of
intensity–difference noise $S_{N_{-}}$ and for the unbalanced case in terms of
gemellity $\mathcal{G}$. In the regime of high amplification previously
described, this model is in good quantitative agreement with the measured
correlations GlorieuxPhD . Exploring the parameters space in this model, we
have found a new region where the 4WM process generates quantum correlations
in the absence of overall amplification. This regime is therefore very similar
to the linear amplifier model followed by a lossy medium described above.
Nevertheless, the microscopic model predicts that in this regime, the
gemellity can be significantly enhanced in contrast to the linear model and
exceeds the -2.8 dB limit discussed previously.
Let us start by presenting the classical behavior of the probe and conjugate
beams in the region of interest of parameter space (theoretical data are
compared to the experimental results). We plot in Fig. 3, the gain for the two
fields as a function the two–photon detuning $\delta$. The main difference
with respect to the high gain parameter region is the choice of the atomic
density (experimentally driven by the temperature). The large gain results of
Fig. 2 were obtained for a temperature of 118 ∘C while the curves in Fig. 3
are obtained for $T=95^{\circ}$C. The approximately one order of magnitude
lower optical density, together with the different choice of $\delta$ and
$\Delta$, explain the drastic reduction of $G_{a}$ and $G_{b}$.
”Beam–splitter” regime is obtained near the two-photon resonance, where
$G_{a}$ goes to zero due to a Raman process involving a probe and a pump
photonGlorieuxPRA10 . Due to the pump–induced AC-Stark shift, this two–photon
resonance is shifted to negative values of $\delta$ and its exact position
depends on the one-photon detuning $\Delta$ and on the pump Rabi frequency
$\Omega$. Within a very narrow region of parameter space, the sum of the two
beams output intensities becomes slightly smaller or almost equal to the input
probe intensity. It is interesting to note that for potential applications
this very narrow feature could be considered as a limitation. Notwithstanding
we have verified numerically, by changing simultaneously $\Delta$, $\Omega$
and the optical depth $\mathcal{N}\alpha L$, that the detuning for which this
system exhibits the behavior of a quantum beam splitter can be tuned over more
than 100 MHz. As already remarked in GlorieuxPRA10 , we note that, despite the
fact that the model is based on a cold atom sample, it yields without any
adjustable parameter a qualitative agreement with the experimental data
obtained in a hot vapor.
Figure 3: (Color Online) Theoretically predicted (left) and experimentally
measured (right) gain for the probe beam ($G_{a}$, ) and conjugate beam
($G_{b}$, black) as a function of the two–photon detuning $\delta$. The
parameters used in the simulations are : , optical depth $\mathcal{N}\alpha
L=500$, Pump Rabi frequency $\Omega$ = 0.42 GHz, single photon detuning
$\Delta/2\pi$=0.8 GHz. Measured parameters are : pump power $P$=0.6 W
($\Omega/2\pi=0.4$ GHz), T=95∘C, single photon detuning $\Delta/2\pi$=0.8 GHz.
Figure 4: (Color Online) Quantum intensity correlations between the probe and
conjugate beams as a function of the analysis frequency with the same
parameters as above and $\delta/2\pi$=-52 MHz.
### II.3 Demonstration of the quantum beam splitter
Motivated by these theoretical predictions, we have experimentally
investigated this original regime. We plot in Fig. 4 the experimentally
measured intensity difference noise as a function of the analysis frequency
$\omega$. We observe significant quantum correlations, down to 1.0$\pm 0.2$ dB
below the SQL around an analysis frequency of 1 MHz. At the same time, the
power of the two beams normalized to the probe input power is measured to be
0.65 and 0.35 for the probe and conjugate respectively. This demonstrates
clearly the behavior of a quantum beam splitter for photons where one laser
beam is split into two beams without gain but generating quantum correlations.
We note that the measured noise reduction is slightly smaller than the one
predicted theoretically (Fig. 4): this discrepancy can be attributed to the
fact that the model is based on a cold atomic sample, far from the
experimental regime.
In this situation, $\mathcal{G}$ can be calculated to compare it to the
theoretical limit of the linear amplifier model. By measuring the noise on the
two individuals beams respectively equal to + 3 and + 2 dB for probe and
conjugate, we obtain a value of the gemellity equal to $\mathcal{G}=-1.8\pm
0.5$ dB. This value does not exceed the maximum limit of -2.8 dB predicted by
the linear amplifier model. As previously noted, the theoretical model does
not take into account the velocity distribution of the atoms, thus time
transit effects and Doppler broadening which are expected to play a
detrimental role. This can explain why the linear amplifier model limit cannot
be reached in this configuration whereas the microscopic model predicts that
gemellities better than $\mathcal{G}$=-3.2 dB can be obtained with the above
parameters and an optical depth, $\mathcal{N}\sigma L=1500$.
We have thus shown firstly that generating quantum correlations does not
require overall amplification and, secondly, that the ideal linear amplifier
is not the ideal device to perform this operation but that 4WM in atomic
vapours presents an interesting avenue in this context. This setup could also
be used as the input beam splitter introducing quantum correlations for an
original version of the Mach-Zender interferometer as well as in a so-called
SU(1,1) interferometer Yurke86 .
## III Conclusion
We have studied the production of quantum correlated beams in four–wave mixing
in a 85Rb cell. First, we have identified and experimentally realized an
optimal regime in the high gain region where intensity–difference noise down
to -9.2 dB below the standard quantum limit (gemellity
$\mathcal{G}=-9.8~{}dB$) have been measured. This result is important in the
domain of quantum communications where both large non–classical effects and
the availability of an atom–based storage media form strong requirements
Sangouard10 ; PL1 ; PL2 .
We have also predicted and observed an original regime where quantum
correlations are present despite significant losses on the probe beam. This
regime is of particular interest, because it can occur in a situation in which
the sum of the two output beam intensities is smaller or equal to the input
probe intensity. Therefore the atomic medium controlled by the pump laser acts
like a beam splitter device that creates quantum correlations (quantum beam
splitter). Although this effect could in principle be observed with an ideal
amplifier, it is to our knowledge the first demonstration of it. In this
context, we have discussed the use of the gemellity criterion more appropriate
in the case of unbalanced beams produced by four–wave mixing. Finally, a
microscopic model allowed us to demonstrate that 4WM in the quantum beam
splitter regime can beat theoretically the limit of quantum correlations
predicted by the model of a linear amplifier followed by a lossy medium.
In our experiment, with a hot atomic vapor, a value of $\mathcal{G}=-1.8\pm
0.5$ dB has been reported. Although the parameter values required to beat the
linear amplifier model limit are presently beyond reach of experiments
performed with cold atoms, our model provides an interesting avenue to surpass
this limit using hot or cold atoms.
## Acknowledgements
We thank E. Arimondo, R. Dubessy and P. Lett for fruitful discussions. We
thank N. Treps for the loan of high–efficiency photodiodes.
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|
arxiv-papers
| 2011-11-17T14:08:45 |
2024-09-04T02:49:24.424992
|
{
"license": "Public Domain",
"authors": "Quentin Glorieux, Luca Guidoni, Samuel Guibal, Jean-Pierre Likforman,\n and Thomas Coudreau",
"submitter": "Quentin Glorieux",
"url": "https://arxiv.org/abs/1111.4103"
}
|
1111.4108
|
Jordan product determined points in matrix algebras
Wenlei Yang 111E-mail address: yang121045760@yahoo.cn, Jun Zhu222E-mail
address: zhu$\\_$gjun@yahoo.com.cn
Institute of Mathematics, Hangzhou Dianzi University, Hangzhou 310018,
People’s Republic of China
Abstract
Let $M_{n}(R)$ be the algebra of all $n\times n$ matrices over a unital
commutative ring $R$ with 6 invertible. We say that $A\in M_{n}(R)$ is a
Jordan product determined point if for every $R$-module $X$ and every
symmetric $R$-bilinear map $\\{\cdot,\cdot\\}$ : $M_{n}(R)\times M_{n}(R)\to
X$ the following two conditions are equivalent: (i) there exists a fixed
element $w\in X$ such that $\\{x,y\\}=w$ whenever $x\circ y=A$, $x,y\in
M_{n}(R)$; (ii) there exists an $R$-linear map $T:M_{n}(R)^{2}\to X$ such that
$\\{x,y\\}=T(x\circ y)$ for all $x,y\in M_{n}(R)$. In this paper, we mainly
prove that all the matrix units are the Jordan product determined points in
$M_{n}(R)$ when $n\geq 3$. In addition, we get some corollaries by applying
the main results.
AMS Classification: 15A04 15A27
Keywords : Jordan product determined point; matrix algebra; Jordan all-
multiplicative point; Jordan all-derivable point
## 1\. Introduction
In this paper, we will mainly discuss Jordan product determined points on
matrix algebras. Before proceeding let us fix some symbols and notations in
this paper. Let $M_{n}(R)$ be the algebra of all $n\times n$ matrices over a
unital commutative ring $R$ with 6 invertible. Matrix units are denoted by
$e_{ij}$ and the Jordan product “$\circ$” is defined as $x\circ y=xy+yx$. The
identity matrix is denoted by $I$. By $M_{n}(R)^{2}$ we denote the $R$-linear
span of all elements of the form $xy$ where $x,y\in M_{n}(R)$.
The concept of zero product (resp. Jordan product, Lie product) determined
algebras was introduced by Brešar et al. [5]. According to [5], $M_{n}(B)$
($n\geq 2$) is zero product determined where $B$ is a unital algebra. If $B$
is a unital algebra with 2 invertible, then $M_{n}(B)$ ($n\geq 3$) is zero
Jordan product determined. From the results above, we can study the linear
maps preserving zero product (resp. Jordan product) and derivable (resp.
Jordan derivable) at zero point respectively. Wang et al. [1, 2] showed that
(1) if a symmetric bilinear map $\\{\cdot,\cdot\\}$ : $M_{n}(R)\times
M_{n}(R)\to X$ satisfies the condition that $\\{u,u\\}=\\{e,u\\}$ whenever
$u^{2}=u$, then there exists a linear map $f$ from $M_{n}(R)$ to $X$ such that
$\\{x,y\\}=f(x\circ y)$ for all $x,y\in M_{n}(R)$; and (2) if an invertible
linear map $\delta$ on $M_{n}(R)$ preserves identity-product, then it is a
Jordan automorphism; and a linear map $\sigma$ on $M_{n}(R)$ is derivable at
the identity matrix if and only if it is an inner derivation. Zhu et al. [3]
showed that for every $G\in M_{n}$, det$G=0$, is an all-multiplicative point
in $M_{n}$. Gong and Zhu [6] considered the case of Jordan all-multiplicative
point in $M_{n}$. Zhu et al. [4] showed that a matrix $G$ is an all-derivable
point in $M_{n}$ if and only if $G\neq 0$. Zhao et al. [7] showed that every
element of the algebra of all upper triangular matrices is a Jordan all-
derivable point.
Motivated by the concepts and results above, we will consider Jordan product
determined points in matrix algebras. For $A\in M_{n}(R)$, we say that $A$ is
a Jordan product determined point if for every $R$-module $X$ and every
symmetric $R$-bilinear map $\\{\cdot,\cdot\\}$ : $M_{n}(R)\times M_{n}(R)\to
X$ the following two conditions are equivalent: (i) there exists a fixed
element $w\in X$ such that $\\{x,y\\}=w$ whenever $x\circ y=A$, $x,y\in
M_{n}(R)$; (ii) there exists an $R$-linear map $T:M_{n}(R)^{2}\to X$ such that
$\\{x,y\\}=T(x\circ y)$ for all $x,y\in M_{n}(R)$. We say that $G\in M_{n}(R)$
is a Jordan all-multiplicative point in $M_{n}(R)$ if for every
$M_{n}(R)$-module $X$ and every Jordan multiplicative $R$-linear map $\varphi$
: $M_{n}(R)\to X$ at $G$ (i.e. $\varphi(S\circ T)=\varphi(S)\circ\varphi(T)$
for any $S,T\in M_{n}(R)$, $S\circ T=G$) with $\varphi(I)=I$ is a
multiplicative mapping in $M_{n}(R)$. We say that $H\in M_{n}(R)$ is a Jordan
all-derivable point in $M_{n}(R)$ if for every $M_{n}(R)$-module $X$ and every
Jordan derivable $R$-linear map $\varphi$ : $M_{n}(R)\to X$ at $H$ (i.e.
$\varphi(S\circ T)=\varphi(S)\circ T+S\circ\varphi(T)$ for any $S,T\in
M_{n}(R)$, $S\circ T=H$) with $\varphi(I)=0$ is a Jordan derivation in
$M_{n}(R)$. The above two definitions are somewhat different from [6] and [7].
In this paper, we will prove that every matrix unit $e_{ij}$ is a Jordan
product determined point in $M_{n}(R)$ when $n\geq 3$. As an application of
the result above, we will show that every matrix unit $e_{ij}$ is a Jordan
all-multiplicative point and a Jordan all-derivable point respectively.
This paper is organized as follows. Section 2 concerns Jordan product
determined points in $M_{n}(R)$, and we obtain the major results Theorem 2.2
and Theorem 2.3 in this paper. In Section 3, we get some corollaries by
applying the main results in section 2.
## 2\. Jordan product determined points in $M_{n}(R)$
According to [5], we give the lemma below.
Lemma 2.1. For $A\in M_{n}(R)$, $A$ is a Jordan product determined point if
and only if for every $R$-module $X$ and every symmetric $R$-bilinear map
$\\{\cdot,\cdot\\}$ satisfy the condition (i), the following condition (iii)
holds true.
(iii) for every $x_{t},y_{t}\in M_{n}(R)$ with $\sum_{t=1}^{l}x_{t}\circ
y_{t}=0$, $t=1,2,\dots,l$, $\sum_{t=1}^{l}\\{x_{t},y_{t}\\}=0$ holds true.
Proof. Obviously, the “only if” part holds true.
Conversely, if the condition (iii) holds true, we can define $R$-linear map
$T$ : $M_{n}(R)^{2}\to X$ as $T(\sum_{t}x_{t}\circ
y_{t})=\sum_{t}\\{x_{t},y_{t}\\}$ according to [5]. Then $T$ satisfies
condition (ii). We only need to prove that $T$ is well defined. Indeed, if
condition (iii) is fulfilled, $T$ is well defined obviously. Hence $A$ is a
Jordan product determined point. $\Box$
Theorem 2.2. $e_{ss}$, $s\in\\{1,2,\dots,n\\}$, is a Jordan product determined
point in $M_{n}(R)$ when $n\geq 3$.
Proof. Let $s$ be a fixed number, $X$ be an $R$-module, $\\{\cdot,\cdot\\}$ :
$M_{n}(R)\times M_{n}(R)\to X$ be a symmetric (i.e. $\\{x,y\\}=\\{y,x\\}$)
$R$-bilinear map. Now we assume that there exists a fixed element $w\in X$
such that $\\{x,y\\}=w$ whenever $x\circ y=e_{ss}$, $x,y\in M_{n}(R)$.
Throughout the proof, $i,j,k,m$ will denote arbitrary indices.
We begin by noticing that $(\frac{1}{2}e_{ss})\circ e_{ss}=e_{ss}$ and as
$\\{\cdot,\cdot\\}$ is symmetric, then
$\\{\frac{1}{2}e_{ss},e_{ss}\\}=w=\\{e_{ss},\frac{1}{2}e_{ss}\\}.$ (1)
Next we suppose $n\geq 3$ and divide the proof into three steps.
Step 1. In this step, we assume $i\neq s,j\neq s,k\neq s\;\mathrm{and}\;m\neq
s$.
Case 1.1. $i\neq m\;\mathrm{and}\;j\neq k$.
Since $(\frac{1}{2}e_{ss}+e_{ij})\circ e_{ss}=e_{ss}$ and as
$\\{\cdot,\cdot\\}$ is symmetric, we have that
$\\{e_{ij},e_{ss}\\}=0=\\{e_{ss},e_{ij}\\}.$ (2)
Noting $(\frac{1}{2}e_{ss}+e_{ij})\circ(e_{ss}+e_{km})=e_{ss}$, it follows
$\\{\frac{1}{2}e_{ss}+e_{ij},e_{ss}+e_{km}\\}=w$. As $\\{\cdot,\cdot\\}$ is
symmetric, applying (1) and (2), this yields
$\\{e_{ij},e_{km}\\}=0.$ (3)
Case 1.2. $i,j,k$ are distinct.
From $(\frac{1}{2}e_{ss}+e_{ik})\circ(e_{ss}+e_{kk}-e_{ii})=e_{ss}$, we obtain
that $\\{\frac{1}{2}e_{ss}+e_{ik},e_{ss}+e_{kk}-e_{ii}\\}=w$. Because
$\\{\cdot,\cdot\\}$ is symmetric, it follows from (1) and (2) that
$\\{e_{ik},e_{kk}\\}=\\{e_{ik},e_{ii}\\}=\\{e_{ii},e_{ik}\\}=\\{e_{kk},e_{ik}\\}.$
(4)
Now we assume $n>3$. As
$(\frac{1}{2}e_{ss}+e_{ij}+e_{ik})\circ(e_{ss}+e_{jk}-e_{kk})=e_{ss}$, we
derive $\\{\frac{1}{2}e_{ss}+e_{ij}+e_{ik},e_{ss}+e_{jk}-e_{kk}\\}=w$. Using
(1), (2) and (3), this can be reduced to
$\\{e_{ij},e_{jk}\\}=\\{e_{ik},e_{kk}\\}.$ (5)
(4) together with (5) yield
$\\{e_{ij},e_{jk}\\}=\\{e_{ik},e_{kk}\\}=\\{e_{ik},e_{ii}\\}=\\{e_{ii},e_{ik}\\}.$
(6)
Since $\\{\cdot,\cdot\\}$ is symmetric, it follows that
$\\{e_{ii},e_{ik}\\}=\\{e_{ik},e_{ii}\\}=\\{e_{kk},e_{ik}\\}=\\{e_{jk},e_{ij}\\}.$
(7)
Case 1.3. $i\neq j$.
By
$(\frac{1}{2}e_{ss}+\frac{1}{2}e_{ii}+e_{ij}-\frac{1}{2}e_{jj})\circ(e_{ss}+e_{ji}-e_{ii}+e_{jj})=e_{ss}$,
it is clear that
$\\{\frac{1}{2}e_{ss}+\frac{1}{2}e_{ii}+e_{ij}-\frac{1}{2}e_{jj},e_{ss}+e_{ji}-e_{ii}+e_{jj}\\}=w$.
Then we can get from (1), (2), (3) and (4) that
$\\{e_{ij},e_{ji}\\}=\frac{1}{2}\\{e_{ii},e_{ii}\\}+\frac{1}{2}\\{e_{jj},e_{jj}\\}.$
(8)
Step 2. In this step, we consider some of the indices of the matrix units
equal to $s$.
Case 2.1. $i\neq s,j\neq s,k\neq s\;\mathrm{and}\;m\neq s$.
For $(\frac{1}{2}e_{ss}-e_{si})\circ(e_{ss}-e_{ii})=e_{ss}$, we have that
$\\{\frac{1}{2}e_{ss}-e_{si},e_{ss}-e_{ii}\\}=w$. Applying (1) and (2), it can
be reduced to $\\{e_{si},e_{ss}\\}=\\{e_{si},e_{ii}\\}$. As
$\\{\cdot,\cdot\\}$ is symmetric, then it follows that
$\\{e_{si},e_{ss}\\}=\\{e_{si},e_{ii}\\}=\\{e_{ss},e_{si}\\}=\\{e_{ii},e_{si}\\}.$
(9)
Since $(\frac{1}{2}e_{ss}-e_{is})\circ(e_{ss}-e_{ii})=e_{ss}$, a similar
discussion as above shows that
$\\{e_{is},e_{ss}\\}=\\{e_{is},e_{ii}\\}=\\{e_{ss},e_{is}\\}=\\{e_{ii},e_{is}\\}.$
(10)
Let $j\neq k$, then we have
$(\frac{1}{2}e_{ss}+e_{sj})\circ(e_{ss}+e_{km}-e_{jj})=e_{ss}$ and it implies
that $\\{\frac{1}{2}e_{ss}+e_{sj},e_{ss}+e_{km}-e_{jj}\\}=w$. As
$\\{\cdot,\cdot\\}$ is symmetric and by (1), (2) and (9), this yields
$\\{e_{sj},e_{km}\\}=0=\\{e_{km},e_{sj}\\}\;\mathrm{if}\;j\neq k.$ (11)
From $(\frac{1}{2}e_{ss}+e_{si})\circ(e_{ss}-2e_{si})=e_{ss}$, we derive
$\\{\frac{1}{2}e_{ss}+e_{si},e_{ss}-2e_{si}\\}=w$. Then it follows from (1)
and (9) that
$\\{e_{si},e_{si}\\}=0.$ (12)
Since
$(\frac{1}{2}e_{ss}+e_{sj}-\frac{1}{2}e_{si})\circ(e_{ss}+e_{si}-e_{jj})=e_{ss}$
when $i\neq j$, then we have
$\\{\frac{1}{2}e_{ss}+e_{sj}-\frac{1}{2}e_{si},e_{ss}+e_{si}-e_{jj}\\}=w$.
Using (1), (2), (9), (11) and (12), it is clear that
$\\{e_{sj},e_{si}\\}=0\;\mathrm{if}\;i\neq j.$ (13)
If $j\neq m$, from
$(\frac{1}{2}e_{ss}+e_{js})\circ(e_{ss}+e_{km}-e_{jj})=e_{ss}$ we can get that
$\\{\frac{1}{2}e_{ss}+e_{js},e_{ss}+e_{km}-e_{jj}\\}=w$. As
$\\{\cdot,\cdot\\}$ is symmetric and by (1), (2) and (10), it follows
$\\{e_{js},e_{km}\\}=0=\\{e_{km},e_{js}\\}\;\mathrm{if}\;j\neq m.$ (14)
For $(\frac{1}{2}e_{ss}+e_{is})\circ(e_{ss}-2e_{is})=e_{ss}$, we have that
$\\{\frac{1}{2}e_{ss}+e_{is},e_{ss}-2e_{is}\\}=w$. Applying (1) and (10), it
can be reduced to
$\\{e_{is},e_{is}\\}=0.$ (15)
Assume $i\neq j$, it follows from
$(\frac{1}{2}e_{ss}+e_{js}-\frac{1}{2}e_{is})\circ(e_{ss}+e_{is}-e_{jj})=e_{ss}$
that
$\\{\frac{1}{2}e_{ss}+e_{js}-\frac{1}{2}e_{is},e_{ss}+e_{is}-e_{jj}\\}=w$.
Using (1), (2), (10), (14) and (15), this yields
$\\{e_{js},e_{is}\\}=0\;\mathrm{if}\;i\neq j.$ (16)
Case 2.2. $i\neq s,k\neq s\;\mathrm{and}\;j\neq s$.
Since
$(\frac{1}{2}e_{ss}+e_{ki}-e_{ks}-\frac{1}{2}e_{ii})\circ(e_{ss}+e_{is})=e_{ss}$
if $k\neq i$, then we have that
$\\{\frac{1}{2}e_{ss}+e_{ki}-e_{ks}-\frac{1}{2}e_{ii},e_{ss}+e_{is}\\}=w$. By
(1), (2), (10) and (16), this yields $\\{e_{ki},e_{is}\\}=\\{e_{ks},e_{ss}\\}$
if $k\neq i$. For our purpose, it is more convenient to rewrite this equation
as $\\{e_{ik},e_{ks}\\}=\\{e_{is},e_{ss}\\}$ if $i\neq k$. As
$\\{\cdot,\cdot\\}$ is symmetric, then we can conclude from the above equation
and (10) that
$\\{e_{ik},e_{ks}\\}=\\{e_{is},e_{ss}\\}=\\{e_{ii},e_{is}\\}=\\{e_{is},e_{ii}\\}=\\{e_{ss},e_{is}\\}=\\{e_{ks},e_{ik}\\}\;\mathrm{if}\;i\neq
k.$ (17)
If $i\neq k$, then we have
$(\frac{1}{2}e_{ss}+e_{ik}-e_{sk}-\frac{1}{2}e_{ii})\circ(e_{ss}+e_{si})=e_{ss}$.
By a similar discussion as above, this yields
$\\{e_{sk},e_{ki}\\}=\\{e_{ss},e_{si}\\}=\\{e_{si},e_{ii}\\}=\\{e_{ii},e_{si}\\}=\\{e_{si},e_{ss}\\}=\\{e_{ki},e_{sk}\\}\;\mathrm{if}\;i\neq
k.$ (18)
If $j\neq k$, from
$(\frac{1}{2}e_{ss}+e_{sk}+e_{jk}-\frac{1}{2}e_{js})\circ(e_{ss}+e_{js}-e_{kk})=e_{ss}$
we have that
$\\{\frac{1}{2}e_{ss}+e_{sk}+e_{jk}-\frac{1}{2}e_{js},e_{ss}+e_{js}-e_{kk}\\}=w$.
For $\\{\cdot,\cdot\\}$ is symmetric, applying (1), (2), (14), (15), (17) and
(18), it follows
$\\{e_{sk},e_{js}\\}=\\{e_{jk},e_{kk}\\}=\\{e_{js},e_{sk}\\}\;\mathrm{if}\;j\neq
k.$ (19)
Case 2.3. $i\neq s$.
Since
$(2e_{si}+\frac{1}{2}e_{is}-e_{ii})\circ(e_{si}+\frac{1}{4}e_{is}+\frac{1}{2}e_{ii})=e_{ss}$,
we obtain
$\\{2e_{si}+\frac{1}{2}e_{is}-e_{ii},e_{si}+\frac{1}{4}e_{is}+\frac{1}{2}e_{ii}\\}=w$.
Using (1), (12), (15), (17) and (18), this can be reduced to
$\\{e_{si},e_{is}\\}+\\{e_{is},e_{si}\\}=\\{e_{ss},e_{ss}\\}+\\{e_{ii},e_{ii}\\}$.
As $\\{e_{si},e_{is}\\}=\\{e_{is},e_{si}\\}$, it leads to
$\\{e_{si},e_{is}\\}=\\{e_{is},e_{si}\\}=\frac{1}{2}\\{e_{ss},e_{ss}\\}+\frac{1}{2}\\{e_{ii},e_{ii}\\}.$
(20)
Step 3. Now concluding from case 1.1 and case 2.1, we can obtain
$\\{e_{ij},e_{km}\\}=0\;\mathrm{for\;every}\;i,j,k,m,\;\mathrm{if}\;i\neq
m\;\mathrm{and}\;j\neq k.$ (21)
If $i,j,k$ are distinct and $n=3$, it follows from (4), (17), (18) and (19)
that
$\\{e_{ij},e_{jk}\\}=\\{e_{ik},e_{kk}\\}=\\{e_{ii},e_{ik}\\}=\\{e_{ik},e_{ii}\\}=\\{e_{kk},e_{ik}\\}=\\{e_{jk},e_{ij}\\}.$
(22)
If $i,j,k$ are distinct and $n>3$, from (6), (7), (17), (18) and (19) we have
the equations above as well. So (22) holds true whenever $n\geq 3$.
From case 1.3 and case 2.3, we have
$\\{e_{ij},e_{ji}\\}=\frac{1}{2}\\{e_{ii},e_{ii}\\}+\frac{1}{2}\\{e_{jj},e_{jj}\\}\;\mathrm{if}\;i\neq
j.$ (23)
Let $\sum_{t=1}^{l}x_{t}\circ y_{t}=0$ where $x_{t},y_{t}\in M_{n}(R)$,
$t=1,2,\ldots,l$. We write $x_{t}$ and $y_{t}$ for
$x_{t}=\sum_{i=1}^{n}\sum_{j=1}^{n}a_{ij}^{t}e_{ij},\;y_{t}=\sum_{k=1}^{n}\sum_{m=1}^{n}b_{km}^{t}e_{km}.$
Then for all $i$ and $m$ we have that
$\sum_{t=1}^{l}\sum_{j=1}^{n}(a_{ij}^{t}b_{jm}^{t}+b_{ij}^{t}a_{jm}^{t})=0.$
(24)
Now we will show $\sum_{t=1}^{l}\\{x_{t},y_{t}\\}=0$. Note that
$\displaystyle\sum_{t=1}^{l}\\{x_{t},y_{t}\\}$ $\displaystyle=$
$\displaystyle\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{k=1}^{n}\sum_{m=1}^{n}\\{a_{ij}^{t}e_{ij},b_{km}^{t}e_{km}\\}$
$\displaystyle=$
$\displaystyle\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{k=1}^{n}\sum_{m=1}^{n}a_{ij}^{t}b_{km}^{t}\\{e_{ij},e_{km}\\}.$
According to our assumptions, it follows from (21), (22), (23) and (24) that
$\displaystyle\sum_{t=1}^{l}\\{x_{t},y_{t}\\}$ $\displaystyle=$
$\displaystyle\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{k=1\atop k\neq
j}^{n}\sum_{m=1\atop m\neq
i}^{n}a_{ij}^{t}b_{km}^{t}\\{e_{ij},e_{km}\\}+\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{m=1\atop
m\neq i}^{n}a_{ij}^{t}b_{jm}^{t}\\{e_{ij},e_{jm}\\}$ $\displaystyle+$
$\displaystyle\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{k=1\atop k\neq
j}^{n}a_{ij}^{t}b_{ki}^{t}\\{e_{ij},e_{ki}\\}+\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}a_{ij}^{t}b_{ji}^{t}\\{e_{ij},e_{ji}\\}$
$\displaystyle=$
$\displaystyle\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}\sum_{m=1\atop m\neq
i}^{n}(a_{ij}^{t}b_{jm}^{t}+b_{ij}^{t}a_{jm}^{t})\\{e_{ij},e_{jm}\\}$
$\displaystyle+$
$\displaystyle\frac{1}{2}\sum_{t=1}^{l}\sum_{i=1}^{n}\sum_{j=1}^{n}a_{ij}^{t}b_{ji}^{t}\left(\\{e_{ii},e_{ii}\\}+\\{e_{jj},e_{jj}\\}\right)$
$\displaystyle=$ $\displaystyle\sum_{i=1}^{n}\sum_{m=1\atop m\neq
i}^{n}\left[\sum_{t=1}^{l}\sum_{j=1}^{n}(a_{ij}^{t}b_{jm}^{t}+b_{ij}^{t}a_{jm}^{t})\\{e_{im},e_{mm}\\}\right]$
$\displaystyle+$
$\displaystyle\frac{1}{2}\sum_{i=1}^{n}\left[\sum_{t=1}^{l}\sum_{j=1}^{n}(a_{ij}^{t}b_{ji}^{t}+b_{ij}^{t}a_{ji}^{t})\\{e_{ii},e_{ii}\\}\right]$
$\displaystyle=$ $\displaystyle 0.$
By Lemma 2.1, $e_{ss}$ is a Jordan product determined point. $\Box$
Theorem 2.3. $e_{pq}$, $p\neq q$, is a Jordan product determined point in
$M_{n}(R)$ when $n\geq 3$.
Proof. Let $p,q$ be the distinct fixed indices, $\\{\cdot,\cdot\\}$ :
$M_{n}(R)\times M_{n}(R)\to X$ be a symmetric (i.e. $\\{x,y\\}=\\{y,x\\}$)
$R$-bilinear map where $X$ is an $R$-module. And we assume that there exists a
fixed element $w\in X$ such that $\\{x,y\\}=w$ whenever $x\circ y=e_{pq}$,
$x,y\in M_{n}(R)$. Throughout the proof, $i,j,k,m$ will denote arbitrary
indices.
According to Lemma 2.1 and Step 3 in the proof of Theorem 2.2, we only need to
verify (21), (22) and (23) hold true when $n\geq 3$. Now we suppose $n\geq 3$
and divide the proof into several steps.
Step 1. In this step, we assume $i\neq q\;\mathrm{and}\;m\neq p$.
Since $e_{ps}\circ e_{sq}=e_{pq}\>(s=1,2,\ldots,n)$ and as $\\{\cdot,\cdot\\}$
is symmetric, hence we have
$\\{e_{ps},e_{sq}\\}=w=\\{e_{sq},e_{ps}\\},\;s=1,2,\ldots,n.$ (25)
Choosing $s\neq j\;\mathrm{and}\;s\neq k$, then we obtain
$e_{ps}\circ(e_{sq}+e_{km})=e_{pq}$ and $(e_{ps}+e_{ij})\circ e_{sq}=e_{pq}$
respectively. For $\\{\cdot,\cdot\\}$ is symmetric, applying (25), it follows
$\\{e_{ps},e_{km}\\}=\\{e_{km},e_{ps}\\}=\\{e_{ij},e_{sq}\\}=\\{e_{sq},e_{ij}\\}=0\;\mathrm{if}\>s\neq
j\;\mathrm{and}\;s\neq k.$ (26)
Given $s\neq j\;\mathrm{and}\;s\neq k$, then if $i\neq m\;\mathrm{and}\;j\neq
k$ we have $(e_{ps}+e_{ij})\circ(e_{sq}+e_{km})=e_{pq}$. It is clear that
$\\{e_{ps}+e_{ij},e_{sq}+e_{km}\\}=w$. As $\\{\cdot,\cdot\\}$ is symmetric,
using (25) and (26), this yields
$\\{e_{ij},e_{km}\\}=0=\\{e_{km},e_{ij}\\}\>\mathrm{if}\>i\neq
m\;\mathrm{and}\;j\neq k.$ (27)
Step 2. In this step, we suppose $i\neq q\;\mathrm{and}\;m\neq p$.
Assuming $s\neq i\;\mathrm{and}\;s\neq m$, then if $i\neq
m\;\mathrm{and}\;i\neq p$ we can verify
$(e_{ps}+e_{im})\circ(e_{sq}+e_{mm}-e_{ii})=e_{pq}$. It follows that
$\\{e_{ps}+e_{im},e_{sq}+e_{mm}-e_{ii}\\}=w$. Since $\\{\cdot,\cdot\\}$ is
symmetric, by (25) and (27), this yields
$\\{e_{im},e_{mm}\\}=\\{e_{im},e_{ii}\\}=\\{e_{ii},e_{im}\\}\>\mathrm{if}\>i\neq
m\;\mathrm{and}\;i\neq p.$ (28)
From $e_{pm}\circ(e_{mq}+e_{mm}-e_{pp})=e_{pq}$, we have that
$\\{e_{pm},e_{mq}+e_{mm}-e_{pp}\\}=w$. As $\\{\cdot,\cdot\\}$ is symmetric, it
follows from (25) that
$\\{e_{pm},e_{mm}\\}=\\{e_{pm},e_{pp}\\}=\\{e_{pp},e_{pm}\\}.$ (29)
Choosing $s\neq j\;\mathrm{and}\;s\neq m$, if $i,j,k$ are distinct we obtain
$(e_{ps}+e_{ij}+e_{im})\circ(e_{sq}+e_{jm}-e_{mm})=e_{pq}$. Then it leads to
$\\{e_{ps}+e_{ij}+e_{im},e_{sq}+e_{jm}-e_{mm}\\}=w$. Applying (25) and (27),
this yields
$\\{e_{ij},e_{jm}\\}=\\{e_{im},e_{mm}\\}\;\mathrm{if}\;i,j,k\;\mathrm{aredistinct}.$
(30)
Since $\\{\cdot,\cdot\\}$ is symmetric, if $i,j,m$ are distinct, we can
conclude from (28), (29) and (30) that
$\\{e_{ij},e_{jm}\\}=\\{e_{im},e_{mm}\\}=\\{e_{ii},e_{im}\\}=\\{e_{im},e_{ii}\\}=\\{e_{mm},e_{im}\\}=\\{e_{jm},e_{ij}\\}.$
(31)
Step 3. In this step, we assume $i\neq q\;\mathrm{and}\;m\neq p$.
Choosing $s\neq i\;\mathrm{and}\;s\neq m$, we suppose $n>3,\>i\neq p,\>i\neq
m\;\mathrm{and}\;m\neq q$. As
$(e_{ps}+\frac{1}{2}e_{ii}+e_{im}-\frac{1}{2}e_{mm})\circ(e_{sq}+e_{mi}-e_{ii}+e_{mm})=e_{pq}$,
it follows that
$\\{e_{ps}+\frac{1}{2}e_{ii}+e_{im}-\frac{1}{2}e_{mm},e_{sq}+e_{mi}-e_{ii}+e_{mm}\\}=w$.
Using (25), (27) and (31), if $n>3,\>i\neq p,\>i\neq m\;\mathrm{and}\;m\neq q$
we have that
$\\{e_{im},e_{mi}\\}=\frac{1}{2}\\{e_{ii},e_{ii}\\}+\frac{1}{2}\\{e_{mm},e_{mm}\\}.$
(32)
Step 4. In this step, we assume $i=q\;\mathrm{and}\;m\neq p$.
Case 4.1. $j\neq p\;\mathrm{and}\;k\neq q$.
By (27), we have $\\{e_{km},e_{qj}\\}=0$ if $m\neq q\;\mathrm{and}\;j\neq k$.
As $\\{\cdot,\cdot\\}$ is symmetric, this yields
$\\{e_{qj},e_{km}\\}=\\{e_{km},e_{qj}\\}=0\;\mathrm{if}\;m\neq
q\;\mathrm{and}\;j\neq k.$ (33)
Case 4.2. $j\neq p\;\mathrm{and}\;k=q$.
Since $(e_{pq}+e_{qj})\circ(e_{qq}-e_{jj})=e_{pq}$ if $j\neq q$, then we have
$\\{e_{pq}+e_{qj},e_{qq}-e_{jj}\\}=w$. For $\\{\cdot,\cdot\\}$ is symmetric,
it follows from (25) and (27) that
$\\{e_{qj},e_{qq}\\}=\\{e_{qj},e_{jj}\\}=\\{e_{jj},e_{qj}\\}=\\{e_{qq},e_{qj}\\}\;\mathrm{if}\;j\neq
q.$ (34)
Noting $(e_{pq}+e_{qj}+e_{pp})\circ(e_{qq}+e_{qm}-e_{jj}-e_{pm})=e_{pq}$ if
$j\neq q\;\mathrm{and}\;m\neq q$, then we obtain
$\\{e_{pq}+e_{qj}+e_{pp},e_{qq}+e_{qm}-e_{jj}-e_{pm}\\}=w$. Applying (25),
(27), (31), (33) and (34), it can be reduced to
$\\{e_{qj},e_{qm}\\}=0\;\mathrm{if}\;j\neq q\;\mathrm{and}\;m\neq q.$ (35)
Case 4.3. $j=p\;\mathrm{and}\;k\neq q$.
From
$(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp})\circ(e_{pq}-2e_{qp}+e_{qq})=e_{pq}$,
we have that
$\\{\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp},e_{pq}-2e_{qp}+e_{qq}\\}=w$. As
$\\{\cdot,\cdot\\}$ is symmetric, it follows from (25) and (27) that
$\\{e_{qp},e_{qq}\\}-\\{e_{pp},e_{qp}\\}-2\\{e_{qp},e_{qp}\\}=0.$ (36)
Since
$(\frac{1}{2}e_{pp}+e_{pq}+2e_{qp})\circ(e_{pq}-2e_{qp}+\frac{1}{2}e_{qq})=e_{pq}$,
a similar discussion shows that
$\\{e_{qp},e_{qq}\\}-\\{e_{pp},e_{qp}\\}-4\\{e_{qp},e_{qp}\\}=0.$ (37)
Because $\\{\cdot,\cdot\\}$ is symmetric, comparing (36) with (37), we get
that
$\\{e_{qp},e_{qp}\\}=0,$ (38)
$\\{e_{qp},e_{qq}\\}=\\{e_{pp},e_{qp}\\}=\\{e_{qp},e_{pp}\\}=\\{e_{qq},e_{qp}\\}.$
(39)
Noting
$(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp})\circ(e_{pq}-2e_{qp}+e_{qq}+e_{km})=e_{pq}$
if $k\neq p\;\mathrm{and}\;m\neq q$, we have that
$\\{\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp},e_{pq}-2e_{qp}+e_{qq}+e_{km}\\}=w$.
As $\\{\cdot,\cdot\\}$ is symmetric, applying (25), (27) and (36), if $k\neq
p\;\mathrm{and}\;m\neq q$ it leads to
$\\{e_{qp},e_{km}\\}=0=\\{e_{km},e_{qp}\\}.$ (40)
Case 4.4. $j=p\;\mathrm{and}\;k=q$.
If $m\neq q$, then we can verify that
$(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}-e_{mm})\circ(e_{pq}-2e_{qp}+e_{qq}+e_{qm}+e_{pm})=e_{pq},$
$(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}+\frac{1}{2}e_{mm})\circ(e_{pq}-2e_{qp}+e_{qq}-2e_{qm}+e_{pm})=e_{pq}.$
Hence we have that
$\\{\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}-e_{mm},e_{pq}-2e_{qp}+e_{qq}+e_{qm}+e_{pm}\\}=w,$
$\\{\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}+\frac{1}{2}e_{mm},e_{pq}-2e_{qp}+e_{qq}-2e_{qm}+e_{pm}\\}=w.$
For $\\{\cdot,\cdot\\}$ is symmetric, it follows from (25), (27), (31), (38),
(39) and (40) that
$\\{e_{qp},e_{pm}\\}+\\{e_{qp},e_{qm}\\}-\\{e_{mm},e_{qm}\\}=0,$
$\\{e_{qp},e_{pm}\\}-2\\{e_{qp},e_{qm}\\}-\\{e_{mm},e_{qm}\\}=0.$
Thus if $m\neq q$ we have that
$\\{e_{qp},e_{qm}\\}=0=\\{e_{qm},e_{qp}\\},$ (41)
$\\{e_{qp},e_{pm}\\}=\\{e_{mm},e_{qm}\\}=\\{e_{qm},e_{mm}\\}=\\{e_{pm},e_{qp}\\}.$
(42)
Step 5. In this step, we assume $i\neq q\;\mathrm{and}\;m=p$.
Case 5.1. $j\neq p\;\mathrm{and}\;k\neq q$.
From (27), if $i\neq p\;\mathrm{and}\;j\neq k$ it follows that
$\\{e_{ij},e_{kp}\\}=0=\\{e_{kp},e_{ij}\\}.$ (43)
Case 5.2. $j=p\;\mathrm{and}\;k\neq q$.
Since $(e_{pq}+e_{ip}-e_{kq})\circ(e_{qq}+e_{kp})=e_{pq}$ if $i\neq
p\;\mathrm{and}\;k\neq p$, it is clear that
$\\{e_{pq}+e_{ip}-e_{kq},e_{qq}+e_{kp}\\}=w$. Applying (25), (27) and (31), it
can be reduced to
$\\{e_{ip},e_{kp}\\}=0.$ (44)
Case 5.3. $j\neq p\;\mathrm{and}\;k=q$.
By (40), if $i\neq p\;\mathrm{and}\;j\neq q$ we have that
$\\{e_{ij},e_{qp}\\}=0.$ (45)
Case 5.4. $j=p\;\mathrm{and}\;k=q$.
If $i\neq p\;\mathrm{and}\;i\neq q$, then we can verify that
$\displaystyle(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}-e_{ii})\circ(e_{pq}-2e_{qp}+e_{qq}+2e_{ip}+e_{iq})=e_{pq},$
$\displaystyle(\frac{1}{2}e_{pp}+\frac{1}{2}e_{pq}+e_{qp}+\frac{1}{2}e_{ii})\circ(e_{pq}-2e_{qp}+e_{qq}-e_{ip}+e_{iq})=e_{pq}.$
According to our assumptions, using (25), (27), (31), (38), (39) and (40), we
have
$\\{e_{pp},e_{ip}\\}+\\{e_{qp},e_{iq}\\}-2\\{e_{ii},e_{ip}\\}+2\\{e_{qp},e_{ip}\\}=0,$
(46)
$\\{e_{pp},e_{ip}\\}-2\\{e_{qp},e_{iq}\\}+\\{e_{ii},e_{ip}\\}+2\\{e_{qp},e_{ip}\\}=0.$
(47)
Comparing the two equations above, as $\\{\cdot,\cdot\\}$ is symmetric, it
follows that
$\\{e_{qp},e_{iq}\\}=\\{e_{ii},e_{ip}\\}=\\{e_{iq},e_{qp}\\}=\\{e_{ip},e_{ii}\\}\;\mathrm{if}\;i\neq
p\;\mathrm{and}\;i\neq q.$ (48)
Substituting (48) into (46), we obtain
$\\{e_{pp},e_{ip}\\}-\\{e_{ii},e_{ip}\\}+2\\{e_{qp},e_{ip}\\}=0\;\mathrm{if}\;i\neq
p\;\mathrm{and}\;i\neq q.$ (49)
Noting $(e_{pq}+e_{ip})\circ(e_{ii}-e_{pp}+2e_{qq})=e_{pq}$ if $i\neq
p\;\mathrm{and}\;i\neq q$, then we have that
$\\{e_{pq}+e_{ip},e_{ii}-e_{pp}+2e_{qq}\\}=w$. Since $\\{\cdot,\cdot\\}$ is
symmetric, it follows from (25) and (27) that
$\\{e_{ip},e_{ii}\\}=\\{e_{ip},e_{pp}\\}=\\{e_{pp},e_{ip}\\}=\\{e_{ii},e_{ip}\\}\;\mathrm{if}\;i\neq
p\;\mathrm{and}\;i\neq q.$ (50)
Substituting (50) into (49), we derive
$\\{e_{qp},e_{ip}\\}=0=\\{e_{ip},e_{qp}\\}\;\mathrm{if}\;i\neq
p\;\mathrm{and}\;i\neq q.$ (51)
Step 6. In this step, we assume $i=q\;\mathrm{and}\;m=p$.
Case 6.1. $j\neq p\;\mathrm{and}\;k\neq q$.
From (27), if $j\neq k$ it is clear that
$\\{e_{qj},e_{kp}\\}=\\{e_{kp},e_{qj}\\}=0.$ (52)
Case 6.2. $j\neq p\;\mathrm{and}\;k=q$.
From (41), we have
$\\{e_{qj},e_{qp}\\}=\\{e_{qp},e_{qj}\\}=0\;\mathrm{if}\;j\neq q.$ (53)
Case 6.3. $j=p\;\mathrm{and}\;k\neq q$.
From (51), it follows
$\\{e_{qp},e_{kp}\\}=\\{e_{kp},e_{qp}\\}=0\;\mathrm{if}\;k\neq p.$ (54)
Case 6.4. $j=p\;\mathrm{and}\;k=q$.
From (38), we obtain that
$\\{e_{qp},e_{qp}\\}=0.$ (55)
Step 7. In this step, we assume $i=q\;\mathrm{and}\;m\neq p$.
If $n>3,\>j\neq p,\>j\neq q,\>j\neq m\;\mathrm{and}\;m\neq q$, choosing $s\neq
q\;\mathrm{and}\;s\neq m$, then we can verify
$(e_{ps}+e_{mm}-e_{jm})\circ(e_{sq}+e_{qm}+e_{qj})=e_{pq}$. It follows that
$\\{e_{ps}+e_{mm}-e_{jm},e_{sq}+e_{qm}+e_{qj}\\}=w$. As $\\{\cdot,\cdot\\}$ is
symmetric, by (25) and (27), if $n>3,\>j\neq p,\>j\neq q,\>j\neq
m\;\mathrm{and}\;m\neq q$ we have that
$\\{e_{mm},e_{qm}\\}=\\{e_{jm},e_{qj}\\}=\\{e_{qj},e_{jm}\\}=\\{e_{qm},e_{mm}\\}.$
(56)
In the case of $n=3$, if $q,j,m$ are distinct we can conclude from (34) and
(42) that
$\\{e_{qj},e_{jm}\\}=\\{e_{qm},e_{mm}\\}=\\{e_{qq},e_{qm}\\}=\\{e_{qm},e_{qq}\\}=\\{e_{mm},e_{qm}\\}=\\{e_{jm},e_{qj}\\}.$
(57)
In the case of $n>3$, if $q,j,m$ are distinct, from (34), (42) and (56) we
have the same equations as well. So (57) holds true whenever $n\geq 3$.
Step 8. In this step, we assume that $i\neq q\;\mathrm{and}\;m=p$.
If $n>3,\>i\neq p,\>i\neq j,\>j\neq p\;\mathrm{and}\;j\neq q$, choosing $s\neq
p\;\mathrm{and}\;s\neq i$, then we have that
$(e_{ps}+e_{jp}-e_{ip})\circ(e_{sq}+e_{ij}+e_{ii})=e_{pq}$. From the
assumptions this yields $\\{e_{ps}+e_{jp}-e_{ip},e_{sq}+e_{ij}+e_{ii}\\}=w$.
Because $\\{\cdot,\cdot\\}$ is symmetric, if $n>3,\>i\neq p,\>i\neq j,\>j\neq
p\;\mathrm{and}\;j\neq q$ it follows from (25) and (27) that
$\\{e_{jp},e_{ij}\\}=\\{e_{ip},e_{ii}\\}=\\{e_{ij},e_{jp}\\}=\\{e_{ii},e_{ip}\\}.$
(58)
If $i,j,p$ are distinct and $n=3$, then we can conclude from (48) and (50)
that
$\\{e_{ij},e_{jp}\\}=\\{e_{ip},e_{pp}\\}=\\{e_{ii},e_{ip}\\}=\\{e_{ip},e_{ii}\\}=\\{e_{pp},e_{ip}\\}=\\{e_{jp},e_{ij}\\}.$
(59)
If $i,j,p$ are distinct and $n>3$, from (48), (50) and (58) we obtain the
equations above as well. Hence, (59) holds true whenever $n\geq 3$.
Step 9. Selecting $s\neq p\;\mathrm{and}\;s\neq m$, then if $m\neq
p\;\mathrm{and}\;m\neq q$ we can verify
$(e_{ps}+\frac{1}{2}e_{pp}+e_{mp}-\frac{1}{2}e_{mm})\circ(e_{sq}+e_{ss}+e_{pm}-e_{pp}+e_{mm})=e_{pq}$.
It implies that
$\\{e_{ps}+\frac{1}{2}e_{pp}+e_{mp}-\frac{1}{2}e_{mm},e_{sq}+e_{ss}+e_{pm}-e_{pp}+e_{mm}\\}=w$.
As $\\{\cdot,\cdot\\}$ is symmetric, by (25), (27), (31) and (59), if $m\neq
p\;\mathrm{and}\;m\neq q$ this yields
$\\{e_{mp},e_{pm}\\}=\frac{1}{2}\\{e_{pp},e_{pp}\\}+\frac{1}{2}\\{e_{mm},e_{mm}\\}=\\{e_{pm},e_{mp}\\}.$
(60)
Choosing $s\neq q\;\mathrm{and}\;s\neq m$, then if $m\neq
p\;\mathrm{and}\;m\neq q$ we have
$(e_{ps}+e_{mq}+\frac{1}{2}e_{qq}-\frac{1}{2}e_{ss}-\frac{1}{2}e_{mm})\circ(e_{sq}+e_{qm}-e_{qq}+e_{mm})=e_{pq}$.
It is clear that
$\\{e_{ps}+e_{mq}+\frac{1}{2}e_{qq}-\frac{1}{2}e_{ss}-\frac{1}{2}e_{mm},e_{sq}+e_{qm}-e_{qq}+e_{mm}\\}=w$.
Because $\\{\cdot,\cdot\\}$ is symmetric, if $m\neq p\;\mathrm{and}\;m\neq q$
it follows from (25), (27), (31) and (57) that
$\\{e_{mq},e_{qm}\\}=\frac{1}{2}\\{e_{mm},e_{mm}\\}+\frac{1}{2}\\{e_{qq},e_{qq}\\}=\\{e_{qm},e_{mq}\\}.$
(61)
Since
$(e_{pp}+\frac{5}{4}e_{pq}+e_{qp}+e_{qq})\circ(e_{pp}-\frac{3}{4}e_{pq}-e_{qp}+e_{qq})=e_{pq}$,
then we have that
$\\{e_{pp}+\frac{5}{4}e_{pq}+e_{qp}+e_{qq},e_{pp}-\frac{3}{4}e_{pq}-e_{qp}+e_{qq}\\}=w$.
As $\\{\cdot,\cdot\\}$ is symmetric, using (25), (27) and (38), it can be
reduced to
$\\{e_{pq},e_{qp}\\}=\frac{1}{2}\\{e_{pp},e_{pp}\\}+\frac{1}{2}\\{e_{qq},e_{qq}\\}=\\{e_{qp},e_{pq}\\}.$
(62)
Step 10. If $j\neq p\;\mathrm{and}\;j\neq q$, we can verify that
$(e_{pp}+e_{pq}-4e_{qp}-3e_{qq}+e_{qj}+e_{jj})\circ(2e_{pp}-4e_{qp}-e_{qq}-4e_{jp}-2e_{jq}+e_{jj})=e_{pq}$.
It follows that
$\\{e_{pp}+e_{pq}-4e_{qp}-3e_{qq}+e_{qj}+e_{jj},2e_{pp}-4e_{qp}-e_{qq}-4e_{jp}-2e_{jq}+e_{jj}\\}=w$.
Since $\\{\cdot,\cdot\\}$ is symmetric, applying (25), (27), (31), (38), (39),
(40), (53), (54), (57), (59), (61) and (62), if $j\neq p\;\mathrm{and}\;j\neq
q$ we have that
$\\{e_{qj},e_{jp}\\}=\\{e_{qp},e_{qq}\\}=\\{e_{qp},e_{pp}\\}=\\{e_{qq},e_{qp}\\}=\\{e_{pp},e_{qp}\\}=\\{e_{jp},e_{qj}\\}.$
(63)
Step 11. Now, if $i\neq m\;\mathrm{and}\;j\neq k$ we can conclude from step 1,
step 4, step 5 and step 6 that
$\\{e_{ij},e_{km}\\}=0.$ (64)
If $i,j,m$ are distinct, it follows from step 2, step 7, step 8 and step 10
that
$\\{e_{ij},e_{jm}\\}=\\{e_{im},e_{mm}\\}=\\{e_{ii},e_{im}\\}=\\{e_{im},e_{ii}\\}=\\{e_{mm},e_{im}\\}=\\{e_{jm},e_{ij}\\}.$
(65)
If $i\neq m$ and $n=3$, from (60), (61) and (62) we have that
$\\{e_{im},e_{mi}\\}=\frac{1}{2}\\{e_{ii},e_{ii}\\}+\frac{1}{2}\\{e_{mm},e_{mm}\\}.$
(66)
If $i\neq m$ and $n>3$, from step 3 and step 9 we have the same equations as
well. So we complete the proof. $\Box$
## 3\. Several applications
Now, we will give two applications of the theorems above.
Definition 3.1. We say that $G\in M_{n}(R)$ is a Jordan all-multiplicative
point in $M_{n}(R)$ if for every $M_{n}(R)$-module $X$ and every Jordan
multiplicative $R$-linear map $\varphi$ : $M_{n}(R)\to X$ at $G$ (i.e.
$\varphi(S\circ T)=\varphi(S)\circ\varphi(T)$ for any $S,T\in
M_{n}(R),\>S\circ T=G$) with $\varphi(I)=I$ is a multiplicative mapping in
$M_{n}(R)$.
Corollary 3.2. Every matrix units $e_{ij}$ in $M_{n}(R)$, $n\geq 3$, is a
Jordan all-multiplicative point.
Proof. Let $X$ be an $M_{n}(R)$-module, $\varphi$ be a Jordan multiplicative
$R$-linear map at $e_{ij}$. Then it follows
$\varphi(I)=I,\;\varphi(S\circ T)=\varphi(S)\circ\varphi(T)\
\mathrm{for\;all}\>S,T\in M_{n}(R),\>S\circ T=e_{ij}.$
Set $\\{S,T\\}=\varphi(S)\circ\varphi(T)$ for any $S,T\in M_{n}(R)$, thus
$\\{\cdot,\cdot\\}$ is a symmetric $R$-bilinear map. Since $\varphi$ is a
multiplicative map at $e_{ij}$, we have
$\\{S,T\\}=\varphi(S)\circ\varphi(T)=\varphi(e_{ij})\
\mathrm{for\;all}\>S,T\in M_{n}(R)\;\mathrm{with}\;S\circ T=e_{ij}.$
As $e_{ij}$ is a Jordan product determined point in $M_{n}(R)$, then there
exists an $R$-linear map $\phi$ : $M_{n}(R)^{2}\to X$ such that
$\\{S,T\\}=\phi(S\circ T)$ for all $S,T\in M_{n}(R)$. So
$\varphi(S)\circ\varphi(T)=\\{S,T\\}=\phi(S\circ T),\;\forall\>S,T\in
M_{n}(R).$
Set $S=I$ in the equation above, then we have $\varphi(T)=\phi(T)$ for every
$T\in M_{n}(R)$. It follows that
$\varphi(S\circ T)=\varphi(S)\circ\varphi(T),\;\forall\;S,T\in M_{n}(R).$
Hence $\varphi$ is a multiplicative mapping in $M_{n}(R)$. The proof is
completed. $\Box$
Definition 3.3. We say that $H\in M_{n}(R)$ is a Jordan all-derivable point in
$M_{n}(R)$ if for every $M_{n}(R)$-module $X$ and every Jordan derivable
$R$-linear map $\varphi$ : $M_{n}(R)\to X$ at $H$ (i.e. $\varphi(S\circ
T)=\varphi(S)\circ T+S\circ\varphi(T)$ for any $S,T\in M_{n}(R),\>S\circ T=H$)
with $\varphi(I)=0$ is a Jordan derivation in $M_{n}(R)$.
Corollary 3.4. Every matrix units $e_{ij}$ in $M_{n}(R)$, $n\geq 3$, is a
Jordan all-derivable point.
Proof. Let $X$ be an $M_{n}(R)$-module, $\tau$ be a Jordan derivable
$R$-linear map at $e_{ij}$. Then we have
$\tau(I)=0,\;\tau(S\circ T)=\tau(S)\circ
T+S\circ\tau(T)\;\mathrm{for\;all}\;S,T\in M_{n}(R),\;S\circ T=e_{ij}.$
Set $\\{S,T\\}=\tau(S)\circ T+S\circ\tau(T)$ for any $S,T\in M_{n}(R)$, so
$\\{\cdot,\cdot\\}$ is a symmetric $R$-bilinear map. Since $\tau$ is a
derivable map at $e_{ij}$, it follows
$\\{S,T\\}=\tau(S)\circ
T+S\circ\tau(T)=\tau(e_{ij})\;\mathrm{for\;all}\;S,T\in
M_{n}(R)\;\mathrm{with}\;S\circ T=e_{ij}.$
As $e_{ij}$ is a Jordan product determined point in $M_{n}(R)$, then there
exists an $R$-linear map $\psi$ : $M_{n}(R)^{2}\to X$ such that
$\\{S,T\\}=\psi(S\circ T)$ for all $S,T\in M_{n}(R)$. Hence
$\tau(S)\circ T+S\circ\tau(T)=\\{S,T\\}=\psi(S\circ T),\;\forall\>S,T\in
M_{n}(R).$
Set $S=I$ in the equation above, then we have $\tau(T)=\psi(T)$ for every
$T\in M_{n}(R)$. It follows that
$\tau(S\circ T)=\tau(S)\circ T+S\circ\tau(T),\;\forall\>S,T\in M_{n}(R).$
Then $\tau$ is a Jordan derivation in $M_{n}(R)$. The proof is completed.
$\Box$
## References
* [1] D. Wang, X. Li, H. Ge, Idempotent elements determined matrix algebras, Linear Algebra Appl. 435 (2011) 2889-2895.
* [2] D. Wang, X. Li, H. Ge, Maps determined by action on identity-product elements, Linear Algebra Appl. 436 (2012) 112-119.
* [3] J. Zhu, C. Xiong, H. Zhu, Multiplicative mappings at some points on matrix algebras, Linear Algebra Appl. 433 (2010) 914-927.
* [4] J. Zhu, C. Xiong, L. Zhang, All-derivable points in matrix algebras, Linear Algebra Appl. 430 (2009) 2070-2079.
* [5] M. Brešar, M. Grašič, J. Ortega, Zero product determined matrix algebras, Linear Algebra Appl. 430 (2009) 1486-1498.
* [6] M. Gong, J. Zhu, Jordan multiplicative mappings at some points on matrix algebras, Journal of Advanced Research in Pure Mathematics. 4 (2010) 84-93.
* [7] S. Zhao, J. Zhu, Jordan all-derivable points in the algebra of all upper triangular matrices, Linear Algebra Appl. 433 (2010) 1922-1938.
|
arxiv-papers
| 2011-11-17T14:23:22 |
2024-09-04T02:49:24.431615
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Yang Wenlei and Zhu Jun",
"submitter": "Jun Zhu Professor",
"url": "https://arxiv.org/abs/1111.4108"
}
|
1111.4122
|
# Reversible electron beam heating for suppression of microbunching
instabilities at free-electron lasers
Christopher Behrens1, Zhirong Huang2, and Dao Xiang2 1 Deutsches Elektronen-
Synchrotron DESY, Notkestr. 85, 22607 Hamburg, Germany
2 SLAC National Accelerator Laboratory, Menlo Park, CA 94025, USA
###### Abstract
The presence of microbunching instabilities due to the compression of high-
brightness electron beams at existing and future X-ray free-electron lasers
(FELs) results in restrictions on the attainable lasing performance and
renders beam imaging with optical transition radiation impossible. The
instability can be suppressed by introducing additional energy spread, i.e.,
“heating” the electron beam, as demonstrated by the successful operation of
the laser heater system at the Linac Coherent Light Source. The increased
energy spread is typically tolerable for self-amplified spontaneous emission
FELs but limits the effectiveness of advanced FEL schemes such as seeding. In
this paper, we present a reversible electron beam heating system based on two
transverse deflecting radio-frequency structures (TDSs) up and downstream of a
magnetic bunch compressor chicane. The additional energy spread is introduced
in the first TDS, which suppresses the microbunching instability, and then is
eliminated in the second TDS. We show the feasibility of the microbunching
gain suppression based on calculations and simulations including the effects
of coherent synchrotron radiation. Acceptable electron beam and radio-
frequency jitter are identified, and inherent options for diagnostics and on-
line monitoring of the electron beam’s longitudinal phase space are discussed.
###### pacs:
29.27.-a, 41.60.Cr, 41.85.Ct
## I Introduction
X-ray free-electron lasers (FELs) provide an outstanding tool for studying
matter at ultrafast time and atomic length scales LCLSnature1 , and have
become a reality with the operation of the Free-Electron Laser in Hamburg
(FLASH) FLASHAnature , the Linac Coherent Light Source (LCLS) LCLSnature2 ,
and the SPring-8 Angstrom Compact Free Electron Laser (SACLA) SACLAnature .
The required high transverse and longitudinal brightness of the X-ray FEL
driving electron bunches may encounter various degradation effects due to
collective effects like coherent synchrotron radiation (CSR) or microbunching
instabilities (e.g., Refs. CSR ; CSR-ub ; lsc-ub ), and need to be preserved
and controlled. In order to suppress a microbunching instability associated
with longitudinal bunch compression that deteriorates the FEL performance, the
LCLS uses a laser heater system to irreversibly increase the uncorrelated
energy spread within the electron bunches, i.e., the slice energy spread, to a
level tolerable for operation of a self-amplified spontaneous emission FEL LH
; LCLSheater . For future X-ray FELs that plan to use external quantum lasers
(seed lasers) to seed the FEL process in order to achieve better temporal
coherence and synchronization for pump-probe experiments, a smaller slice
energy spread is required to leave room for the additional energy modulation
imprinted by the seed laser. Thus, the amount of tolerable beam heating is
more restrictive and the longitudinal phase space control becomes more
critical (e.g., Refs. NLS ; flatflat ). The same strict requirement on small
slice energy spreads is valid for optical klystron enhanced self-amplified
spontaneous emission free-electron lasers okly .
Originally designed for high-energy particle separation by radio-frequency
(rf) fields LOLA , transverse deflecting rf structures (TDSs) are routinely
used for high-resolution temporal electron beam diagnostics at present X-ray
FELs (e.g., Refs. Kick ; Roehrs ; Filippetto ; xtcav ) and are proposed to use
for novel beam manipulation methods (e.g., Refs. psxray ; exchange1 ;
exchange2 ; mapping ; ramp ; psex ). Recently, a TDS was used to increase the
slice energy spread in an echo-enabled harmonic generation FEL experiment
Xiang ; Xiang2 . In this paper, we present a reversible electron beam heating
system that uses two TDSs located up and downstream of a magnetic bunch
compressor chicane. The additional slice energy spread is introduced in the
first TDS, which suppresses the microbunching instability, and then is
eliminated in the second TDS.
The method of reversible beam heating is shown in Sec. II by means of linear
beam optics and a corresponding matrix formalism. In Sec. III, we show the
feasibility of this scheme to preserve both the transverse and longitudinal
brightness of the electron beam, and discuss the impact of coherent
synchrotron radiation. Section IV covers the gain suppression of microbunching
instabilities by analytical calculations and numerical simulations, and in
Sec. V we discuss the impact of beam and rf jitter, and show inherent options
for diagnosis and on-line monitoring of the electron beam’s longitudinal phase
space. The results and conclusions are summarized in Sec. VI.
## II Method
Figure 1: Layout of a reversible electron beam heater system including two
transverse deflecting rf structures located up and downstream of a magnetic
bunch compressor (BC) chicane, and longitudinal phase space diagnostics using
screens and synchrotron radiation monitors (SRM). Parameters related to the
reversible beam heater system are denoted in curly brackets.
In this and the following sections, we consider a linear accelerator (linac)
employing a single bunch compressor for a soft X-ray FEL, such as the proposed
linac configuration for the Next Generation Light Source (NGLS) at LBNL
Corlett . The choice of a single magnetic bunch compressor simplifies our
consideration and analysis, although the concept is also applicable for
typical bunch compressor arrangements with multiple stages. We note that a
single bunch compressor arrangement has also been considered for the
FERMI@Elettra FEL in order to minimize the impact of microbunching
instabilities Venturini .
The generic layout of the reversible electron beam heater system is depicted
in Fig. 1. It consists of linac sections providing and accelerating high-
brightness electron beams, a magnetic bunch compressor chicane in order to
achieve sufficient peak currents to drive the FEL process, and two transverse
deflecting rf structures located up and downstream of the bunch compressor. An
additional higher-harmonic rf linearizer system (Linearizer), like at the LCLS
or FLASH Lin , can be used to achieve uniform bunch compression by means of
longitudinal phase space linearization upstream of the bunch compressor. The
whole system can be supplemented by dedicated longitudinal phase space
diagnostics (see Sec. V), and except for the two TDSs, the layout is commonly
used for bunch compression at present and future X-ray FELs.
The principle of the reversible electron beam heater relies on the physics of
TDSs arising from the Panofsky-Wenzel theorem Panowsky ; Bro , which states
that the transverse momentum gain $\Delta\vec{p}_{\perp}$ of a relativistic
electron imprinted by a TDS is related to the transverse gradient of the
longitudinal electric field $\nabla_{\perp}\mathcal{E}_{z}$ inside the TDS,
and yields
$\Delta\vec{p}_{\perp}=-i\frac{e}{\omega}\int_{0}^{L}\nabla_{\perp}\mathcal{E}_{z}d\tilde{z}\,,$
(1)
where $\omega/(2\pi)$ is the operating rf frequency, $e$ is the elementary
charge, $L$ is the structure length, and $\tilde{z}$ is the longitudinal
position inside the structure (not to be confused with the beamline
coordinate, which is given by $s$ in the following). Operating a TDS with
vertical deflection, i.e., in y-direction, near the zero-crossing rf phase
$\psi=\omega/c\,z$, electrons experience transverse kicks Kick
$\Delta y^{\prime}=\frac{e\omega V_{y}}{cE}z=K_{y}z$ (2)
and relative energy deviations ($\delta=\Delta E/E$) IES1 ; IES2
$\Delta\delta=K_{y}\frac{1}{L}\int_{0}^{L}{y(s)ds}=K_{y}\overline{y}\,,$ (3)
where $K_{y}=e\omega V_{y}/(cE)$ is the vertical kick strength, $V_{y}$ is the
peak deflection voltage in the TDS, $c$ is the speed of light in vacuum, $E$
is the electron energy, and $\overline{y}$ is the mean vertical position over
the structure length $L$ along the beamline relative to the central axis
inside the finite TDS. Here, $z$ is the internal bunch length coordinate of
the electron relative to the zero-crossing rf phase. Both the additional
transverse kicks and relative energy deviations are induced by the TDS
operation itself and generate correlations within an electron bunch. In fact,
near the zero-crossing rf phase (see Eq. (2)), the induced transverse kick
correlates linearly with the internal bunch length coordinate ($z=ct$) and
enables high-resolution temporal diagnostics (e.g., Refs. Kick ; Roehrs ;
Filippetto ), whereas the induced relative energy deviation correlates with
the vertical offset inside the TDS and results in an induced relative energy
spread $\Delta\sigma_{\delta}=K_{y}\sigma_{y}$. Here, the symbol $\sigma$
denotes the root mean square (r.m.s.) value, and $\sigma_{y}$ is the vertical
r.m.s. beam size. This additional energy spread (cf. laser heater LH ;
LCLSheater ), in combination with the momentum compaction $R_{56}$ of a bunch
compressor chicane, is able to smear microbunch structures, and
correspondingly suppresses the associated instability as is shown in Sec. IV.
The effect of induced energy spread (“beam heating”) is generated by off-axis
longitudinal electric fields, related to the principle of a TDS by the
Panofsky-Wenzel theorem, and has been observed experimentally at FLASH IES3
and the LCLS emma . The induced energy spread is uncorrelated in the
longitudinal phase space $(z,\delta)$, but shows correlations in the phase
space $(y,\delta)$, which is the reason that it can be eliminated (“beam
cooling”) with a second TDS in a reversible mode as is shown in the following
by two different approaches.
### II.1 Linear beam optics
The transverse betatron motion of an electron passing through a TDS with
vertical deflection (in $y$) is given by
$y(s)=y_{0}(s)+S_{y}(s,s_{0})z$ (4)
with the vertical shear function (e.g., Refs. Kick ; Roehrs ; IES2 )
$S_{y}(s,s_{0})=R_{34}K_{y}=\sqrt{\beta_{y}(s)\beta_{y}(s_{0})}\mathrm{sin}(\Delta\phi_{y}(s,s_{0}))\frac{e\omega
V_{y}}{cE}\,,$ (5)
where $R_{34}$ is the angular-to-spatial element of the vertical beam transfer
matrix from the TDS at $s_{0}$ to any position $s$, $\beta_{y}$ is the
vertical beta function, $\Delta\phi_{y}$ is the vertical phase advance between
$s_{0}$ and $s$, and $y_{0}$ describes the vertical beam offset independent of
any TDS shearing effect. Referring to the layout depicted in Fig. 1 and taking
bunch compression into account, the induced vertical beam offset ($\Delta
y=y-y_{0}$) downstream of the second TDS becomes (omitting the subscript $y$
in $S_{y}$)
$\displaystyle\Delta y(s)=$ $\displaystyle
S_{1}(s,s_{1})z_{1}+S_{2}(s,s_{2})z_{2}$ $\displaystyle=$
$\displaystyle(CS_{1}(s,s_{1})+S_{2}(s,s_{2}))z_{2}$ (6)
with the bunch compression factor $C=z_{1}/z_{2}$ and the shear functions
$S_{1,2}(s,s_{1,2})$ of the corresponding TDSs. Here, $S_{1}(s,s_{1})z_{1}$
describes the vertical beam offset induced by the first TDS located at $s_{1}$
that is independent of the second TDS. In order to cancel the spatial chirp
induced by the combined TDS operation, the beam offset $\Delta y$ in Eq. (6)
must vanish for any $z_{2}$. Hence, using Eq. (5) for $S_{1,2}$ in Eq. (6) and
taking acceleration from $E_{1}$ to $E_{2}$ in Linac2 into account by making
the replacement
$\beta_{y}(s)\beta_{y}(s_{1})\rightarrow\beta_{y}(s)\beta_{y}(s_{1})E_{1}/E_{2}$
Kick , we get
$\displaystyle
C\sqrt{\beta_{y}(s_{1})}\mathrm{sin}(\Delta\phi_{y}(s,s_{1}))\sqrt{E_{1}}K_{1}$
$\displaystyle+\sqrt{\beta_{y}(s_{2})}\mathrm{sin}(\Delta\phi_{y}(s,s_{2}))\sqrt{E_{2}}K_{2}=0\,,$
(7)
where $K_{1,2}$ are the vertical kick strengths of the corresponding TDSs, and
$\Delta\phi_{y}(s,s_{1,2})$ describes the vertical phase advances between
$s_{1,2}$ and $s$, respectively. As a consequence, the phase advance between
both TDSs is
$\Delta\phi_{y}(s_{2},s_{1})=\Delta\phi_{y}(s,s_{1})-\Delta\phi_{y}(s,s_{2})$.
A general solution, valid for any position $s$ downstream of the second TDS,
is only possible for a phase advance difference of
$\Delta\phi_{y}(s_{2},s_{1})=n\cdot\pi$ (8)
with $n$ being integer, and the kick strength
$K_{2}=\pm
C\sqrt{\frac{\beta_{y}(s_{1})}{\beta_{y}(s_{2})}}\sqrt{\frac{E_{1}}{E_{2}}}K_{1}\,,$
(9)
where the sign depends on the actual phase advance, i.e.,
$\Delta\phi_{y}(s_{2},s_{1})=\pi+n\cdot 2\pi$ for ($+$) and
$\Delta\phi_{y}(s_{2},s_{1})=n\cdot 2\pi$ for ($-$). The different sign of $K$
can technically be achieved by changing the rf phase in the TDS by
$180^{\circ}$ which results in a zero-crossing rf phase with opposite slope
and deflection. Besides cancelation of the induced spatial chirps, the induced
energy spread of the first TDS needs to be eliminated in the second structure
in order to have a fully reversible electron beam heater. Applying Eq. (3)
similar to Eq. (6), the relative energy deviation downstream of the second TDS
for finite structure lengths become (omitting the argument in $y(s)$ and
$S(s)$)
$\Delta\delta=K_{1}\overline{y_{1}}C\frac{E_{1}}{E_{2}}+K_{2}\overline{(y_{2}+S_{1}z_{1})}$
(10)
with the mean vertical offsets $\overline{y_{1}}$ and
$\overline{(y_{2}+S_{1}z_{1})}$ inside the TDSs. For constant vertical offsets
inside the TDSs or short structure lengths, the mean vertical offsets can be
replaced by the actual offsets, i.e., $\overline{y_{1}}\rightarrow y_{1}$ and
$\overline{(y_{2}+S_{1}z_{1})}\rightarrow(y_{2}+S_{1}z_{1})$. The latter
describes the offset in the second TDS and involves the spatial chirp induced
by the first TDS with $S_{1}\sim\mathrm{sin}(\Delta\phi_{y}(s_{2},s_{1}))$,
which vanishes in the case of spatial chirp cancelation given by Eq. (8). In
order to cancel the relative energy spread induced by the combined TDS
operation, it follows
$K_{1}\overline{y_{1}}C\frac{E_{1}}{E_{2}}+K_{2}\overline{y_{2}}=0\,.$ (11)
The general transverse beam transport optics with the vertical phase advance
condition in Eq. (8) gives
$\overline{y_{2}}=\pm\overline{y_{1}}\sqrt{\beta_{y}(s_{2})/\beta_{y}(s_{1})}$,
and taking $\beta_{y}(s_{2})\rightarrow\beta_{y}(s_{2})E_{1}/E_{2}$ (see prior
Eq. (7)) into account yields exactly the same condition as in Eq. (9).
Simultaneous spatial chirp and energy spread cancelation in the second TDS is
the basic principle for reversible electron beam heating and enables local
increase of slice energy spread. The additional energy spread in the bunch
compressor, which is added in quadrature by the first TDS, can be controlled
by the kick strength $K_{1}$ and the vertical beam size $\sigma_{y}(s_{1})$.
In the following, a complementary approach to discuss the reversible beam
heating system is shown. It uses the matrix formalism for beam transport and
provides an analytical way to show microbunching gain suppression and to
discuss the impact of beam and rf jitter.
### II.2 Matrix formalism
We adopt the beam transport matrix notation of a 6x6 matrix for
$(x,x^{\prime},y,y^{\prime},z,\delta)$ but leaves $(x,x^{\prime})$ out for
simplicity, i.e., $(y,y^{\prime},z,\delta)$ is used in the following. The 4x4
beam transport matrix for a vertical deflecting TDS in thin-lens approximation
reads (e.g., Refs. exchange1 ; IES1 ; IES3 )
${\mathbf{R}}_{T}^{thin}=\begin{pmatrix}1&0&0&0\\\ 0&1&K&0\\\ 0&0&1&0\\\
K&0&0&1\end{pmatrix}.$ (12)
As discussed above, the main components of the given reversible heater system
shown in Fig. 1 consist of TDS1 with the kick strength $K_{1}$, a bunch
compressor with the momentum compaction factor $R_{56}$, and TDS2 with the
kick strength $K_{2}$. Including the momentum compaction factor $R_{56}$ and
acceleration in Linac2 ($E_{1}\rightarrow E_{2}$), the 4x4 beam matrix between
the two TDSs is given by
${\mathbf{R}}_{C}=\begin{pmatrix}R_{33}&R_{34}&0&0\\\ R_{43}&R_{44}&0&0\\\
0&0&1&R_{56}\\\ 0&0&0&\frac{E_{1}}{E_{2}}\end{pmatrix}\,.$ (13)
In order to allow the energy change in the first TDS to be compensated for in
the second TDS, we require the point-to-point imaging from TDS1 to TDS2 (i.e.,
$R_{34}=0$), which corresponds to an equivalent vertical phase advance of
$\Delta\phi_{y}(s_{2},s_{1})=n\cdot\pi$ with $n$ being integer (see Eq. (8)).
Then we get the magnification factor
$R_{33}=\pm\sqrt{\beta_{y}(s_{2})/\beta_{y}(s_{1})}$ and $R_{44}=1/R_{33}$.
The linear accelerator section with higher-harmonic rf linearizer (Linac1 and
Linearizer) upstream of the first TDS introduces an appropriate energy chirp
$h$ for uniform bunch compression. Without loss of generality, we neglect
acceleration between the two TDSs, i.e., we do not consider Linac2 anymore.
Including Linac2 would simply result in a correction term $\sqrt{E_{1}/E_{2}}$
(cf. Eqs. (9) and (15) below) but would leave the physics unchanged. Then the
entire 4x4 beam transport matrix from the beginning of TDS1 to the end of TDS2
becomes
$\begin{pmatrix}R_{33}&0&0&0\\\
R_{43}+K_{1}K_{2}R_{56}&\frac{1}{R_{33}}&\frac{K_{1}}{R_{33}}+K_{2}(1+hR_{56})&K_{2}R_{56}\\\
K_{1}R_{56}&0&1+hR_{56}&R_{56}\\\
{K_{1}}+R_{33}K_{2}&0&0&\frac{1}{1+hR_{56}}\end{pmatrix}\,.$ (14)
Cancelation of the induced spatial chirp ($\Delta y^{\prime}\sim z_{0}$, cf.
Eq. (2)) requires $R_{45}=0$ (6x6-matrix notation), i.e.,
$K_{1}/R_{33}+K_{2}(1+hR_{56})=0\,,$ (15)
where $R_{45}$ describes the coupling between $y^{\prime}$ and $z_{0}$. We
note that the coupling between $\delta$ and $y_{0}$ (i.e., $R_{63}$ element)
is nonzero in Eq. (14) because the bunch is energy-chirped after compression
($\delta\sim z\sim y_{0}$), which can be removed by Linac3 downstream of TDS2.
For uniform bunch compression with $C^{-1}=(1+hR_{56})$, no acceleration in
Linac2, i.e., $E_{2}=E_{1}$, and taking into account that
$R_{33}=\pm\sqrt{\beta_{y}(s_{2})/\beta_{y}(s_{1})}$, Eq. (15) is identical to
Eq. (9). Thus, both formalisms yield the same result.
Since the kick strength of the first TDS is very weak, it can be implemented
by means of a short rf structure and the thin-lens approximation is still
valid. However, the kick strength of the second TDS is usually stronger, and
the effect of the finite structure length should be taken into account. The
symplectic beam transport matrix of a finite TDS with the length $L_{2}$ is
given in Ref. exchange1 by
${\mathbf{R}}_{T}^{thick}=\begin{pmatrix}1&L_{2}&K_{2}L_{2}/2&0\\\
0&1&K_{2}&0\\\ 0&0&1&0\\\ K_{2}&K_{2}L_{2}/2&K_{2}^{2}L_{2}/6&1\end{pmatrix}.$
(16)
In this case, we require the point-to-point imaging is from the first TDS to
the middle of the second TDS in order to have a complete cancellation. The
overall matrix from TDS1 to the end of TDS2, when Eq. (15) is fulfilled,
becomes more complicated. A few correction terms containing the length $L_{2}$
of TDS2 appear, which however does not change the working principle of the
reversible beam heater system. It should be pointed out that downstream of the
reversible heater system, the beam is slightly coupled in
$y^{\prime}-\delta_{0}$ and $y-z_{0}$, which results in a small growth of the
projected emittance given by
$\epsilon_{y,z}^{2}=\epsilon_{y0,z0}^{2}+\epsilon_{y0}\epsilon_{z0}\frac{\beta_{y0}\gamma_{z0}K_{1}^{2}R_{56}^{2}}{(1+hR_{56})^{2}}\,,$
(17)
where $\epsilon_{y0,z0}$ is the initial vertical (longitudinal) emittance, and
$\beta_{y0}$ and $\gamma_{z0}$ are the initial Twiss parameters. As is shown
in the following section, this projected emittance growth is typically
negligible.
## III Reversible heating and emittance preservation
We demonstrate the feasibility of the reversible beam heater system by
numerical simulations using the particle tracking code elegant Elegant , and
the simulations in the following include $5\,\times\,10^{5}$ particles. Table
1 summarizes the main parameters used in the simulations, and the adopted
accelerator optics model, including the positions of the TDSs, is shown in
Fig. 2. The magnetic bunch compressor chicane is assumed to bend in the
horizontal plane, and the TDSs are oriented perpendicularly with vertical
deflection. In the previous section,
Table 1: Parameters of the electron beam, of the bunch compressor system, and of the transverse deflecting rf structures. Parameter | Symbol | Value | Unit
---|---|---|---
Beam energy at TDS1/2 | $E$ | 350 | MeV
Lorentz factor at TDS1/2 | $\gamma$ | 685 |
Initial transverse emittance | $\gamma\epsilon_{x,y}$ | 0.6 | $\mu$m
Initial slice energy spread | $\sigma_{E}$ | $\sim$ 1 | keV
Momentum compaction factor | $R_{56}$ | $-138$ | mm
Compression factor | $C$ | $\sim$ 13 |
Final bunch current | $I_{f}$ | $\sim$ 520 | A
TDS1/2 rf frequency | $\omega/2\pi$ | 3.9 | GHz
Voltage of TDS1 | $V_{1}$ | 0.415 | MV
Voltage of TDS2 (without CSR) | $V_{2}$ | 5.440 | MV
Length of TDS1 | $L_{1}$ | 0.1 | m
Length of TDS2 | $L_{2}$ | 0.5 | m
Figure 2: Relevant accelerator optics (Twiss parameters) and positions of the
transverse deflecting rf structures used to numerically demonstrate the
reversible beam heater system.
we included Linac2 for a general derivation of the method, but in practice,
due to wakefield concerns, we recommend putting TDS2 right after the bunch
compressor. In order to show numerical examples based on this approach, Linac2
is not considered anymore throughout the rest of this paper. Except for the
TDSs, the parameters are similar to the magnetic bunch compressor system
discussed for the Next Generation Light Source at LBNL Corlett ; Venturini2 .
The initial longitudinal electron bunch profile is assumed to be flat-top with
a peak current of $\sim$ 40 A and a slice energy spread of $\sim$ 1 keV
(r.m.s.). The initial linear and quadratic chirp is set for a uniform
compression factor $C$ of about 13 across the entire bunch length. This is
possible even with bunch compressor nonlinearities by using a higher-harmonic
rf linearizer upstream of the bunch compressor Lin and needed to achieve
uniform cancelation of the induced energy spread downstream of TDS2.
Figure 3 shows the principle of the reversible beam heater system by means of
simulation of the longitudinal phase space at different positions along the
beamline.
(a) Upstream of TDS1.
(b) Downstream of TDS1.
(c) Upstream of TDS2.
(d) Downstream of TDS2.
Figure 3: Simulation of the longitudinal phase space after removing the
correlated energy chirp: (a) upstream of the first TDS, (b) directly
downstream of the first TDS, (c) directly downstream of the bunch compressor
and upstream of the second TDS, and (d) downstream of the second TDS. The axes
scales change from (b) to (c) when bunch compression takes place. The bunch
head is on the left, i.e., where $z/c<0$.
The impact of CSR is not taken into account (cf. next subsection for CSR
effects). The initial slice energy spread is heated up to $\sim$ 10 keV
(r.m.s.) in the first TDS, increased by the compression factor in the bunch
compressor to $\sim$ 130 keV (r.m.s.), and finally cooled down to $\sim$ 13
keV (r.m.s.) by the second TDS (see Figs. 3(a)-3(d)). The plot in Fig. 4(b)
shows that the heating induced by the first TDS is perfectly reversible, and
the final slice energy spread is simply the initial slice energy spread scaled
with the compression factor, which would be exactly the same like in the case
without using the reversible beam heater system. Figure 4(a) shows the heater
system impact on both the projected emittance (horizontal and vertical) and
the core energy spread, i.e., the slice energy spread in the center of the
bunch, for different voltages in the second TDS. The minimum of the vertical
emittance is related to the cancelation of the spatial chirp and energy spread
induced by the first TDS. The horizontal emittance is not affected at all, and
the small projected emittance growth (6 %) in the vertical plane at the
minimum is due to residual coupling generated by the system that is described
by Eq. (17). Nevertheless, as shown in Sec. III.1, even in the case with CSR
effects, the horizontal slice emittance stays unaffected at all and the
vertical slice emittance exhibits only deviations in the bunch head ($z/c<0$)
and tail ($z/c>0$).
(a) Projected emittances and core energy spread.
(b) Slice energy spread.
Figure 4: Simulations without CSR effects on the impact of the reversible
heater system on projected emittances, core energy spread, and slice energy
spread : (a) Projected emittances (normalized) and core energy spread, and (b)
slice energy spread for $V_{\mathrm{2}}$ at minimum emittance (see Fig. 4(a)).
The longitudinal coordinate is normalized to the bunch length.
### III.1 Impact of coherent synchrotron radiation
The previous results undergo small modifications when including CSR effects,
which is shown in Fig. 5. The voltage of the second TDS for minimum projected
emittance in the vertical is shifted by about 0.2 MV to lower values which is
due the additional energy chirp induced by CSR. In comparison to the case
without any CSR effects (cf. Fig. 4), the projected emittance in the vertical
plane is slightly increased and the slice energy spread is not perfectly
canceled in the head and tail. The slice energy spread in the core part of the
bunch is also slightly increased to 17.5 keV (r.m.s.) (instead of 13.5 keV
(r.m.s.) in the absence of CSR). The projected emittance in the horizontal is
about 1.7 larger which is independent of the reversible beam heater operation.
This horizontal emittance growth can further be reduced by minimizing the
horizontal beta function in the last dipole of the chicane where the bunch
length becomes the shortest. This optimization is independent of the relevant
motion in the vertical and does not affect the results of the reversible
heater system.
(a) Projected emittances and core energy spread.
(b) Slice energy spread.
Figure 5: Simulation on the impact of the reversible beam heater system on
projected emittances, core energy spread, and slice energy spread: (a)
Projected emittances (normalized) and core energy spread, and (b) slice energy
spread for $V_{\mathrm{2}}$ at minimum emittance (see Fig. 5(a)). CSR effects
are included by means of the 1-dimensional model in elegant Elegant .
Albeit the fact that the projected emittances are increased, the horizontal
slice emittance stays unaffected and the vertical slice emittance exhibits
only deviations in the head and tail due to CSR effects as is shown in Fig. 6.
Thus, the core emittances are well preserved. We note that vertically streaked
bunches in the bunch compressor chicane may change the impact of CSR effects
but require a 3-dimensional “point-to-point” tracking which is not available
neither in elegant nor in CSRtrack CSRtrack , and is beyond the scope of this
paper.
Figure 6: Simulation of the normalized slice emittance for both the vertical
and horizontal upstream of the first and downstream of the second TDS. CSR
effects are included.
## IV Microbunching gain suppression
The principle of the microbunching gain suppression with the reversible beam
heater system is shown by an analytical treatment following Ref. LCLSheater
and by using the beam transport matrix in Eq. (14). Then we show the
feasibility of the reversible heater system to suppress microbunching
instabilities by means of particle tracking simulations with initial density
and energy modulations.
### IV.1 Analytical calculations
Using the vector notation $(y_{0},y_{0}^{\prime},z_{0},\delta_{0})$ for
particles in the first linac upstream of the first TDS, the longitudinal
position downstream of the second TDS is given by
$z=K_{1}R_{56}y_{0}+(1+hR_{56})z_{0}+R_{56}\delta_{0}\,.$ (18)
Suppose that $\delta_{0}=\delta_{u}+\delta_{m}$, where $\delta_{u}$ is the
uncorrelated relative energy deviation, and $\delta_{m}(z_{0})$ is the
relative energy modulation accumulated before and in the first linac (Linac1).
Following Ref. LCLSheater , the initial energy modulation at the wavenumber
$k_{0}$ is converted into additional density modulation at a compressed
wavenumber $k$. For a 4-dimensional (4-D) distribution function
$F(y,y^{\prime},z,\delta)$, the bunching factor $b(k)$ is given by
$\displaystyle b(k)=$ $\displaystyle\int dydy^{\prime}dzd\delta
e^{-ikz}F(y,y^{\prime},z,\delta)$ $\displaystyle=$ $\displaystyle\int
dy_{0}dy_{0}^{\prime}dz_{0}d\delta_{u}e^{-ikK_{1}R_{56}y_{0}-ik(1+hR_{56})z_{0}}$
$\displaystyle
e^{-ikR_{56}(\delta_{u}+\delta_{m}(z_{0}))}F_{0}(y_{0},y_{0}^{\prime},z_{0},\delta_{u})\,,$
(19)
where $F_{0}(y_{0},y_{0}^{\prime},z_{0},\delta_{u})$ is the initial 4-D
distribution. If the induced energy modulation is small such that
$|kR_{56}\delta_{m}|\ll 1$, we can expand the exponent of Eq. (19) up to the
first order in $\delta_{m}$ to obtain
$\displaystyle b(k)\approx$ $\displaystyle~{}b_{0}(k_{0})-ikR_{56}\int
dz_{0}\delta_{m}(z_{0})e^{-ik_{0}z_{0}}$ $\displaystyle\times$
$\displaystyle\int
dy_{0}d\delta_{u}e^{-ikK_{1}R_{56}y_{0}-ikR_{56}\delta_{u}}U(y_{0})V(\delta_{u})\,,$
(20)
where $k=Ck_{0}$, $C=1/(1+hR_{56})$, $U(y_{0})$ describes the transverse
profile, and $V(\delta_{u})$ is the initial energy distribution. For both
Gaussian profiles ($U$ and $V$), we have
$\displaystyle b(k)=b_{0}(k_{0})-$ $\displaystyle
ikR_{56}\delta_{m}(k_{0})\exp\left[-(k^{2}R_{56}^{2}K_{1}^{2}\sigma_{y1}^{2}/2)\right]$
$\displaystyle\times$ $\displaystyle\exp\left[-(k^{2}R_{56}^{2}\sigma_{\delta
u}^{2}/2)\right]\,.$ (21)
Here, we denote the Fourier transform of $\delta_{m}(z_{0})$ as
$\delta_{m}(k_{0})$, which is the accumulated energy modulation at the
wavenumber $k_{0}$ in the first linac due to longitudinal space charge and
other collective effects. The initial energy spread is given by
$\sigma_{\delta u}$, and $\sigma_{y1}$ is the vertical beam size in the first
TDS. We see that $K_{1}\sigma_{y1}$ acts like effective energy spread for
microbunching gain suppression.
### IV.2 Numerical simulations
(a) Downstream of TDS2: Heater system off.
(b) Downstream of TDS2: Heater system on.
Figure 7: Simulation on suppression of microbunching instabilities due to an
initial density modulation, i.e., simulating CSR-driven microbunching. The
entire longitudinal phase space, after removing the correlated energy chirp,
is shown.
Suppression of microbunching instabilities is demonstrated by using both a
pure initial density modulation with 5 $\%$ peak amplitude and
$100\,\mathrm{\mu m}$ modulation wavelength ($\lambda_{m}$), and a pure
initial energy modulation with 3 keV peak amplitude and
$\lambda_{m}=50\,\mathrm{\mu m}$. Whereas the case with initial energy
modulation is immediately consistent with the previous analytical treatment
and describes the longitudinal space charge driven microbunching instability
lsc-ub , the initial density modulations need to be converted into energy
modulations by longitudinal CSR-impedance which expresses the consistency and
describes the CSR-driven microbunching instability CSR-ub . The simulations
were performed using the code elegant with $1\,\times\,10^{6}$ particles.
Figure 7 shows the longitudinal phase space downstream of the second TDS,
after removing the correlated energy chirp (linear and quadratic chirp), for
both the reversible beam heater system switched off (Fig. 7(a)) and on (Fig.
7(b)). In the case without reversible beam heater, energy and density
modulations at the compressed modulation wavelength $\lambda_{m}/C$ appear,
i.e., CSR-driven microbunching becomes visible. When switching the reversible
beam heater on, the microbunching instability disappears and the resulting
longitudinal phase space remains smooth. The reason is that the microbunches
at the compressed wavelength are smeared due to $R_{56}K_{1}\sigma_{y1}$ (cf.
Eq. (21)), and accordingly, the modulations appear as correlations in the
phase spaces $(y,z)$ and $(y^{\prime},\delta)$. The same effect of
microbunching suppression is given for initial energy modulations as shown in
Fig. 8.
(a) Downstream of TDS2: Heater system off.
(b) Downstream of TDS2: Heater system on.
Figure 8: Simulation on suppression of microbunching instabilities due to an
initial energy modulation, i.e., simulating longitudinal space charge driven
microbunching. For the sake of clarity, only the core of the longitudinal
phase space, after removing the correlated energy chirp, is shown.
The effect of the microbunching instability appears even stronger compared to
the simulations case with initial density modulations, but the performance of
the reversible heater system is the same with a smooth residual longitudinal
phase space when the reversible beam heater is switched on (see Fig. 8(b)).
Figures 7 and 8 are obtained for a magnetic bunch compressor system as shown
in Fig. 1. The electron bunch will be further accelerated and transported
throughout the rest of the accelerator to reach the final beam energy and peak
current in order to drive an X-ray FEL (not studied in this paper). A
microbunched electron beam as illustrated in Figs. 7(a) and 8(a), i.e., when
the reversible beam heater system is switched off, will accumulate additional
energy and density modulations, which would lead to unacceptable longitudinal
phase space properties for an X-ray FEL such as a large slice energy spread.
## V Practical considerations
The previous sections covered the principle of reversible electron beam
heating and microbunching gain suppression by means of analytical calculations
and numerical simulations. In real accelerators, we also have to deal with
imperfections, jitter and drifts of various parameters, and accordingly
supplementary studies with respect to sensitivity on jitter sources and
tolerances have to be performed. In the following, we discuss the impact of
beam and rf jitter on the reversible beam heater system, and also point out
the inherent possibility of longitudinal phase space diagnostics and on-line
monitoring.
### V.1 Jitter and tolerances
The impact of beam and rf jitter on the reversible beam heater method can
effectively be discussed using the Eqs. (2) and (14) with the condition in Eq.
(15). Deviations from the conditions in Eq. (15) can appear from jitter of the
individual peak deflection voltages $V_{1}$ and $V_{2}$ of the TDSs, and lead
to growth of the projected vertical emittance as is shown in Fig. 5(a), where
the voltage of the second TDS is varied. Even in the case of a large TDS
voltage jitter of 1 %, the vertical projected emittance growth is less than 2
$\%$ (see Fig. 5(a)). In the case of acceleration between the first and second
TDS, also energy jitter, which is similar or smaller than TDS voltage jitter,
due to this intermediate acceleration leads to deviation of the condition in
Eq. (15). The choice of superconducting accelerator technology even provide
much better rf stability HS ; CS . Pure arrival time jitter upstream of the
first TDS has no impact as long as the condition in Eq. (15), which describes
the coupling between $y^{\prime}$ and $t=z/c$, is fulfilled. In the case that
Eq. (15) is not exactly fulfilled, e.g., due to TDS voltage jitter which is on
the percent-level, the impact of typical arrival time jitter well below 100
fs, like at the LCLS LCLSnature2 or FLASH CS , is negligible. The most
critical jitter sources arise from energy jitter upstream of the bunch
compressor chicane and from rf phase jitter in the TDSs. The momentum
compaction factor translates energy jitter into arrival time jitter, which
leads to vertical kicks in the second TDS. The same effect of additional
vertical kicks is generated by rf phase jitter in the TDSs. In order to have
small impact of vertical kicks on the remaining beam transport, we demand
$\Delta\sigma_{y^{\prime}}\ll\sigma_{y^{\prime}}$ directly downstream of the
second TDS with the induced vertical r.m.s. kick $\Delta\sigma_{y^{\prime}}$
and the intrinsic beam divergence $\sigma_{y^{\prime}}$. The relevant total
vertical r.m.s. kick is given by
$\displaystyle\Delta\sigma_{y^{\prime}}=$
$\displaystyle\sqrt{\left(K_{2}R_{56}\frac{\sigma_{E}}{E}\right)^{2}+\left(K_{2}\frac{c}{\omega}\sigma_{\varphi_{2}}\right)^{2}+\left(\frac{K_{1}}{R_{33}}\frac{c}{\omega}\sigma_{\varphi_{1}}\right)^{2}}$
$\displaystyle=$
$\displaystyle\sqrt{\left(K_{2}R_{56}\frac{\sigma_{E}}{E}\right)^{2}+\left(K_{2}\frac{c}{\omega}\right)^{2}\left(\sigma_{\varphi_{2}}^{2}+\frac{1}{C^{2}}\sigma_{\varphi_{1}}^{2}\right)}$
$\displaystyle\approx$
$\displaystyle\sqrt{\left(K_{2}R_{56}\frac{\sigma_{E}}{E}\right)^{2}+\left(K_{2}\frac{c}{\omega}\right)^{2}\sigma_{\varphi_{2}}^{2}}$
(22)
with the energy jitter $\sigma_{E}/E$ upstream of the bunch compressor, the rf
phase jitter $\sigma_{\varphi_{1,2}}$ of the TDSs, the magnification factor
$R_{33}$ from the first to the second TDS (see Eq. (13)), and using Eq. (15)
with the compression factor $C=(1+hR_{56})^{-1}$. We see that the vertical
r.m.s. kick due to rf phase jitter in the first TDS scales with $C^{-2}$ and
can be neglected compared to the vertical r.m.s. kick induced by the second
TDS when we assume the same amount of rf phase jitter in both TDSs. The
condition for trajectory stability
$\Delta\sigma_{y^{\prime}}\ll\sigma_{y^{\prime}_{2}}=\sqrt{\epsilon_{y_{2}}/\beta_{y_{2}}}$
with the intrinsic (uncorrelated) r.m.s. beam divergence
$\sigma_{y^{\prime}_{2}}$ downstream of the bunch compressor at TDS2, where
$\epsilon_{y_{2}}$ is the geometrical emittance, can be restated as
$\sqrt{\left(R_{56}\frac{\sigma_{E}}{E}\right)^{2}+\left(\frac{c}{\omega}\sigma_{\varphi_{2}}\right)^{2}}\ll\frac{\epsilon_{y_{2}}}{K_{2}\sqrt{\beta_{y_{2}}\epsilon_{y_{2}}}}=\frac{\sqrt{\epsilon_{y_{2}}\epsilon_{y_{1}}}}{C\Delta\sigma_{\delta_{1}}}\,.$
(23)
Here, $\Delta\sigma_{\delta_{1}}$ is the additional relative energy spread
induced by the first TDS for suppression of microbunching instabilities, and
$\epsilon_{y_{1}}$ denotes the geometrical emittance upstream of the bunch
compressor at TDS1.
For the example parameters discussed throughout this paper (see also Table 1),
i.e., $C=13$, $\gamma\epsilon_{y_{1}}=0.6\,\mathrm{\mu m}$,
$\gamma\epsilon_{y_{2}}=0.72\,\mathrm{\mu m}$ (see Fig. 5(a)), and
$\Delta\sigma_{\delta_{1}}E\approx 10\,\mathrm{keV}$ with
$E=350\,\mathrm{MeV}$ ($\gamma=685$), the stability condition in Eq. (23)
yields pure relative energy jitter (neglecting rf phase jitter) of
$\sigma_{E}/E\ll 1.9\cdot 10^{-5}$ or pure rf phase jitter (neglecting energy
jitter) of $\sigma_{\varphi_{2}}\ll 0.012^{\circ}$. A combination of both will
obviously tighten the acceptable jitter. This level of rf stability is
difficult to achieve in normal conducting linacs with single bunch operation,
but might be achieved with superconducting accelerator technology like at
FLASH or as planned for NGLS, where many bunches can be accelerated in a long
rf pulse, i.e., in a bunch train. Currently, several rf feedforward and
feedback controls are able to stabilize the bunches at FLASH to
$\sigma_{E}/E=3.0\cdot 10^{-5}$ and $\sigma_{\varphi}=0.007^{\circ}$ at 150
MeV HS ; SP , and further improvements towards $\sigma_{E}/E\leq 1.0\cdot
10^{-5}$ are planned using a fast normal conducting cavity upstream of the
bunch compressors CS ; HS . With perfect scaling of rf jitter from several
independent rf power stations that adds uncorrelated, we would expect an
improvement of $\sqrt{150\,\mathrm{MeV}/350\,\mathrm{MeV}}\approx 0.66$
compared to the results at FLASH with $150\,\mathrm{MeV}$ and assuming the
beam energy of $350\,\mathrm{MeV}$ in the bunch compressor of the NGLS design.
Continuous-wave rf operation, as planned for the NGLS design Corlett , and a
proper choice of rf working points for FEL operation might improve the
stability further.
### V.2 Integrated longitudinal phase space diagnostics
A practical spin-off of the reversible beam heater system is the availability
of longitudinal phase space diagnostics. The vertical betatron motion of
electrons passing through a TDS is described by Eq. (4), which enables a
mapping from time (longitudinal coordinate) to the vertical Kick ; Roehrs ;
Filippetto , and finally a possibility to obtain temporal bunch information by
means of transverse beam diagnostics. In a similar manner, the relative energy
deviation is mapped to the horizontal in the presence of horizontal momentum
dispersion, like in a magnetic bunch compressor chicane (see, e.g., Refs.
Roehrs ; Filippetto ). The combined operation makes single-shot measurements
of the longitudinal phase space possible, and in the case of the generic
layout of a reversible electron beam heater system as depicted in Fig. 1,
longitudinal phase space measurements become feasible using the first TDS and
observation screens in the dispersive section of the bunch compressor chicane.
In order to get information of the bunch length after the bunch compression,
the second TDS can be used with downstream observation screens (not shown in
Fig. 1).
In addition to invasive longitudinal phase space measurements of a single
bunch using observation screens, even fully noninvasive measurements utilizing
incoherent synchrotron radiation, emitted in the bunch compressor bending
magnets, are possible (see, e.g., Ref. Gerth ). When using a fast gated
camera, the implication will be the possibility of on-line monitoring the
longitudinal phase space of individual bunches in multi-bunch accelerators.
## VI Summary and conclusions
Our studies show that the reversible beam heater system proposed here can
suppress microbunching instabilities and preserve the high beam brightness at
the same time. Due to CSR effects, some vertical emittance degradation in the
head and tail region of the bunch occurs, but the core emittances are well
preserved. In the numerical demonstrations using the code elegant, the first
TDS generates about 10 keV (r.m.s.) slice energy spread, which is similar to
the laser heater but with a more Gaussian energy distribution (cf. laser
heater). The bunch compression process increases the slice energy spread to
$\sim$ 130 keV (r.m.s.), which is then reversed to $\sim$ 17 keV (r.m.s.)
after the second TDS in the presence of CSR effects. Without CSR effects, the
slice energy spread is reversed to $\sim$ 13 keV (r.m.s.), which demonstrates
perfect cancelation. The simulations also show that initial bunching in energy
and density in the beam can be smeared out during the process in the
reversible beam heater system, i.e., microbunching instabilities can be
suppressed. The resulting smooth beam can then propagate through the remaining
accelerator without further generation of much additional energy spread and is
advantageous for any kind of laser seeding manipulations and experiments. For
example, this scheme significantly loosen the required laser power for short-
wavelength HHG seeding NLS and may strongly impact the design of future
seeded FELs. In addition, the reversible beam heater system exhibits
integrated options for diagnosis and on-line monitoring of the longitudinal
phase space applicable for multi-bunch machines, which is also the preferred
type of accelerator for the reversible heater system due to large
sensitivities on energy and rf jitter. Linear accelerators based on
superconducting rf technology might be able to match the strict tolerances in
order to keep vertical kicks small and to achieve a sufficient trajectory
stability in the downstream undulators.
###### Acknowledgements.
We would like to thank P. Emma, Ch. Gerth, A. Lumpkin, H. Schlarb, and J.
Thangaraj for useful discussions and suggestions. This work was supported by
Department of Energy Contract No. DE-AC02-76SF00515.
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|
arxiv-papers
| 2011-11-17T15:03:00 |
2024-09-04T02:49:24.439895
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Christopher Behrens, Zhirong Huang, Dao Xiang",
"submitter": "Christopher Behrens",
"url": "https://arxiv.org/abs/1111.4122"
}
|
1111.4424
|
# Initial conditions for star formation in clusters: physical and kinematical
structure of the starless core OphA-N6
Tyler L. Bourke11affiliation: Harvard-Smithsonian Center for Astrophysics, 60
Garden Street, Cambridge, MA 02138; email tbourke@cfa.harvard.edu , Philip C.
Myers11affiliation: Harvard-Smithsonian Center for Astrophysics, 60 Garden
Street, Cambridge, MA 02138; email tbourke@cfa.harvard.edu , Paola
Caselli22affiliation: School of Physics & Astronomy, E.C. Stoner Building, The
University of Leeds, Leeds, LS2 9JT, UK , James Di Francesco33affiliation:
National Research Council Canada, Herzberg Institute of Astrophysics,
Victoria, BC, Canada , Arnaud Belloche44affiliation: Department of Physics and
Astronomy, University of Calgary, Calgary, AB, Canada , René
Plume55affiliation: Max-Planck-Institut für Radioastronomie, Auf dem Hügel 69,
D-53121 Bonn, Germany , David J. Wilner11affiliation: Harvard-Smithsonian
Center for Astrophysics, 60 Garden Street, Cambridge, MA 02138; email
tbourke@cfa.harvard.edu
###### Abstract
We present high spatial ($<$300 AU) and spectral (0.07 km s-1) resolution
Submillimeter Array observations of the dense starless cluster core Oph A N6,
in the 1 mm dust continuum and the 3-2 line of N2H+ and N2D+. The dust
continuum observations reveal a compact source not seen in single-dish
observations, of size $\sim$1000 AU and mass 0.005-0.01 ${\rm M}_{\sun}$. The
combined line and single-dish observations reveal a core of size 3000 $\times$
1400 AU elongated in a NW-SE direction, with almost no variation in either
line width or line center velocity across the map, and very small non-thermal
motions. The deuterium fraction has a peak value of $\sim$0.15 and is $>$ 0.05
over much of the core. The N2H+ column density profile across the major axis
of Oph A-N6 is well represented by an isothermal cylinder, with temperature 20
K, peak density $7.1\times 10^{6}$ cm-3, and N2H+ abundance $2.7\times
10^{-10}$. The mass of Oph A-N6 is estimated to be 0.29 ${\rm M}_{\sun}$,
compared to a value of 0.18 ${\rm M}_{\sun}$ from the isothermal cylinder
analysis, and 0.63 ${\rm M}_{\sun}$ for the critical mass for fragmentation of
an isothermal cylinder. Compared to isolated low-mass cores, Oph A-N6 shows
similar narrow line widths and small velocity variation, with a deuterium
fraction similar to “evolved” dense cores. It is significantly smaller than
isolated cores, with larger peak column and volume density. The available
evidence suggests Oph A-N6 has formed through the fragmentation of the Oph A
filament and is the precursor to a low-mass star. The dust continuum emission
suggests it may already have begun to form a star.
ISM: individual (Oph-A N6) – stars: formation – stars: low-mass
## 1 Introduction
It is now well established that low-mass (i.e., solar-like) stars form from
the collapse of dense cores within molecular clouds (e.g., Larson 2003). The
initial conditions of isolated star formation are known from continuum and
line observations of many tens of dense cores in molecular cloud complexes
such as Taurus, with relatively sparse concentrations of young stars (Di
Francesco et al. 2007). The best commonly observed molecular line tracers of
the central few thousand AU of centrally condensed cores on the verge of star
formation are NH3 (ammonia), N2H+, and N2D+ (Caselli et al. 2002c; Tafalla et
al. 2002; Crapsi et al. 2005, 2007). Observations of these lines and the dust
continuum have enabled the properties of dense isolated starless cores thought
to be near to or at the point of star formation to be well determined in
recent years (Di Francesco et al. 2007; Bergin & Tafalla 2007). These
properties include (i) a high degree of deuterium fractionation, (i.e., large
N(N2D+)/N(N2H+)) (ii) a large central density ($\sim 10^{6}$ cm-3), (iii)
depletion of CO and other C-bearing species, (iv) cold central regions ($<$10
K), and (v) line asymmetries indicative of infall motions.
In this paper, we present the first detailed observational study of the
internal structure of a starless cluster core, to explore how star formation
in clusters compares to isolated star formation. Most stars are believed to
form in close proximity to other stars within embedded clusters (Lada & Lada
2003; Allen et al. 2007; Bressert et al. 2010). As a result, a fundamental
problem in astrophysics is whether star formation in cluster environments is
similar to better-understood isolated star formation, or whether cluster star
formation is more turbulent and dynamic. An important question is to what
degree do stars in clusters form in the same way as in isolation, i.e., from
cores whose properties strongly influence the stellar properties (Shu et al.
2004; Larson 2005; Tan et al. 2006; see reviews by Shu et al. 1987; Larson
2003), and to what degree do they form in a different, more dynamic way, where
external forces and interactions matter more than initial conditions (e.g.,
Bonnell et al. 2001; Bonnell & Bate 2006).
A number of observational studies made over recent years have provided
information on the global properties of dense cores within nearby cluster-
forming regions (Ward-Thompson et al. 2007). Dust continuum observations
tracing high column densities have been more numerous, due to their faster
mapping speeds compared to molecular line observations. As a result, large
maps have been made of molecular clouds containing embedded clusters in
Ophiuchus (L1688; Motte et al. 1998; Johnstone et al. 2000; Stanke et al.
2006; Young et al. 2006), Perseus (NGC 1333, IC348-SW; Sandell & Knee 2001;
Hatchell et al. 2005; Enoch et al. 2006; H. Kirk et al. 2006), Corona
Australis (Chini et al. 2003; Nutter et al. 2005), Serpens (Davis et al. 1999;
Enoch et al. 2007), and the more distant and massive Orion cluster (Chini et
al. 1997; Lis et al. 1998; Johnstone & Bally 1999; Li et al. 2007). Some
progress has been made in large area mapping of molecular line dense gas
tracers with resolution approaching that of dust continuum observations of
$\leq 30\arcsec$, in nearby low-mass regions (Williams & Myers 2000; André et
al. 2007; Walsh et al. 2007; Friesen et al. 2009, 2010a,b), and Orion (Ikeda
et al. 2007, 2009; Tatematsu et al. 2008). Although some studies have been
able to “resolve” starless cores in clusters in molecular spectral lines
(Williams & Myers 2000; Walsh et al 2007), in the sense that the measured core
size is larger than the beam, no study has examined the internal structure of
any cluster core in detail. Observations of cores in cluster-forming regions
that resolve individual cores are thus needed to understand how these regions
can make stars so much more efficiently than in isolation. Progress toward
this goal has been slower than for isolated regions. Most young clusters are
more crowded so their cores suffer more confusion; their cores are smaller and
more distant so they are harder to resolve, and they form stars more
frequently so that starless cores are less common (Jijina et al. 1999).
### 1.1 Ophiuchus A N6
The Ophiuchus molecular cloud, at a distance of only 125 pc (de Geus et al.
1989; Knude & Hog 1998; Loinard et al. 2008; Lombardi et al. 2008), is the
nearest example of cluster formation embedded within dense gas and dust, and
is an ideal region in which to study the initial conditions for cluster
formation (see review by Wilking, Gagné and Allen 2008; Motte et al. 1998;
Johnstone et al. 2000, 2004; André et al. 2007; Enoch et al. 2008; Jørgensen
et al. 2008; Simpson et al. 2008; Padgett et al. 2008; Friesen et al. 2009,
2010a,b; van Kempen et al. 2009; Gutermuth et al. 2009; Maruta et al. 2010).
Within Ophiuchus, the Oph A ridge is the brightest of the clumps in both dust
emission and N2H+ 1-0, containing 8 dust cores identified throught 1.3 mm
continuum emission (Motte et al. 1998) and 6 local maxima of integrated N2H+
1-0 emission (Di Francesco et al. 2004, hereafter DAM04; André et al. 2007).
The N2H+ cores in the main part of the ridge, where the dust emission is
strongest, show line widths that are not significantly different from those
observed toward isolated starless cores (Di Francesco et al. 2004; Crapsi et
al. 2005). In this region, N6 is the best core for studies of internal
structure as it is isolated from the other Oph A cores (thereby suffering less
from confusion), is bright in molecular lines, and is larger and thus better
resolved than the other Oph A cores (Di Francesco et al. 2004; Pon et al.
2009). For these reasons, we have undertaken high angular resolution
submillimeter observations of N6 using the Submillimeter Array in the high
density gas tracers N2H+ and N2D+ 3-2, and combined those data with single
dish observations using the James Clerk Maxwell Telescope (JCMT) and Institut
de Radioastronomie Millimétrique (IRAM) 30-m telescope.
This paper is divided into the following sections. In §2 we present our
observations using the Submillimeter Array (SMA), JCMT, and IRAM 30-m, in §3
we present the line and continuum results, §4 presents the analysis, including
the column densities, deuterium fraction, and structure of N6, §5 presents a
discussion on the structure and evolution of N6, and compares it to both
isolated cores and cores within cluster-forming regions, while a summary is
presented in §6.
## 2 Observations
### 2.1 Submillimeter Array
Observations in the N2D+ 3-2 line toward Oph A-N6 were undertaken with the SMA
on 2006 July 22. The array was in its compact configuration and zenith
opacities at 225 GHz were typically 0.1-0.13. The SMA correlator was
configured with 2048 channels over 104 MHz for the N2D+ 3-2 line at 231.321
GHz, providing a channel spacing of 0.066 km s-1. This high resolution mode
decreases the available bandwidth for continuum observations, resulting in 1
GHz of continuum bandwidth in the upper sideband and 1.15 GHz in the lower.
The effective continuum frequency is 226 GHz (1.3 mm). The observations of Oph
A N6 were interleaved with the quasar J1626-298 for complex gain calibration.
Uranus and 3c454.3 were used for bandpass calibration, and Uranus was also
used for flux calibration. The data were calibrated and edited in the SMA’s
MIR software package.
Observations in the N2H+ 3-2 line toward Oph A-N6 were undertaken with the SMA
on 2007 March 30. The array was in its compact configuration and zenith
opacities at 225 GHz were typically $\sim 0.06$. The SMA correlator was
configured with 2048 channels over 104 MHz for the N2H+ 3-2 line at 279.512
GHz, providing a channel spacing of 0.055 km s-1. This high resolution mode
decreases the available bandwidth for continuum observations, resulting in 1.3
GHz of continuum bandwidth in both the upper and lower sidebands. The
effective continuum frequency is 276 GHz (1.1 mm). The observations of Oph
A-N6 were interleaved with the quasars J1626-298 and J1517-243 for complex
gain calibration. The quasar 3c279 was used for bandpass calibration, and
Titan and 3c279 were used for flux calibration. The data were calibrated and
edited in the SMA’s MIR software package.
Further observations in the N2H+ 3-2 line toward Oph A-N6 were undertaken with
the SMA on 2007 May 2, with the array in its subcompact configuration. Zenith
opacities at 225 GHz were typically 0.06-0.08. The correlator setup,
resolution, available continuum bandwidth, and complex gain calibrators were
the same as for the 2007 March 30 observations. Neptune and 3c273 were used
for bandpass calibration, and Neptune was also used for flux calibration. The
data were calibrated and edited in the SMA’s MIR software package.
Based upon independent observations of the gain and passband quasars at
similar frequencies ($\pm 10$ GHz) and times (within one month), and upon
independent calibration of the observations presented here, as part of the
SMA’s ongoing monitoring of quasar fluxes111see
http://sma1.sma.hawaii.edu/callist/callist.html, we estimate the flux
calibration to be good to 20% for all our observations.
### 2.2 JCMT
Observations in the N2D+ 3-2 line toward Oph A-N6 were undertaken with the
James Clerk Maxwell Telescope (JCMT) between February and July 2005. All
observations were made on a $5\times 5$ grid with 10″ spacing, with angular
resolution of 22″, and a spectral resolution of 0.1 km s-1. Details of the
observations have been presented in Pon et al. (2009), and the reader is
referred to that paper for further information.
### 2.3 IRAM 30-m
Observations of Oph A-N6 with the IRAM 30-m telescope were undertaken in May
2007 as part of a larger mapping program of the Oph A ridge. Although four SIS
heterodyne receivers were used simultaneously in the 3, 2, and 1.2 mm
atmospheric windows, here we only focus on the observations of N2H+ 3-2 at
279.511 GHz that are used in this paper. The autocorrelation spectrometer
VESPA was used as backend with a channel spacings of 40 kHz and bandwidth of
80 MHz. The system temperature ranged from 650 to 3620 K and the pointing was
checked every 1–2 hours on bright quasars and found to be good to 2–3″ (rms).
The telescope focus was optimized on Saturn and Jupiter every 3–4 hours. At
the frequency of N2H+ 3-2 the telescope beam size (full-width at half-power)
is 9″. The observations were performed in position-switching mode with the OFF
position offset by ($\Delta\alpha$,$\Delta\delta$) = (-900″,0″) from the
nominal map center of 16h26m26$\fs$46, -24°24′30.8″, located at VLA 1623. No
emission was found at the OFF position down to an rms noise level of 0.42 K in
$T_{a}^{*}$ scale. Mapping was done in on-the-fly scanning mode with a step of
4′′, providing fully sampled maps. We scanned alternately in right ascension
and declination to avoid striping artefacts. The data were reduced using the
CLASS software in its Fortran 90 version222see
http://www.iram.fr/IRAMFR/GILDAS
### 2.4 Combined Interferometric and Single-Dish Observations
The procedure used to combine the interferometric and single dish data sets is
similar to that described by Zhang et al. (2000) and Takakuwa et al. (2007),
and is based on the methods described in Vogel et al. (1984) and Wilner &
Welch (1994). The MIRIAD software package (Sault et al. 1995) was used for the
combination and subsequent imaging.
#### 2.4.1 N2D+ 3-2: SMA + JCMT
For the N2D+ 3-2 observations, the data sets were resampled along the velocity
axis to a channel spacing of 0.07 km s-1. The JCMT data were converted to Jy
using a conversion factor of $S(Jy)=27.4\times T_{A}^{*}(K)$, and deconvolved
with a 22″ FWHM Gaussian used to represent the JCMT beam at 231 GHz. Next, the
JCMT data were convolved by a 55″ Gaussian representing the SMA primary beam
(full width at half power). Side lobe effects are not well known and are thus
ignored. Then the JCMT image cube was fourier transformed into a visibility
data set, with a sampling density in the (u,v) plane chosen to closely match
that of the SMA in their overlap region. Finally, the JCMT and SMA visibility
data sets were fourier transformed together back into the image plane. Because
of the extended nature of the emission, a correction for the SMA primary beam
attenuation away from the phase center was applied. With a robust weighting of
0 applied during the transform, the resultant image cube has a resolution of
$5\farcs 1\times 3\farcs 4$ (synthesised beam full-width at half power) with a
1$\sigma$ rms sensitivity of 0.43 Jy beam-1 channel-1.
#### 2.4.2 N2H+ 3-2: SMA + 30-m
The procedure for combining the SMA and 30-m N2H+ 3-2 data sets is the same as
that used for the SMA and JCMT data sets. The data sets were first resampled
along the velocity axis to a channel spacing of 0.07 km s-1, for direct
comparison with the N2D+ 3-2 data. A conversion factor of $S(Jy)=9.3\times
T_{A}^{*}(K)$ was used for the 30-m data, and the assumed Gaussian beams sizes
used were 8$\farcs$8 and 45″ for the 30-m and SMA respectively at 279.5 GHz.
The 30-m and SMA visibility data sets were fourier transformed together using
a robust weighting of 0 with a Gaussian taper of 2″. Because of the extended
nature of the emission, a correction for the SMA primary beam attenuation away
from the phase center was applied. The resultant image cube has a resolution
of $5\farcs 6\times 3\farcs 7$, similar to that of the N2D+ 3-2 cube, with a
1$\sigma$ rms sensitivity of 0.72 Jy beam-1 channel-1. A comparison of the
central spectra of the combined SMA + 30-m dataset with the 30-m only dataset,
after smoothing the SMA + 30-m dataset to the angular resolution of the 30-m
dataset, shows that essentially all of the single-dish flux is recovered, and
the line-shapes are very similar.
### 2.5 Continuum imaging
The MIRIAD software package was used to fourier transform and produce images
from the interferometric-only continuum data. At 226 GHz (1.3 mm) a robust
weighting of 2 was used with a 3″ taper, to improve the sensitivity, resulting
in a resolution of $5\farcs 2\times 4\farcs 3$ with a 1$\sigma$ rms
sensitivity of 2.6 mJy beam-1. At 276 GHz (1.1 mm) a robust weighting of 0 was
used with a 3″ taper (note that with both compact and subcompact array data, a
robust weighting of 2 would excessively down-weight the longer baselines),
resulting in a resolution of $4\farcs 6\times 3\farcs 5$ with a 1$\sigma$ rms
sensitivity of 3.6 mJy beam-1.
## 3 Results
### 3.1 Molecular line maps
Figure 1 shows the integrated line maps of N2D+ 3-2 and N2H+ 3-2, compared to
the distribution of integrated N2H+ 1-0 emission within the entire Oph A ridge
(DAM04). The integration is performed over the hyperfine structure,
corresponding to velocity ranges of 3.23–4.56 km s-1 (N2D+) and 0.22–6.18 km
s-1 (N2H+). The general agreement between the N2D+ and N2H+ emission is good,
but their peaks are offset by 9″ (Figure 2). Similarly, the N2H+ 1-0 map has
its peak offset from the N2D+ 3-2 map, and shows good positional coincidence
with N2H+ 3-2. The offsets between N2H+ and N2D+ could be due to optical depth
effects (see §4.1), or chemical differentiation (Pon et al. 2009). Figure 2
also shows that the 3-2 maps have slightly different position angles. The N2D+
3-2 map is closely aligned with the N2H+ 1-0 map, although their peaks do not
coincide.
The N2H+ 3-2 spectrum is composed of three groups of hyperfine (hf) features.
Integrated maps of these hyperfine groups are shown in Figure 3, labelled
according to their velocity offsets relative to the line frequency as “low-V”,
“main-V”, and “high-V”. The low-V and high-V groups are sometimes referred to
as the satellite hyperfine groups. Figure 3(a) shows the spectrum at the
position of peak integrated emission in the main-V map, with the positions and
relative intensities of the hyperfine components in the optically thin case
indicated. The main-V group shows considerable saturation for the components
in the range 3.5–4.0 km s-1, indicating large optical depths (confirmed
through fitting of the hyperfine structure, see §4.1). As a result, the
integrated map of the main-V group is larger than that of the other hyperfine
groups and the region of peak emission is more extended. The peaks of the
satellite hyperfine groups, low-V and high-V, are not coincident, as might be
expected as they both account for the same relative line strength. Instead,
the high-V group peaks at the same position as the main-V group, while the
low-V group peaks at the position of peak N2H+ 1-0 emission (DAM04). The
reason for this is not clear; perhaps it suggests that non-LTE excitation
anomalies, as seen in the 1-0 line (Caselli, Myers & Thaddeus 1995; Daniel,
Cernicharo, & Dubernet 2006) are present in the 3-2 line. Modelling of the
current observations with a non-LTE line code that does not assume the
hyperfines are in statistical equilibrium (Keto & Caselli 2010) could
potentially address this question, but is beyond the scope of this paper. The
satellite hyperfines are particularly strong in N6, both in absolute
intensity, and in their relative intensity compared to the main-V group. Their
relative strength compared to the main-V group is mostly due to the high total
optical depth, as noted. However, their absolute intensity is about an order-
of-magnitude brighter than has been seen toward any other low-mass starless
core, for example by comparison to L1544, using just the 30-m data for each
core (Caselli et al. 2002a,b; Daniel et al. 2007).
The map sizes as traced by the molecular line emission have been estimated
through two dimensional Gaussian fitting to the integrated maps, and through
approximate measurements by eye using the contour level tracing 50% of the
peak emission in the same maps. The results are similar for both methods for
each line, suggesting the integrated emission can be approximated by a
Gaussian. For N2H+ 3-2, we only used the maps of the low-V and high-V emission
to estimate the size, so including N2D+ we used three maps in all. All maps
and methods give consistent results, with the half-maximum diameter measured
to be $\sim 3100\times 1600$ AU, with uncertainties of a few hundred AU for
each axis. The geometric mean diameter is $\sim 2200$ AU, which is smaller
than the geometric mean diameter of $\sim 3400$ AU determined by DAM04 through
N2H+ 1-0 observations, likely due to the finer resolution of the observations
presented here (5″ cf. 10″). The ratio of major-to-minor axes is about 2:1.
The core is well-resolved, as each axis is significantly greater than the beam
diameter ($640\times 425$ AU for N2D+, $700\times 460$ AU for N2H+).
### 3.2 Dust continuum emission
Weak dust continuum emission is detected at both 1.3 mm and 1.1 mm (Fig. 4).
Unlike the emission in the single-dish map (Fig. 4a; Motte et al. 1998), N6 is
a local peak of emission at these frequencies. The interferometer has
effectively filtered out the larger scale bright emission to reveal the weak
emission associated with N6. More extended emission is missing from the 1.3 mm
map due to the array configurations used, and this may be enough to cause the
emission in this map to look smaller than the 1.1 mm emission. The dust
emission is similar in size and orientation at both wavelengths, with a small
offset (3$\farcs$5) between their peaks. However, this offset is at the
1$\sigma$ flux level and so is probably not significant. The orientation of
the dust emission is very similar to that of the molecular line emission
(Figure 5), and to that of the large scale dust emission.
The mass of the region traced by the dust emission can be determined in the
standard manner (Hildebrand 1983) using the flux density of the region. For
these calculations we use the peak dust temperature of 20 K (Pon et al. 2009),
a gas-to-dust ratio of 100, and assume the dust opacity is given by the
commonly used “OH5” opacities that are believed to best represent the dust
properties within cold dense regions (Ossenkopf & Henning 1994; Evans et al.
2001). Using the flux density above the 2 sigma level leads to a total mass
estimate of 0.011 ${\rm M}_{\sun}$ at both 1.3 mm and 1.1 mm. Using the 3
sigma level as the cutoff, the mass estimates are 0.006 ${\rm M}_{\sun}$ at
1.3 mm and 0.005 ${\rm M}_{\sun}$ at 1.1 mm. Thus the values at 1.3 mm and 1.1
mm are in agreement. At the 3 sigma intensity level the core size is of order
1000 AU at both wavelengths, much smaller than the size of the line emitting
regions. This small size may be due to large scale structure being resolved
out by the interferometer.
### 3.3 Kinematic Structure
Figures 6 and 7 show the velocity and linewidth maps of N2D+ and N2H+ 3-2. The
line velocity was determined through fits of the hyperfine structure of each
line using the hfs method in CLASS333http://www.iram.fr/IRAMFR/GILDAS, which
performs a simultaneous fitting of all hyperfine components using CERN’s
“Function Minimization and Error Analysis” package
MINUIT444http://wwwasdoc.web.cern.ch/wwwasdoc/minuit/minmain.html. MINUIT has
been shown to produce accurate values for line velocities and widths, even in
the case of severe line overlap (Pon et al. 2009). Typical uncertainties
reported by CLASS for the fits reported here are $\lesssim$ 0.01 km s-1 in
velocity and $\lesssim$ 0.02 km s-1 in line-width, for N2H+, and $\lesssim$
0.01 km s-1 in velocity and $\lesssim$ 0.03 km s-1 in line-width, for N2D+.
There is very little variation in line center velocity across the N2D+ 3-2 map
(Fig. 6a), with a possible hint of a gradient across the long axis in the
southern part of the core. The variation in line center velocity is mostly
$<0.1$ km s-1, larger than the typical uncertainty from hyperfine fitting, and
of the same order as the channel separation. The two integrated line peaks
have velocities that are separated by about a line width ($\sim 0.25$ km s-1).
The results are similar for N2H+ 3-2, in that there is very little variation
in line center velocity across the map, with the largest variation occurring
along the western edge.
DAM04 found that the linewidth of N2H+ 1-0 over N6 is generally $\leq$ 0.3 km
s-1 with a mean value of 0.25 km s-1, with typical uncertainties of 0.005-0.01
km s-1. Our observations with higher angular resolution confirm this result,
and suggest that the linewidth varies by less than the channel width of the
observations over most of the map where significant emission is present (Fig.
7). There is a suggestion of large line widths on the western edge, which may
be due to the nearby dust continuum source SM2, or simply due to lower S/N in
the edges of the map. There appears to be an increase in N2H+ line width of
order $\sim 0.03-0.04$ km s-1 near the N2D+ SE peak, which is also near the
position of the continuum source (Fig. 7). This increase is seen in the N2D+
data (Figure 8), as the line width of the SE peak of N2D+ ($0.279\pm 0.018$ km
s-1) is significantly larger than that of the NW peak ($0.217\pm 0.015$ km
s-1) As in the case of N2H+ 1-0 (DAM04), the observed line width across N6 in
the 3-2 lines of N2H+ and N2D+ is $\sim 0.25\pm 0.02$ km s-1. Observations of
the (1,1) and (2,2) lines of NH3 indicate a gas temperature of $20\pm 2$ K
(Pon et al. 2009). Using this gas temperature, the thermal velocity dispersion
$\sigma_{\rm T}$ is 0.26 km s-1, implying that the non-thermal velocity
dispersion $\sigma_{\rm NT}$ is $\sim 0.08$ km s-1. Thus the non-thermal
motions within N6 are highly subsonic, with $\sigma_{\rm NT}/\sigma_{\rm
T}\sim 0.3$.
## 4 Analysis
### 4.1 N2H+ and N2D+ Column Density
Line optical depths ($\tau$) and excitation temperatures ($T_{\rm ex}$) were
determined from fits to the hyperfine components of each transition, using the
fitting routines in CLASS (DAM04). CLASS provides an estimate of the optical
depth, and the product [$J_{\nu}(T_{\rm ex})-J_{\nu}(T_{\rm bg})$]$\tau$,
where $J_{\nu}(T_{\rm ex})$ and $J_{\nu}(T_{\rm bg})$ are the equivalent
Rayleigh-Jeans excitation and background temperatures. The method of Caselli
et al. (2002b; their appendix, in particular equation (A4)) was used to
determine the column density of N2H+ and N2D+ by integrating over the
hyperfine features. This method assumes that the line emission is optically
thin. Equation (A4) from Caselli et al. (2002b) is repeated here, as it is
important for the following discussion,
$N_{\rm tot}=\frac{8\pi
W}{\lambda^{3}A}\frac{g_{l}}{g_{u}}\frac{1}{J_{\nu}(T_{\rm ex})-J_{\nu}(T_{\rm
bg})}\frac{1}{1-\exp(-h\nu/kT_{\rm ex})}\frac{Q_{\rm
rot}}{g_{l}\exp(-E_{l}/kT_{\rm ex})}\,\,,$ (1)
where $N_{\rm tot}$ is the total column density, $W$ is the integrated line
emission, $\lambda$ and $\nu$ are the wavelength and frequency of the
observations, $A$ is the Einstein coefficient, $g_{l}$ and $g_{u}$ are the
statistical weights of the lower and upper levels, and $Q_{\rm rot}$ is the
partition function.
Fits to the hyperfine structure of N2D+ 3-2 show that this line is optically
thin at most positions (but with CLASS usually reporting values $>$ 0.1), with
the largest total opacities measured near to the map center at $\tau\sim 2-3$.
Positions where the total optical depth was $\geq$ 1 were used to estimate a
single excitation temperature for the whole map, for which we find
$\mbox{$T_{\rm ex}$}=10.0\pm 3.3$ K, so we assume a constant $T_{\rm ex}$ of
10 K for N2D+. For N2D+, the total integrated line emission and $T_{\rm ex}$
were used to calculate the column density $N$ at each map position. The
results are shown in Figure 9(a). The column density of N2D+ ranges from
$9.8\times 10^{11}$ cm-2 to $4.7\times 10^{12}$ cm-2, with values greater than
$2\times 10^{12}$ cm-2 over much of the map.
The N2H+ 3-2 emission was found to be very optically thick over much of N6,
making it difficult to estimate $\tau$, and thus determine $T_{\rm ex}$. In
addition, the saturated lines means that the observed integrated line emission
is only a lower limit of the true emission, and equation (1) is not valid for
optically thick emission. To overcome this problem, a multistep approach was
used to obtain an estimate of $T_{\rm ex}$and measure the integrated line
emission, so that the column density could be determined. Most of the optical
depth of N2H+ 3-2 is due to the main-V hyperfine group. The low-V and high-V
groups only account for 0.0742 of the total line strength (normalized to 1.0;
Daniel et al. 2006; Pagani, Daniel, & Dubernet 2009). The total integrated
line emission was thus found by integrating only over the low-V ($0.22-1.58$
km s-1) and high-V ($5.07-6.18$ km s-1) hf groups, and scaling by the inverse
of their relative line strength. In the outer part of the N2H+ map, the total
optical depth drops to reasonable values ($<15$), allowing the total line
emission to be measured, and compared to the value obtained using only the
low-V and high-V hf groups scaled by 1/0.0742. At these positions, the results
were found to be in general agreement (better than 20%). Although the total
optical depth is high, the individual hyperfine features are optically thin
(39 hyperfine features in total, 17 in the low-V and high-V hyperfine groups),
and the total optical depth of the low-V and high-V hf groups together are
also thin in these data, or at most $\tau\sim 2$ with and uncertainty of
similar size.
While it was possible to obtain good fits to essentially every map position
using only the low-V and high-V hf groups, in most cases this resulted in an
optically thin fit ($\tau$ = 0.1 in CLASS), so that $T_{\rm ex}$ is
unconstrained. In order to obtain an estimate of $T_{\rm ex}$, full hf fitting
is needed. Using only those positions away from the map center where the full
hf fit gives $\tau$ $<$ 20, we obtain $\mbox{$T_{\rm ex}$}=10.0\pm 2.2$ K. A
full hf fit to a spectrum generated from the inner $8\times 10$ positions,
gives a similar result. As we were unable to obtain a reliable estimate for
each individual map position, we assume that $\mbox{$T_{\rm ex}$}=10\pm 2$ K
across the whole map. This value is significantly lower than the value of 17 K
determined by DAM04 for N2H+ 1-0, and the value of the kinetic temperature of
20 K. This difference could suggest that while the 1-0 line is thermalized,
the 3-2 line is not. Alternatively, the denser interior of the core, better
traced by the 3-2 line, could be colder. However, $T_{\rm ex}$ is fairly
constant over the region mapped in N2H+ 3-2, and the temperature derived from
dust observations is closer to 20 K, so this alternative is the less likely of
the two possibilities.
To determine the N2H+ 3-2 column density, we assumed that the total column
density $N_{\rm tot}=N_{\rm hf}/0.0742$, where $N_{\rm hf}$ is the column
density of the outer hyperfines, calculated using the integrated line emission
of the low-V and high-V hyperfine groups, and assuming a constant $T_{\rm ex}$
of 10 K. As shown in equation (1), the column density $N$ (whether $N_{\rm
tot}$ or $N_{\rm hf}$) is simply a function of $T_{\rm ex}$, $f(\mbox{$T_{\rm
ex}$})$, times the integrated line intensity, $W$, so that $N=W\times
f(\mbox{$T_{\rm ex}$})$. The column density of N2H+ determined in this manner
ranges from $3.5\times 10^{12}$ cm-2 to $4.6\times 10^{13}$ cm-2, with most
values being greater than $10^{13}$ cm-2, and with a significant fraction of
the inner map region having values $>2.5\times 10^{13}$ cm-2 (Figure 9(b)).
The typical uncertainty in a particular measurement of the column density is
$N^{+100\%}_{-50\%}$. Similar values for the N2H+ column density were found by
DAM04.
We have checked our results, using the outer hyperfine satellite groups and
assuming optically thin emission, against N2H+ column densities determined
from hyperfine fits to the full hyperfine spectra (Caselli et al. 2002b; Di
Francesco et al. 2004; Friesen et al. 2010a). We find that the results are
consistent, in that the values from the full fit are within the uncertainties
of the method we have used. However, $N$ determined from the full fit case are
typically, but not systematically, higher (but are sometimes lower) by up to
50%. Because the total optical depth is so high its actual value is not well
constrained by the full fit at any particular position, so we prefer the
method we have used for estimating $N$.
### 4.2 Deuterium Fraction
The ratio of N2H+ and N2D+ column densities can be used to estimate the
deuteration fraction within N6. This is shown in Figure 10, where the ratio
$N$(N2D+)/$N$(N2H+) is shown, compared to the integrated intensity maps of
each molecule. From this Figure it can be seen that the D/H ratio is of order
0.05 over a large fraction of the map, reaching higher values toward the
western side, of order 0.15. These values are larger than those determined by
Pon et al. (2009), from lower resolution observations. Figure 10 also shows
that the NW N2D+ peak has a higher D/H ratio than the SE peak, as might be
expected from Pon et al. (2009), where only the NW peak is clearly detected in
the JCMT data. This result shows that Oph A-N6 has a high central degree of
deuteration, and is similar to values found for isolated low-mass starless
cores (Crapsi et al. 2005). In some map locations the D/H value is close to
the dividing line of 0.1 used to characterize the isolated cores as prestellar
or starless, with the idea that prestellar cores are those closest to star
formation (Crapsi et al. 2005). Of the prestellar cores identified by Crapsi
et al. (2005), all but one, like OphA-N6, have $N$(N2H+) $>10^{13}$ cm-2. It
is notable that even though the kinetic temperatures are near to 20 K, where
the D/H ratio should decrease dramatically (Caselli et al. 2008), and
significantly higher than in isolated cores, the D/H ratio is as high as in
most starless cores, if not higher.
### 4.3 Structure & Mass
N6 is elongated and may represent a fragment of a filament. The simplest model
of a filament is a self-gravitating isothermal cylinder, whose radial density
profile is (Ostriker 1964; Johnstone et al. 2003),
$n(r)=\frac{n_{0}}{\left[1+\left(\frac{r^{2}}{8H^{2}}\right)\right]^{2}}\,,$
(2)
where $n_{0}$ is the peak number density, $r$ is the radial offset, and the
scale length $H$ is
$H^{2}\equiv\frac{c^{2}}{4\pi G\rho_{0}}\,,$ (3)
where $c$ is the sound speed, $\rho_{0}$ the peak density, and $G$ is the
gravitational constant.
If N6 is viewed perpendicular to its axis, then the column density along the
line-of-sight is
$\displaystyle N(r)$ $\displaystyle=$
$\displaystyle\frac{\pi}{2}\frac{n_{0}H}{\left[1+\left(\frac{r^{2}}{8H^{2}}\right)^{2}\right]^{3/2}}$
(4) $\displaystyle=$ $\displaystyle
N_{0}\frac{\pi}{4R}\frac{H}{\left[1+\left(\frac{r^{2}}{8H^{2}}\right)^{2}\right]^{3/2}}$
(5)
where $N_{0}$ is the peak column density and $R$ is the radius.
Figure 12 shows the radial column density profile across the minor axis of N6
derived from N2H+ 3-2 compared to the profile of an isothermal cylinder (dark
continuous curve). This profile was constructed from N2H+ 3-2 data imaged with
a 2$\farcs$4 beam and 1$\farcs$2 pixels (Nyquist sampling), using the method
of “super-resolution” (Briggs 1994; Chandler et al. 2005), in order to better
sample the radial profile. Eighteen independent, consecutive profiles were
extracted across the major axis at 1$\farcs$2 intervals along the major axis.
The region over which the profiles were extracted is shown in Figure 11. Each
profile was normalized to its peak values, and the normalized profiles
averaged together to form the composite profile shown in Figure 12. This
figure shows that the column density profile of N6 is very well represented by
an isothermal cylinder, as the model matches the data within its 1$\sigma$
uncertainties at eight consecutive positions across the peal of the profile.
The model allows the peak density and hence abundance of N2H+ to be estimated,
keeping other parameters fixed at their previously determined values; radius
$R$ = 800 AU, temperature of 20 K (Pon et al. 2009), and peak N2H+ column
density of 4.6 $\times 10^{13}$ cm-2. Using these values, we find a good match
to the data, as shown in Figure 12, assuming a constant N2H+ abundance
$X_{N_{2}H^{+}}=2.7\pm 0.2\times 10^{-10}$, resulting in values of peak
density $n_{0}=7.1^{+0.6}_{-0.5}\times 10^{6}$ cm-3, and scale length
$H=362^{+12}_{-14}$ AU. Allowing for a 5 pc uncertainty in the distance does
not change these values.
Even though N6 is not a local dust emission peak, DAM04 estimated the column
density, $N$(H2), to be $3\times 10^{23}$ cm-2 using the dust continuum
emission, assuming isothermal dust at a temperature of 20 K. From this and
their value for $N$(N2H+) they infer an abundance $X$(N2H+) of $3\times
10^{-10}$, in very close agreement with the value used here that provides an
excellent match between the isothermal cylinder model and the data. This
abundance is in good agreement with values inferred for isolated low-mass
cores, including the evolved prestellar cores discussed above (Benson, Caselli
& Myers 1998; Caselli et al. 2002c; Crapsi et al. 2005).
The mass per unit length of an isothermal cylinder is
$M(r)=2\pi\rho_{0}\int_{0}^{R}rdr\left[1+\left(\frac{r^{2}}{8H^{2}}\right)\right]^{-2}$
(6)
After integrating, the mass of a cylinder of length $L$ can be written:
$M=L\,\frac{2c^{2}}{G}\left[1+\left(\frac{2c^{2}}{\pi
G\rho_{0}R^{2}}\right)\right]^{-1}\,.$ (7)
For $T=20$ K, $n_{0}=7.1\times 10^{6}$ cm-2, $R$ = 800 AU, and $L$ = 3100 AU,
the mass is $M=0.18\pm 0.02$ ${\rm M}_{\sun}$, where the uncertainty is due to
the uncertainties in $n_{0}$ given above and the distance uncertainty.
We can determine the total mass traced by N2H+, using the N2H+ column density
map (Fig. 9(b)), with the result for the N2H+ abundance. The map gives the
column density per pixel, from which the mass per pixel ($M_{p}$) can be
determined, and hence the total mass, using
$M_{p}=X\,\mu_{m}\,A_{p}\,N_{X}\,\,,$ (8)
where $\mu_{m}$ is the mean particle mass (2.37 amu; Stahler & Palla 2005;
Kauffmann et al. 2008), $A_{p}$ is the area per pixel, $X$ is the abundance of
the molecule used, and $N_{X}$ is its column density. In a Nyquist sampled
map, the total mass is then just the sum over all pixels. For $T_{\rm ex}$ =
10 K and $X_{N_{2}H^{+}}$ of $2.7\times 10^{-10}$, we measure
$M=0.29^{+0.05}_{-0.04}$ ${\rm M}_{\sun}$ for positions within the half-power
level of the column density map. The uncertainties come from the uncertainties
in $X$ and the distance. The change in mass by assuming $T_{\rm ex}$ = 9 or 11
K is much smaller than either of these.
The critical mass is the mass of a condensation whose radius is equal to the
shortest wavelength of a periodic perturbation that will grow. Larson (1985)
has studied the critical mass for fragmentation of a number of geometries, and
for an isothermal filament (i.e., a cylinder) finds (Larson 1985, equation 21)
$\displaystyle M_{c}$ $\displaystyle=$
$\displaystyle\frac{7.88c^{4}}{G^{2}\mu_{m}N}$ (9) $\displaystyle=$
$\displaystyle 1.1\left(\frac{T}{20\,{\rm
K}}\right)^{2}\left(\frac{10^{23}\,\,\mbox{cm${}^{-2}$}}{N}\right)[\mbox{${\rm
M}_{\sun}$}].$ (10)
With $T=20$ K and $N=1.7\times 10^{23}\mbox{cm${}^{-2}$}$ (from the peak N2H+
column density and $X$), $M_{c}=0.63^{+0.05}_{-0.04}$ ${\rm M}_{\sun}$. This
value is within about a factor of 2 of the mass computed for N6 of 0.29 ${\rm
M}_{\sun}$ from eqn. (8). Given the uncertainties in computing masses, such as
determining the “size” of a core, and our method of measuring $N$ at each
position, this result suggests that N6 is consistent with having formed from
the fragmentation of an isothermal filament, in this case Oph A.
## 5 Discussion
### 5.1 Kinematics
Internally, Oph A-N6 is rather quiescent. It shows very narrow N2H+ and N2D+
line-widths of about 0.25 km s-1 across its extent, barely more than the
thermal line width for the measured gas temperature of 20 K, of 0.18 km s-1.
Its non-thermal motions are very sub-sonic, but the surrounding gas shows
significantly larger line-widths (DAM04; André et al. 2007; Pon et al. 2009)
suggesting that N6 has lost any turbulent motions it may have had. The lack of
significant variation in line centroid velocity and line-width over the core
indicate that N6 is an example of a coherent core, as has been seen in more
isolated cores (Barranco & Goodman 1998; Goodman et al. 1998; Caselli et al.
2002a; Tafalla et al. 2004; Pineda et al. 2010). This result suggests that
small non-thermal motions typical of isolated cores are found in some starless
cores within turbulent molecular clouds.
Observations of HCO+ and DCO+ 3-2 show the expected signature of inward
motions (Evans 1999; Pon et al. 2009), but the complex hyperfine structure of
N2H+ and N2D+ 3-2 makes identifying any similar signature in these lines
impossible. In addition, the very narrow line widths of N2H+ 1-0 together with
the spectral resolution and signal-to-noise of the data make it difficult to
identify any signature of inward motions (DAM04), regardless of the hyperfine
structure. Data with finer spectral resolution and improved signal-to-noise
are required to search for inward motions in N2H+. However, the very narrow
line-widths already suggest that any inward motions on the size scales probed
by N2H+ ($\sim 300$ AU) must be small. The ratio of non-thermal-to-thermal
line-width in N6 is about 0.3, which is lower than observed in most starless
dense cores in Perseus (Walsh et al. 2007; H. Kirk et al. 2007; Rosolowsky et
al. 2008), for dense cores elsewhere in Ophiuchus (André et al. 2007; Friesen
et al. 2009), or for most isolated low-mass dense cores (Myers 1983). Further,
the absence of line-broadening toward the center of N6 suggests a lack of a
central source. The motions observed in HCO+ may be infall onto the core,
rather than core collapse (Pon et al. 2009).
### 5.2 Dust Emission
Starless cores are usually defined through observations of the dust continuum
or molecular lines in single-dish observations at millimeter wavelengths, with
angular resolution 10-20″ (typically the line observations are of lower
resolution than the continuum). As a result, by definition they generally only
show a single peak of emission, and fairly simple structures, being round or
elongated with small aspect ratios (less than 2). When observed with an
interferometer, which acts as a spatial filter, many such cores are not
detected, or still only appear as single peaks of emission, due to their
smooth large-scale structure and lack of significant sub-structure (Williams &
Myers 1999; Williams et al. 1999, 2006; Harvey et al. 2003a; Olmi et al. 2005;
Schnee et al. 2010). Combining the single-dish and interferometer line data,
as we have done here, allows the small scale structure to be studied, without
concerns about missing flux. These studies usually show that starless cores do
not break up into sub-cores on small scales. One exception is L183, which is
composed of 3 sub-cores in N2H+ 1-0 (J. Kirk et al. 2009), but it shows a very
elongated structure in single dish maps, so perhaps this is not too
surprising.
The nature of the compact dust emission detected toward the peak of integrated
N2H+ emission is unclear, given that N6 is not a local maximum in single-dish
continuum observations between 1300 and 450 µm, with 10-15″ resolution (Motte
et al. 1998; Wilson et al. 1999; Johnstone et al. 2000). However, a dust
temperature map derived from the ratio of 450-to-850 µm flux, assuming a
constant dust emissivity, shows a similar structure to the N2H+ maps, although
with lower resolution (Pon et al. 2009). The dust temperature map shows a peak
of 20 K at the N2H+ peak, and is elongated in the NW-SE direction. It is not
yet known if the gas temperature varies in a similar manner on similar scales,
as the NH3 observations only have a resolution of about 30″. However, NH3 may
not probe the highest densities toward the center of N6, and so determining
the gas temperature there with confidence will be difficult. Supporting
evidence for a relatively constant gas temperature within N6 comes from the
comparison of its column density profile with that of an isothermal cylinder
(Fig. 12), and from the almost constant line-widths.
N6 is embedded within the Oph A ridge, and the large column of dust due to the
ridge may make it difficult to distinguish a compact core within it as a
separate entity. The dust temperature map suggests that it is a local
temperature maximum, at about 20 K (Pon et al. 2009). This result is unlike
those in detailed studies of isolated starless cores, which show flat
temperature profiles in low-resolution observations (Jijina et al. 1999;
Tafalla et al. 2004), but a drop in temperature toward the core center in
observations with finer resolution (Crapsi et al. 2007). The mass of dust seen
in N6 is very low, only of order 0.005-0.01 ${\rm M}_{\sun}$, and the inferred
peak column density of $\sim 1.3\times 10^{22}$ cm-2 is an order of magnitude
below that found from N2H+ observations, and from single-dish continuum
observations.
Recently, interferometers have detected compact millimeter dust emission
toward three “starless” cores (Chen et al. 2010; Enoch et al. 2010; Pineda et
al. 2011; Dunham et al. 2011). Supporting evidence, in the form of CO outflows
or faint, compact 70 µm emission, and SED modeling, suggests that in all cases
the emission is due to an internal heating source of very low temperature
($>$100 K), and the inferred luminosities are very low ($<$0.1 $L_{\sun}$).
These cores are all candidates to be the long theorized first hydrostatic core
(FHSC; Larson 1969; Boss & Yorke 1995; Omukai 2007; Tomida et al. 2010),
although none are in complete agreement with theoretical predictions. They
have outflows that are too fast (L1448-IRS2E – Chen et al. 2010), too
collimated ((L1448-IRS2E – Chen et al. 2010; Per-Bolo 58 – Dunham et al.
2011), or are detected at too short a wavelength (Per-Bolo 58 – Enoch et al.
2011), to be consistent with current models. For one source, L1451-mm, the
observations are in better agreement with theoretical models, but a model of a
protostar plus a disk provides an equally good fit to its SED and continuum
interferometric visibilities (Pineda et al. 2011). Theoretically, a FHSC is
expected to be of low mass, with a maximum value of order 0.05 ${\rm
M}_{\sun}$, essentially undetectable at $<100$ µm, and with an observed SED
resembling a blackbody at 30 K (Omukai 2007; Saigo & Tomisaka 2011). A low-
velocity, compact outflow may also be present. Due to crowding, it is
difficult to determine if 70 µm Spitzer emission is present toward N6, and as
stated it does not show a local maximum in 450-1300 µm emission in single-dish
observations. Observations to search for the presence of an outflow in CO 2-1
are compromised by the bright outflow from VLA 1623 (André et al. 1990), that
passes close to N6 and is detected in the SMA observations presented here. The
compact 1 mm dust emission detected with the SMA has the right mass to be
considered a candidate FHSC, but further evidence is needed, particularly
detections at other wavelengths. At present there is insufficient information
to suggest whether N6 is a better FHSC candidate than the other candidates.
### 5.3 Comparison to other Starless Cores
Studies of starless cores with resolution similar to the one presented here
are rare, particularly in molecular lines. This is a greater problem for cores
in clusters, which are typically smaller than more isolated cores, such as
those in Taurus (Ward-Thompson et al. 2007, and references therein). Studies
with 10-15″ resolution have been made for the cluster-forming regions in
Ophiuchus (125 pc; Motte et al. 1998; Johnstone et al. 2000; DAM04; Simpson et
al. 2008; Friesen et al. 2009), Perseus ($\sim 235$ pc; Hatchell et al. 2005;
Walsh et al. 2007) and Serpens ($>$300 pc; Testi & Sargent 1998; Williams &
Myers 1999), and for isolated cores in Taurus (140 pc; Ward-Thompson et al.
1994; Caselli et al. 2002a,b; Tafalla et al. 2002; J. Kirk et al. 2005). A
comparison of the properties of N6 with the cores in these studies shows that
N6 is denser than most starless cores, by about an order-of-magnitude
($10^{7}$ cm-3 cf. $10^{6}$ cm-3). These observations typically do not have
the resolution of our observations of N6, and derived average densities of
$10^{6}$ cm-3 over the central 1000 AU could be consistent with peak densities
of $\sim 10^{7}$ cm-3 (Keto & Caselli 2010).
N6 is smaller than most cores in all three cluster-forming regions listed
above, (Walsh et al. 2007; Friesen et al. 2009), but this could be partly an
effect of resolution, as this study has finer resolution by at least a factor
of two over previous work. The small size could also be partly due to the
molecular transition used, as we have used higher J transitions that
preferentially trace higher column density material. Walsh et al. (2007) list
a few cores with sizes comparable to N6, but these are at the limit of their
resolution. N6 is significantly smaller than all isolated cores that have been
well resolved, whether studied in line-emission, dust continuum, or extinction
(Ward-Thompson et al. 1999; Bacmann et al. 2000; Crapsi et al. 2005; Kandori
et al. 2005; Kauffmann et al. 2008).
As noted earlier, the N2H+ linewidths in N6 are narrower than almost all other
cluster cores (André et al. 2007; Walsh et al. 2007; H. Kirk et al. 2007;
Friesen et al. 2010a), and are almost totally due to thermal motions (as
discussed in detail in DAM04). The linewidth barely varies across N6, and it
is an excellent example of a coherent core, a core where the non-thermal
motions are subsonic, and constant, so that it appears to be cut-off from the
surrounding turbulent gas (Mouschovias 1991; Myers 1998; Barranco & Goodman
1998; Goodman et al. 1998; Caselli et al. 2002c; Pineda et al. 2010).
The N2H+ column density in N6 is larger than most cores elsewhere in Ophiuchus
(Friesen et al. 2010a) and Perseus (H. Kirk et al. 2007), with a peak value of
$\sim 5\times 10^{13}$ cm-2 compared to values of $\sim 10^{13}$ cm-2. This
peak value is about three times greater than the peak value observed in a
sample of 28 isolated starless cores, and about eight times greater than the
sample mean ($\sim 8\times 10^{12}$ cm-2; Crapsi et al. 2005; see also Daniel
et al. 2007). Similarly, the N2D+ column density of N6, $\sim 5\times 10^{12}$
cm-2, is greater than that of cores in Oph B, where the peak value is $\sim
7\times 10^{11}$ cm-2, and greater than the mean value of 25 isolated starless
cores, of $<10^{12}$ cm-2 (Friesen et al. 2010b; Crapsi et al. 2005; see also
Daniel et al. 2007). However, the peak value of N(N2D+) for isolated starless
cores, observed toward L1544, L429 and L694-2, is similar to N6 (Crapsi et al.
2005).
The deuterium fraction, ranging from a mean near 0.05 to a maximum value of
about 0.15, is similar to that seen in 28 isolated cores (range 0.03 – 0.44),
where 22 of the cores have D/H $<$ 0.1 (Crapsi et al. 2005). The range of
values in N6 is also similar to that of the cluster-core Oph B2, which has a
peak of 0.16 but with most of the core showing values around 0.03 (Friesen et
al. 2010b). The mean temperature of the isolated cores is about 10 K, while
the mean in Oph B2 is higher at around 13-14 K. The deuterium fraction is
expected to be significantly higher in the cold ($<$20 K) dense interiors of
starless cores than the cosmic D/H ratio. This is due to two main factors.
First, the main pathway for the formation of H2D+, the parent molecule of
deuterated molecules, is exothermic by 230 K, and the backward reaction is not
available. Second, CO will freeze out at temperatures below 20 K, removing the
main destroyer of H2D+. However, N2 should also freeze out at essentially the
same rate at CO at $<$20 K, so the picture is not so simple. A number of more
subtle factors that affect the D/H ratio, such as grain size (surface
chemistry), ionization rate, ortho-to-para H2 value, and the CO depletion
factor are examined in detail by Caselli et al. (2008). They show that the D/H
ratio can still be relatively high near 20 K, but drops sharply at $<$15 K.
One obvious explanation for the relatively high values of D/H in N6 may be
temperatures a few degrees lower than 20 K, but the complete reason for the
high D/H values is likely to be more complicated. The D/H ratio is largest
away from the dust temperature peak, to the NW, and Pon et al. (2009) do infer
a radial temperature drop.
Thus, N6 appears to be denser and smaller than starless cores in both cluster-
forming and isolated environments. While a very small number of cluster-cores
have a similar size, no isolated core does, and no starless core has a mean
density as high.
### 5.4 Structure and Evolution
Starless cores have typically been modeled with spherically symmetric
geometries, with a radial density profile that is almost constant at small
radii but decrease as a power law at larger radii (Ward-Thompson et al. 1994,
1999; Bacmann et al. 2000; Evans et al. 2001; Kandori et al. 2005; J. Kirk et
al. 2005). However, N6 is clearly elongated, with an aspect ratio of at least
2:1, as observed in many cores (Myers et al. 1991; Ryden 1996), suggesting
that it is prolate or filamentary in nature. Over its half-maximum size its
dust temperature is fairly uniform at about 20 K, to within 1-2 K, decreasing
at larger distances (Pon et al. 2009). We have thus compared the column
density profile of N6 to that of an isothermal cylinder (Ostriker 1964; Curry
2000), finding an extremely good match between the data and the model (Fig.
12).
The mass of the isothermal model cylinder is $\sim 0.2$ ${\rm M}_{\sun}$,
similar to the observationally derived mass of 0.3 ${\rm M}_{\sun}$, given the
uncertainties in the input values (cylinder size, column density per pixel,
temperature). Similarly, the critical mass for fragmentation of an isothermal
filament with properties similar to those of N6 is $\sim 0.6$ ${\rm M}_{\sun}$
(Larson 1985), also close to the observational value given the uncertainties
(in particular as the observed value of mass depends on column density,
whereas the critical cylinder mass depends on its inverse). The excellent
match between the column density profile of N6 and the isothermal cylinder
model, and the similarity of the observed mass to the mass of an isothermal
filament with the properties of N6, strongly suggest that N6 has formed via
fragmentation of the Oph A filament, and is in a critical state at the
beginning of star formation, or has already started the star-formation process
(as evidenced by the low-mass compact dust continuum emission). N6 is aligned
with its parent core, Oph A, supporting the view that it has formed as the
result of fragmentation of Oph A along its axis.
Although non-spherically symmetric models are rarely used, they have been
quite successful in explaining the observed elongated structure of dense cores
and filamentary nature of molecular clouds (Harvey et al. 2003b; Johnstone et
al. 2003). It has also been shown that the density profiles of collapsing
centrally condensed (“Bonner-Ebert”) spheres and cylinders are remarkably
similar, and distinguishing between these cases observationally, using only
column density maps, may not be possible (Myers 2005). Resorting to the ease
of using spherical models without considering alternatives should be avoided,
a message that is often given without heed (Hartmann 2004).
Harvey et al. (2003b) performed a detailed study of the isolated starless core
L694-2, through near-infrared extinction mapping. They found that spherical
models where the radial density profile is described by a power-law, or a
Bonnor-Ebert sphere, did not provide accurate matches to the data. Instead
they found that a cylindrical model, like that described here, provided an
excellent fit to their data. In the case of L694-2, the cylinder is inclined
to the line-of-sight, so initially it was not obvious that a cylindrical model
was needed. The steep power-law index that best matched the data,
significantly steeper than predicted by inside-out collapse, motivated the
cylindrical model.
Many distant infrared dark clouds and nearby star-forming complexes have
filamentary or elongated structures (Schneider & Elmegreen 1979; Myers 2009),
and in at least one case, that of the infrared dark cloud G11.11-0.12, the
radial profile of 850 µm emission closely matches that of an isothermal
filament of radius $\approx$ 0.1 pc over a length of more than 10 pc
(Johnstone et al. 2003). While many molecular clouds have been compared to
models of infinite or finite sheets, models of individual clouds as cylinders
are almost absent, although theoretical studies suggest that such models could
provide good matches to available data (Curry 2000). New results from the
Herschel Space Observatory show filamentary structures in nearby molecular
clouds (André et al. 2010; Arzoumanian et al. 2011), and modelling suggests
radial profiles much shallower than isothermal cylinders, at radii $>$ 0.1 pc.
At such large radii the assumptions of constant temperature and/or hydrostatic
equilibrium are unlikely, so this result is not surprising. The observations
do not have the resolution to probe scales similar to that of N6. The profiles
could be consistent with collapsing polytropic cylinders (Arzoumanian et al.
2011), or magnetized filaments in virial equilibrium (Fiege & Pudritz 2000).
It will be interesting to observe cores within these filaments with finer
resolution to study their structure.
The physical, kinematic, and chemical properties of dense cores have been used
to assess their evolutionary state (Ward-Thompson et al. 1999; Crapsi et al.
2005; Di Francesco et al. 2007). Crapsi et al. (2005) searched for
evolutionary indicators in a sample of 31 isolated cores, using observations
of N2H+, N2D+, CO and the dust continuum. They proposed eight chemical and
kinematic evolutionary indicators, and identified as “evolved” those cores
that met at least four of the conditions. These conditions include large
column densities of N2H+ and N2D+, a large deuterium fraction, large CO
depletion and central density, broad linewidths, infall asymmetry in the line
profiles, and compact central regions. All of these conditions, with the
exception of CO depletion, can be tested in N6 with our data.
The peak column densities of N2H+ ($4.6\times 10^{13}$ cm-2) and N2D+
($4.7\times 10^{12}$ cm-2) are greater than the dividing values given by
Crapsi et al. (2005), of $8.5\times 10^{12}$ cm-2 and $1.0\times 10^{12}$
cm-2, respectively. Similarly, the peak of the deuterium fraction, 0.15, is
greater than the value of 0.1 used by Crapsi et al. to separate candidate
prestellar cores from starless cores. The peak density we determine,
$7.0\times 10^{6}$ cm-3, is far greater than the separator of $5.1\times
10^{5}$ cm-3 used by Crapsi et al., as is the degree of central concentration
(this is not a function of resolution, as the single dish data have been
included in our results). The narrow line widths, and hyperfine structure,
makes the search for line skewness difficult. As mentioned, Pon et al. (2009)
observe infall asymmetry in HCO+, so there is some evidence of contraction in
N6. However, the N2H+ and N2D+ lines widths are of similar size to the
separating value of 0.25 km s-1 used by Crapsi et al. (2005). So of the seven
evolutionary indicators we can measure, five clearly indicate that N6 is an
evolved prestellar core, according to the analysis of Crapsi et al. (2005),
while the other two indicators are borderline (line widths) or unclear (infall
asymmetry or line skewness). These results strongly suggest that N6 is at an
advanced stage of prestellar evolution. The compact dust continuum emission
further supports this idea.
## 6 Summary & Conclusions
We have presented high spatial ($\sim$500 AU) and spectral (0.07 km s-1)
resolution observations in N2H+ 3-2, N2D+ 3-2, and the dust continuum at 1-mm,
of the starless core Oph A-N6, embedded within the Oph A molecular ridge in
the Ophuichus cluster-forming molecular cloud. These are the highest
resolution observations of a starless dense core presented to date. Such a
small condensation as N6 would be hard to detect with coarser angular
resolution observations, especially in clouds much more distant than
Ophiuchus.
The major results of this study are summarized below:
1. 1.
The observations reveal a compact dust continuum peak, of size $\sim$1000 AU
and mass 0.005-0.01 ${\rm M}_{\sun}$, not seen in single dish observations,
except in a dust temperature map (Pon et al. 2009). The small size and mass
suggests it might be the first indication of collapse, either representing a
temperature increase, or possibly a first hydrostatic core.
2. 2.
The size of N6 from the line observations is larger, with a projected half-
power diameter of $3100\times 1600$ AU, and thus an aspect ratio of 2:1. The
N2H+ and N2D+ integrated line maps show slightly different position angles,
probably representing chemical variations, or the very high optical depth in
N2H+.
3. 3.
Very little variation is seen in either linewidth or line center velocity in
either line across their maps. The variations are so small that N6 appears to
be a coherent core, with very small non-thermal motions.
4. 4.
The peak column densities are $4.6\times 10^{13}$ cm-2 for N2H+, and
$4.7\times 10^{12}$ cm-2 for N2D+. The positions of peak column density are
offset, with the N2D+ peak located to the NW of the N2H+ peak.
5. 5.
The deuterium fraction has a peak value of 0.15, and is greater than or about
equal to 0.05 over much of the mapped area. The maximum value of deuteration
lies in the NW, and not at the position of the dust continuum peak, nor the
N2H+ peak.
6. 6.
The column density profile of N6 across its minor axis, as determined from the
N2H+ observations, is very well represented by an isothermal cylinder (at 20
K), of peak density $7.1\times 10^{6}$ cm-3, and N2H+ abundance $2.7\times
10^{-10}$.
7. 7.
The mass of N6, determined from the mapped positions, lies in the range
0.25-0.34 ${\rm M}_{\sun}$, depending strongly on the assumed N2H+ abundance
($2.7\pm 0.2\times 10^{-10}$). The low value compares favourably to the mass
determined from the cylindrical analysis, of $\sim 0.2$ ${\rm M}_{\sun}$,
while the high value compares well to the critical mass for fragmentation of
an isothermal filament with similar properties, of $\sim 0.6$ ${\rm
M}_{\sun}$.
8. 8.
Compared to isolated low-mass cores, Oph A-N6 shows similar narrow line widths
and small velocity variation, with a deuterium fraction that is similar to
“evolved” dense cores. It is significantly smaller than isolated cores, with
larger peak column density and volume density, while the previously measured
kinetic temperature is significantly higher than isolated starless cores.
These results strongly suggest Oph A-N6 has formed from the fragmentation of
the Oph A filament, and is a precursor to a low-mass star in a cluster-forming
region. The results also suggest Oph A-N6 has completed almost all of its
prestellar evolution, and may even have begun to form a star.
This research is supported in part by the National Science Foundation under
grant number 0708158 (T.L.B.). We thank Andy Pon and Rachel Friesen for
sharing results in advance of publication, and Rachel Friesen and Chris de
Vries for checking column density estimates for optically thick N2H+. We thank
Mark Gurwell for his diligent maintenance of the “Submillimeter Calibrator
List”. The Submillimeter Array is a joint project between the Smithsonian
Astrophysical Observatory and the Academia Sinica Institute of Astronomy and
Astrophysics and is funded by the Smithsonian Institution and the Academia
Sinica. The James Clerk Maxwell Telescope is operated by The Joint Astronomy
Centre on behalf of the Science and Technology Facilities Council of the
United Kingdom, the Netherlands Organisation for Scientific Research, and the
National Research Council of Canada. Based on observations carried out with
the IRAM 30 m telescope. IRAM is supported by INSU/CNRS (France), MPG
(Germany) and IGN (Spain). This research has made use of NASA’s Astrophysics
Data System Bibliographic Services
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Figure 1: Integrated combined single-dish + interferometer line maps of (a)
N2H+ 1-0 (from DAM04), (b) N2D+ 3-2, and (c) N2H+ 3-2. The contour levels are
(a) 6, 9, 12 … times the $1\sigma$ sensitivity of 0.6 Jy/beam km s-1 for N2H+
1-0, (b) 9, 12, 15 … times the $1\sigma$ sensitivity of 0.12 Jy/beam km s-1
for N2D+ 3-2, and (c) 12, 15, 18 … times the $1\sigma$ sensitivity of 0.47
Jy/beam km s-1 for N2H+ 3-2. The white cross is the position of the peak
integrated N2H+ 1-0 emission from DAM04. The large dashed circles in (b) and
(c) indicate the size of the SMA primary beam (full-width at half-maximum
sensitivity). The small grey ovals at lower right in each panel indicate the
synthesised beam sizes (full-width at half-maximum sensitivity).
Figure 2: Comparison of the integrated line maps of N2D+ 3-2 (black contours)
and N2H+ 3-2 (grey contours). The black cross is the position of the peak
integrated N2H+ 1-0 from DAM04. The synthesised beam size for each observation
is shown in the corresponding colours at lower right, and the primary beam
size for the N2D+ observations is shown as the dashed circle. Contour levels
are 10, 14, 18 … times the $1\sigma$ sensitivity of 0.12 Jy/beam km s-1 for
N2D+ 3-2, and 12, 16, 20 … times the $1\sigma$ sensitivity of 0.47 Jy/beam km
s-1 for N2H+ 3-2.
Figure 3: Spectrum of N2H+ 3-2, panel (a), and maps of integrated N2H+ 3-2
emission over the 3 hyperfine groups, panels (b)-(d). Data is combined 30-m +
SMA (compact+subcompact). In (a), the location of the hyperfine components,
and their relative weights, are indicated by the light grey lines. The limits
of integration are indicated by “low-V”, “main-V”, and “high-V”, for panels
(b)-(d). The rms noise in the spectrum is 0.5 Jy/beam. Contour levels are (b)
3, 4, 5, … times the $1\sigma$ sensitivity of 0.26 Jy/beam km s-1, (c) 15, 18,
21, … times the $1\sigma$ sensitivity of 0.30 Jy/beam km s-1, and (c) 3, 4, 5,
… times the $1\sigma$ sensitivity of 0.21 Jy/beam. The white cross is the
position of the peak integrated N2H+ 1-0 emission (DAM04), while the white
square marks the position of the spectrum shown in (a). The large dashed
circle indicates the size of the SMA primary beam. The small grey ovals at
lower right in panels (b)-(d) indicate the synthesised beam size.
Figure 4: Maps of (a) 1.3 mm emission (30-m), (b) 1.3 mm emission (SMA) and
(c) 1.1 mm emission (SMA). Contours for (a) the 30-m observations are 10, 20,
30, … times the $1\sigma$ sensitivity of 10 mJy/beam (Motte et al. 1998;
DAM04). Contours for the SMA observations are 2, 3, 4, … times the $1\sigma$
sensitivity of (b) 2.6 mJy/beam (1.3 mm) and (c) 3.6 mJy/beam (1.1 mm). Dotted
contours indicate negative levels. The cross is the position of the peak
integrated N2H+ 1-0 emission (DAM04). The large dashed circles indicates the
size of the SMA primary beam. The small grey ovals at lower right each panel
indicate the synthesised beam size.
Figure 5: Comparison of continuum and integrated line maps. (a) Contours of
1.3 mm continuum emission (white; SMA) over contours of integrated N2D+ 3-2
emission (SMA+JCMT). The greyscale is integrated N2D+ 3-2 emission. (b)
Contours of 1.1 mm emission (SMA) over contours of integrated N2H+ 3-2
emission (SMA+30-m). Contour levels for N2D+ are 10, 14, 18, … times the
$1\sigma$ sensitivity of 0.12 Jy/beam km s-1, while contour levels for N2H+
are 12, 16, 20, … times $1\sigma$ sensitivity of 0.47 0.12 Jy/beam km s-1.
Continuum contour levels are the same as shown in Figures 4(a)–(b). The white
cross is the position of the peak integrated N2H+ 1-0 emission (DAM04). The
primary beam size of the SMA for each observation is shown as the dashed
circle.
Figure 6: (a) Velocity map from hyperfine fits to the N2D+ 3-2 data. The data
have been resampled onto a 2$\farcs$5 grid with a 5″ circular beam (shown at
lower right). The positions where a fit was performed are indicated by
crosses. Typical uncertainties in the fitted velocities are 0.01-0.02 km s-1.
The contours represent the integrated intensity map of N2D+ and are the same
as shown in Figure 1. (b) Velocity map from hyperfine fits to the N2H+ 3-2
data. Typical uncertainties in the fitted velocities are 0.005-0.02 km s-1.
Other details are the same as in (a), except that the contours represent the
integrated intensity map of N2H+ and are the same as shown in Figure 1. The
origin of the maps is the position of the peak integrated N2H+ 1-0 emission
(DAM04). The open circle indicates the beam size. [COLOUR FIGURE]
Figure 7: (a) Linewidth map from hyperfine fits to the N2D+ 3-2 data. The
data have been resampled onto a 2$\farcs$5 grid with a 5″ circular beam (shown
at lower right). Typical uncertainties in the fitted linewidths are 0.02-0.04
km s-1. The positions where a fit was performed are indicated by crosses. The
contours represent the integrated intensity map of N2D+ and are the same as
shown in Figure 1. (b) Linewidth map from hyperfine fits to the N2D+ 3-2 data.
Typical uncertainties in the fitted linewidths are 0.01-0.03 km s-1. Other
details are the same as in (a), except that the contours represent the
integrated intensity map of N2H+ and are the same as shown in Figure 1. The
origin of the maps is the position of the peak integrated N2H+ 1-0 emission
(DAM04). [COLOUR FIGURE]
Figure 8: Spectra at the integrated line map peaks of N2D+ 3-2. (a) The
integrated map as shown in Figure 1, with the locations of the two spectra
indicated. The large dashed circle indicates the SMA primary beam size, while
the small grey oval indicates the synthesised beam size. The white cross is
the position of the peak integrated N2H+ 1-0 emission (DAM04). (b) The
spectrum at position NW (histogram), with a model fit of the N2D+ 3-2
hyperfine structure (continuous line in grey). (c) The spectrum at position SE
(histogram), with a model fit of the N2D+ 3-2 hyperfine structure (continuous
line in grey).
Figure 9: (a) Column density of N2D+, assuming a constant excitation
temperature of 10 K. The data have been resampled onto a 2$\farcs$5 grid with
a 5″ circular beam (shown at lower right). The positions where the column
density was determined are indicated by crosses. The contours represent the
integrated intensity map of N2D+ and are the same as shown in Figure 1. (b)
Column density of N2H+, assuming a constant excitation temperature of 10 K.
Other details are the same as in (a), except that the contours represent the
integrated intensity map of N2H+ and are the same as shown in Figure 1. The
origin of the maps is the position of the peak integrated N2H+ 1-0 emission
(DAM04). [COLOUR FIGURE]
Figure 10: (a) Ratio of N2D+ to N2H+ column densities, using the results
presented in Figure 9. The contours are integrated intensity of N2D+ and are
the same as shown in Figure 1. (b) Ratio of N2D+ to N2H+ column densities,
using the results presented in Figure 9. The contours are integrated intensity
of N2H+and are the same as shown in Figure 1. The resolution of the
observations is indicated by the circular beam (shown at lower right). The
origin of the maps is the position of the peak integrated N2H+ 1-0 emission
(DAM04). [COLOUR FIGURE]
Figure 11: Location of the N2H+ 3-2 radial column density cuts used to
construct the radial profile shown in Figure 12. The box indicates the area
from which the cuts were extracted, with the vector indicating the direction
of the cuts. The image is the low resolution column density map, with contour
levels of 1.5, 2.5, 3.6, and 4.6 $\times 10^{13}$ cm-2.
Figure 12: Normalized mean N2H+ 3-2 radial column density cut parallel to the
minor axis of N6 (histogram), constructed from a series of cuts across N6 that
were averaged and then normalized to the peak value. The dark continuous line
represents a model of an isothermal cylinder whose parameters are described in
§4.3. The model shown here has the values of radius (800 AU), kinetic
temperature (20 K) and peak N2H+ column density ($4.6\times 10^{13}$ cm-2)
fixed prior to comparison with the data, an N2H+ abundance ($2.75\times
10^{-10}$), determined by matching the model to the data (by eye), with a
resultant peak density of $7.0\times 10^{6}$ cm-3 and scale length of 365 AU.
The shaded grey areas indicate where the data and model differ significantly
and the core begins to merge into the background emission.
|
arxiv-papers
| 2011-11-18T16:55:58 |
2024-09-04T02:49:24.457193
|
{
"license": "Public Domain",
"authors": "Tyler L. Bourke, Philip C. Myers, Paola Caselli, James Di Francesco,\n Arnaud Belloche, Ren\\'e Plume, David J. Wilner",
"submitter": "Tyler Bourke",
"url": "https://arxiv.org/abs/1111.4424"
}
|
1111.4526
|
# Signal Propagation in Feedforward Neuronal Networks with Unreliable Synapses
Daqing Guo1, Chunguang Li2 Email: dqguo07@gmail.comAuthor for correspondence,
Email: cgli@zju.edu.cn
(1 School of Electronic Engineering, University of Electronic Science and
Technology of China, Chengdu 610054, People’s Republic of China
2 Department of Information Science and Electronic Engineering, Zhejiang
University, Hangzhou 310027, People’s Republic of China
)
###### Abstract
In this paper, we systematically investigate both the synfire propagation and
firing rate propagation in feedforward neuronal network coupled in an all-to-
all fashion. In contrast to most earlier work, where only reliable synaptic
connections are considered, we mainly examine the effects of unreliable
synapses on both types of neural activity propagation in this work. We first
study networks composed of purely excitatory neurons. Our results show that
both the successful transmission probability and excitatory synaptic strength
largely influence the propagation of these two types of neural activities, and
better tuning of these synaptic parameters makes the considered network
support stable signal propagation. It is also found that noise has significant
but different impacts on these two types of propagation. The additive Gaussian
white noise has the tendency to reduce the precision of the synfire activity,
whereas noise with appropriate intensity can enhance the performance of firing
rate propagation. Further simulations indicate that the propagation dynamics
of the considered neuronal network is not simply determined by the average
amount of received neurotransmitter for each neuron in a time instant, but
also largely influenced by the stochastic effect of neurotransmitter release.
Second, we compare our results with those obtained in corresponding
feedforward neuronal networks connected with reliable synapses but in a random
coupling fashion. We confirm that some differences can be observed in these
two different feedforward neuronal network models. Finally, we study the
signal propagation in feedforward neuronal networks consisting of both
excitatory and inhibitory neurons, and demonstrate that inhibition also plays
an important role in signal propagation in the considered networks.
Keywords: Feedforward neuronal network, unreliable synapse, signal
propagation, synfire chain, firing rate
## 1 Introduction
A major challenge in neuroscience is to understand how the neural activities
are propagated through different brain regions, since many cognitive tasks are
believed to involve this process (Vogels and Abbott, 2005). The feedforward
neuronal network is the most used model in investigating this issue, because
it is simple enough yet can explain propagation activities observed in
experiments. In recent years, two different modes of neural activity
propagation have been intensively studied. It has been found that both the
synchronous spike packet (synfire), and the firing rate, can be transmitted
across deeply layered networks (Abeles 1991; Aertsen et al. 1996; Diesmann et
al. 1999; Diesmann et al. 2001; C$\hat{\text{a}}$teau and Fukai 2001; Gewaltig
et al. 2001; Tetzlaff et al. 2002; Tetzlaff et al. 2003; van Rossum et al.
2002; Vogels and Abbott 2005; Wang et al. 2006; Aviel et al. 2003; Kumar et
al. 2008; Kumar et al. 2010; Shinozaki et al. 2007; Shinozaki et al. 2010).
Although these two propagation modes are quite different, the previous results
demonstrated that a single network with different system parameters can
support stable and robust signal propagation in both of the two modes, for
example, they can be bridged by the background noise and synaptic strength
(van Rossum et al. 2002; Masuda and Aihara 2002; Masuda and Aihara 2003).
Neurons and synapses are fundamental components of the brain. By sensing
outside signals, neurons continually fire discrete electrical signals known as
action potentials or so-called spikes, and then transmit them to postsynaptic
neurons through synapses (Dayan and Abbott 2001). The spike generating
mechanism of cortical neurons is generally highly reliable. However, many
studies have shown that the communication between neurons is, by contrast,
more or less unreliable (Abeles 1991; Raastad et al. 1992; Smetters and Zador
1996). Theoretically, the synaptic unreliability can be explained by the
phenomenon of probabilistic transmitter release (Branco and Staras 2009; Katz
1966; Katz 1969; Trommershäuser et al. 1999), i.e., synapses release
neurotransmitter in a stochastic fashion, which has been confirmed by well-
designed biological experiments (Allen and Stevens 1994). In most cases, the
transmission failure rate at a given synapse tends to exceed the fraction of
successful transmission (Rosenmund et al. 1993; Stevens and Wang 1995). In
some special cases, the synaptic transmission failure rate can be as high as
0.9 or even higher (Allen and Stevens 1994). Further computational studies
have revealed that the unreliability of synaptic transmission might be a part
of information processing of the brain and possibly has functional roles in
neural computation. For instance, it has been reported that the unreliable
synapses provide a useful mechanism for reliable analog computation in space-
rate coding (Maass and Natschl$\ddot{\text{a}}$ger 2000); and it has been
found that suitable synaptic successful transmission probability can improve
the information transmission efficiency of synapses (Goldman 2004) and can
filter the redundancy information by removing autocorrelations in spike trains
(Goldman et al. 2002). Furthermore, it has also been demonstrated that
unreliable synapses largely influence both the emergence and dynamical
behaviors of clusters in an all-to-all pulse-coupled neuronal network, and can
make the whole network relax to clusters of identical size (Friedrich and
Kinzel 2009).
Although the signal propagation in multilayered feedforward neuronal networks
has been extensively studied, to the best of our knowledge the effects of
unreliable synapses on the propagation of neural activity have not been widely
discussed and the relevant questions still remain unclear (but see the
footnote111An anonymous reviewer kindly reminded us that there might be a
relevant abstract (Trommershäuser and Diesmann 2001) discussing the effect of
synaptic variability on the synchronization dynamics in feedforward cortical
neural networks, but the abstract itself does not contain the results
presumably presented on the poster and also the follow-up publications do not
exist.). In this paper, we address these questions and provide insights by
computational modeling. For this purpose, we examine both the synfire
propagation and firing rate propagation in feedforward neuronal networks. We
mainly investigate the signal propagation in feedforward neuronal networks
composed of purely excitatory neurons connected with unreliable synapses in an
all-to-all coupling fashion (abbr. URE feedforward neuronal network) in this
work. We also compare our results with the corresponding feedforward neuronal
networks (we will clarify the meaning of “corresponding” later) composed of
purely excitatory neurons connected with reliable synapses in a random
coupling fashion (abbr. RRE feedforward neuronal network). Moreover, we study
feedforward neuronal networks consisting of both excitatory and inhibitory
neurons connected with unreliable synapses in an all-to-all coupling fashion
(abbr. UREI feedforward neuronal network).
The rest of this paper is organized as follows. The network architecture,
neuron model, and synapse model used in this paper are described in Sec. 2.
Besides these, the measures to evaluate the performance of synfire propagation
and firing rate propagation, as well as the numerical simulation method are
also introduced in this section. The main results of the present work are
presented in Sec. 3. Finally, a detailed conclusion and discussion of our work
are given in Sec. 4.
## 2 Model and method
### 2.1 Network architecture
In this subsection, we introduce the network topology used in this paper. Here
we only describe how to construct the URE feedforward neuronal network. The
methods about how to build the corresponding RRE feedforward neuronal network
and the UREI feedforward neuronal network will be briefly given in Secs. 3.4
and 3.5, respectively. The architecture of the URE feedforward neuronal
network is schematically shown in Figure 1. The network totally contains
$L=10$ layers, and each layer is composed of $N_{s}=100$ excitatory neurons.
Since neurons in the first layer are responsible for receiving and encoding
the external input signal, we therefore call this layer sensory layer and
neurons in this layer are called sensory neurons. In contrast, the function of
neurons in the other layers is to propagate neural activities. Based on this
reason, we call these layers transmission layers and the corresponding neurons
cortical neurons. Because the considered neuronal network is purely
feedforward, there is no feedback connection from neurons in downstream layers
to neurons in upstream layers, and there is also no connection among neurons
within the same layer. For simplicity, we call the $i$-th neuron in the $j$-th
layer neuron $(i,j)$ in the following.
Figure 1: Network architecture of the URE feedforward neuronal network. The
network totally contains 10 layers. The first layer is the sensory layer and
the others are the transmission layers. Each layer consists of 100 excitatory
neurons. For clarity, only 6 neurons are shown in each layer.
### 2.2 Neuron model
We now introduce the neuron model used in the present work. Each cortical
neuron is modeled by using the integrate-and-fire (IF) model (Nordlie et al.
2009), which is a minimal spiking neuron model to mimic the action potential
firing dynamics of biological neurons. The subthreshold dynamics of a single
IF neuron obeys the following differential equation:
$\begin{split}\tau_{m}\frac{dV_{ij}}{dt}=V_{\text{rest}}-V_{ij}+RI_{ij},\end{split}$
(2.1)
with the total input current
$\begin{split}I_{ij}=I_{ij}^{\text{syn}}+I_{ij}^{\text{noise}}.\end{split}$
(2.2)
Here $i=1,2,\ldots,N_{s}$ and $j=2,3,\ldots,L$, $V_{ij}$ represents the
membrane potential of neuron $(i,j)$, $\tau_{m}=20$ ms is the membrane time
constant, $V_{\text{rest}}=-60$ mV is the resting membrane potential, $R=20$
M$\Omega$ denotes the membrane resistance, and $I_{ij}^{\text{syn}}$ is the
total synaptic current. The noise current
$I_{ij}^{\text{noise}}=\sqrt{2D_{t}}\xi_{ij}(t)$ represents the external or
intrinsic fluctuations of the neuron, where $\xi_{ij}(t)$ is a Gaussian white
noise with zero mean and unit variance, and $D_{t}$ is referred to as the
noise intensity of the cortical neurons. In this work, a deterministic
threshold-reset mechanism is implemented for spike generation. Whenever the
membrane potential of a neuron reaches a fixed threshold at
$V_{\text{th}}=-50$ mV, the neuron fires a spike, and then the membrane
potential is reset according to the resting potential, where it remains
clamped for a 5-ms refractory period.
On the other hand, we use different models to simulate the sensory neurons
depending on different tasks. To study the synfire propagation, we assume that
each sensory neuron is a simple spike generator, and control their firing
behaviors by ourselves. While studying the firing rate propagation, the
sensory neuron is modeled by using the IF neuron model with the same
expression (see Eq. (2.1)) and the same parameter settings as those for
cortical neurons. For each sensory neuron, the total input current is given by
$\begin{split}I_{i1}=I(t)+I_{i1}^{\text{noise}},\end{split}$ (2.3)
where $i=1,2,...,N_{s}$ index neurons. The noise current
$I_{i1}^{\text{noise}}$ has the same form as that for cortical neurons but
with the noise intensity $D_{s}$. $I(t)$ is a time-varying external input
current which is injected to all sensory neurons. For each run of the
simulation, the external input current is constructed by the following
process. Let $\eta(t)$ denote an Ornstein-Uhlenbeck process, which is
described by
$\begin{split}\tau_{c}\frac{d\eta(t)}{dt}=-\eta(t)+\sqrt{2A}\xi(t),\end{split}$
(2.4)
where $\xi(t)$ is a Gaussian white noise with zero mean and unit variance,
$\tau_{c}$ is a correlation time constant, and $A$ is a diffusion coefficient.
The external input current $I(t)$ is defined as
$\begin{split}I(t)=\begin{cases}\eta(t)&\text{if $\eta(t)\geq 0$},\\\
0&\text{if $\eta(t)<0$}.\end{cases}\end{split}$ (2.5)
Parameter $A$ can be used to denote the intensity of the external input signal
$I(t)$. In this work, we choose $A=200$ $\text{nA}^{2}$ and $\tau_{c}=80$ ms.
By its definition, the external input current $I(t)$ corresponds to a
Gaussian-distributed white noise low-pass filtered at 80 ms and half-wave
rectified. It should be noted that this type of external input current is
widely used in the literature, in particular in the research papers which
study the firing rate propagation (van Rossum et al. 2002; Vogels and Abbott
2005; Wang and Zhou 2009).
### 2.3 Synapse model
The synaptic interactions between neurons are implemented by using the
modified conductance-based model. Our modeling methodology is inspired by the
phenomenon of probabilistic transmitter release of the real biological
synapses. Here we only introduce the model of unreliable excitatory synapses,
because the propagation of neural activity is mainly examined in URE
feedforward neuronal networks in this work. The methods about how to model
reliable excitatory synapses and unreliable inhibitory synapses will be
briefly introduced in Secs. 3.4 and 3.5, respectively.
The total synaptic current onto neuron $(i,j)$ is the linear sum of the
currents from all incoming synapses,
$\begin{split}I_{ij}^{\text{syn}}=\sum_{k=1}^{N_{s}}G(i,j;k,j-1)\cdot(E_{{\text{syn}}}-V_{ij}).\end{split}$
(2.6)
In this equation, the outer sum runs over all synapses onto this particular
neuron, $G(i,j;k,j-1)$ is the conductance from neuron $(k,j-1)$ to neuron
$(i,j)$, and $E_{\text{syn}}=0$ mV is the reversal potential of the excitatory
synapse. Whenever the neuron $(k,j-1)$ emits a spike, an increment is assigned
to the corresponding synaptic conductances according to the synaptic
reliability parameter, which process is given by
$\begin{split}G(i,j;k,j-1)\leftarrow G(i,j;k,j-1)+J(i,j;k,j-1)\cdot
h(i,j;k,j-1),\end{split}$ (2.7)
where $h(i,j;k,j-1)$ denotes the synaptic reliability parameter of the synapse
from neuron $(k,j-1)$ to neuron $(i,j)$, and $J(i,j;k,j-1)$ stands for the
relative peak conductance of this particular excitatory synapse which is used
to determine its strength. For simplicity, we assume that $J(i,j;k,j-1)=g$,
that is, the synaptic strength is identical for all excitatory connections.
Parameter $p$ is defined as the successful transmission probability of spikes.
When a presynaptic neuron $(k,j-1)$ fires a spike, we let the corresponding
synaptic reliability variables $h(i,j;k,j-1)=1$ with probability $p$ and
$h(i,j;k,j-1)=0$ with probability $1-p$. That is to say, whether the
neurotransmitter is successfully released or not is in essence controlled by a
Bernoulli on-off process in the present work. In other time, the synaptic
conductance decays by an exponential law:
$\begin{split}\frac{d}{dt}G(i,j;k,j-1)=\frac{1}{\tau_{s}}G(i,j;k,j-1),\end{split}$
(2.8)
with a fixed synaptic time constant $\tau_{s}$. In the case of synfire
propagation, we choose $\tau_{s}=2$ ms, and in the case of firing rate
propagation, we choose $\tau_{s}=5$ ms.
### 2.4 Measures of the synfire and firing rate propagation
We now introduce several useful measures used to quantitatively evaluate the
performance of the two different propagation modes: the synfire mode and
firing rate mode. The propagation of synfire activity is measured by the
survival rate and the standard deviation of the spiking times of the synfire
packet (Gewaltig et al. 2001). Let us first introduce how to calculate the
survival rate for the synfire propagation. In our simulations, we find that
the synfire propagation can be divided into three types: the failed synfire
propagation, the stable synfire propagation, as well as the synfire
instability propagation (for detail, see Sec. 3.1). For neurons in each layer,
a threshold method is developed to detect the local highest “energy” region.
To this end, we use a 5 ms moving time window with 0.1 ms sliding step to
count the number of spikes within each window. Here a high energy region means
that the number of spikes within the window is larger than a threshold
$\theta=50$. Since we use a moving time window with small sliding step, there
might be a continuous series of windows contain more than 50 spikes around a
group of synchronous spikes. In this work, we only select the first window
which covers the largest number of spikes around a group of synchronous spikes
as the local highest energy region. We use the number of local highest energy
region to determine which type of synfire propagation occurs. If there is no
local highest energy region detected in the final layer of the network, we
consider it as the failed synfire propagation. When two or more separated
local highest energy regions are detected in one layer, we consider it as the
synfire instability propagation. Otherwise, it means the occurrence of the
stable synfire propagation. For each experimental setting, we carry out the
simulation many times. The survival rate of the synfire propagation is defined
as the ratio of the number of occurrence of the stable synfire propagation to
the total number of simulations. In additional simulations, it turns out that
the threshold value $\theta$ can vary in a wide range without altering the
results. Under certain conditions, noise can help the feedforward neuronal
network produce the spontaneous spike packets, which promotes the occurrence
of synfire instability propagation and therefore decreases the survival rate.
For stable synfire propagation, there exists only one highest energy region
for neurons in each layer. Spikes within this region are considered as the
candidate synfire packet, which might also contain a few spontaneous spikes
caused by noise and other factors. In this work, an adaptive algorithm is
introduced to eliminate spontaneous spikes from the candidate synfire packet.
Suppose now that there is a candidate synfire packet in the $i$-th layer with
the number of spikes it contains $\alpha_{i}$ and the corresponding spiking
times $\\{t_{1},t_{2},\ldots,t_{\alpha_{i}}\\}$. The average spiking time of
the candidate synfire packet is therefore given by
$\begin{split}\bar{t}_{i}=\frac{1}{\alpha_{i}}\sum_{k=1}^{\alpha_{i}}t_{k}.\end{split}$
(2.9)
Thus the standard deviation of the spiking times in the $i$-th layer can be
calculated as follows:
$\begin{split}\sigma_{i}=\sqrt{\frac{1}{\alpha_{i}}\sum_{k=1}^{\alpha_{i}}[t_{k}-\bar{t}_{i}]^{2}}.\end{split}$
(2.10)
We remove the $j$-th spike from the candidate synfire packet if it satisfies:
$|t_{j}-\bar{t}_{i}|>\mu\sigma_{i}$, where $\mu$ is a parameter of our
algorithm. We recompute the average spiking time as well as the standard
deviation of the spiking times for the new candidate synfire packet, and
repeat the above eliminating process, until no spike is removed from the new
candidate synfire packet anymore. We define the remaining spikes as the
synfire packet, which is characterized by the final values of $\alpha_{i}$ and
$\sigma_{i}$. Parameter $\mu$ determines the performance of the proposed
algorithm. If $\mu$ is too large, the synfire packet will lose several useful
spikes at its borders, and if $\mu$ is too small, the synfire packet will
contain some noise data. In our simulations, we found that $\mu=4$ can result
in a good compromise between these two extremes. It should be emphasized that
our algorithm is inspired by the method given in (Gewaltig et al. 2001). Next,
we introduce how to measure the performance of the firing rate propagation.
The performance of firing rate propagation is evaluated by combining it with a
population code. Specifically, we compute how similar the population firing
rates in different layers to the external input current $I(t)$ (van Rossum et
al. 2002; Vogels and Abbott 2005). To do this, a 5 ms moving time window with
1 ms sliding step is also used to estimate the population firing rates
${r_{i}(t)}$ for different layers as well as the smooth version of the
external input current $I_{s}(t)$. The correlation coefficient between the
population firing rate of the $i$-th layer and external input current is
calculated by
$\begin{split}C_{i}(\tau)=\frac{\left\langle\left[I_{s}(k+\tau)-\overline{I}_{s}\right]\left[r_{i}(k)-\overline{r}_{i}\right]\right\rangle_{t}}{\sqrt{\left\langle\left[I_{s}(k+\tau)-\overline{I}_{s}\right]^{2}\right\rangle_{t}\left\langle\left[r_{i}(k)-\overline{r}_{i}\right]^{2}\right\rangle_{t}}},\end{split}$
(2.11)
where $\langle\cdot\rangle_{t}$ denotes the average over time. Here we use the
maximum cross-correlation coefficient $Q_{i}=\max\\{C_{i}(\tau)\\}$ to
quantify the performance of the firing rate propagation in the $i$-th layer.
Note that $Q_{i}$ is a normalization measure and a larger value corresponds to
a better performance.
### 2.5 Numerical simulation method
In all numerical simulations, we use the standard Euler-Maruyama integration
scheme to numerically calculate the aforementioned stochastic differential
Eqs. (2.1)-(2.8) (Kloeden et al. 1994). The temporal resolution of integration
is fixed at 0.02 ms for calculating the measures of the synfire mode and at
0.05 ms for calculating the measures of the firing rate mode, as the
measurement of the synfire needs higher precise. In additional simulations, we
have found that further reducing the integration time step does not change our
numerical results in a significant way. For the synfire mode, all simulations
are executed at least 100 ms to ensure that the synfire packet can be
successfully propagated to the final layer of the considered network. While
studying the firing rate mode, we perform all simulations up to 5000 ms to
collect enough spikes for statistical analysis. It should be noted that, to
obtain convincing results, we carry out several times of simulations (at least
200 times for the synfire mode and 50 times for the firing rate mode) for each
experimental setting to compute the corresponding measures.
## 3 Simulation results
In this section, we report the main results obtained in the simulation. We
first systematically investigate the signal propagation in the URE feedforward
neuronal networks. Then, we compare these results with those for the
corresponding RRE feedforward neuronal networks. Finally, we further study the
signal propagation in the UREI feedforward neuronal networks.
### 3.1 Synfire propagation in URE feedforward neuronal networks
Here we study the role of unreliable synapses on the propagation of synfire
packet in the URE feedforward neuronal networks. In the absence of noise, we
artificially let each sensory neuron fire and only fire an action potential at
the same time ($\alpha_{1}=100$ and $\sigma_{1}=0$ ms). Without loss of
generality, we let all sensory neurons fire spikes at $t=10$ ms. Figure 2
shows four typical spike raster diagrams of propagating synfire activity. Note
that the time scales in Figs. 2-2 are different. The URE feedforward neuronal
network with both small successful transmission probability and small
excitatory synaptic strength badly supports the synfire propagation. In this
case, due to high synaptic unreliability and weak excitatory synaptic
interaction between neurons, the propagation of synfire packet cannot reach
the final layer of the whole network (see Fig. 2). For suitable values of $p$
and $g$, we find that the synfire packet can be stably transmitted in the URE
feedforward neuronal network. Moreover, it is obvious that the width of the
synfire packet at any layer for $p=0.8$ is much narrower than that of the
corresponding synfire packet for $p=0.25$ (see Figs. 2 and 2). At the same
time, the transmission speed is also enhanced with the increasing of $p$.
These results indicate that the neuronal response of the considered network is
much more precise and faster for suitable large successful transmission
probability. However, our simulation results also reveal that a strong
excitatory synaptic strength with large value of $p$ might destroy the
propagation of synfire activity. As we see from Fig. 2, the initial tight
synfire packet splits into several different synfire packets during the
transmission process. Such phenomenon is called the “synfire instability”
(Tetzlaff et al. 2002; Tetzlaff et al. 2003), which mainly results from the
burst firings of several neurons caused by the strong excitatory synaptic
interaction as well as the stochastic fluctuation of the synaptic connections.
Figure 2: (Color online) Several typical spike raster diagrams for different
values of successful transmission probability $p$ and excitatory synaptic
strength $g$. Shown are samples of (a) failed synfire propagation, (b) and (c)
stable synfire propagation, and (d) synfire instability. System parameters are
$p=0.19$ and $g=1$ nS (a), $p=0.25$ and $g=1$ nS (b), $p=0.8$ and $g=1$ nS
(c), and $p=0.7$ and $g=2$ nS (d), respectively. As we see, the synfire packet
can reach the final layer of the network successfully only for appropriate
values of $p$ and $g$. It should be noted that the time scales in these four
subfigures are different.
Figure 3: (Color online) Effects of the successful transmission probability
and excitatory synaptic strength on the synfire propagation in URE feedforward
neuronal networks. (a) The survival rate of the synfire propagation versus $p$
for different values of $g$. (b) Schematic of three different synfire
propagation regimes in the $(g,p)$ panel ($41\times 41=1681$ points). Regime
I: the failed synfire propagation region; regime II: the stable synfire
propagation region; and regime III: the synfire instability propagation
region. (c) The value of $\sigma_{10}$ as a function of $p$ for different
values of $g$. (d) The value of $\sigma_{10}$ as a function of $g$ for
different values of $p$. In all cases, the noise intensity $D_{t}=0$. Each
data point shown here is computed based on 200 independent simulations with
different random seeds.
In Fig. 3, we depict the survival rate of synfire propagation as a function of
the successful transmission probability $p$ for different values of excitatory
synaptic strength $g$, with the noise intensity $D_{t}=0$. We find that each
survival rate curve can be at least characterized by one corresponding
critical probability $p_{\text{on}}$. For small $p$, due to low synaptic
reliability, any synfire packet cannot reach the final layer of the URE
feedforward neuronal network. Once the successful transmission probability $p$
exceeds the critical probability $p_{\text{on}}$, the survival rate rapidly
transits from 0 to 1, suggesting that the propagation of synfire activity
becomes stable for a suitable high synaptic reliability. On the other hand,
besides the critical probability $p_{\text{on}}$, we find that the survival
rate curve should be also characterized by another critical probability
$p_{\text{off}}$ if the excitatory synaptic strength is sufficiently strong
(for example, $g=3.5$ nS in Fig. 3). In this case, when $p\geq
p_{\text{off}}$, our simulation results show that the survival rate rapidly
decays from 1 to 0, indicating that the network fails to propagate the stable
synfire packet again. However, it should be noted that this does not mean that
the synfire packet cannot reach the final layer of the network in this
situation, but because the excitatory synapses with both high reliability and
strong strength lead to the occurrence of the redundant synfire instability in
transmission layers.
To systematically establish the limits for the appearance of stable synfire
propagation as well as to check that whether our previous results can be
generalized within a certain range of parameters, we further calculate the
survival rate of synfire propagation in the $(g,p)$ panel, which is shown in
Fig. 3. As we see, the whole $(g,p)$ panel can be clearly divided into three
regimes. These regimes include the failed synfire propagation regime (regime
I), the stable synfire propagation regime (regime II), and the synfire
instability propagation regime (regime III). Our simulation results reveal
that transitions between these regimes are normally very fast and therefore
can be described as a sharp transition. The data shown in Fig. 3 further
demonstrate that synfire propagation is determined by the combination of both
the successful transmission probability and excitatory synaptic strength. For
a lower synaptic reliability, the URE feedforward neuronal network might need
a larger $g$ to support the stable propagation of synfire packet.
In reality, not only the survival rate of the synfire propagation but also its
performance is largely influenced by the successful transmission probability
and the strength of the excitatory synapses. In Figs. 3 and 3, we present the
standard deviation of the spiking times of the output synfire packet
$\sigma_{10}$ for different values of $p$ and $g$, respectively. Note that
here we only consider parameters $p$ and $g$ within the stable synfire
propagation regime. The results illustrated in Fig. 3 clearly demonstrate that
the propagation of synfire packet shows a better performance for a suitable
higher synaptic reliability. For the ideal case $p=1$, the URE feedforward
neuronal network even has the capability to propagate the perfect synfire
packet ($\alpha_{i}=100$ and $\sigma_{i}=0$ ms) in the absence of noise. On
the other hand, it is also found that for a fixed $p$ the performance of
synfire propagation becomes better and better as the value of $g$ is increased
(see Fig. 3). The above results indicate that both high synaptic reliability
and strong excitatory synaptic strength are able to help the URE feedforward
neuronal network maintain the precision of neuronal response in the stable
synfire propagation regime.
Figure 4: (Color online) Dependence of the synfire propagation on the
parameters of the initial spike packet. Here we display the survival rate of
the synfire propagation versus $\alpha_{1}$ (a) and $\sigma_{1}$ (b) for
different successful transmission probabilities, and the values of $n_{i}$ (c)
and $\sigma_{i}$ (d) as a function of the layer number for different initial
spike packets and different intrinsic parameters of the network, respectively.
Note that the vertical axis in (d) is a log scale. In all cases, the noise
intensity $D_{t}=0$. Other parameters are $g=1$ nS and $\sigma_{1}=3$ ms (a),
$g=1$ nS and $\alpha_{1}=70$ (b). Each data shown in (a) and (b) is computed
based on 600 independent simulations, whereas each data shown in (c) and (d)
is calculated based on 500 independent successful synfire propagation
simulations.
Up to now, we only use the perfect initial spike packet ($\alpha_{1}=100$ and
$\sigma_{1}=0$ ms) to evoke the synfire propagation. This is a special case
which is simplified for analysis, but it is not necessary to restrict this
condition. To understand how a generalized spike packet is propagated through
the URE feedforward neuronal network, we randomly choose $\alpha_{1}$ neurons
from the sensory layer, and let each of these neurons fire and only fire a
spike at any moment according to a Gaussian distribution with the standard
deviation $\sigma_{1}$. In Figs. 4 and 4, we plot the survival rate of the
synfire propagation as a function of $\alpha_{1}$ and $\sigma_{1}$ for four
different values of successful transmission probability, respectively. When
the successful transmission probability is not too large (for example,
$p=0.24$, 0.27, and 0.3 in Figs. 4 and 4), the synfire activity is well build
up after several initial layers for sufficiently strong initial spike packet
(large $\alpha_{1}$ and small $\sigma_{1}$), and then this activity can be
successfully transmitted along the entire network with high survival rate. In
this case, too weak initial spike packet (small $\alpha_{1}$ and large
$\sigma_{1}$) leads to the propagation of the neural activities becoming
weaker and weaker with the increasing of layer number. Finally, the neural
activities are stopped before they reach the final layer of the network.
Moreover, with the increasing of the successful transmission probability,
neurons in the downstream layers will share more common synaptic currents from
neurons in the corresponding upstream layers. This means that neurons in the
considered network have the tendency to fire more synchronously for suitable
larger $p$ (not too large). On the other hand, for sufficiently high synaptic
reliability (for instance, $p=0.6$ in Figs. 4 and 4), a large $\alpha_{1}$ or
a suitable large $\sigma_{1}$ may result in the occurrence of synfire
instability, which also reduces the survival rate of the synfire propagation.
Therefore, for a fixed $g$, the URE feedforward neuronal network with suitable
higher synaptic reliability has the ability to build up stable synfire
propagation from a slightly weaker initial spike packet (see Figs. 4 and 4).
Figures 4 and 4 illustrate the values of $\alpha_{i}$ and $\sigma_{i}$ versus
the layer number for different initial spike packets and several different
intrinsic system parameters of the network (the successful transmission
probability $p$ and excitatory synaptic strength $g$). For each case shown in
Figs. 4 and 4, once the synfire propagation is successfully established,
$n_{i}$ converges fast to the saturated value 100 and $\sigma_{i}$ approaches
to an asymptotic value. Although the initial spike packet indeed determines
whether the synfire propagation can be established or not as well as
influences the performance of synfire propagation in the first several layers,
but it does not determine the value of $\sigma_{i}$ in deep layers provided
that the synfire propagation is successfully evoked. For the same intrinsic
system parameters, if we use different initial spike packets to evoke the
synfire propagation, the value of $\sigma_{i}$ in deep layers is almost the
same for different initial spike packets (see Fig. 4). The above results
indicate that the performance of synfire propagation in deep layers of the URE
feedforward neuronal network is quite stubborn, which is mainly determined by
the intrinsic parameters of the network but not the parameters of the initial
spike packet. In fact, many studies have revealed that the synfire activity is
governed by a stable attractor in the $(\alpha,\sigma)$ space (Diesmann et al.
1999; Diesmann et al. 2001; Diesmann 2002; Gewaltig et al. 2001). Our above
finding is a signature that the stable attractor of synfire propagation does
also exist for the feedforward neuronal networks with unreliable synapses.
Figure 5: (Color online) Effects of the noise intensity $D$ on the synfire
propagation. (a) The survival rate versus successful transmission probability
$p$ for different excitatory coupling strength and noise intensities. (b) The
value of $\sigma_{10}$ versus the noise intensity $D$ for different successful
transmission probabilities, with the excitatory synaptic strength $g=1.5$ nS.
In all cases, the parameters of the initial spike packet are $\alpha_{1}=100$
and $\sigma_{1}=0$ ms. Each data shown here is computed based on 200
independent simulations with different random seeds.
Next, we study the dependence of synfire propagation on neuronal noise. It is
found that both the survival rate of synfire propagation and its performance
are largely influenced by the noise intensity. There is no significant
qualitative difference between the corresponding survival rate curves in low
synaptic reliable regime. However, we find important differences between these
curves for small $g$ in high synaptic reliable regime, as well as for large
$g$ in intermediate synaptic reliable regime, that is, during the transition
from the successful synfire regime to the synfire instability propagation
regime (see from Fig. 5). For each case, it is obvious that the top region of
the survival rate becomes smaller with the increasing of noise intensity. This
is at least due to the following two reasons: (i) noise makes neurons
desynchronize, thus leading to a more dispersed synfire packet in each layer.
For relatively high synaptic reliability, a dispersed synfire packet has the
tendency to increase the occurrence rate of the synfire instability. (ii)
Noise with large enough intensity results in several spontaneous neural firing
activities at random moments, which also promote the occurrence of the synfire
instability. Figure 5 presents the value of $\sigma_{10}$ as a function of the
noise intensity $D_{t}$ for different values of successful transmission
probability $p$. As we see, the value of $\sigma_{10}$ becomes larger and
larger as the noise intensity is increased from 0 to 0.1 (weak noise regime).
This is also due to the fact that the existence of noise makes neurons
desynchronize in each layer. However, although noise tends to reduce the
synchrony of synfire packet, the variability of $\sigma_{i}$ in deep layers is
quite low (data not shown). The results suggest that, in weak noise regime,
the synfire packet can be stably transmitted through the feedforward neuronal
network with small fluctuation in deep layers, but displays slightly worse
performance compared to the case of $D_{t}=0$. Further increase of noise will
cause many spontaneous neural firing activities which might significantly
deteriorate the performance of synfire propagation. However, it should be
emphasized that, although the temporal spread of synfire packet tends to
increase as the noise intensity grows, several studies have suggested that
under certain conditions the basin of attraction of synfire activity reaches a
maximum extent (Diesmann 2002; Postma et al. 1996; Boven and Aertsen 1990).
Such positive effect of noise can be compared to a well known phenomenon
called aperiodic stochastic resonance (Collins et al. 1995b; Collins et al.
1996; Diesmann 2002).
### 3.2 Firing rate propagation in URE feedforward neuronal networks
In this subsection, we examine the firing rate propagation in URE feedforward
neuronal networks. To this end, we assume that all sensory neurons are
injected to a same time-varying external current $I(t)$ (see Sec. 2.2 for
detail). Note now that the sensory neurons are modeled by using the integrate-
and-fire neuron model in the study of the firing rate propagation.
Figure 6: The maximum cross-correlation coefficient between the smooth version
of external input current $I_{s}(t)$ and the population firing rate of sensory
neurons $r_{1}(t)$ for different noise intensities.
Figure 7: (Color online) Impacts of noise on the encoding performance (the
firing rate mode) of sensory neurons. For each case, the smooth version of the
external input current $I_{s}(t)$ (top panel), spike raster diagram of sensory
neurons (middle panel), and population firing rate of sensory neurons (bottom
panel) are shown. Noise intensities are $D_{s}=0$ (a), $D_{s}=0.1$ (b),
$D_{s}=0.6$ (c), $D_{s}=0.8$ (d), $D_{s}=2$ (e), and $D_{s}=5$ (f),
respectively.
Before we present the results of the firing rate propagation, let us first
investigate how noise influences the encoding capability of sensory neurons by
the population firing rate. This is an important preliminary step, because how
much input information represented by sensory neurons will directly influence
the performance of firing rate propagation. The corresponding results are
plotted in Figs. 6 and 7, respectively. When the noise is too weak, the
dynamics of sensory neurons is mainly controlled by the same external input
current, which causes neurons to fire spikes almost at the same time (see
Figs. 7 and 7). In this case, the information of the external input current is
poorly encoded by the population firing rate since the synchronous neurons
have the tendency to redundantly encode the same aspect of the external input
signal. When the noise intensity falls within a special intermediate range
(about 0.5-10), neuronal firing is driven by both the external input current
and noise. With the help of noise, the firing rate is able to reflect the
temporal structural information (i.e., temporal waveform) of the external
input current to a certain degree (see Figs. 7 to 7), and therefore $Q_{1}$
has large value in this situation. For too large noise intensity, the external
input current is almost drowned in noise, thus resulting that the input
information cannot be well read from the population firing rate of sensory
neurons again. On the other hand, sensory neurons can fire “false” spikes
provided that they are driven by sufficiently strong noise (as for example at
$t\approx 2800$ ms in Fig. 7). Although the encoding performance of the
sensory neurons might be good enough in this case, our numerical simulations
reveal that such false spikes will seriously reduce the performance of the
firing rate propagation in deep layers, which will be discussed in detail in
the later part of this section. By taking these factors into account, we
consider the noise intensity of sensory neurons to be within the range of 0.5
to 1 in the present work.
Figure 8: An example of the firing rate propagation in the URE feedforward
network. Here we show the smooth version of the external input current
$I_{s}(t)$, as and the population firing rates of layers 1, 2, 4, 6, 8, and
10, respectively. System parameters are $g=0.4$ nS, $p=0.2$, and
$D_{s}=D_{t}=0.7$.
Figure 8 shows a typical example of the firing rate propagation. In view of
the overall situation, the firing rate can be propagated rapidly and basically
linearly in the URE feedforward neuronal network. However, it should be noted
that, although the firing rates of neurons from the downstream layers tend to
track those from the upstream layers, there are still several differences
between the firing rates for neurons in two adjacent layers. For example, it
is obvious that some low firing rates may disappear or be slightly amplified
in the first several layers, as well as some high firing rates are weakened to
a certain degree during the whole transmission process. Therefore, as the
neural activities are propagated across the network, the firing rate has the
tendency to lose a part of local detailed neural information but can maintain
a certain amount of global neural information. As a result, the maximum cross-
correlation coefficient between $I_{s}(t)$ and $r_{i}(t)$ basically drops with
the increasing of the layer number.
Figure 9: (Color online) Effects of the unreliable synapses on the
performance of firing rate propagation. (a) The value of $Q_{10}$ as a
function of $p$ for different values of excitatory synaptic strength. (b) The
value of $Q_{10}$ as a function of $g$ for different values of successful
transmission probability. Noise intensities are $D_{s}=D_{t}=0.5$ in all
cases. Here each data point is computed based on 50 different independent
simulations with different random seeds.
Let us now assess the impacts of the unreliable synapses on the performance of
firing rate propagation in the URE feedforward neuronal network. Figure 9
presents the value of $Q_{10}$ versus the success transmission probability $p$
for various excitatory synaptic strengths. For a fixed value of $g$, a bell-
shaped $Q_{10}$ curve is clearly seen by changing the value of successful
transmission probability, indicating that the firing rate propagation shows
the best performance at an optimal synaptic reliability level. This is
because, for each value of $g$, a very small $p$ will result in the
insufficient firing rate propagation due to low synaptic reliability, whereas
a sufficiently large $p$ can lead to the excessive propagation of firing rate
caused by burst firings. Based on above reasons, the firing rate can be well
transmitted to the final layer of the URE feedforward neuronal network only
for suitable intermediate successful transmission probabilities. Moreover,
with the increasing of $g$, the considered network needs a relatively small
$p$ to support the optimal firing rate propagation. In Fig. 9, we plot the
value of $Q_{10}$ as a function of the excitatory synaptic strength $g$ for
different values of $p$. Here the similar results as those shown in Fig. 9 can
be observed. This is due to the fact that increasing $g$ and fixing the value
of $p$ is equivalent to increasing $p$ and fixing the value of $g$ to a
certain degree. According to the aforementioned results, we conclude that both
the successful transmission probability and excitatory synaptic strength are
critical for firing rate propagation in URE feedforward networks, and better
choosing of these two unreliable synaptic parameters can help the cortical
neurons encode neural information more accurately.
Figure 10: (Color online) Impacts of noise on the performance of firing rate
propagation. (a) The value of $Q_{10}$ as a function of the noise intensity of
cortical neurons $D_{t}$ for different values of $D_{s}$. (b) The performance
of firing rate propagation in each layer at $D_{t}=0.6$ for two different
noise intensities of sensory neurons. In all cases, $p=0.2$ and $g=0.4$ nS.
Here each data point is computed based on 50 different independent simulations
with different random seeds.
Next, we examine the dependence of the firing rate propagation on neuronal
noise. The corresponding results are plotted in Figs. 10 and 11, respectively.
Figure 10 demonstrates that the noise of cortical neurons plays an important
role in firing rate propagation. Noise of cortical neurons with appropriate
intensity is able to enhance their encoding accuracy. It is because
appropriate intermediate noise, on the one hand, prohibits synchronous firings
of cortical neurons in deep layers, and on the other hand, ensures that the
useful neural information does not drown in noise. However, the level of
enhancement is largely influenced by the noise intensity of sensory neurons.
As we see, for a large value of $D_{s}$, such enhancement is weakened to a
great extent. This is because slightly strong noise intensity of sensory
neurons will cause these neurons to fire several false spikes and a part of
these spikes can be propagated to the transmission layers. If enough false
spikes appear around the weak components of the external input current, these
spikes will help the network abnormally amplify these weak components during
the whole transmission process. The aforementioned process can be seen clearly
from an example shown in Fig. 11. As a result, the performance of the firing
rate propagation might be seriously deteriorated in deep layers. However, it
should be noted that this kind of influence typically needs the accumulation
of several layers. Our simulation results show that the performance of firing
rate propagation can be well maintained or even becomes slightly better
(depending on the noise intensity of sensory neurons, see Fig. 6) in the first
several layers for large $D_{s}$ (see Fig. 10). In fact, the above results are
based on the assumption that each cortical neuron is driven by independent
noise current with the same intensity. Our results can be generalized from the
sensory layer to the transmission layers if we suppose that noise intensities
for neurons in different transmission layers are different. All these results
imply that better tuning of the noise intensities of both the sensory and
cortical neurons can enhance the performance of firing rate propagation in the
URE feedforward neuronal network.
Figure 11: (Color online) An example of weak external input signal
amplification. System parameters are the successful transmission probability
$p=0.2$, excitatory synaptic strength $g=0.4$ nS, and noise intensities
$D_{s}=2.5$ and $D_{t}=0.6$, respectively.
### 3.3 Stochastic effect of neurotransmitter release
From the numerical results depicted in Secs. 3.1 and 3.2, we find that
increasing $g$ with $p$ fixed has similar effects as increasing $p$ while
keeping $g$ fixed for both the synfire mode and firing rate mode. Some persons
might therefore postulate that the signal propagation dynamics in feedforward
neuronal networks with unreliable synapses can be simply determined by the
average amount of received neurotransmitter for each neuron in a time instant,
which can be reflected by the product of $g\cdot p$. To check whether this is
true, we calculate the measures of these two signal propagation modes as a
function of $g\cdot p$ for different successful transmission probabilities. If
this postulate is true, the URE feedforward neuronal network will show the
same propagation performance for different values of $p$ at a fixed $g\cdot
p$. Our results shown in Figs. 12-12 clearly demonstrate that the signal
propagation dynamics in the considered network can not be simply determined by
the product $g\cdot p$ or, equivalently, by the average amount of received
neurotransmitter for each neuron in a time instant. For both the synfire
propagation and firing rate propagation, although the propagation performance
exhibits the similar trend with the increasing of $g\cdot p$, the
corresponding measure curves do not superpose in most parameter region for
each case, and in some parameter region the differences are somewhat
significant (see Figs. 12 and 12). This is because of the stochastic effect of
neurotransmitter release, that is, the unreliability of neurotransmitter
release will add randomness to the system. Different successful transmission
probabilities may introduce different levels of randomness, which will further
affect the nonlinear spiking dynamics of neurons. Therefore, the URE
feedforward neuronal network might display different propagation performance
for different values of $p$ even at a fixed $g\cdot p$. If we set the value of
$g\cdot p$ constant, a low synaptic reliability will introduce large
fluctuations in the synaptic inputs. For small $p$, according to the above
reason, some neurons will fire spikes more than once in the large $g\cdot p$
regime. This mechanism increases the occurrence rate of the synfire
instability. Thus, the URE feedforward neuronal network has the tendency to
stop the stable synfire propagation for a small synaptic transmission
probability (see Fig. 12). On the other hand, a high synaptic reliability will
introduce small fluctuations in the synaptic inputs for a fixed $g\cdot p$.
This makes neurons in the considered network fire spikes almost synchronously
for a large $p$, thus resulting the worse performance for the firing rate
propagation in large $g\cdot p$ regime (see Fig. 12). Our above results
suggest that the performance of the signal propagation in feedforward neuronal
networks with unreliable synapses is not only purely determined by the change
of synaptic parameters, but also largely influenced by the stochastic effect
of neurotransmitter release.
Figure 12: (Color online) Dependence of signal propagation dynamics on the
product of $g\cdot p$ in the URE feedforward neuronal network. _Synfire_ mode:
survival rate (a) and $\sigma_{10}$ versus $g\cdot p$ with $D_{t}=0$. The
parameters of the initial spike packet are $\alpha_{1}=100$ and $\sigma_{1}=0$
ms. _Firing rate_ mode: $Q_{10}$ as a function of $g\cdot p$ with
$D_{s}=D_{t}=0.5$. Here each data point shown in (a) and (b) is calculated
based on 200 different independent simulations, whereas each data point shown
in (c) are based on 50 different independent simulations.
### 3.4 Comparison with corresponding RRE feedforward neuronal networks
In this subsection, we make comparisons on the propagation dynamics between
the URE and the RRE feedforward networks. We first introduce how to generate a
corresponding RRE feedforward neuronal network for a given URE feedforward
neuronal network. Suppose now that there is a URE feedforward neuronal network
with successful transmission probability $p$. A corresponding RRE feedforward
neuronal network is constructed by using the connection density $p$ (on the
whole), that is, a synapse from one neuron in the upstream layer to one neuron
in the corresponding downstream layer exists with probability $p$. As in the
URE feedforward neuronal network given in Sec. 2.1, there is no feedback
connection from downstream neurons to upstream neurons and also no connection
among neurons within the same layer in the RRE feedforward neuronal network.
It is obvious that parameter $p$ has different meanings in these two different
feedforward neuronal network models. The synaptic interactions between neurons
in the RRE feedforward neuronal network are also implemented by using the
conductance-based model (see Eqs. (2.6) and (2.7) for detail). However, here
we remove the constraint of the synaptic reliability parameter for the RRE
feedforward neuronal network, e.g., $h(i,j;k,j-1)=1$ in all cases. A naturally
arising question is what are the differences, if have, between the synfire
propagation and firing propagation in URE feedforward neuronal networks and
those in RRE feedforward neuronal networks, although the numbers of active
synaptic connections that taking part in transmitting spikes in a time instant
are the same from the viewpoint of mathematical expectation.
Figure 13: (Color online) The difference between the synfire propagation in
the URE feedforward neuronal network and the RRE feedforward neuronal network.
Here we show the value of survival rate as a function of $\sigma_{1}$ for
different network models. In all cases, $D_{t}=0$ and $\alpha_{1}=70$. Other
system parameters are $g=4$ nS and $p=0.15$ (dot: “$\bullet$”, and circle:
“$\circ$” ), and $g=3.5$ nS and $p=0.12$ (square: “$\square$”, and asterisk:
“$\ast$”). Each data point is calculated based on 500 different independent
simulations with different random seeds.
For the synfire propagation, our simulation results indicate that, compared to
the RRE feedforward neuronal network, the URE feedforward neuronal network is
able to suppress the occurrence of synfire instability to a certain degree,
which can be seen clearly in Fig. 13. Typically, this phenomenon can be
observed in strong excitatory synaptic strength regime. Due to the
heterogeneity of connectivity, some neurons in the RRE feedforward neuronal
network will have more input synaptic connections than the other neurons in
the same network. For large value of $g$, these neurons tend to fire spikes
very rapidly after they received synaptic currents. If the width of the
initial spike packet is large enough, these neurons might fire spikes again
after their refractory periods, which are induced by a few spikes from the
posterior part of the dispersed initial spike packet. These spikes may
increase the occurrence rate of the synfire instability. While in the case of
URE feedforward neuronal network, the averaging effect of unreliable synapses
tends to prohibit neurons fire spikes too quickly. Therefore, under the
equivalent parameter conditions, less neurons can fire two or more spikes in
the URE feedforeard neuronal network. As a result, the survival rate of the
synfire propagation for the URE feedforeard neuronal network is larger than
that for the RRE feedforward neuronal network (see Fig. 13), though not so
significant.
In further simulations, we find interesting results in small $p$ regime for
the firing rate propagation. Compared to the case of the URE feedforward
neuronal network, the RRE feedforward neuronal network can better support the
firing rate propagation in this small $p$ regime for strong excitatory
synaptic strength (see Fig. 14). It is because the long-time averaging effect
of unreliable synapses at small $p$ tends to make neurons fire more
synchronous spikes in the URE feedforward neuronal network through the
homogenization process of synaptic currents. However, with the increasing of
$p$, neurons in the downstream layers have the tendency to share more common
synaptic currents from neurons in the corresponding upstream layers for both
types of feedforward neuronal networks. The aforementioned factor makes the
difference of the performance of firing rate propagation between these two
types of feedforward neuronal networks become small so that the $Q_{10}$
curves almost coincide with each other for the case of $p=0.6$ (see Fig. 14).
Although from the above results we can not conclude that unreliable synapses
have advantages and play specific functional roles in signal propagation, not
like those results shown in the previous studies (Goldman et al. 2002; Goldman
2004), at least it is shown that the signal propagation activities are
different in URE and RRE to certain degrees. We should be cautioned when using
random connections to replace unreliable synapses in modelling research.
However, it should be noted that the RRE feedforward neuronal network
considered here is just one type of diluted feedforward neuronal networks.
There exists several other possibilities to construct the corresponding
diluted feedforward neuronal networks (Hehl et al. 2001). The similar
treatments for these types of diluted feedforward neuronal networks require
further investigation.
Figure 14: (Color online) Firing rate propagation in URE feedforward neuronal
network and RRE feedforward neuronal network. The value of $Q_{10}$ as a
function of excitatory synaptic strength $g$ for $p=0.2$ (a) and $p=0.6$ (b),
respectively. Noise intensities are $D_{t}=D_{s}=0.5$. Each data point is
calculated based on 50 different independent simulations with different random
seeds.
### 3.5 Signal propagation in UREI feedforward neuronal networks
In this subsection, we further study the signal propagation in the feedforward
neuronal networks composed of both excitatory and inhibitory neurons connected
in an all-to-all coupling fashion (i.e., the UREI feedforward neuronal
networks). This study is necessary because real biological neuronal networks,
especially mammalian neocortex, consist not only of excitatory neurons but
also of inhibitory neurons. The UREI feedforward neuronal network studied in
this subsection has the same topology as that shown in Fig. 1. In simulations,
we randomly choose 80 neurons in each layer as excitatory and the rest of them
as inhibitory, as the ratio of excitatory to inhibitory neurons is about $4:1$
in mammalian neocortex. The dynamics of the unreliable inhibitory synapse is
also modeled by using Eqs. (2.6) and (2.7). The reversal potential of the
inhibitory synapse is fixed at -75 mV, and its strength is set as $J=K\cdot
g$, where $K$ is a scale factor used to control the relative strength of
inhibitory and excitatory synapses. Since the results of the signal
propagation in UREI feedforward neuronal networks are quite similar to those
in URE feedforward neuronal networks, we omit most of them and only discuss
the effects of inhibition in detail.
Figure 15: (Color online) Partition of three different synfire propagation
regimes in the $(K,p)$ panel ($41\times 41=1681$ points). Regime I: the failed
synfire propagation region; regime II: the stable synfire propagation region;
and regime III: the synfire instability propagation region. we set $g=2.5$ nS
(a), $g=3$ nS (b), and $g=3.5$ nS (c). In all cases, the parameters of the
initial spike packet are $\alpha_{1}=100$ and $\sigma_{1}=0$ ms, and the noise
intensity is $D_{t}=0$. Each data point shown here is calculated based on 200
different independent simulations with different random seeds.
Figure 15 shows the survival rate of synfire propagation in the $(K,p)$ panel
for three different excitatory synaptic strengths. Depending on whether the
synfire packet can be successfully and stably transmitted to the final layer
of the UREI feedforward neuronal network, the whole $(K,p)$ panel can also be
divided into three regimes. For each considered case, the network with both
small successful transmission probability and strong relative strength of
inhibitory and excitatory synapses (failed synfire regime) prohibits the
stable propagation of the synfire activity. While in the case of high synaptic
reliability and small $K$ (synfire instability propagation regime), the
synfire packet also cannot be stably transmitted across the whole network due
to the occurrence of synfire instability. Therefore, the UREI feedforward
neuronal network is able to propagate the synfire activity successfully in a
stable way only for suitable combination of parameters $p$ and $K$. Moreover,
due to the competition between excitation and inhibition, the transitions
between these different regimes cannot be described as a sharp transition
anymore, in particular, for large scale factor $K$. Our results suggest that
such non-sharp character is strengthen with the increasing of $g$. On the
other hand, the partition of these different propagation regimes depends not
only on parameters $p$ and $K$ but also on the excitatory synaptic strength
$g$. As the value of $g$ is decreased, both the synfire instability
propagation regime and stable synfire propagation regime are shifted to the
upper left of the $(K,p)$ panel at first, and then disappear one by one (data
not shown). In contrast, a strong excitatory synaptic strength has the
tendency to extend the areas of the synfire instability propagation regime,
and meanwhile makes the stable synfire propagation regime move to the lower
right of the $(K,p)$ panel.
Figure 16: (Color online) Effect of inhibition on firing rate propagation.
Here we show the value of $Q_{10}$ as a function of scale factor $K$ for
different excitatory synaptic strengths. System parameters are $p=0.2$, and
$D_{s}=D_{t}=0.6$ in all cases. Each data point is calculated based on 50
different independent simulations with different random seeds.
For the case of firing rate propagation, we plot the value of $Q_{10}$ versus
the scale factor $K$ for different excitatory synaptic strengths in Fig. 16,
with a fixed successful transmission probability $p=0.2$. When the excitatory
synaptic strength is small (for instance $g=0.4$ nS), due to weak excitatory
synaptic interaction between neurons the UREI feedforward neuronal network
cannot transmit the firing rate sufficiently even for $K=0$. In this case,
less and less neural information can be propagated to the final layer of the
considered network with the increasing of $K$. Therefore, $Q_{10}$
monotonically decreases with the scale factor $K$ at first and finally
approaches to a low steady state value. Note that here the low steady state
value is purely induced by the spontaneous neural firing activities, which are
caused by the additive Gaussian white noise. As the excitatory synaptic
strength grows, more neural information can be successfully transmitted for
small value of $K$. When $g$ is increased to a rather large value, such as
$g=0.6$ nS, the coupling is so strong that too small scale factor will lead to
the excessive propagation of firing rate. However, in this case, the
propagation of firing rate can still be suppressed provided that the relative
strength of inhibitory and excitatory synapses is strong enough. As a result,
there always exists an optimal scale factor to best support the firing rate
propagation for each large excitatory synaptic strength (see Fig. 16). If we
fix the value of $g$ (not too small), then the similar results can also be
observed by changing the scale factor for a large successful transmission
probability (data not shown). Once again, this is due to the fact that
increasing $g$ and fixing $p$ is equivalent to increasing $p$ and fixing $g$
to a certain degree.
## 4 Conclusion and discussion
The feedforward neuronal network provides us an effective way to examine the
neural activity propagation through multiple brain regions. Although
biological experiments suggested that the communication between neurons is
more or less unreliable (Abeles 1991; Raastad et al. 1992; Smetters and Zador
1996), so far most relevant computational studies only considered that neurons
transmit spikes based on reliable synaptic models. In contrast to these
previous work, we took a different approach in this work. Here we first built
a 10-layer feedforward neuronal network by using purely excitatory neurons,
which are connected with unreliable synapses in an all-to-all coupling
fashion, that is, the so-called URE feedforward neuronal network in this
paper. The goal of this work was to explore the dependence of both the synfire
propagation and firing rate propagation on unreliable synapses in the URE
neuronal network, but was not limited this type of feedforward neuronal
network. Our modelling methodology was motivated by experimental results
showing the probabilistic transmitter release of biological synapses (Branco
and Staras 2009; Katz 1966; Katz 1969; Trommershäuser et al. 1999).
In the study of synfire mode, it was observed that the synfire propagation can
be divided into three types (i.e., the failed synfire propagation, the stable
synfire propagation, and the synfire instability propagation) depending on
whether the synfire packet can be successfully and stably transmitted to the
final layer of the considered network. We found that the stable synfire
propagation only occurs in the suitable region of system parameters (such as
the successful transmission probability and excitatory synaptic strength). For
system parameters within the stable synfire propagation regime, it was found
that both high synaptic reliability and strong excitatory synaptic strength
are able to support the synfire propagation in feedforward neuronal networks
with better performance and faster transmission speed. Further simulation
results indicated that the performance of synfire packet in deep layers is
mainly influenced by the intrinsic parameters of the considered network but
not the parameters of the initial spike packet, although the initial spike
packet determines whether the synfire propagation can be evoked to a great
degree. In fact, this is a signature that the synfire activity is governed by
a stable attractor, which is in agreement with the results given in (Diesmann
et al. 1999; Diesmann et al. 2001; Diesmann 2002; Gewaltig et al. 2001).
In the study of firing rate propagation, our simulation results demonstrated
that both the successful transmission probability and the excitatory synaptic
strength are critical for firing rate propagation. Too small successful
transmission probability or too weak excitatory synaptic strength results in
the insufficient firing rate propagation, whereas too large successful
transmission probability or too strong excitatory synaptic strength has the
tendency to lead to the excessive propagation of firing rate. Theoretically
speaking, better tuning of these two synaptic parameters can help neurons
encode the neural information more accurately.
On the other hand, neuronal noise is ubiquitous in the brain. There are many
examples confirmed that noise is able to induce many counterintuitive
phenomena, such as stochastic resonance (Collins et al. 1995a; Collins et al.
1995b; Collins et al. 1996; Chialvo et al. 1997; Guo and Li 2009) and
coherence resonance (Pikovsky and Kurths 1996; Lindner and Schimansky-Geier
2002; Guo and Li 2009). Here we also investigated how the noise influences the
performance of signal propagation in URE feedforward neuronal networks. The
numerical simulations revealed that noise tends to reduce the performance of
synfire propagation because it makes neurons desynchronized and causes some
spontaneous neural firing activities. Further studies demonstrated that the
survival rate of synfire propagation is also largely influenced by the noise.
In contrast to the synfire propagation, our simulation results showed that
noise with appropriate intensity is able to enhance the performance of firing
rate propagation in URE feedforward neuronal networks. In essence, it is
because suitable noise can help neurons in each layer maintain more accurate
temporal structural information of the the external input signal. Note that
the relevant mechanisms about noise have also been discussed in several
previous work (van Rossum et al. 2002; Masuda and Aihara 2002; Masuda and
Aihara 2003), and our results are consistent with the findings given in these
work.
Furthermore, we have also investigated the stochastic effect of
neurotransmitter release on the performance of signal propagation in the URE
feedforward neuronal networks. For both the synfire propagation and firing
rate propagation, we found that the URE feedforward neuronal networks might
display different propagation performance, even when their average amount of
received neurotransmitter for each neuron in a time instant remains unchanged.
This is because the unreliability of neurotransmitter release will add
randomness to the system. Different synaptic transmission probabilities will
introduce different levels of stochastic effect, and thus might lead to
different spiking dynamics and propagation performance. These findings
revealed that the signal propagation dynamics in feedforward neuronal networks
with unreliable synapses is also largely influenced by the stochastic effect
of neurotransmitter release.
Finally, two supplemental work has been also performed in this paper. In the
first work, we compared both the synfire propagation and firing rate
propagation in URE feedforward neuronal networks with the results in
corresponding feedforward neuronal networks composed of purely excitatory
neurons but connected with reliable synapses in an random coupling fashion
(RRE feedforward neuronal network). Our simulations showed that several
different results exist for both the synfire propagation and firing rate
propagation in these two different feedforward neuronal network models. These
results tell us that we should be cautioned when using random connections to
replace unreliable synapses in modelling research. In the second work, we
extended our results to more generalized feedforward neuronal networks, which
consist not only of the excitatory neurons but also of inhibitory neurons
(UREI feedforward neuronal network). The simulation results demonstrated that
inhibition also plays an important role in both types of neural activity
propagation, and better choosing of the relative strength of inhibitory and
excitatory synapses can enhance the transmission capability of the considered
network.
The results presented in this work might be more realistic than those obtained
based on reliable synaptic models. This is because the communication between
biological neurons indeed displays the unreliable properties. In real neural
systems, neurons may make full use of the characteristics of unreliable
synapses to transmit neural information in an adaptive way, that is, switching
between different signal propagation modes freely as required. Further work on
this topic includes at least the following two aspects: (i) since all our
results are derived from numerical simulations, an analytic description of the
synfire propagation and firing rate propagation in our considered feedforward
neuronal networks requires investigation. (ii) Intensive studies on signal
propagation in the feedforward neuronal network with other types of
connectivity, such as the Mexican-hat-type connectivity (Hamaguchi et al.
2004; Hamaguchi and Aihara 2005) and the Gaussian-type connectivity (van
Rossum et al. 2002), as well as in the feedforward neuronal network imbedded
into a recurrent network (Aviel et al. 2003; Vogels and Abbott 2005; Kumar et
al. 2008), from the unreliable synapses point of view are needed as well.
## Acknowledgement
We thank Feng Chen, Yuke Li, Qiuyuan Miao, Xin Wei and Qunxian Zheng for
valuable discussions on an early version of this manuscript. We gratefully
acknowledge the anonymous reviewers for providing useful comments and
suggestion, which greatly improved our paper. We also sincerely thank one
reviewer for reminding us of a critical citation (Trommershäuser and Diesmann
2001). This work is supposed by the National Natural Science Foundation of
China (Grant No. 60871094), the Foundation for the Author of National
Excellent Doctoral Dissertation of PR China, and the Fundamental Research
Funds for the Central Universities (Grant No. 1A5000-172210126). Daqing Guo
would also like to thank the award of the ongoing best PhD thesis support from
the University of Electronic Science and Technology of China.
## References
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|
arxiv-papers
| 2011-11-19T02:36:39 |
2024-09-04T02:49:24.472642
|
{
"license": "Public Domain",
"authors": "Daqing Guo and Chunguang Li",
"submitter": "Daqing Guo",
"url": "https://arxiv.org/abs/1111.4526"
}
|
1111.4537
|
# Some common fixed points results on metric spaces over topological modules
Ion Marian Olaru
###### Abstract
In this paper, we replace the real numbers by a topological R-module and
define R-metric spaces $(X,d)$. Also, we prove some common fixed point
theorems on R-module metric spaces. We obtain, as a particular case the Perov
theorem (see [3])
2010 Mathematical Subject Classification: 47H10
Keywords: R-metric spaces, fixed point theory, topological rings, topological
modules
## 1 R-metric spaces
In this section we shall define $R$-metric spaces and prove some properties.
All axioms for an ordinary metric space can be meaningfully formulated for an
abstract metric space, where the abstract metric takes values in a partially
ordered topological module of a certain type which will be defined below. Such
a space will be called $R$-metric space.
We begin this section by recalling a few facts concerning topological rings,
topological modules and partially ordered rings. Unless explicitly stated
otherwise all rings will be assumed to possess an identity element, denoted by
$1$.
###### Definition 1.1.
(see [5]) A topology $\tau$ on a ring $(R,+,\cdot)$ is a ring topology and
$R$, furnished with $\tau$, is a topological ring if the following conditions
hold:
* (TR 1)
$(x,y)\rightarrow x+y$ is continuous from $R\times R$ to $R$;
* (TR 2)
$x\rightarrow-x$ is continuous from $R$ to $R$;
* (TR 3)
$(x,y)\rightarrow x\cdot y$ is continuous from $R\times R$ to $R$,
where $R$ is given topology $\tau$ and $R\times R$ the cartesian product
determined by topology $\tau$.
###### Definition 1.2.
(see [5]) Let $R$ be a topological ring, $E$ an R-module. A topology
$\mathcal{T}$ on $E$ is a R-module topology and E, furnished with
$\mathcal{T}$, is a topological R-module if the following conditions hold:
* (TM 1)
$(x,y)\rightarrow x+y$ is continuous from $E\times E$ to $E$;
* (TM 2)
$x\rightarrow-x$ is continuous from $E$ to $E$;
* (TM 3)
$(a,x)\rightarrow a\cdot x$ is continuous from $R\times E$ to $E$,
where $E$ is given topology $\mathcal{T}$, $E\times E$ the cartesian product
determined by topology $\mathcal{T}$ and $A\times E$ the cartesian product
determined by topology of R and E.
By a partially ordered ring is meant a pair consisting of a ring and a
compatible partial order, denoted by $\preceq$(see [4]).
In the following we always suppose that $R$ is an ordered topological ring
such that $0\preceq 1$ and $E$ is a topological $R$-module.
###### Definition 1.3.
A subset $P$ of $E$ is called a cone if:
* (i)
$P$ is closed, nonempty and $P\neq\\{0_{E}\\}$;
* (ii)
$a,b\in R$, $0\preceq a$, $0\preceq b$ and $x,y\in P$ implies $a\cdot x+b\cdot
y\in P$;
* (iii)
$P\cap-P=\\{0\\}.$
Given a cone $P\subset E$, we define on E the partial ordering $\leq_{P}$ with
respect to $P$ by
(1.1) $x\leq_{P}y\ if\ and\ only\ if\ y-x\in P.$
We shall write $x<_{P}y$ to indicate that $x\leq_{P}y$ but $x\neq y$, while
$x\ll y$ will stand for $y-x\in intP$(interior of $P$).
###### Example 1.1.
Let $R=\mathcal{M}_{n\times n}(\mathbb{R})$ be the ring of all matrices with n
rows and n columns with entries in $\mathbb{R}$ and $E=\mathbb{R}^{n}$. We
define the partial order $\preceq$ on $M_{n\times n}(\mathbb{R})$ as follows
$A\preceq B\ if\ and\ only\ if\ for\ each\ i,j=\overline{1,n}\ we\ have\
a_{ij}\leq b_{ij}.$
Then
* (a)
the topology $\tau$, generated by matrix norm
$N:M_{n\times n}(\mathbb{R})\rightarrow\mathbb{R},$
$N(A)=\max\limits_{i=\overline{1,n}}\sum\limits_{j=1}^{n}|a_{ij}|,$
is a ring topology;
* (b)
the standard topology $\mathcal{D}$ is a R-module topology on
$\mathbb{R}^{n}$;
* (c)
$P=\\{(x_{1},x_{2},\cdots,x_{n})\in\mathbb{R}^{n}\mid x_{i}\geq
0,(\forall)i=\overline{1,n}\\}$ is a cone in E.
Indeed, Theorem 1.3 pp 2 of [5] leads us to $(a)$.
It is obvious that $(TM\ 1)$ and $(TM\ 2)$ are satisfied. Now we consider
$A_{n}\stackrel{{\scriptstyle\tau}}{{\rightarrow}}A$ and
$x_{n}\stackrel{{\scriptstyle\mathcal{D}}}{{\rightarrow}}x$ as
$n\rightarrow\infty$. Then
$\|A_{n}\cdot x_{n}-A\cdot x\|_{\mathbb{R}^{n}}=\|A_{n}\cdot x_{n}-A_{n}\cdot
x+A_{n}\cdot x-A\cdot x\|_{\mathbb{R}^{n}}\leq$
$\|A_{n}\cdot(x_{n}-x)\|_{\mathbb{R}^{n}}+\|(A_{n}-A)\cdot
x\|_{\mathbb{R}^{n}}\leq
N(A_{n})\|x_{n}-x\|_{\mathbb{R}^{n}}+N(A_{n}-A)\|x\|_{\mathbb{R}^{n}}\stackrel{{\scriptstyle
n\rightarrow\infty}}{{\rightarrow}}0.$
Hence, $(TM\ 3)$ holds. Thus, we have obtained $(b)$. Finally, it easy to see
that $P$ is a cone in $E$.
In the following we always suppose that $E$ is a topological $R$-module, $P$
is a cone in $E$ with $intP\neq\emptyset$ and $\leq_{P}$ is a partial ordering
with respect to $P$.
###### Definition 1.4.
Let $X$ be a nonempty set. Suppose that a mapping
$d:X\times X\rightarrow E$
satisfies:
* $(d_{1})$
$0_{E}\leq_{P}d(x,y)$ for all $x,y\in X$ and $d(x,y)=0_{E}$ if and only if
$x=y$;
* $(d_{2})$
$d(x,y)=d(y,x)$, for all $x,y\in X$ ;
* $(d_{3})$
$d(x,y)\leq_{P}d(x,z)+d(z,y)$, for all $x,y,z\in X$.
Then $d$ is called a R-metric on $X$ and $(X,d)$ is called a R-metric space.
###### Example 1.2.
Any cone metric space is a R-metric space.
###### Example 1.3.
Let $R=M_{n\times n}(\mathbb{R})$ be the ring of all matrices with n rows and
n columns with entries in $\mathbb{R}$, $E=\mathbb{R}^{n}$, $X=\mathbb{R}^{n}$
and
$P=\\{(x_{1},x_{2},\cdots,x_{n})\in\mathbb{R}^{n}\mid x_{i}\geq
0,(\forall)i=\overline{1,n}\\}$
a cone in E.
We define the partial order $\preceq$ on $M_{n\times n}(\mathbb{R})$ as
follows
$A\preceq B\ if\ and\ only\ if\ for\ each\ i,j=\overline{1,n}\ we\ have\
a_{ij}\leq b_{ij}.$
Then for all $A=(a_{ij})$, $a_{ij}>0$ we have that
$d:\mathbb{R}^{n}\times\mathbb{R}^{n}\rightarrow\mathbb{R}^{n},$
$d(x,y)=(\sum\limits_{j=1}^{n}a_{1j}|x_{j}-y_{j}|,\cdots,\sum\limits_{j=1}^{n}a_{ij}|x_{j}-y_{j}|,\cdots,\sum\limits_{j=1}^{n}a_{nj}|x_{j}-y_{j}|)$
is a $R$-metric on $X$.
Indeed,
* $(d_{1})$
Since $\sum\limits_{j=1}^{n}a_{ij}|x_{j}-y_{j}|\geq 0$ for all
$i=\overline{1,n}$, we have that $0\leq_{P}d(x,y)$ for all
$x,y\in\mathbb{R}^{n}$. Also, $d(x,y)=0$ involve that
$\sum\limits_{j=1}^{n}a_{ij}|x_{j}-y_{j}|=0$ which means that $x_{j}=y_{j}$
for all $j=\overline{1,n}$. It follows that $x=y$.
* $(d_{2})$
It is obvious that $d(x,y)=d(y,x)$, for all $x,y\in\mathbb{R}^{n}$.
* $(d_{3})$
Let be $x,y,z\in\mathbb{R}^{n}$. Since
$\sum\limits_{j=1}^{n}a_{ij}|x_{j}-y_{j}|\leq\sum\limits_{j=1}^{n}a_{ij}|x_{j}-z_{j}|+\sum\limits_{j=1}^{n}a_{ij}|z_{j}-y_{j}|$,
we have that $d(x,y)\leq_{P}d(x,z)+d(z,y)$, for all $x,y,z\in\mathbb{R}^{n}$.
In the following, we shall write $x\prec y$ to indicate that $x\preceq y$ but
$x\neq y$.
###### Remark 1.1.
We have that:
* (i)
$intP+intP\subseteq intP$;
* (ii)
$\lambda\cdot intP\subseteq intP$, where $\lambda$ is a invertible element of
the ring $R$ such that $0\prec\lambda$;
* (iii)
if $x\leq_{P}y$ and $0\preceq\alpha$, then $\alpha\cdot x\preceq\alpha\cdot
y$.
Proof:
$(i)$ Let be $x\in intP+intP$. Then there exists $x_{1},x_{2}\in intP$ such
that $x=x_{1}+x_{2}$. It follows that there exists $V_{1}$ neighborhood of
$x_{1}$ and $V_{2}$ neighborhood of $x_{2}$ such that
$x_{1}\in V_{1}\subset P\ and\ x_{2}\in V_{2}\subset P.$
Since for each $x_{0}\in E$, the mapping $x\rightarrow x+x_{0}$ is a
homeomorphism of $E$ onto itself,we have that $V_{1}+V_{2}$ is a neighborhood
of $x$ with respect to topology $\mathcal{T}$. Thus, $x\in intP$.
$(ii)$ Let $0\prec\lambda$ be an invertible element of the ring $R$ and
$x=\lambda\cdot c$, $c\in intP$. It follows that there exists a neighborhood
$V$ of $c$ such that $c\in V\subset P$.
Since the mapping $x\rightarrow\lambda x$ is a homeomorphism of $E$ onto
itself, we have that $\lambda\cdot V$ is a neighborhood of $x$ with respect to
the topology $\mathcal{T}$. Thus, $x\in intP$.
$(iii)$ If $x\leq_{P}y$, then $y-x\in P$. It follows that for all
$0\preceq\alpha$ we have that $\alpha\cdot(y-x)\in P$ i.e. $\alpha\cdot
x\preceq\alpha\cdot y$.
In the following, unless otherwise specified, we always suppose that there
exists at least one a sequence $\\{\alpha_{n}\\}\subset R$ of invertible
elements such that $0\prec\alpha_{n}$ and $\alpha_{n}\rightarrow 0$ as
$n\rightarrow\infty$.
###### Remark 1.2.
Let $E$ be a $R$-topological module and $P\subset E$ a cone. We have that:
* (i)
If $u\leq_{P}v$ and $v\ll w$, then $u\ll w$;
* (ii)
If $u\ll v$ and $v\leq_{P}w$, then $u\ll w$;
* (iii)
If $u\ll v$ and $v\ll w$, then $u\ll w$;
* (iv)
If $0\leq_{P}u\ll c$ for each $c\in IntP$, then $u=0$;
* (v)
If $a\leq_{P}b+c$ for each $c\in IntP$, then $a\leq_{P}b$;
* (vi)
If $0\ll c$, $0\leq_{P}a_{n}$ and $a_{n}\rightarrow 0$ then there exists
$n_{0}\in\mathbb{N}$ such that $a_{n}\ll c$ for all $n\geq n_{0}$.
Proof:
* (i)
We have to prove that $w-u\in intP$ if $v-u\in P$ and $w-v\in intP$. The
condition $(TM_{1})$ implies that there exists a neighborhood $V$ of $0$ such
that $w-v+V\subset P$. It follows that $w-u+V=(w-v)+V+(v-u)\subset P+P\subset
P$. Since for each $x_{0}\in E$ the mapping $x\rightarrow x+x_{0}$ is a
homeomorphism of $E$ onto itself we have that $w-u+V$ is a neighborhood of
$w-u$ with respect to the topology $\mathcal{T}$. Thus, $w-u\in int\ P$.
* (ii)
Analogous with $(i)$.
* (iii)
We have to prove that $w-u\in intP$ if $v-u\in intP$ and $w-v\in intP$. Since
we have $intP+intP\subset intP$, it easy to see that the above assertion is
satisfied.
* (iv)
Let $\\{\alpha_{n}\\}_{n\in\mathbb{N}}\subset R$ be a sequence of invertible
elements such that for each $n\in\mathbb{N}$ we have $0\prec\alpha_{n}$ and
$\alpha_{n}\rightarrow 0$ as $n\rightarrow\infty$.
Since for each invertible element $\lambda_{0}\in R$ the mapping
$x\rightarrow\lambda_{0}\cdot x$ is a homeomorphism of $E$ onto itself, we
have that if $V$ is a neighborhood of zero then $\lambda_{0}\cdot V$ is a
neighborhood of zero. Hence, $\alpha_{n}\cdot c\in intP$ for each
$n\in\mathbb{N}$.
Then $\alpha_{n}\cdot c-u\in intP$. It follows that
$\lim\limits_{n\rightarrow\infty}\alpha_{n}\cdot c-u=-u\in\overline{P}=P$.
Thus, $u\in P\cap-P=\\{0\\}$.
* (v)
Analogous with $(iv)$.
* (vi)
Let be $0\ll c$, $0\leq_{P}a_{n}$ and $a_{n}\rightarrow 0$. Since for each
invertible element $\lambda_{0}\in R$ the mapping
$x\rightarrow\lambda_{0}\cdot x$ is a homeomorphism of $E$ onto itself, it
follows that for all neighborhood $V$ of zero we have that $-V$ is a
neighborhood of zero. Now, $0\ll c$ implies that there exists a neighborhood
$U$ of zero such that $c+U\subset P$. Let $V=U\cap-U$ be a neighborhood of
zero. Since $a_{n}$ converges to zero, it follows that there exists
$n_{0}\in\mathbb{N}$ such that $a_{n}\in V$ for all $n\geq n_{0}$. Then we
have that $c-a_{n}\in c+V\subset c+U\subset P$ for all $n\geq n_{0}$. Thus,
$a_{n}\ll c$ for all $n\geq n_{0}$.
###### Definition 1.5.
Let $\\{x_{n}\\}$ be a sequence in a R-metric space $(X,d)$ and $x\in X$. We
say that:
* (i)
the sequence $\\{x_{n}\\}$ converges to $x$ and is denoted by
$\lim\limits_{n\rightarrow\infty}x_{n}=x$ if for every $0\ll c$ there exists
$N\in\mathbb{N}$ such that $d(x_{n},x)\ll c$, for all $n>N$;
* (ii)
the sequence $\\{x_{n}\\}$ is a Cauchy sequence if for every $c\in E$, $0\ll
c$ there exists $N\in\mathbb{N}$ such that $d(x_{m},x_{n})\ll c$ for all
$m,n>N$;
The R-metric space $(X,d)$ is called complete if every Cauchy sequence is
convergent.
From the above remark we obtain
###### Remark 1.3.
Let $(X,d)$ be a R-metric space and $\\{x_{n}\\}$ be a sequence in $X$. If
$\\{x_{n}\\}$ converges to $x$ and $\\{x_{n}\\}$ converges to $y$, then $x=y$.
Indeed, for all $0\ll c$
$d(x,y)\leq_{P}d(x,x_{n})+d(x_{n},y)\ll 2\cdot c.$
Hence, $d(x,y)=0$ i.e. $x=y$.
## 2 Common fixed points theorems
In this section we obtain several coincidence and common fixed point theorems
for mappings defined on a $R$-metric space.
###### Definition 2.1.
(see [1]) Let $f$ and $g$ be self maps of a set $X$. If $w=fx=gx$ for some
$x\in X$, then $x$ is called a coincidence point of $f$ and $g$, and $w$ is
called a point of coincidence of $f$ and $g$.
Jungck [2], defined a pair of self mappings to be weakly compatible if they
commute at their coincidence points.
###### Proposition 2.1.
(see [1]) Let $f$ and $g$ be weakly compatible self maps of a set $X$. If $f$
and $g$ have a unique point of coincidence $w=fx=gx$, then $w$ is the unique
common fixed point of $f$ and $g$.
Let $\mathcal{K}$ be the set of all $k\in R$, $0\preceq k$ which have the
property that there exists a unique $S\in R$ such that
$S=\lim\limits_{n\rightarrow\infty}(1+k+\cdots+k^{n})$.
###### Example 2.1.
Let be $A\in\mathcal{M}_{n\times n}(\mathbb{R}_{+})$ such that $\rho(A)<1$.
Then $A\in\mathcal{K}$.
###### Theorem 2.1.
Let $(X,d)$ be a R-metric space and suppose that the mappings
$f,g:X\rightarrow X$ satisfy:
* (i)
the range of $g$ contains the range of $f$ and $g(X)$ is a complete subspace
of $X$;
* (ii)
there exists $k\in\mathcal{K}$ such that $d(fx,fy)\leq_{P}k\cdot d(gx,gy)$ for
all $x,y\in X$.
Then $f$ and $g$ have a unique point of coincidence in $X$. Moreover, if $f$
and $g$ are weakly compatible then, $f$ and $g$ have a unique common fixed
point.
Proof: Let $x_{0}$ be an arbitrary point in $X$. We choose a point $x_{1}\in
X$ such that $f(x_{0})=g(x_{1}).$ Continuing this process, having chosen
$x_{n}\in X$, we obtain $x_{n+1}\in X$ such that $f(x_{n})=g(x_{n+1})$. Then
$d(gx_{n+1},gx_{n})=d(fx_{n},fx_{n-1})\leq_{P}k\cdot
d(gx_{n},gx_{n-1})\leq_{P}$ $\leq_{P}k^{2}\cdot
d(gx_{n-1},gx_{n-2})\leq_{P}\cdots\leq_{P}k^{n}\cdot d(gx_{1},gx_{0}).$
We denote $S_{n}=1+k+\cdots+k^{n}$and we get that
$d(gx_{n},gx_{n+p})\leq_{P}d(gx_{n},gx_{n+1})+d(gx_{n+1},gx_{n+2})+\cdots+d(gx_{n+p-1},gx_{n+p})\leq_{P}$
$\leq_{P}(k^{n}+k^{n+1}+\cdots+k^{n+p-1})\cdot
d(gx_{1},gx_{0})=(S_{n+p-1}-S_{n-1})\cdot d(gx_{1},gx_{0}),$
for all $p\geq 1$. Thus, via Remark 1.2 $(vi)$, we obtain that $gx_{n}$ is a
Cauchy sequence. Since $g(X)$ is complete, there exists $q\in g(X)$ such that
$gx_{n}\rightarrow q$ as $n\rightarrow\infty$. Consequently, we can find $p\in
X$ such that $g(p)=q$. Further, for each $0\ll c$ there exists
$n_{0}\in\mathbb{N}$ such that for all $n\geq n_{0}$
$d(gx_{n},fp)=d(fx_{n-1},f(p))\leq_{P}k\cdot d(gx_{n-1},gp)\ll c.$
It follows that $gx_{n}\rightarrow fp$ as $n\rightarrow\infty$. The uniqueness
of the limit implies that $fp=gp=q$. Next we show that $f$ and $g$ have a
unique point of coincidence. For this, we assume that there exists another
point $p_{1}\in X$ such that $fp_{1}=gp_{1}$. We have
$d(gp_{1},gp)=d(fp_{1},fp)\leq_{P}k\cdot d(gp_{1},gp)=k\cdot
d(fp_{1},fp)\leq_{P}k^{2}\cdot d(gp_{1},gp)\leq_{P}\cdots\leq_{P}k^{n}\cdot
d(gp_{1},gp).$
Let be $0\ll c$. Since $k^{n}\rightarrow 0$ as $n\rightarrow\infty$ it follows
that there exists $n_{0}\in\mathbb{N}$ such that $k^{n}\cdot d(gp_{1},gp)\ll
c$ for all $n\geq n_{0}$. Then $d(gp_{1},gp)\ll c$ for each $0\ll c$. Thus
$d(gp_{1},gp)=0$ i.e. $gp_{1}=gp$. From Proposition 2.1 $f$ and $g$ have a
unique common fixed point.
###### Remark 2.1.
The above theorem generalizes Theorem 2.1 of Abbas and Jungck [1], which
itself is a generalization of Banach fixed point theorem.
###### Corollary 2.1.
Let $(X,d)$ be a complete $R$-metric space and we suppose that the mapping
$f:X\rightarrow X$ satisfies:
* (i)
there exists $k\in\mathcal{K}$ such that $d(fx,fy)\leq_{P}k\cdot d(x,y)$ for
all $x,y\in X$.
Then $f$ has in $X$ a unique fixed point point.
Proof: The proof uses Theorem 3.1 by replacing $g$ with the identity mapping.
From the above corollary using Example 1.1, we obtain the Perov fixed point
theorem (see [3])
###### Corollary 2.2.
Let $(X,d)$ be a complete $\mathcal{M}_{n\times n}(\mathbb{R}_{+})-$ metric
space and $E=\mathbb{R}^{n}$ and we suppose the mapping $f:X\rightarrow X$
satisfies:
* (i)
there exists $A\in\mathcal{M}_{n\times n}(\mathbb{R}_{+})$ with $\rho(A)<1$
such that
$d(fx,fy)\leq_{P}A\cdot d(x,y),$
for all $x,y\in X$.
Then $f$ has in $X$ a unique fixed point point.
## 3 Comparison function
###### Definition 3.1.
Let P be a cone in a topological R-module E. A function $\varphi:P\rightarrow
P$ is called a comparison function if
* (i)
$\varphi(0)=0$ and $\varphi(t)<_{P}t$ for all $t\in P-\\{0\\}$;
* (ii)
$t_{1}\leq_{P}t_{2}$ implies $\varphi(t_{1})\leq_{P}\varphi(t_{2})$;
* (iii)
$t\in intP$ implies $t-\varphi(t)\in intP$;
* (iv)
if $t\in P-\\{0\\}$ and $0\ll c$, then there exists $n_{0}\in\mathbb{N}$ such
that $\varphi^{n}(t)\ll c$ for each $n\geq n_{0}$.
###### Example 3.1.
Let P be a cone in a topological R-module E and $\lambda\in\mathcal{K}$ such
that $0\prec 1-\lambda$. Then $\varphi:P\rightarrow P$, defined by
$\varphi(t)=\lambda\cdot t$ is a comparison function.
Indeed,
$(i)$ It is obvious that $\varphi(0)=0$ and $\varphi(t)<_{P}t$ for all $t\in
P-\\{0\\}$.
$(ii)$ if $t_{1}\leq_{P}t_{2}$ and $\lambda\in\mathcal{K}$ then
$\lambda\cdot(t_{2}-t_{1})\in P$. Thus $\varphi(t_{1})\leq_{P}\varphi(t_{2})$.
$(iii)$ We remark that if $\lambda\in\mathcal{K}$, then $1-\lambda$ is an
invertible element of the ring $R$. Now, let be $t\in intP$. Then
$(1-\lambda)\cdot t\in(1-\lambda)intP\subset intP$.
$(iv)$ Let be $t\in P-\\{0\\}$ and $0\ll c$. Then
$\varphi^{n}(t)=\lambda^{n}\cdot t\stackrel{{\scriptstyle
n\rightarrow\infty}}{{\rightarrow}}0.$
We obtain, via Remark 1.2, that there exists $n_{0}\in\mathbb{N}$ such that
$\varphi^{n}(t)\ll c$ for each $n\geq n_{0}$.
Let $(X,d)$ be a $R$-metric space and let $\varphi:K\rightarrow K$ be a
comparison function. For a pair $(f,g)$ of self-mappings on $X$ consider the
following condition
* (C)
for arbitrary $x,y\in X$ there exists $u\in\\{d(gx,gy),d(gx,fx),d(gy,fy)\\}$
such that $d(fx,fy)\leq_{P}\varphi(u)$.
###### Theorem 3.1.
Let $(X,d)$ be a R-metric space and let $f,g:X\rightarrow X$ such that
* (i)
the pair $(f,g)$ satisfies the condition (C) for some comparison function
$\varphi$;
* (ii)
$f(X)\subset g(X)$ and f(X) or g(X) is a complete subspace of X.
Then f and g have a unique point of coincidence in X. Moreover if $f$ and $g$
are weakly compatible, then $f$ and $g$ have a unique common fixed point.
Proof: Let $x_{0}$ be an arbitrary point in $X$. We choose a point $x_{1}\in
X$ such that $fx_{0}=gx_{1}.$ Continuing this process, having chosen $x_{n}\in
X$, we obtain $x_{n+1}\in X$ such that $fx_{n}=gx_{n+1}$.
$\ulcorner$ We shall prove that the sequence $\\{y_{n}\\}$, where
$y_{n}=fx_{n}=gx_{n+1}$( the so-called Jungck sequence ) is a Cauchy sequence
in $R$-metric space $(X,d)$.
If $y_{N}=y_{N+1}$ for some $N\in\mathbb{N}$ then $y_{m}=y_{N}$ for each $m>N$
and the conclusion follows. Indeed, we prove by induction arguments that
(3.1) $y_{N+k}=y_{N+k+1},(\forall)k\in\mathbb{N}.$
For $k=0$ we have $y_{N}=y_{N+1}$. Let (3.1) hold for all $k=\overline{0,i}$.
Then
$d(y_{N+i+1},y_{N+i+2})=d(fx_{N+i},fx_{N+i+1})\leq_{P}\varphi(u),$
where
$u\in\\{d(gx_{N+i},gx_{N+i+1}),d(gx_{N+i},fx_{N+i}),d(gx_{N+i+1},fx_{N+i+1})\\}=$
$\\{d(y_{N+i-1},y_{N+i}),d(y_{N+i-1},y_{N+i}),d(y_{N+i},y_{N+i+1})\\}=\\{0\\}.$
Hence, $d(y_{N+i+1},y_{N+i+2})\leq_{P}\varphi(u)=0$ i.e.
$y_{N+i+1}=y_{N+i+2}$. Q.E.D
Suppose that $y_{n}\neq y_{n+1}$ for each $n\in\mathbb{N}$. The condition (C)
implies that
$d(y_{n},y_{n+1})=d(fx_{n},fx_{n+1})\leq_{P}\varphi(u),$
where
$u\in\\{d(gx_{n},gx_{n+1}),d(fx_{n},gx_{n}),d(fx_{n+1},gx_{n+1})\\}=\\{d(y_{n-1},y_{n}),d(y_{n},y_{n-1}),d(y_{n+1},y_{n})\\}.$
The case $u=d(y_{n+1},y_{n})$ is impossible, since this would imply
$d(y_{n+1},y_{n})\leq_{P}\varphi(d(y_{n+1},y_{n}))<_{P}d(y_{n+1},y_{n}).$
Thus, $u=d(y_{n},y_{n-1})$ and
$d(y_{n+1},y_{n})\leq_{P}\varphi(d(y_{n},y_{n-1}))\leq_{P}\cdots\leq_{P}\varphi^{n}(d(y_{1},y_{0})).$
Using Remark 1.2 (i) and property $(iv)$ of the comparison function we obtain
that for all $0\ll\varepsilon$ there exists $n_{0}\in\mathbb{N}$ such that
(3.2) $d(y_{n},y_{n+1})\ll\varepsilon,\ (\forall)n\geq n_{0}.$
Now, let be $0\ll c$. Then, using property $(iii)$ of the comparison function,
we get that
(3.3) $d(y_{n},y_{n+1})\ll c-\varphi(c),\ (\forall)n\geq n_{0}.$
Let us fix now $n\geq n_{0}$ and let us prove that
(3.4) $d(y_{n},y_{k+1})\ll c,\ (\forall)k\geq n.$
Indeed, for $k=n$ we have
$d(y_{n},y_{n+1})\ll c-\varphi(c)\leq_{P}c.$
Hence,
$d(y_{n},y_{n+1})\ll c.$
Let (3.4 ) hold for some $k\geq n$. Then we have
$d(y_{n},y_{k+2})\leq_{P}d(y_{n},y_{n+1})+d(y_{n+1},y_{k+2})\ll$
$c-\varphi(c)+d(fx_{n+1},fx_{k+2})\leq_{P}c-\varphi(c)+\varphi(u),$
where
$u\in\\{d(gx_{n+1},gx_{k+2}),d(gx_{n+1},fx_{n+1}),d(gx_{k+2},fx_{k+2})\\}.$
Consider now the following three possible cases:
* Case 1:
$u=d(gx_{n+1},gx_{k+2})$. Then
$\varphi(u)=\varphi(d(gx_{n+1},gx_{k+2}))=\varphi(d(y_{n},y_{k+1}))\leq_{P}\varphi(c).$
From the above relation it follows that,
$d(y_{n},y_{k+2})\ll
c-\varphi(c)+\varphi(u)\leq_{P}c-\varphi(c)+\varphi(c)=c.$
Hence, $d(y_{n},y_{k+2})\ll c.$
* Case 2:
$u=d(gx_{n+1},fx_{n+1})=d(y_{n},y_{n+1})$. Then
$\varphi(u)\leq_{P}\varphi(d(y_{n},y_{n+1}))\leq_{P}\varphi(c-\varphi(c))\leq_{P}\varphi(c).$
From the above relation we get that,
$d(y_{n},y_{k+2})\ll
c-\varphi(c)+\varphi(u)\leq_{P}c-\varphi(c)+\varphi(c)=c.$
Hence, $d(y_{n},y_{k+2})\ll c.$
* Case 3:
$u=d(gx_{k+2},fx_{k+2})$. Then
$\varphi(u)=\varphi(d(gx_{k+2},fx_{k+2}))=\varphi(d(y_{k+1},y_{k+2}))\leq_{P}\varphi(c-\varphi(c))\leq_{P}\varphi(c).$
From the above relation we get that,
$d(y_{n},y_{k+2})\ll
c-\varphi(c)+\varphi(u)\leq_{P}c-\varphi(c)+\varphi(c)=c.$
Hence, $d(y_{n},y_{k+2})\ll c.$
So, it has been proved by induction that $\\{y_{n}\\}$ is a Cauchy
sequence.$\lrcorner$
Since, by assumption, $f(X)$ or $g(X)$ is a complete subspace of $X$, we
conclude that there exists $q\in g(X)$ such that
$y_{n}=fx_{n}=gx_{n+1}\rightarrow q$ as $n\rightarrow\infty$ and there exists
$p\in X$ such that $q=gp$.
$\ulcorner$ We claim that $q=fp$. Indeed, if we suppose that $d(fp,q)\neq 0$,
then we have
$d(fp,q)\leq_{P}d(fp,fx_{n})+d(fx_{n},q)\leq_{P}\varphi(u)+d(y_{n},q),$
where
$u\in\\{d(gp,gx_{n}),d(gp,fp),d(gx_{n},fx_{n})\\}$
Let $0\ll c$. At least one of the following three cases holds for infinitely
many $n\in\mathbb{N}$:
* Case 1:
$u=d(gp,gx_{n})$. Then, there exists $n_{0}(c)\in\mathbb{N}$ such that for all
$n\geq n_{0}(c)$
$d(fp,q)\leq_{P}\varphi(d(gp,gx_{n}))+d(y_{n},q)<_{P}d(q,y_{n-1})+d(y_{n},q)\ll
2\cdot c.$
It follows that $d(fp,q)=0$, which is a contradiction.
* Case 2:
$u=d(gp,fp)=d(fp,q)$ . Then we have
$d(fp,q)\leq_{P}\varphi(d(q,fp))+d(y_{n},q)\ll\varphi(d(q,fp))+c.$
Thus, $d(fp,q)\leq_{P}\varphi(d(q,fp))$. So, by using of the properties $(ii)$
and $(iv)$ of the comparison function, we obtain that there exists
$n_{0}\in\mathbb{N}$ such that $d(fp,q)\leq_{P}\varphi^{n}(d(fp,q))\ll c$ i.e.
$d(fp,q)=0$, which is a contradiction.
* Case 3:
$u=d(gx_{n},fx_{n})=d(y_{n-1},y_{n})$. Then, there exists
$n_{0}(c)\in\mathbb{N}$ such that for all $n\geq n_{0}(c)$ we have
$d(fp,q)\leq_{P}\varphi(d(y_{n-1},y_{n}))+d(y_{n},q)<_{P}d(y_{n-1},y_{n})+d(y_{n},q)\ll
2\cdot c,$
i.e. $d(fp,q)=0$, which is a contradiction.
It follows that $fp=gp=q$ i.e. $p$ is a coincidence point of the pair $(f,g)$
and $q$ is a point of coincidence.$\lrcorner$
$\ulcorner$ Next we show that $f$ and $g$ have a unique point of coincidence.
For this we assume that there exists another point $p_{1}\in X$ such that
$fp_{1}=gp_{1}$. If we suppose that $d(fp_{1},fp)\neq 0$ we get that
$d(fp_{1},fp)\leq_{P}\varphi(u)$, where
$u\in\\{d(gp_{1},gp),d(gp_{1},fp_{1}),d(gp,fp)\\}=\\{d(gp_{1},gp),0\\}.$
In both possible cases a contradiction easily follows :
$d(fp_{1},fp)\leq_{P}\varphi(d(fp_{1},fp))<_{P}d(fp_{1},fp)$ or
$d(fp_{1},fp)\leq_{P}\varphi(0)=0$. We conclude that the mappings $f$ and $g$
have a unique point of coincidence. From Proposition 2.1 $f$ and $g$ have a
unique common fixed point.$\lrcorner$
## References
* [1] M. Abbas, G. Jungck,Common fixed point results for noncommuting mappings without continuity in cone metric spaces, J. Math.Anal.Appl. 341 (2008)
* [2] G. Jungck, Common fixed points for noncontinuous nonself maps on non-metric spaces, Far East J. Math. Sci.(FJMS) 4(1996) 199-215.
* [3] A. I. Perov, On the Cauchy problem for a system of ordinary differential equations, Pviblizhen. Met. Reshen. Differ. Uvavn., vol. 2, pp. 115 134, 1964.
* [4] Stuart A. Steinberg, Lattice-ordered rings and modules, Springer Science + Business Media, LLC 2010, DOI 10.1007/978-1-4419-1721-8-1.
* [5] Seth Warner, Topological rings, North-Holland Math. Studies 178, 1993.
Department of Mathematics,
Faculty of Science,
University ”Lucian Blaga” of Sibiu,
Dr. Ion Ratiu 5-7, Sibiu, 550012, Romania
E-mail: marian.olaru@ulbsibiu.ro
|
arxiv-papers
| 2011-11-19T07:46:08 |
2024-09-04T02:49:24.485558
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Ion Olaru",
"submitter": "Ion Olaru",
"url": "https://arxiv.org/abs/1111.4537"
}
|
1111.4579
|
# Canalization of the evolutionary trajectory of the human influenza virus
Trevor Bedford Department of Ecology and Evolutionary Biology, University of
Michigan, Ann Arbor, MI, USA. Howard Hughes Medical Institute, University of
Michigan, Ann Arbor, MI, USA. Andrew Rambaut Institute of Evolutionary
Biology, University of Edinburgh, Edinburgh, UK. Fogarty International
Center, National Institutes of Health, Bethesda, MD, USA. Mercedes Pascual
Department of Ecology and Evolutionary Biology, University of Michigan, Ann
Arbor, MI, USA. Howard Hughes Medical Institute, University of Michigan, Ann
Arbor, MI, USA.
Abstract
Since its emergence in 1968, influenza A (H3N2) has evolved extensively in
genotype and antigenic phenotype. Antigenic evolution occurs in the context of
a two-dimensional ‘antigenic map’, while genetic evolution shows a
characteristic ladder-like genealogical tree. Here, we use a large-scale
individual-based model to show that evolution in a Euclidean antigenic space
provides a remarkable correspondence between model behavior and the
epidemiological, antigenic, genealogical and geographic patterns observed in
influenza virus. We find that evolution away from existing human immunity
results in rapid population turnover in the influenza virus and that this
population turnover occurs primarily along a single antigenic axis. Thus,
selective dynamics induce a canalized evolutionary trajectory, in which the
evolutionary fate of the influenza population is surprisingly repeatable and
hence, in theory, predictable.
## Author Summary
Each year, hundreds of millions become sick with the influenza virus and
hundreds of thousands die from the infection. After recovery from the flu, a
person gains permanent immunity specific to the infecting strain. However,
owing to its RNA makeup, mutations occur rapidly to the virus genome. Some of
these mutations change the shape of proteins visible to the human immune
system, and thus alter the virus’s antigenic phenotype. These mutations allow
the virus to re-infect those who have previously recovered from an earlier
strain, and thus will quickly spread through the virus population. It is
because of influenza’s rapid antigenic evolution that the flu vaccine needs
frequent updating. However, despite strong pressure to evolve away from human
immunity and to diversify in antigenic phenotype, influenza, and especially
influenza A (H3N2), shows paradoxically limited genetic and antigenic
diversity present at any one time. Here, we propose a simple model of
influenza that displays rapid evolution but low standing diversity and
simultaneously accounts for the epidemiological, genetic, antigenic and
geographical patterns displayed by the virus. We find that in evolving
directly away from existing human immunity, the virus severely limits its
future evolutionary potential.
## Introduction
Epidemic influenza is responsible for between 250,000 and 500,000 global
deaths annually, with influenza A (and in particular, A/H3N2) having caused
the bulk of human mortality and morbidity [1]. Influenza A (H3N2) has
continually circulated within the human population since its introduction in
1968, exhibiting recurrent seasonal epidemics in temperate regions and less
periodic transmission in the tropics. During this time, H3N2 influenza has
continually evolved both genetically and antigenically. Most antigenic drift
is thought to be driven by changes to epitopes in the hemagglutinin (HA)
protein [2]. Phylogenetic analysis of the genetic relationships among HA
sequences has revealed a distinctive genealogical tree showing a single
predominant trunk lineage and side branches that persist for only 1–5 years
before going extinct [3]. This tree shape is indicative of serial replacement
of strains over time; H3N2 influenza shows rapid evolution, but low standing
genetic diversity.
This observation has remained puzzling from an epidemiological standpoint.
Antigenic evolution occurs rapidly and strong diversifying selection exists to
escape from human immunity; why then do we see serial replacement of strains
rather than continual genetic and antigenic diversification? Indeed, simple
epidemiological models show explosive diversity of genotype and phenotype over
time [4, 5]. Previous work has sought model-based explanations of the limited
diversity of influenza, relying on short-lived strain-transcending immunity
[4, 5], complex genotype-to-phenotype maps [6] or a limited repertoire of
antigenic phenotypes [7].
Experimental characterization of antigenic phenotype is possible through the
hemagglutination inhibition (HI) assay, which measures the cross-reactivity of
HA from one virus strain to serum raised against another strain [8]. The
results of many HI assays can be combined to yield a two-dimensional map,
representing antigenic similarity and distance between strains as an easily
visualized and quantified measure [9]. The path traced across this map by
influenza A (H3N2) from 1968 until present is largely linear, showing serial
replacement of one strain by another; there are no major bifurcations of
antigenic phenotype [9].
Herein, we seek to simultaneously model the genetic and antigenic evolution of
the influenza virus. We represent antigenic phenotypes as points in a
$N$-dimensional Euclidean space. Based on the finding that a two-dimensional
map adequately explains observed antigenic distance between strains [9], we
begin with antigenic phenotypes as points on a plane, but relax this
assumption later on in the analysis. After exposure to a virus, a host’s risk
of infection is proportional to the Euclidean distance between the infecting
phenotype and the closest phenotype in the host’s immune history. Mutations
perturb antigenic phenotype, moving phenotype in a random radial direction and
for a randomly distributed distance. We implemented this geometrical model in
a large-scale individual-based simulation intended to directly model the
antigenic map and genealogical tree of the global influenza population. The
simulation includes multiple host populations with different seasonal forcing,
hosts with complete immune histories of infection, and viruses with antigenic
phenotypes. As the simulation proceeds, infections are tracked and a complete
genealogy connecting virus samples is constructed. Results shown here are for
a single representative simulation of 40 years of virus evolution in a
population of 90 million hosts.
## Results
The virus persists over the course of the 40-year simulation, infecting a
significant fraction of the host population through annual winter epidemics in
temperate regions and through less periodic epidemics in the tropics (Figure
1A). Across replicate simulations, we observe average yearly attack rates of
6.8% in temperate regions and rates of 7.1% in the tropics, comparable with
estimated attack rates of influenza A (H3N2) of 3–8% per year [10, 11]. Over
the course of the simulation, the virus population evolves in antigenic
phenotype exhibiting, at any point, a handful of highly abundant phenotypes
sampled repeatedly and a large number of phenotypes appearing at low abundance
(Figure 1B). By including measurement noise on antigenic locations, we
approximate an experimental antigenic map of H3N2 influenza (Figure 1D). The
appearance of clusters in the antigenic map comes from the regular spacing of
high abundance phenotypes combined with measurement noise. Over time, clusters
of antigenically similar strains are replaced by novel clusters of more
advanced strains (Figure S1A). Across replicate simulations, clusters persist
for an average of 5.0 years measured as the time it takes for a new cluster to
reach 10% frequency, peak and decline to 10% frequency. The transition between
clusters occurs quickly, taking an average of 1.8 years.
Remarkably, although antigenic phenotype is free to mutate in any direction in
the two-dimensional space, selection pressures force the virus population to
move in nearly a straight line in antigenic space (Figure 1B). Across
replicate simulations, 94% of the variance of antigenic phenotype can be
explained by a single dimension of variation. This mirrors the empirical
results showing a largely linear antigenic map for H3N2 influenza isolates
from 1968 to 2003 [9]. Because of the primarily one-dimensional movement,
antigenic distance from the original phenotype increases nearly linearly with
time (Figure S1B). Antigenic evolution occurs in a punctuated fashion; periods
of relative stasis are interspersed with more rapid antigenic change (Figure
S1B). Antigenic and epidemiological dynamics show a fundamental linkage, so
that large jumps of antigenic phenotype result in increased rates of infection
(Figure 1, Figure S2). In general, evolution via many smaller mutations
results in a smoother antigenic map and less variation in yearly epidemics
(Figure S3), while evolution via rare mutations of large effect exhibits a
more clustered antigenic map and wider variation in seasonal incidence (Figure
S4).
The genealogical tree connecting the evolving virus population appears
characteristically sparse with pronounced trunk and short-lived side branches
(Figure 1C). This tree shape is reflected in low levels of standing diversity;
across replicates, an average of 5.68 years of evolution separate two randomly
sampled viruses in the population. This level of diversity matches what is
observed in phylogenies of influenza A (H3N2) [12]. A spindly genealogical
tree is indicative of population turnover, wherein novel antigenic phenotypes
continually replace more primitive ‘spent’ phenotypes, purging their
genealogical diversity. In general, natural selection reduces effective
population size and genealogical diversity [13]. Here, by comparing mutations
occurring on trunk branches vs. mutations occurring on side branches, we find
evidence for pervasive positive selection for antigenic change (Table 1).
Trunk mutations tend to push antigenic phenotype forward along the line of
primary antigenic variation (Figure S5). We find a roughly linear relationship
between the antigenic effect of a mutation and the likelihood that this
mutation becomes incorporated into the trunk (Figure S6). Additionally, we
find that trunk mutations occur at strikingly regular intervals, with less
variation of waiting times than expected under a simple random process (Figure
S7). There is a relative scarcity of mutation events occurring in intervals
under 1 year and a relative excess of a mutation events occurring in 2–3 year
intervals (Figure S7).
The genealogical tree also contains detailed information on the history of
migration between regions. We find that, consistent with empirical estimates
[14, 15], the trunk resides primarily within the tropics, where seasonal
dynamics are less prevalent (Figure 2A). Across replicate simulations, we
observe 72% of the trunk’s history within the tropics and 28% within temperate
regions. With symmetric host contact rates and equal host population sizes,
and without seasonal forcing, we would expect trunk proportions of one third
for each region. We calculated rates of migration based on observed event
counts across replicate simulations, separating region-specific rates on side
branches from region-specific rates on trunk branches. We find that migration
patterns on side branches are close to symmetric, with similar rates between
all regions, while migration patterns on trunk branches are highly asymmetric,
with high rates of movement between temperate regions and from temperate
regions into the tropics (Figure 2B). Extrapolating from these rates, we
arrive at an expected stationary distribution of 76% tropics and 24% temperate
regions, in line with the observed residency patterns of the trunk. It may at
first seem counter-intuitive to see higher rates of movement from the
temperate regions into the tropics along trunk branches, but it makes sense
when thought of in terms of conditional probability. Only those lineages that
migrate into the tropics or those lineages which rapidly migrate between the
north and south have a chance at becoming the trunk lineage, while lineages
that remain within the temperate regions are doomed to extinction.
These findings suggest that persistence and migration are fundamentally
connected and have important implications for future phylogeographic analyses.
Russell et al. [14] emphasize a source-sink model of movement of the HA
protein of influenza A (H3N2) based on their finding of a trunk lineage
residing within China and the Southeast Asian tropics. Whereas, Bedford et al.
[15] emphasize a global metapopulation model based on phylogenetic inference
of migration rates across the entire tree. Our results suggest that both
scenarios are simultaneously possible; side branches may be highly volatile
moving rapidly and symmetrically between regions, while the trunk lineage may
be more stable remaining within a region (or within a highly connected network
of regions) that has more persistent transmission. In light of these results,
we suggest that future work on the phylogeography of influenza take into
account trunk vs. side branch differences in migration patterns.
## Discussion
### Correspondence between model and data
Although multiple epidemiological/evolutionary mechanisms have been proposed
to explain the restricted genetic diversity and rapid population turnover of
influenza A (H3N2) [4, 5, 6, 7], our results show that a simple model coupling
antigenic and genealogical evolution exhibits broad explanatory power. We find
a strong correspondence between the antigenic and genealogical patterns
generated by our model (Figure 1) and patterns of genetic and antigenic
evolution exhibited by influenza A (H3N2) [3, 9]. Our model suggests that
punctuated antigenic evolution need only be explained by a lack of mutational
opportunity and predicts that more detailed classification of influenza
strains will support a relatively small number of predominant phenotypes
(Figure 1B). We suggest that a large proportion of intra-cluster variation in
the observed antigenic map is due to experimental noise, rather than each
strain possessing a unique antigenic location. Additionally, our model
accurately predicts the contrasting dynamics of other types/subtypes of
influenza. We find that lowering mutation size/effect or lowering intrinsic
$R_{0}$ results in decreased incidence, slower antigenic movement and greater
genealogical diversity, all distinguishing characteristics of H1N1 influenza
and influenza B (Figure S8, Figure S9).
In our model, when antigenic phenotype remains static, there may be multiple
consecutive seasons without appreciable incidence (Figure 1A), a pattern
apparently absent from H3N2 influenza [16]. We suggest that any model
exhibiting punctuated evolution broadly consistent with the punctuated change
seen in the experimental antigenic map will show similar patterns of
incidence. We can ‘fix’ the incidence patterns, but at the cost of too smooth
an antigenic map (Figure S3). Evolutionary patterns of the neuraminidase (NA)
protein may provide an explanation. Epitopes in the HA and NA proteins are
jointly responsible for determining antigenicity [2], and it is now clear that
levels of adaptive evolution are similar between HA and NA [17]. Thus, changes
in NA may be driving incidence patterns as well, resulting in an observed
timeseries of incidence partially divorced from the antigenic map of HA.
It remains a central question as to the extent that short-lived strain-
transcending immunity is responsible for influenza’s limited diversity and
spindly genealogical tree [4, 5]. Our findings suggest a possible resolution.
Although lacking short-lived immunity, our model shows a detailed
correspondence to both the antigenic map and genealogical tree of H3N2
influenza. If an antigenic map were to show a deep bifurcation, where two
viral lineages move in different antigenic directions, then we would expect
the same bifurcation to be evident in the genealogical tree. Short-lived
strain-transcending immunity provides a mechanism by which lineages may
diverge in antigenic phenotype, but still show epidemiological interference.
This mechanism would explain a situation where bifurcations emerge in the
antigenic map, but competition results in the extinction of divergent
antigenic lineages. The empirical antigenic map [9] suggests that this is not
the case; one cluster leads to another cluster in orderly succession and there
is never competition between antigenically distant clusters. This supports the
hypothesis that antigenic evolution is primarily limited by a lack of
mutational availability. This is not to say that short-lived strain-
transcending immunity is not present; observed interference between subtypes
[4, 18] and evolution at CTL epitopes [19] provides substantial evidence for
its existence. Instead, we suggest that short-lived strain-transcending
immunity does not automatically generate antigenic maps and genealogical trees
consistent with empirical evidence.
### Linear antigenic movement
It would seem possible for one viral lineage to move in one antigenic
direction, while another lineage moves tangentially, eventually resulting in
two non-interacting viral lineages. Instead, we find that only movement in a
single antigenic direction is favored. The origins of this pattern can be seen
in the interaction between virus evolution and host immunity (Figure 3). As
the virus population evolves forward it leaves a wake of immunity in the host
population, and evolution away from this immunity results in the canalization
of the antigenic phenotype; mutations that continue along the line of primary
antigenic variation will show a transmission advantage compared to more
tangential mutations.
Following the work of Smith et al. [9], it remained an open question of why a
two-dimensional map should explain the antigenic variation of H3N2 influenza.
Although the authors astutely speculated that “there is a selective advantage
for clusters that move away linearly from previous clusters as they most
effectively escape existing population-level immunity, and this is a plausible
explanation for the somewhat linear antigenic evolution in regions of the
antigenic map.” This hypothesis remained to be tested. Here, we show from a
simple model of epidemiology and evolution that a linear trajectory of
antigenic evolution dynamically emerges due to basic selective pressures. This
result simultaneously explains the linear pattern of antigenic drift [9] and
the characteristically spindly genealogical tree [3] exhibited by influenza A
(H3N2).
To consider to what extent these results were contingent on the dimensionality
of the underlying antigenic model, we further implemented our model in a
10-dimensional antigenic space. Here, mutations occur as 10-spheres, but the
distance moved by a mutation is the same as in the previous two-dimensional
formulation. We arrive at nearly the same results with this model; principal
components analysis shows that the first and second dimensions of variation
account for 87% and 7%, respectively, of the total variance (Figure S10).
Thus, our model predicts that future work probing mutational effects will
support an underlying high-dimensional antigenic space, even though a two-
dimensional map is sufficient to explain observed antigenic relationships
among strains.
### Winding back the tape
It seems clear that, in our model, selection reduces the degrees of freedom of
antigenic evolution. In light of this, we wanted to examine the degree of
stochasticity in replicated evolutionary trajectories, and thereby test what
happens when we “wind back the tape” [20] on the evolution of the virus. We
ran 100 replicate simulations, each starting from the endpoint of the original
40-year simulation (Figure 4, Figure 5). Initially, we find a great detail of
repeatability; during the first year of evolution, every replicate virus
population undergoes a similar antigenic transition (Figure 4), resulting in a
repeatable peak in northern hemisphere incidence (Figure 5). After three
years, repeatability has mostly disappeared, with antigenic phenotype and
incidence appearing highly variable across replicates (Figure 4, Figure 5).
The 1–2 year timescale of repeatability can be explained by the presence of
standing antigenic variation. In the initial virus population, there are
several novel antigenic variants present at low frequency (Figure 3), one of
which, without fail, comes to predominate the virus population.
We see that the initial evolutionary trajectory, during which time standing
variation plays out, is highly repeatable, and thus predictable given enough
information and the right methods of analysis. However, prediction of longer-
term evolutionary scenarios will necessarily be difficult or impossible except
in a vague sense. Through careful surveillance efforts and genetic and
antigenic characterization of influenza strains, the World Health Organization
makes twice-yearly vaccine strain recommendations [21]. It should be possible
to combine these sorts of modeling approaches with surveillance data to gauge
the likelihood that a sampled variant will spread through the population.
Recent work on empirical fitness landscapes has shown that natural selection
follows few mutational paths [22]. The spindly genealogical tree and the
almost linear serial replacement of influenza strains has remained a puzzling
phenomenon. We suggest that the evolutionary and epidemiological dynamics
displayed by the influenza virus may simply be explained as an outgrowth of
selection to avoid host immunity. Natural selection can only ‘see’ one step
ahead, and so sacrifices long-term gains for short-term advantages. The result
is a canalized evolutionary trajectory lacking antigenic diversification.
## Materials and Methods
### Transmission model
To characterize the joint epidemiological, genealogical, antigenic and spatial
patterns of influenza, we implemented a large-scale individual-based model.
This model consists of daily time steps, in which the states of hosts and
viruses are updated. Hosts may be born, may die, may contact other hosts
allowing viral transmission, or may recover from infection. Viruses may mutate
in antigenic phenotype. Each simulation ran for 40 years of model time.
Hosts in this model are divided between three regions: North, South and
Tropics. There are 30 million hosts within each of the three regions, giving
$N=9\times 10^{7}$ hosts. Host population size remains fixed at this number,
but vital dynamics cause births and deaths of hosts at a rate of $1/30$ years
$=9.1\times 10^{-5}$ per host per day. Within each region, transmission
proceeds through mass-action with contacts between hosts occurring at a rate
of $\beta=0.36$ per host per day. Regional transmission rates in temperate
regions vary according to sinusoidal seasonal forcing with amplitude
$\epsilon=0.15$ and opposite phase in the North and in the South. Transmission
rate does not vary over time in the Tropics. Transmission between region $i$
and region $j$ occurs at rate $m\,\beta_{i}$, where $m=0.001$ and is the same
between each pair of regions and $\beta_{i}$ is the within-region contact
rate. Hosts recover from infection at rate $\nu=0.2$ per host per day, so that
$R_{0}$ in a naive host population is 1.8. There is no super-infection in the
model.
Each virus possesses an antigenic phenotype, represented as a location in
Euclidean space. Here, we primarily use a two-dimensional antigenic location.
After recovery, a host ‘remembers’ the phenotype of its infecting virus as
part of its immune history. When a contact event occurs and a virus attempts
to infect a host, the Euclidean distance from infecting phenotype $\phi_{v}$
is calculated to each of the phenotypes in the host immune history
$\phi_{h_{1}},\dots,\phi_{h_{n}}$. Here, one unit of antigenic distance is
designed to correspond to a twofold dilution of antiserum in a
hemagglutination inhibition (HI) assay [9]. The probability that infection
occurs after exposure is proportional to the distance $d$ to the closest
phenotype in the host immune history. Risk of infection follows the form
$\rho=\textrm{min}\\{d\,s,1\\}$, where $s=0.07$. Cross-immunity $\sigma$
equals $1-\rho$. The initial host population begins with enough immunity to
slow down the initial virus upswing and place the dynamics closer to their
equilibrium state; initial $R$ was 1.28.
Our model follows Gog and Grenfell [23] in representing antigenic distance as
distance between points in a geometric space. By forcing one-dimension to
directly modulate $\beta$, Gog and Grenfell find an orderly replacement of
strains. Kryazhimskiy et al. [24] use a two-dimensional strain-space, but
enforce a cross-immunity kernel that directly favors moving along a diagonal
line away from previous strains. Our model does not ‘build in’ the one-
dimensional direction of antigenic drift, which instead emerges dynamically
from competition among strains.
The initial virus population consisted of 10 infections each with the
identical antigenic phenotype of $\\{0,0\\}$. Over time viruses evolve in
antigenic phenotype. Each day there is a chance $\mu=10^{-4}$ that an
infection mutates to a new phenotype. This mutation rate represents a
phenotypic rate, rather than genetic mutation rate, and can be thought of as
arising from multiple genetic sources. When a mutation occurs, the virus’s
phenotype is moved in a completely random direction
$\sim\textrm{Uniform}(0,360)$ degrees. Mutation size is sampled from the
distribution $\sim\textrm{Gamma}(\alpha,\beta)$, where $\alpha$ and $\beta$
are chosen to give a mean mutation size of 0.6 units and a standard deviation
of 0.4 units. This distribution is parameterized so that mutation usually has
little effect on antigenic phenotype, but occasionally has a large effect.
This is similar to the neutral networks implemented by Koelle et al. [6],
wherein most amino acid changes result in little decrease to cross-immunity
between strains, but some changes result in large jumps in cross-immunity.
### Model output
Daily incidence and prevalence are recorded for each region. During the course
of the simulation, samples of current infections are taken from the evolving
virus population at a rate proportional to prevalence. Each viral infection is
assigned a unique ID, and in addition, infections have their phenotypes,
locations and dates of infection recorded. In this model, viruses lack
sequences and so standard phylogenetic inference of the evolutionary
relationships among strains is impossible. Instead, the viral genealogy is
directly recorded. This is made possible by tracking transmission events
connecting infections during the simulation; infections record the ID of their
‘parent’ infection. Proceeding from a sample of infections, their genealogical
history can be reconstructed by following consecutive links to parental
infections. During this procedure, lineages coalesce to the ancestral lineages
shared by the sampled infections, eventually arriving at the initial infection
introduced at the beginning of the simulation. Commonly, phylodynamic
simulations generate sequences that are subsequently analyzed with
phylogenetic software to produce an estimated genealogy [4, 6, 25]. This step
of phylogenetic inference is imperfect and computationally intensive, and by
side-stepping phylogenetic reconstruction we arrive at genealogies quickly and
accurately. Other authors have implemented similar tracking of infection trees
[26, 27]. This genealogy-centric approach makes many otherwise difficult
calculations transparent, such as calculating lineage-specific region-specific
migration rates (Figure 2) and lineage-specific mutation effects (Figure S5).
Infections are sampled at a rate designed to give approximately 6000 samples
over the course of the simulation, with genealogies constructed from a
subsample of approximately 300 samples. The results presented in Figure 1
represent a single representative model output; one hundred replicate
simulations were conducted to arrive at statistical estimates.
### Parameter selection and sensitivity analysis
Estimating what the basic reproductive number $R_{0}$ for seasonal influenza
would be in a naive population is notoriously difficult. Season-to-season
estimates of effective reproductive number $R$ for the USA and France gathered
from mortality timeseries display an interquartile range of 0.9–1.8 [28].
Geographic spread within the USA suggests an average seasonal $R$ of 1.35
[29]. These estimates of $R$ will be lower than the $R_{0}$ of influenza due
to the effects of human immunity. We assumed $R_{0}$ of 1.8, consistent with
the upper range of seasonal estimates. Duration of infection was chosen based
on patterns of viral shedding shown during challenge studies [30]. The linear
form of the risk of infection and its increase as a function of antigenic
distance $s$ was chosen as 0.07 based on experimental work on equine influenza
[31] and from studies of vaccine effectiveness [32]. Between-region contact
rate $m$ was chosen to yield a trunk lineage that resides predominantly in the
tropics. With much higher rates of mixing, the trunk lineage ceases to show a
preference the tropics, and with much lower rates of mixing, particular
seasons in the north and the south will often be skipped. The amplitude of
seasonal forcing $\epsilon$ was chosen to be just large enough to get
consistent fade-outs in the summer months and is consistent with empirical
estimates [33].
Mutational parameters were based, in part, on model behavior. We assumed 10
amino acid sites involved in antigenicity, each mutating at a rate of
$10^{-5}$ [12] to give a phenotypic mutation rate $\mu=10^{-4}$ per infection
per day. We chose mutational effect parameters ($\textrm{mean}=0.6$,
$\textrm{sd}=0.4$) that would give suitably fast rates of antigenic evolution
corresponding to approximately 1.2 units of antigenic change per year, while
simultaneously giving clustered patterns of antigenic evolution [9]. Similar
outcomes are possible under a variety of parameterizations. If mutations are
more common ($\mu=3\times 10^{-4}$) and show less variation in effect size
($\textrm{mean}=0.6$, $\textrm{sd}=0.2$), then antigenic drift occurs in a
more continuous fashion, resulting in less variation in seasonal incidence and
a smoother distribution of antigenic phenotypes (Figure S3). If mutations are
less common ($\mu=5\times 10^{-5}$) and show more variance in effect
($\textrm{mean}=0.7$, $\textrm{sd}=0.5$), then antigenic change occurs in a
more punctuated fashion (Figure S4). Basic reproductive number $R_{0}$ can be
traded off with mutational parameters to some extent. Less mutational input
and higher $R_{0}$ will give similar patterns of antigenic drift and seasonal
incidence. Similarly, Kucharski and Gog [34] find that increasing $R_{0}$
results in increased rates of emergence of antigenically novel strains.
In 20 out of the 100 replicate simulations, we observed a major bifurcation of
antigenic phenotype and a consequent increase in incidence and genealogical
diversity. These simulations were removed from the analysis. Similar to Koelle
et al. [35], we assume that although the historical evolution of H3N2
influenza followed the path of a single lineage, it could have included a
major bifurcation. Further work in these directions will help to determine the
likelihoods of single lineage vs. bifurcating scenarios.
### Antigenic map
Antigenic phenotypes are modeled as discrete entities on the Euclidean plane;
multiple samples have the same antigenic location. However, in the empirical
antigenic map of influenza A (H3N2), each strain appears in a unique location
[9]. We would argue that some of this pattern comes from experimental noise.
Indeed, Smith et al. [9] find that observed measurements and measurements
predicted from the map differ by an average of 0.83 antigenic units with a
standard deviation of 0.67 antigenic units. We take this as a proxy for
experimental noise and add jitter to each sampled antigenic phenotype by
moving it in a random direction for an exponentially distributed distance with
mean of 0.53 antigenic units. If two samples with the same underlying
antigenic phenotype are jittered in this fashion, the distance between them
averages 0.83 antigenic units with a standard deviation of 0.64 units.
We added noise to each of the 5943 sampled viruses in this fashion resulting
in an approximated antigenic map (Figure 1D). Virus samples were then
clustering following standard clustering algorithms. We tried clustering by
the $k$-means algorithm and also agglomerative hierarchical clustering with a
variety of linkage criterion. We found that clustering by Ward’s criterion
consistently outperformed other methods, when measured in terms of within-
cluster and between-cluster variances. However, the exact clustering algorithm
had little effect on our overall results.
### Acknowledgments
We would like to thank Sarah Cobey, Aaron King, Pejman Rohani and the
attendees of the 2011 RAPIDD Workshop on Phylodynamics for helpful discussion.
We would also like to thank Ed Baskerville and Daniel Zinder for programming
advice. The term ‘canalization’ was originally suggested by Micaela Martinez-
Bakker.
### Funding
TB is supported by the Howard Hughes Medical Institute and by the European
Molecular Biology Organization. AR works as a part of the Interdisciplinary
Centre for Human and Avian Influenza Research (ICHAIR) and the University of
Edinburgh’s Centre for Immunity, Infection and Evolution (CIIE). MP is an
investigator of the Howard Hughes Medical Institute.
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## Tables
Table Table 1: Rates of mutation and phenotypic change on trunk and side branches and mutational expectation. | Baseline | Side branch | Trunk | Trunk / side branch
---|---|---|---|---
Mutation size (AG units) | 0.60 | 0.79 | 1.58 | 1.99$\times$
Mutation rate (mut per year) | 0.04 | 0.06 | 0.81 | 13.23$\times$
Antigenic flux (AG units per year) | 0.02 | 0.05 | 1.27 | 26.25$\times$
## Figures
Figure Figure 1: Simulation results showing epidemiological, antigenic and
genealogical dynamics. (A) Weekly timeseries of incidence of viral infection
in north and tropics regions. (B) Two-dimensional antigenic phenotypes of 5943
viruses sampled over the course of the simulation. Each discrete virus
phenotype is shown as a bubble, with bubble area proportional to the number of
times this phenotype was sampled. (C) Genealogical tree depicting the
infection history of 376 samples from the virus population. Parent/offspring
relationships were tracked over the course of the simulation, giving a direct
observation of the genealogy rather than a phylogenetic inference. (D)
Antigenic map depicting phenotypes of 5943 viruses sampled over the course of
the simulation. To construct the map, noise was added to each sample and the
resulting observations grouped into 11 clusters and colored accordingly. Grid
lines show single units of antigenic distance. Cluster assignments were used
to color all panels in a consistent fashion.
Figure Figure 2: Patterns of geographic movement of virus lineages. (A)
Evolutionary relationships among 376 viruses sampled evenly through time
colored by geographic location. Lineages residing in the north (N), south (S)
and tropics (T) are colored yellow, red and blue respectively. (B) Observed
migration rates between regions on side branch lineages (left) and on trunk
lineages (right). Arrows denote movement of lineages and arrow width is
proportional to migration rate. Circle area is proportional to the expected
stationary frequency of a region given the observed migration rates. In both
cases, migration rates are calculated across 80 replicate simulations.
Figure Figure 3: Host immunity and antigenic history of the virus population.
Contour lines represent the state of host immunity at the end of the 40-year
simulation. They show the mean risk of infection (as a percentage) after a
random host in the population encounters a virus bearing a particular
antigenic phenotype. Contour lines are spaced in intervals of 2.5%. Bubbles
represent a sample of antigenic phenotypes present at the end of the 40-year
simulation. The area of each bubble is proportional to the number of samples
with this phenotype. Lines leading into these bubbles show past antigenic
history. The current phenotypes rapidly coalesce to a trunk phenotype. The
movement of the virus population from the left to the center of the figure can
be seen from the antigenic history. At the end of the simulation several virus
phenotypes exist with similar antigenic locations; all of these phenotypes lie
significantly ahead of the peak of host immunity.
Figure Figure 4: Antigenic phenotypes over the course of 4 years of evolution
across 100 replicate simulations starting from identical initial conditions.
Replicate simulations were initialized with the end state of the 40-year
simulation shown in Figure 1. Each panel shows an additional year of
evolution, with black points representing the mean antigenic phenotypes of the
100 replicate simulations and gray lines representing the history of each mean
antigenic phenotype.
Figure Figure 5: Timeseries of incidence across 100 replicate simulations with
identical initial conditions. Panels show incidence in the North, Tropics and
South regions over the course of 6 years. Solid black lines represent the
median weekly incidence across the 100 replicate simulations, while gray
intervals represent the interquartile range across simulations. There is
little variability for the first year of replicate simulations. Replicate
simulations were initialized with the end state of the 40-year simulation
shown in Figure 1.
## Supporting Information
Figure Figure S1: Antigenic evolution over the course of the 40-year
simulation. (A) Proportion of virus population comprised of each antigenic
cluster through time. (B) Antigenic distance from initial phenotype ($x=0$,
$y=0$) for each of 5943 virus samples relative to time of virus sampling.
Viruses were sampled at a constant rate proportional to prevalence and
coloring was determined from the antigenic map in Figure 1D.
Figure Figure S2: Correlation between antigenic drift and attack rate.
Antigenic drift is measured as the distance between the centroid of phenotypes
of one year and the centroid of phenotypes of the following year. Measurements
were taken across 80 replicate simulations. Individual pairs of measurements
are shown as gray points and a locally-linear regression (LOESS) is shown as a
black dashed line.
Figure Figure S3: Simulation results showing epidemiological, antigenic and
genealogical dynamics for ‘smoother’ mutation model. (A) Weekly timeseries of
incidence of viral infection in north and tropics regions. (B) Antigenic map
depicting phenotypes of viruses sampled over the course of the simulation.
Grid lines show single units of antigenic distance. (C) Genealogical tree
depicting the infection history of samples from the virus population. Cluster
assignments were used to color panels (A), (B) and (C) in a consistent
fashion. Alternative mutational parameters are $\mu=3\times 10^{-4}$, mean
mutation size of 0.6 units and standard deviation of mutation size of 0.2
units.
Figure Figure S4: Simulation results showing epidemiological, antigenic and
genealogical dynamics for ‘rougher’ mutation model. (A) Weekly timeseries of
incidence of viral infection in north and tropics regions. (B) Antigenic map
depicting phenotypes of viruses sampled over the course of the simulation.
Grid lines show single units of antigenic distance. (C) Genealogical tree
depicting the infection history of samples from the virus population. Cluster
assignments were used to color panels (A), (B) and (C) in a consistent
fashion. Alternative mutational parameters are $\mu=5\times 10^{-5}$, mean
mutation size of 0.7 units and standard deviation of mutation size of 0.5
units.
Figure Figure S5: Mutation spectrum in two-dimensional antigenic space of side
branch mutations and trunk mutations. (A) Histogram of mutation effects along
the axis of primary antigenic variation across 80 replicate simulations. The
left panel shows the distribution of effects of side branch mutations and the
right panel shows the distribution of effects of trunk mutations. (B) Smoothed
two-dimensional histogram of mutation effects along the primary and secondary
axes of antigenic variation across 80 replicate simulations. Histograms were
constructed from 21,405 side branch mutations and 1584 trunk mutations.
Figure Figure S6: Relationship between a mutation’s phenotypic effect and its
likelihood of being part of the trunk. The $x$-axis represents the effect of a
mutation along the line of primary antigenic variation, and the $y$-axis
represents the probability that the mutation is part of the trunk. Mutations
of large effect are increasingly rare, but when they do occur are increasingly
likely to be part of the trunk.
Figure Figure S7: Observed vs. expected distributions of waiting times between
phenotypic mutations along genealogy trunk. (A) Histogram bins show the
observed distribution of waiting times in years across 80 replicate
simulations representing 1584 mutations. The mean of this distribution is 1.76
years. The dashed line shows the Poisson process expectation of exponentially
distributed waiting times. (B) The density distribution of waiting times is
transformed into a hazard function, representing the rate of trunk mutation
after a specific waiting time. The dashed line shows the memoryless hazard
function of the Poisson process expectation.
Figure Figure S8: Simulation results showing epidemiological, antigenic and
genealogical dynamics with weaker mutation. (A) Weekly timeseries of incidence
of viral infection in north and tropics regions. (B) Antigenic map depicting
phenotypes of viruses sampled over the course of the simulation. Grid lines
show single units of antigenic distance. (C) Genealogical tree depicting the
infection history of samples from the virus population. Cluster assignments
were used to color panels (A), (B) and (C) in a consistent fashion. Here,
$\mu=5\times 10^{-5}$, mean mutation size is 0.42 units and standard deviation
of mutation size is 0.28 units.
Figure Figure S9: Simulation results showing epidemiological, antigenic and
genealogical dynamics with lower intrinsic $R_{0}$. (A) Weekly timeseries of
incidence of viral infection in north and tropics regions. (B) Antigenic map
depicting phenotypes of viruses sampled over the course of the simulation.
Grid lines show single units of antigenic distance. (C) Genealogical tree
depicting the infection history of samples from the virus population. Cluster
assignments were used to color panels (A), (B) and (C) in a consistent
fashion. Here, $\beta=0.3$, giving $R_{0}=1.5$.
Figure Figure S10: Principal components of antigenic variation under a
10-sphere mutation model. Each panel shows 5991 samples of antigenic phenotype
over the course of a 40-year simulation. Each phenotype is represented as a
bubble, with bubble area proportional to the number of samples with this
phenotype. Bubbles are colored based on clustering the 10-dimensional
antigenic phenotypes. The original 10-dimensional space was rotated using
principal components analysis to give orthogonal axes in the order of their
contribution to antigenic variation. Each panel shows a two-dimensional slice
of the this rotated space. Principal components 7–10 were left out of the
figure for clarity.
|
arxiv-papers
| 2011-11-19T19:41:26 |
2024-09-04T02:49:24.493026
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Trevor Bedford, Andrew Rambaut and Mercedes Pascual",
"submitter": "Trevor Bedford",
"url": "https://arxiv.org/abs/1111.4579"
}
|
1111.4593
|
footnote
# A graph counterexample to davies’ conjecture
Gady Kozma
###### Abstract.
There exists a graph with two vertices $x$ and $y$ such that the ratio of the
heat kernels $p(x,x;t)/p(y,y;t)$ does not converge as $t\to\infty$.
This paper is concerned with a conjecture of Brian Davies from 1997 on the
heat kernel of Riemannian manifolds [D97, §5]. We will not disprove the
conjecture as stated, but rather transform it to the realm of graphs using a
well-known (though informal) “dictionary” between these two categories, and
build a graph that will serve as a counterexample. We will make some remarks
on how the construction might be carried over back to the category of
manifolds in the end, but we will not give all details. The bulk of this paper
is about graphs.
We start by describing the conjecture in its original setting. Let $M$ be a
connected Riemannian manifold, and let $p$ be the _heat kernel_ associated
with the _Laplace-Beltrami operator_ on $M$. Then Davies conjecture states
that for any $M$ and any 3 points $x,y,z\in M$ the limit
$\lim_{t\to\infty}\frac{p(x,y;t)}{p(z,z;t)}$ (1)
exists and is positive. Here $p(x,y;t)$ is the value of the heat kernel at
time $t$ and at points $x$ and $y$. This property is known as the “strong
ratio limit property” (where the “weak” version is an averaged result due to
Döblin, [D38]) or SRLP for short. So Davies’ conjecture is that in these
settings SRLP always holds. SRLP holds for manifolds with one end [D82] and
for strongly Liouville manifolds (i.e. manifolds where any positive harmonic
function is constant), see [ABJ02, corollary 2.7] who also make interesting
connections between these properties and the infinite Brownian loop.
Ratio limit properties were considered for Markov chains even earlier. If $M$
is any Markov chain on a countable state space, then we say that $M$ satisfies
SRLP if (1) holds for any three states $x$, $y$ and $z$, where $p(x,y;t)$ is
the probability that the Markov chain started at $x$ will be at $y$ at time
$t$. For general Markov chains there are a few examples where SRLP does not
hold. Clearly it does not hold when the Markov chain has some kind of
periodicity. F. J. Dyson constructed an example of an aperiodic recurrent
Markov chain which does not satisfy SRLP [C60, part I, §10]. That example
utilizes long chains of states with only one outgoing edge, which the walker
must traverse sequentially. In particular it is not _reversible_.
Now, the Laplace-Beltrami operator is self-adjoint so a proper analog of
Davies conjecture needs to assume that the Markov chain is reversible.
Reversible Markov chains are also known as random walks on weighted graphs.
The issue of periodicity can be dealt with by looking at random walk in
continuous time or at lazy random walk. Lazy random walk is a walk where the
walker, at every step, chooses with probability $\frac{1}{2}$ to stay where it
is, and with probability $\frac{1}{2}$ moves to one of the neighbours (with
probability proportional to the weights).
The result is this paper is
###### Theorem.
There exists a connected graph $G$ with bounded weights and two vertices
$x,y\in G$ such that the heat kernel of the lazy random walk satisfies
$\frac{p(x,x;t)}{p(y,y;t)}\nrightarrow$ (2)
as $t\to\infty$.
Let us remark on the “bounded weights” clause. When doing analogies between
manifolds and graphs, it is often assumed that the manifold has bounded
geometry and the graph has bounded weights. Davies, however, explicitly does
not make the assumption of bounded geometry. Thus one might wonder what is
exactly the graph analog. All this is mute, of course, since the
counterexample does has bounded weights (and hence a manifold example
constructed along the same line should have bounded geometry).
It is easy to see that in these setting (reversible, irreducible, lazy) this
ratio must be bounded between two constants independent of $n$. Hence if it
does not converge then it must fluctuate between two values. The proof
constructs the graph with two halves, denoted $H^{e}$ and $H^{o}$ ($e$ and $o$
standing for even and odd, $H$ standing for half), which are connected by one
edge, $(x,y)$. On the “odd scales”, $H^{e}$ will “look like $\mathbb{Z}^{22}$”
while $H^{o}$ will “look like $\mathbb{Z}^{3}$”. This means that to get from
$x$ to $x$ (where $x$ is on the $H^{e}$ side), the most beneficial strategy is
to move to $y$ as fast as possible, spend most of your time on the $H^{o}$
side and return to the $H^{e}$ side only in the last minute. Clearly this
would mean that $p(x,x;t)$ is smaller than $p(y,y;t)$ as the random walk
starting from $y$ can stay on its side at all time, not losing the constant
that $x$ needs for the maneuver. On the “even scales” the picture is reversed
and $y$ is at a disadvantage. See figure 1 — drawing in 22 dimensions might
have distracted the reader, so the figure demonstrates the construction in 1
and 2 dimensions. The smallest two braces in the figure are the first scale,
in which $H^{o}$ is really one dimensional and $H^{e}$ is really two-
dimensional. The larger two braces indicate the second scale. This time
$H^{o}$ is a network of lines so it should be thought about as two
dimensional, while $H^{e}$ is a thick column, so it should be thought about as
one dimensional. The third scale is only hinted at in this figure, but one can
imagine that $H^{o}$ now becomes a thick band, so it is again one-dimensional,
while $H^{e}$ becomes a network of these thick columns and bands, so it is
back to being two dimensional.
$x$$y$$H^{e}$$H^{o}$
Figure 1. The graphs $H^{o}$ and $H^{e}$
As one might expect, the numbers $3$ and $22$ have no particular significance.
They both have to be $>2$, since otherwise our graph would be recurrent and
recurrent graphs always satisfy SRLP [O61, theorem 3]. And of course they have
to be different. We took here the large value $22$ in order to be able to be
wasteful in various points (sum over times and such stuff), but the proof
could proceed with any value larger than $3$.
This paper was first written in 2006. I wish to take this opportunity to
apologize to all those who has to wait so long for it to appear, with no real
reason. My intentions were good but my time management was abysmal. I wish to
thank Yehuda Pinchover for telling me about the problem and for reading early
drafts. Partially supported by the Israel Science Foundation.
## 1\. Proof
The construction uses $\mathbb{Z}^{d}$-like graphs as building blocks, so we
start with quoting a few results on these. We first recall the notion of rough
isometry [K85].
###### Definition.
Let $X$ and $Y$ be two metric spaces. We say that $X$ and $Y$ are roughly
isometric if there is a constant $C$ and a map $\varphi:X\to Y$ with the
following properties:
1. (i)
For all $x$ and $y$ in $X$,
$\frac{1}{C}d(x,y)-C\leq d(\varphi(x),\varphi(y))\leq Cd(x,y)+C$
2. (ii)
The image of $\varphi$ is roughly dense, i.e. for all $y\in Y$ there is an
$x\in X$ such that $d(y,\varphi(x))\leq C$
If $G$ and $H$ are graphs we say that they are roughly isometric if they are
roughly isometric when considered with the metric $d$ being the graph
distance, namely $d(x,y)$ is the length of the shortest path between $x$ and
$y$, or $\infty$ if no such path exists.
With this definition we can state the following standard result, essentially
due to Delmotte.
###### Lemma 1.
Let $G$ be a graph roughly isometric to $\mathbb{Z}^{d}$. Then the heat kernel
$p$ for the lazy walk on $G$ satisfies, for all $t\geq 1$,
$ct^{-d/2}\leq p(x,x;t)\leq Ct^{-d/2}.$
Here $G$ is a simple graph — we do not allow weights or multiple edges. $C$
and $c$ are constant which do not depend on $t$. In general we will use $c$
for constants which are small enough and $C$ for constants which are large
enough, and different appearance of $c$ and $C$ might relate to different
constants.
###### Proof.
By Delmotte’s theorem [D99] any $G$ which satisfies volume-doubling and the
Poincaré inequality, satisfies
$p(x,x;t)\approx\frac{1}{|B(x,\sqrt{t})|}$
where $B(x,r)$ is the ball around $x$ with radius $r$ (again with the graph
distance), and $|B(x,r)|$ is the sum of the degrees of the vertices in $B$.
The notation $X\approx Y$ is short for $cY\leq X\leq CY$. The fact that $G$ is
roughly-isometric to $\mathbb{Z}^{d}$ gives
$|B(x,r)|\approx r^{d}$ (3)
so we would get $p(x,x;t)\approx t^{-d/2}$, as needed. So we need only show
that $G$ satisfies volume doubling and Poincaré inequality.
Now, the definition of volume doubling is that for every $x$ a vertex of $G$
and every $r\geq 1$
$|B(x,2r)|\leq C|B(x,r)|$
and it follows immediately from (3). The Poincaré inequality is not much more
complicated. The definition is: for every vertex $x$, for every $r$ and for
every function $f:B(x,2r)\to\mathbb{R}$,
$\sum_{y\in B(x,r)}\deg(y)|f(y)-\overline{f}|^{2}\leq Cr^{2}\sum_{(y,z)\in
E(B(x,2r))}(f(y)-f(z))^{2}$ (4)
where
$\overline{f}=\frac{1}{|B(x,r)|}\sum_{y\in B(x,r)}\deg(y)f(y).$
and where $\deg(y)$ is the degree of $y$, and $E(B)$ is the set of edges both
whose vertices are in $B$. Now, $\mathbb{Z}^{d}$ satisfies the Poincaré
inequality (see e.g. [PSC, §4.1.1]). It is well-known and not difficult to see
that the Poincaré inequality is preserved by rough isometries (it uses the
fact that $\sum\deg(y)|f(y)-a|^{2}$ is minimized when $a=\overline{f}$). This
finishes the proof.∎
###### Lemma 2.
Let $G$ be a graph roughly-isometric to $\mathbb{Z}^{d}$, $d\geq 3$ and let
$x$ be some vertex. Let $p$ be the probability that lazy random walk starting
from $x$ returns to $x$ for the first time at time $t$. Then $p\geq
ct^{-d/2}$.
###### Proof.
Let $p_{1}$ be the same probability but without the restriction that this is
the first return to $x$. This is exactly $p(x,x;t)$ and by theorem 1 we have
$p_{1}\geq ct^{-d/2}$. Fix some $K$ and examine the event that the random walk
returns to $x$ at $t$ and also at some time $s\in[K,t-K]$. Let $p_{2}$ be its
probability. Using the other direction in theorem 1 we can write
$p_{2}\leq\sum_{s=K}^{t-K}p(x,x;s)p(x,x;t-s)\leq
C\sum_{s=K}^{t-K}s^{-d/2}(t-s)^{-d/2}\leq CK^{1-d/2}t^{-d/2}.$
Since $d\geq 3$ we can choose $K$ sufficiently large such that
$p_{2}\leq\frac{1}{2}p_{1}$. So we know that with probability $p_{1}-p_{2}\geq
ct^{-d/2}$ the walk does not return to $x$ between $K$ and $t-K$. If it does
reach $x$ before time $K$, do some local modification so that it does not. For
example, if the original walker reached $x$ at some time $s<K$ and on the next
step went to some neighbour $y$ of $x$, modify it to walk to $y$ in the first
step and stay there for $s$ steps (remember that our walk is lazy) and then
continue like the original walker. Clearly this “costs” only a constant and
ensures our walker does not visit $x$ in the interval $[1,K]$. Do the same for
the interval $[t-K,t-1]$, losing another constant. The finishes the lemma.∎
###### Lemma 3.
Let $G$ be a graph and let $p$ be the heat kernel for the lazy walk on $G$.
Let $t$ and $s$ satisfy that $|t-s|\leq\sqrt{t}$. Then
$|p(x,x;t)-p(x,x;s)|\leq
C\frac{|t-s|\log^{3}t}{\sqrt{t}}p(x,x;t)+Ce^{-c\log^{2}t}.$
###### Proof.
Denote by $q(x,y;t)$ the heat kernel for the _simple_ random walk on $G$. Then
by definition
$p(x,x;t)=\sum_{i=0}^{t}q(x,x;i)\binom{t}{i}2^{-t}.$
Writing the same formula for $p(x,x;s)$ and subtracting we get
$|p(x,x;t)-p(x,x;s)|\leq\sum_{i}q(x,x;i)\left(\binom{t}{i}2^{-t}-\binom{s}{i}2^{-s}\right)=\Sigma_{1}+\Sigma_{2}$
where $\Sigma_{1}$ is the sum over all $|t-2i|\leq\sqrt{t}\log t$ and
$\Sigma_{2}$ is the rest. A simple calculation with Stirling’s formula shows
that
$2^{-t}\binom{t}{i}=\sqrt{\frac{2}{\pi
t}}\exp\left(-\frac{(t-2i)^{2}}{2t}\left(1+O\left(\frac{|t-2i|+1}{t}\right)\right)\right)$
And with some more calculations
$\Sigma_{1}\leq C\sum_{|t-2i|\leq\sqrt{t}\log t}q(x,x;i)\sqrt{\frac{2}{\pi
t}}e^{-(t-2i)^{2}/2t}\frac{|t-s|\log^{3}t}{\sqrt{t}}\leq
C\frac{|t-s|\log^{3}t}{\sqrt{t}}p(x,x;t)$
while
$\Sigma_{2}\leq C\sum_{|t-2i|>\sqrt{t}\log t}e^{-c(t-2i)^{2}/t}\leq
Ce^{-c\log^{2}t}$
proving the lemma. ∎
###### Proof of the theorem.
Abusing notations, for subsets $H\subset\mathbb{Z}^{d}$ we will not
distinguish between $H$ as a set and as an induced subgraph of
$\mathbb{Z}^{d}$ ($d$ will be $22$). For the construction we need a
sufficiently fast increasing sequence $a_{1}<a_{2}<\dotsb$. We further assume
that $a_{k}$ are all even and that $a_{k-1}$ divides $\frac{1}{2}a_{k}$. It
would have probably been enough to choose $a_{k}=2^{a_{k-1}}$, but it turns
out simpler to choose the $a_{k}$ inductively, and we perform this as follows.
Let $a_{1}=2$. Assume now $a_{1},\dotsc,a_{k-1}$ have been defined. Define,
for integers $m<\frac{1}{2}l$ and $i\in\\{1,\dotsc,22\\}$,
$\displaystyle Q_{l,m,i}$
$\displaystyle:=\left\\{\vec{n}\in\mathbb{Z}^{22}:\left|n_{i}\textrm{ mod
}l\right|\leq m\right\\}$ $\displaystyle Q_{l,m}$
$\displaystyle:=\bigcup_{\begin{subarray}{c}I\subset\\{1,\dotsc,22\\}\\\
|I|=19\end{subarray}}\bigcap_{i\in I}Q_{l,m,i}.$
Here $n\textrm{ mod
}l\in\\{-\lfloor\frac{l-1}{2}\rfloor,\dotsc,\lfloor\frac{l}{2}\rfloor\\}$. In
words, $Q_{l,m,i}$ is a $21$-dimensional subspace of $\mathbb{Z}^{22}$
orthogonal to one of the axes, fattened up by $2m+1$ (a “slab”) and repeated
periodically with period $l$. $Q_{l,m}$ is the collection of all
$3$-dimensional subspaces, fattened and repeated similarly. The particular
point $\vec{0}$ is in fact contained in all $\binom{22}{3}$ of these
$3$-dimensional slabs which will be a little inconvenient, so let us shift
$Q_{l,m}$ by
$v(m)=\Big{(}\underbrace{{\textstyle\frac{1}{2}}m,\dotsc,{\textstyle\frac{1}{2}}m}_{3\textrm{
times}},\underbrace{\vphantom{{\textstyle\frac{1}{2}}m}0,\dotsc,0}_{19\textrm{
times}}\Big{)}.$
In the shifted set $Q_{l,m}+v(m)$ the geometry of the neighbourhood of
$\vec{0}$ is simpler, it is contained in just one slab. Compare to the figure
on page 1. The point $x$ is in the _middle_ of a fat column and not at the
intersection of a column and a band.
We want to use these graphs with $l=a_{j}$ and $m$ a little larger than
$a_{j-1}$. Precisely, define
$b_{j}=\sum_{k=1}^{j}a_{k}.$
With this choice of $b_{j}$, $Q(a_{j},b_{j-1},i)$ contains only complete
components of $Q(a_{l},\linebreak[4]b_{l-1},i)$ for each $l<j$. Each such
component is either contained in $Q(a_{j},b_{j-1},i)$ or disjoint from it. The
same holds for the translations $Q(a_{l},b_{l-1},i)+v(a_{l})$ (we need here
that $a_{l}>4a_{l-1}$ so let us assume this from now on). For brevity, define
$v_{j}=v(a_{j})$.
We may now define two graphs, denoted by $H_{k-1}^{\textrm{e}}$ and
$H_{k-1}^{o}$ (“e” and “o” standing for even and odd) by
$H_{k-1}^{\textrm{e/o}}:=\bigcap_{\begin{subarray}{c}2\leq j\leq k-1\\\
j\textrm{ even/odd}\end{subarray}}(Q_{a_{j},b_{j-1}}+v_{j}).$
We shall usually suppress the $k-1$ from the notation. It is not difficult to
check that $H^{\textrm{e/o}}$ are both roughly isometric to $\mathbb{Z}^{22}$
(the rough isometry constant depends on the “past” $a_{1},\dotsc,a_{k-1})$.
Therefore by lemma 1 we see that there exists an $\alpha$ (again, depending on
the past) such that
$p_{H^{\textrm{e/o}}}(x,x;t)\leq\alpha t^{-11}.$ (5)
Examine now the graphs
$F_{k-1}^{\textrm{e/o}}:=H_{k-1}^{\textrm{e/o}}\cap\left\\{\vec{n}\in\mathbb{Z}^{22}:\left|n_{i}\right|\leq
b_{k-1}\>\forall i=4,\dotsc,22\right\\}.$
$F^{\textrm{e/o}}$ are both roughly isometric to $\mathbb{Z}^{3}$ so by lemma
2 there exists some $\beta$ such that
$\mathbb{P}_{F^{\textrm{e/o}}}(\mbox{the walk returns to }\vec{0}\mbox{ for
the first time at }t)\geq\frac{1}{\beta}t^{-3/2}.$ (6)
Define $\gamma_{k}:=\left\lceil\max\\{\alpha,\beta\\}\right\rceil$ (as usual,
$\left\lceil\cdot\right\rceil$ stands for the upper integer value). With these
we can define $a_{k}$ to be any even number satisfying
$a_{k}>2\gamma_{k}^{4}+4a_{k-1}$ and such that $a_{k-1}$ divides
$\frac{1}{2}a_{k}$. This completes the description of the induction, and we
define
$H_{\infty}^{\textrm{e/o}}:=\bigcap_{\begin{subarray}{c}2\leq j\\\ j\textrm{
even/odd}\end{subarray}}(Q_{a_{j},a_{j-1}}+v_{j}).$
These graphs will be the two halves of our target graph $G$.
Before continuing, let us collect some simple facts about
$H_{\infty}^{\textrm{e/o}}$:
1. (i)
$H_{\infty}^{\textrm{e/o}}$ is connected — in fact we used this indirectly
when we claimed $H_{k}^{\textrm{e/o}}$ are roughly isometric to
$\mathbb{Z}^{22}$.
2. (ii)
$H_{\infty}^{\textrm{e/o}}$ are transient — this follows because each contains
a copy of $\mathbb{Z}^{3}$ (namely $\\{n_{4}=\dotsb=n_{22}=0\\}$) and
transience is preserved on adding edges. This last fact follows from
conductance arguments, see e.g. [DS84].
Define therefore the escape probabilities
$\varepsilon^{\textrm{e/o}}:=\mathbb{P}_{H_{\infty}^{\textrm{e/o}}}^{\vec{0}}(R(t)\neq\vec{0}\,\forall
t>0)$
($R$ being the random walk on the graph) and let
$\delta:=\frac{1}{2}\min\\{\varepsilon^{\textrm{e}},\varepsilon^{\textrm{o}}\\}$.
Define the graph $G$ by connecting $H_{\infty}^{\textrm{e}}$ to
$H_{\infty}^{\textrm{o}}$ with a single edge between the two $\vec{0}$ with
weight $\delta$. Define $x:=\vec{0}^{\textrm{e}}$ and
$y=\vec{0}^{\textrm{o}}$. This is our construction and we need to show (2),
which will follow if we show that, for $k$ sufficiently large,
$\left.\begin{aligned} p(x,x;t_{2k})&\geq 3p(y,y;t_{2k})\\\
p(x,x;t_{2k+1})&\leq{\textstyle\frac{1}{3}}p(y,y;t_{2k+1})\end{aligned}\right\\}\quad
t_{k}:=\gamma_{k}^{4}.$ (7)
We will only prove the even case, the odd will follow similarly.
Examine therefore $p(x,x;t_{2k})$. Since $a_{2k}>t_{2k}$ we get that
$H_{\infty}^{\textrm{e/o}}\cap[-t_{2k},t_{2k}]^{22}=H_{2k}^{\textrm{e/o}}\cap[-t_{2k},t_{2k}]^{22}$
or in other words, the steps after $2k$ do not effect us at all. Similarly it
is possible to simplify the last stage namely
$\displaystyle H_{2k}^{\textrm{e}}\cap[-t_{2k},t_{2k}]^{22}$
$\displaystyle=H_{2k-1}^{\textrm{e}}\cap(Q_{a_{2k},b_{2k-1}}+v_{2k})\cap[-t_{2k},t_{2k}]^{22}=$
$\displaystyle=H_{2k-1}^{\textrm{e}}\cap\left\\{\vec{n}\in\mathbb{Z}^{22}:\left|n_{i}\right|\leq
b_{2k-1}\>\forall i=4,\dotsc,22\right\\}=F_{2k-1}^{\textrm{e}}$
(here is where these translations by $v_{j}$ are used). By (6),
$\displaystyle p_{G}(x,x;t_{2k})$
$\displaystyle\geq\frac{1}{2}\mathbb{P}_{H_{2k}^{e}}(R\mbox{ returns to
}x\mbox{ for the first time at }t)\geq$
$\displaystyle\stackrel{{\scriptstyle\textrm{(\ref{eq:beta})}}}{{\geq}}\frac{1}{2\gamma_{2k}}t_{2k}^{-3/2}=\frac{1}{2}t_{2k}^{-7/4}$
(8)
(the $\frac{1}{2}$ comes from the first step).
To estimate $p(y,y;t_{2k})$ we divide the event $\\{R(t_{2k})=y\\}$ according
to whether $R$ “essentially goes through $x$” or not. Formally, denote by
$T_{1}$ and $T_{2}$ the first and last time before $t_{2k}$ when $R(T)=x$ (if
this does not happen, denote $T_{1}=\infty$ and $T_{2}=-\infty$). Then we
define
$\displaystyle p_{1}:=\mathbb{P}(T_{1}>\gamma_{2k},\,R(t_{2k})=y)\qquad
p_{2}:=\mathbb{P}(T_{2}<t_{2k}-\gamma_{2k},\,R(t_{2k})=y)$ $\displaystyle
p_{3}:=\mathbb{P}(T_{1}\leq\gamma_{2k},\,T_{2}\geq
t_{2k}-\gamma_{2k},\,R(t_{2k})=y)$
so that $p(y,y;t_{2k})\leq p_{1}+p_{2}+p_{3}$.
Now, $p_{1}$ and $p_{2}$ are easy to estimate. As above we have
$H_{2k}^{\textrm{o}}\cap[-t_{2k},t_{2k}]^{22}=H_{2k-1}^{\textrm{o}}\cap[-t_{2k},t_{2k}]^{22}$
so (5) applies and we get
$\mathbb{P}_{H_{\infty}^{\textrm{o}}}^{y}(R(t)=y)\leq\gamma_{2k}t^{-11}\quad\forall
t\leq t_{2k}.$ (9)
Therefore
$\displaystyle p_{1}$
$\displaystyle\leq\sum_{t=\gamma_{2k}}^{t_{2k}-1}\mathbb{P}(T_{1}=t,\,R(t_{2k})=y)+\mathbb{P}(T_{1}=\infty,\,R(t_{2k})=y)\leq$
$\displaystyle\leq\sum_{t=\gamma_{2k}-1}^{t_{2k}-2}\mathbb{P}_{H_{\infty}^{\textrm{o}}}(R(t)=y)+\mathbb{P}_{H_{\infty}^{\textrm{o}}}(R(t_{2k})=y)\leq$
$\displaystyle\stackrel{{\scriptstyle\textrm{(\ref{eq:yy})}}}{{\leq}}\sum_{t=\gamma_{2k}-1}^{t_{2k}-2}\gamma_{2k}\cdot
t^{-11}+\gamma_{2k}\cdot t_{2k}^{-11}\leq
C\gamma_{2k}^{-9}=Ct_{2k}^{-9/4}\stackrel{{\scriptstyle\textrm{(\ref{eq:pxxle})}}}{{=}}o(p(x,x;t))$
(10)
and similarly for $p_{2}$. As for $p_{3}$, we have
$\mathbb{P}(T_{1}\leq\gamma_{2k})\leq\Big{(}\sum_{i=0}^{\infty}\mathbb{P}_{H_{\infty}^{\textrm{o}}}(r\textrm{
visits }y\textrm{ }i\textrm{ times before
}\gamma_{2k})\Big{)}\cdot\delta\leq\frac{\delta}{\epsilon^{\textrm{o}}}\leq\frac{1}{2}$
and similarly (using time reversal) for $\mathbb{P}(T_{2}\geq
t_{2k}-\gamma_{2k})$. Hence we get
$p_{3}\leq\frac{1}{4}\max_{t_{2k}-2\gamma_{2k}\leq s\leq t_{2k}}p(x,x;s)$
and by lemma 3,
$\displaystyle p_{3}$
$\displaystyle\leq\frac{1}{4}p(x,x;t_{2k})\left(1+O\left(\frac{\gamma_{2k}\log^{3}t_{2k}}{\sqrt{t_{2k}}}\right)\right)+O(e^{-c\log^{2}t_{2k}})\leq$
$\displaystyle\stackrel{{\scriptstyle\textrm{(\ref{eq:pxxle})}}}{{\leq}}\frac{1}{4}p(x,x;t_{2k})(1+o(1)).$
With the estimate (10) for $p_{1}$ and the corresponding estimate for $p_{2}$
we get
$p(y,y;t_{2k})\leq p(x,x;t_{2k})\left(\frac{1}{4}+o(1)\right).$
A completely symmetric argument shows that at $t_{2k+1}$ the opposite occurs:
$p(x,x;t_{2k+1})\leq p(y,y;t_{2k+1})\left(\frac{1}{4}+o(1)\right)$
proving the theorem.∎
###### Remark.
If you want an example with unweighted graphs, this is not a problem. $H^{e}$
and $H^{o}$ are already unweighted, so the only thing needed is to connect
them, instead of with an edge of weight $\delta$, with a segment sufficiently
long such that the probability to traverse it is $\leq\delta$. The proof
remains essentially the same.
## 2\. Manifolds
We would like to show a Manifold $M$ and two points $x,y\in M$ such that the
heat kernel on $M$ satisfies
$\frac{p(x,x;t)}{p(y,y;t)}\nrightarrow$
as $t\to\infty$. Here is how one might translate the construction of our
theorem to the settings of manifolds. The dimension of the manifold plays
little role, so we might as well construct a surface.
For a subset $H\subset\mathbb{Z}^{22}$ one can associate a manifold $H^{*}$ by
replacing each vertex $v\in H$ with a sphere $v^{*}$ and every edge with a
empty, baseless cylinder. Since the degree of every vertex in $H$ is $\leq 44$
we may simply designate 44 disjoint circles on $\mathbb{S}^{2}$ and attach the
cylinders to the spheres at these circles. This is reminiscent of the well-
known “infinite jungle gym” construction, see some lovely pictures in [ON].
The exact method of doing so is unimportant since anyway the manifold that we
get is roughly isometric to $H$, considered as an induced subgraph of
$\mathbb{Z}^{22}$ (one of the nice features of rough isometry is that
continuous and discrete objects may be roughly isometric, rough isometry
inspects only the large scale geometry). Clearly $H^{*}$ can be made
$C^{\infty}$.
One can then construct a (possibly different) sequence $a_{k}$ and two
manifolds $\big{(}H_{\infty}^{\textrm{e/o}}\big{)}^{*}$ with the only
difference is that the $\alpha$ and $\beta$ must satisfy (5) and (6) for our
choice of the ∗ operation. This should be possible since
$\big{(}H_{k}^{\textrm{e/o}}\big{)}^{*}$ and
$\big{(}F_{k}^{\textrm{e/o}}\big{)}^{*}$ are roughly isometric to
$\mathbb{Z}^{22}$ and $\mathbb{Z}^{3}$ respectively. Instead of Delmotte one
can use the manifold version [SC95] (or rather, [D99] is the graph version of
earlier results for manifolds, see [SC95] for historical remarks).
The argument for the transience of
$\big{(}H_{\infty}^{\textrm{e/o}}\big{)}^{*}$ should also be direct
translation. Each contains a submanifold (with boundary) which is roughly
isometric to $\mathbb{Z}^{3}$ and therefore is transient. Since transience is
equivalent to the fact that for some $c>0$ every function which is $1$ at $x$
and $0$ at infinity satisfies that the Dirichlet form $\langle\nabla f,\nabla
f\rangle>c$, and since restricting to a submanifold only decreases the
Dirichlet form, we see that $\big{(}H_{\infty}^{\textrm{e/o}}\big{)}^{*}$ are
transient. Denote by
$\epsilon^{\textrm{e/o}}=\inf_{x\in
v^{*},v\sim\vec{0}^{\textrm{e/o}}}\mathbb{P}^{x}\left(W[0,\infty)\cap\big{(}\vec{0}^{\textrm{e/o}}\big{)}^{*}=\emptyset\right)$
where $W$ here is the Brownian motion on the manifold
$\big{(}H_{\infty}^{\textrm{e/o}}\big{)}^{*}$; and where the infimum is taken
over all $x$ belonging to a sphere $v^{*}$ where $v$ is some neighbor of
$\vec{0}^{\textrm{e/o}}$ in $H_{\infty}^{\textrm{e/o}}$. One can now define
$\delta=\frac{1}{2}\min(\varepsilon^{\textrm{e}},\varepsilon^{\textrm{o}})$
and connect $\vec{0}^{\textrm{e}}$ to $\vec{0}^{\textrm{o}}$ by a cylinder
sufficiently thin (or sufficiently long) such that the probability to traverse
it in either direction before reaching a neighboring sphere is $\leq\delta$.
This concludes a possible construction of a manifold $M$, and one may take $x$
to be an arbitrary point in $\big{(}\vec{0}^{\textrm{e}}\big{)}^{*}$ and $y$
and arbitrary point in $\big{(}\vec{0}^{\textrm{o}}\big{)}^{*}$.
The proof that $p(x,x;t)/p(y,y;t)$ does not converge should not require
significant changes. We note that in our case it is possible for a Brownian
motion at time $t$ to exit the box $[-t,t]^{22}$, but it is exponentially
difficult to do so. Hence, for example, instead of (8) we get
$p(x,x;t_{2K})\geq\frac{1}{2\gamma_{2k}^{2}}t_{2k}^{-3/2}-Ce^{-ct_{2k}}\geq\frac{1}{4}t_{2k}^{-7/4}$
for $k$ sufficiently large. Another point to note is that lemma 3 needs to be
replaced with an appropriate analog.
## References
* [ABJ02] Jean-Philippe Anker, Philippe Bougerol and Thierry Jeulin, _The infinite Brownian loop on a symmetric space_. Rev. Mat. Iberoamericana 18:1 (2002), 41–97. Available at: projecteuclid.org
* [C60] Kai Lai Chung, _Markov chains with stationary transition probabilities_. Die Grundlehren der mathematischen Wissenschaften 104, Springer-Verlag, 1960.
* [D82] Brian E. Davies, _Metastable states of symmetric Markov semigroups. II_. J. London Math. Soc. 26:3 (1982), 541–556. Available at: oxfordjournals.org
* [D97] Brian E. Davies, _Non-Gaussian aspects of heat kernel behaviour_. J. London Math. Soc. 55:1 (1997), 105–125. Available at: oxfordjournals.org
* [D38] Wolfgang Doeblin, _Sur deux problèmes de M. Kolmogoroff concernant les chaînes dénombrables_. Bull. de la Soc. Math. de France, 66 (1938), 210–220. Available at: numdam.org
* [DS84] Peter G. Doyle and Laurie J. Snell, _Random walks and electric networks_. Carus Mathematical Monographs, 22. Mathematical Association of America, Washington DC, 1984. Available at: dartmouth.edu/~doyle
* [D99] Thierry Delmotte, _Parabolic Harnack inequality and estimates of Markov chains on graphs_ , Rev. Mat. Iberoamericana 15:1 (1999), 181–232. Available at: rsme.es
* [K85] Masahiko Kanai, _Rough isometries, and combinatorial approximations of geometries of non-compact Riemannian manifolds_. J. Math. Soc. Japan 37:3 (1985), 391–413. projecteuclid.org
* [ON] Barrett O’Neill, _An infinite jungle gym_. Available at: ucla.edu/~bon
* [O61] Steven Orey, _Strong ratio limit property_. Bull. Amer. Math. Soc. 67:5 (1961), 571–574. Available at: ams.org
* [PSC] Christophe Pittet and Laurent Saloff-Coste, _A survey on the relationships between volume growth, isoperimetry, and the behavior of simple random walk on Cayley graphs, with examples_. Available from: cornell.edu/~lsc
* [SC95] Laurent Saloff-Coste, _Parabolic Harnack inequality for divergence-form second-order differential operators_. Potential Anal. 4:4 (special issue, 1995), 429–467. Available at: springerlink.com
|
arxiv-papers
| 2011-11-19T23:28:07 |
2024-09-04T02:49:24.502353
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Gady Kozma",
"submitter": "Gady Kozma",
"url": "https://arxiv.org/abs/1111.4593"
}
|
1111.4712
|
# An $L_{p}$-theory of stochastic parabolic equations with the random
fractional Laplacian driven by Lévy processes
Kyeong-Hun Kim111Department of Mathematics, Korea University, 1 Anam-dong,
Sungbuk-gu, Seoul 136-701, Republic of Korea. E-mail: kyeonghun@korea.ac.kr.
The research of this author was supported by Basic Science Research Program
through the National Research Foundation of Korea(NRF) funded by the Ministry
of Education, Science and Technology (20090087117). and Panki Kim222
Department of Mathematical Sciences and Research Institute of Mathematics,
Seoul National University, San56-1 Shinrim-dong Kwanak-gu, Seoul 151-747,
Republic of Korea. E-mail: pkim@snu.ac.kr. This work was supported by Basic
Science Research Program through the National Research Foundation of
Korea(NRF) grant funded by the Korea government(MEST)(2010-0001984).
###### Abstract
In this paper we give an $L_{p}$-theory for stochastic parabolic equations
with random fractional Laplacian operator. The driving noises are general Lévy
processes.
Keywords: Fractional Laplacian, Stochastic partial differential equations,
Lévy processes, $L_{p}$-theory.
AMS 2000 subject classifications: 60H15, 35R60.
## 1 Introduction
Let $d,m\geq 1$ be positive integers, $p\in[2,\infty)$ and $\alpha\in(0,2)$.
As usual $\mathbb{R}^{d}$ stands for the Euclidean space of points
$x=(x^{1},...,x^{d})$. We will use $dx$ to denote the Lebesgue measure in
either $\mathbb{R}^{d}$ or $\mathbb{R}^{m}$, which is clear in each context.
In this article we are dealing with $L_{p}$-theory of the stochastic partial
differential equations of the type
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f(u)\right)dt\,+\sum_{k=1}^{\infty}g^{k}(u)\cdot
dZ^{k}_{t},\quad u(0)=u_{0}$ (1.1)
given for $\omega\in\Omega,t\geq 0$ and $x\in\mathbb{R}^{d}$. Here $\Omega$ is
a probability space, $\Delta^{\alpha/2}$ is the fractional Laplacian defined
in (2.2), $Z^{k}_{t}$ are independent $m$-dimensional Lévy processes, and the
functions $f$ and vector-valued function $g^{k}=(g^{k,1},\dots g^{k,m})$
depend on $(\omega,t,x,u)$ satisfying certain continuity conditions. Our
result will cover the case
$\displaystyle f(u)$
$\displaystyle=b(\omega,t,x)\Delta^{\beta_{1}/2}u+c^{i}(\omega,t,x)u_{x^{i}}I_{\alpha>1}+d(\omega,t,x)u+f_{0},$
$\displaystyle g^{k,j}(u)$
$\displaystyle=\sigma^{k,j}(\omega,t,x)\Delta^{\beta^{j}_{2}/2}u+\nu^{k,j}(\omega,t,x)u+g^{k,j}_{0},\quad
j=1,\dots m.$
where $\beta_{1}<\alpha$ and $\beta^{j}_{2}<\alpha/2$ (see Assumptions 2.13
and 3.7).
An $L_{p}$-theory of (1.1) is introduced in [7] for the case that
$a(\omega,t)=1$ and are only finitely many Winer processes appear in the
equation. The approach in [7] cannot cover the case when there are infinitely
many Wiener processes, and the assumptions on $g$ in [7] are stronger than
conditions in our paper (See Remark 2.12 below). Moreover equations driven by
jump processes are not considered in [7]. A Hölder space theory for more
general (but non-random) integro-differential equations driven by Hilbert
space-valued Wiener process is given in [19] (also see [18] for a
deterministic equation). Even though the main result in [19] provide a nice
Hölder regularity of the solution to such problem, due to the Hölder-type
function spaces defined there, assumptions on $f$ and $g$ are quite strong.
Furthermore in [19] the equations with discontinuous Lévy processes are not
considered. We emphasize that the approach of this paper, based on $L_{p}$
theory in [14], is different from [19]. Our results include the case when $f$
and $g$ are only distributions and the number of derivatives of $f$ and $g$
are negative and fractional. On the other hand if $f$ and $g$ are sufficiently
smooth in $x$ then Sobolev embedding theorem combined with our $L_{p}$-theory
gives pointwise Hölder continuity of the solution even when $Z^{k}$ are
general Lévy processes.
$L_{p}$-theory for second-order stochastic parabolic equations driven by
Wiener processes was first established by Krylov [14]. Recently in [6] $L_{p}$
regularity theory for second-order stochastic parabolic equations driven by
Lévy processes is discussed.
In this paper, we establish an $L_{p}$-theory for stochastic parabolic
equations with the random fractional Laplacian driven by arbitrary Lévy
processes. Our result includes the case when the equation is driven by Lévy
space-time white noise (see Theorems 4.3 and 4.4). Among main tools used in
the article to study $L_{p}$-regularity theory are Burkholder-Davis-Gundy
inequality and a parabolic version of Littlewood-Paley inequality for the
fractional Laplacian operator introduced [11].
The organization of this article is as follows. First, in section 2, we prove
uniqueness and existence results of equation (1.1) driven by Wiener processes
in the space $L_{p}(\Omega\times[0,T],H^{\gamma+\alpha/2}_{p})$ (Theorem
2.15). Here $p\in[2,\infty)$ and $\gamma\in\mathbb{R}$. In section 3 we extend
Theorem 2.15 for the case when $Z^{k}_{t}$ are Lévy processes and $Z^{k}_{t}$
have finite $p$-th moments (see condition (3.2)). In section 4, the condition
(3.2) is weakened, and the uniqueness and existence results are proved in the
space $L_{p,\text{loc}}(\Omega\times[0,T],H^{\gamma+\alpha/2}_{p})$. The
condition (3.2) can be completely dropped if only finitely many Lévy processes
appear in the equation.
If we write $c=c(...)$, this means that the constant $c$ depends only on what
are in parenthesis. The constant $c$ stands for constants whose values are
unimportant and which may change from one appearance to another. The
dependence of the lower case constants on the dimensions $d,m$ may not be
mentioned explicitly. We will use “$:=$” to denote a definition, which is read
as “is defined to be”. For $a,b\in\mathbb{R}$, $a\wedge b:=\min\\{a,b\\}$ and
$a\vee b:=\max\\{a,b\\}$. Let $C^{\infty}_{0}(\mathbb{R}^{d})$ be the
collection of smooth functions with compact supports in $\mathbb{R}^{d}$. Most
of functions we discuss in this paper are random (depend on
$\omega\in\Omega$). For notational convenience, we suppress the dependency on
$\omega$ in most of expressions
## 2 Stochastic Parabolic equations with the random fractional Laplacian
driven by Wiener processes
Let $(\Omega,\mathcal{F},P)$ be a complete probability space,
$\\{\mathcal{F}_{t},t\geq 0\\}$ be an increasing filtration of $\sigma$-fields
$\mathcal{F}_{t}\subset\mathcal{F}$, each of which contains all
$(\mathcal{F},P)$-null sets. We assume that on $\Omega$ we are given
independent one-dimensional Wiener processes $W^{1}_{t},W^{2}_{t},...$
relative to $\\{\mathcal{F}_{t},t\geq 0\\}$. Let $\mathcal{P}$ be the
predictable $\sigma$-field generated by $\\{\mathcal{F}_{t},t\geq 0\\}$.
Let $p(t,x)$, where $t>0$, denote the inverse Fourier transform of
$e^{-|\xi|^{\alpha}t}$ in $\mathbb{R}^{d}$, that is,
$p(t,x):=\frac{1}{(2\pi)^{d/2}}\int_{\mathbb{R}^{d}}e^{i\xi\cdot
x}e^{-|\xi|^{\alpha}t}d\xi.$
For a suitable function $g$ and $t>0$, define the corresponding convolution
operator
$T_{t}g(x):=(p(t,\cdot)*g(\cdot))(x):=\int_{\mathbb{R}^{d}}p(t,x-y)g(y)dy,$
(2.1)
and define
$\partial^{\alpha}_{x}g(x)={\Delta}^{\frac{\alpha}{2}}g(x)=-(-{\Delta})^{\frac{\alpha}{2}}g(x):=\mathcal{F}^{-1}(-|\xi|^{\alpha}\mathcal{F}(g)(\xi))(x),$
(2.2)
where
$\mathcal{F}(g)(\xi)=\hat{g}(\xi):=\frac{1}{(2\pi)^{d/2}}\int_{\mathbb{R}^{d}}e^{-i\xi\cdot
x}g(x)dx$ is the Fourier transform of $g$ in $\mathbb{R}^{d}$.
In this section we study the nonlinear equations of the type
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f(u)\right)\,dt+\sum_{k=1}^{\infty}g^{k}(u)dW^{k}_{t},\quad
u(0)=u_{0},$ (2.3)
where $a(\omega,t)\in(\delta,\delta^{-1})$ for some $\delta>0$, and
$f(u)=f(\omega,t,x,u)$ and $g^{k}(u)=g^{k}(\omega,t,x,u)$ satisfy certain
continuity conditions, which we will put below.
First we introduce some stochastic Banach spaces. Let
$(\phi,\psi):=\int_{\mathbb{R}^{d}}\phi(x)\psi(x)dx$ and for $p\geq 1$,
$L_{p}=L_{p}(\mathbb{R}^{d}):=\\{\phi:\mathbb{R}^{d}\to\mathbb{R},\|\phi\|^{p}_{p}:=\int_{\mathbb{R}^{d}}|\phi(x)|^{p}dx<\infty\\}.$
For $n=0,1,2,...$, define
$H^{n}_{p}=H^{n}_{p}(\mathbb{R}^{d}):=\left\\{u:u,Du,...,D^{n}u\in
L_{p}(\mathbb{R}^{d})\right\\}.$
In general, for $\gamma\in\mathbb{R}$ define the space
$H^{\gamma}_{p}=H^{\gamma}_{p}(\mathbb{R}^{d})=(1-\Delta)^{-\gamma/2}L_{p}$
(called the space of Bessel potentials or the Sobolev space with fractional
derivatives) as the set of all distributions $u$ on $\mathbb{R}^{d}$ such that
$(1-\Delta)^{\gamma/2}u\in L_{p}$. For $u\in H^{\gamma}_{p}$, we define
$\|u\|_{H^{\gamma}_{p}}:=\|(1-\Delta)^{\gamma/2}u\|_{p}:=\|\mathcal{F}^{-1}[(1+|\xi|^{2})^{\gamma/2}\mathcal{F}(u)(\xi)]\|_{p},$
(2.4)
where $\mathcal{F}$ is the Fourier transform in $\mathbb{R}^{d}$. For
$\ell_{2}$-valued $g=(g^{1},g^{2},\dots)$, we define
$\|g\|_{H^{\gamma}_{p}(\ell_{2})}:=\|(1-\Delta)^{\gamma/2}g|_{\ell_{2}}\|_{p}:=\|\mathcal{F}^{-1}[(1+|\xi|^{2})^{\gamma/2}\mathcal{F}(g)(\xi)]|_{\ell_{2}}\|_{p}.$
Let $\overline{\mathcal{P}}$ be the completion of $\mathcal{P}$ with respect
to $dP\times dt$, and
$\mathbb{H}^{\gamma}_{p}(T):=L_{p}(\Omega\times[0,T],\overline{\mathcal{P}},H^{\gamma}_{p})$,
that is, $\mathbb{H}^{\gamma}_{p}(T)$ is the set of all
$\overline{\mathcal{P}}$-measurable processes $u:\Omega\times[0,T]\to
H^{\gamma}_{p}$ so that
$\|u\|_{\mathbb{H}^{\gamma}_{p}(T)}:=\left({\mathbb{E}}\left[\int^{T}_{0}\,\|u(\omega,t)\|^{p}_{H^{\gamma}_{p}}\,dt\right]\right)^{1/p}<\infty.$
###### Lemma 2.1
For any $\beta>0$,
$\eta^{1}_{\beta}(\xi):=\frac{(1+|\xi|^{2})^{\beta/2}}{1+|\xi|^{\beta}}$,
$\eta^{2}_{\beta}=(\eta^{1}_{\beta})^{-1}$,
$\eta^{3}:=\frac{|\xi|^{\beta}}{1+|\xi|^{\beta}}$ and
$\eta^{4}:=\frac{|\xi|^{\beta}}{(1+|\xi|^{2})^{\beta/2}}$ are
$L^{p}(\mathbb{R}^{d})$-multipliers, that is,
$\|\mathcal{F}^{-1}\left(\eta^{i}_{\beta}(\xi)(\mathcal{F}u)(\xi)\right)\|_{L_{p}}\leq
c(p,\beta)\|u\|_{L_{p}},\quad\quad i=1,2,3,4.$
Proof. See Theorem 0.2.6 of [24] (also see the remark below the theorem).
$\Box$
Lemma 2.36 easily yields the following results.
###### Corollary 2.2
(i) Let $\gamma\geq 0$. There exists a constant $c=c(\gamma)>0$ so that
$c\|u\|_{H^{\gamma}_{p}}\leq(\|u\|_{L_{p}}+\|\partial^{\gamma/2}_{x}u\|_{L_{p}})\leq
c^{-1}\|u\|_{H^{\gamma}_{p}}.$
(ii) For any $\beta\in\mathbb{R}$,
$\|\Delta^{\alpha/2}u\|_{H^{\beta}_{p}}\leq
c(\alpha,\beta)\|u\|_{H^{\beta+\alpha}_{p}}.$
###### Remark 2.3
Let $\gamma,\beta\geq 0$. Then due to the well-known inequality
$\|u\|_{L_{p}}\leq\varepsilon\|u\|_{H^{\gamma}_{p}}+c(\varepsilon,\gamma,\beta)\|u\|_{H^{-\beta}_{p}},$
it also follows
$\|u\|_{H^{\gamma}_{p}}\leq
c(\gamma,\beta,p)(\|u\|_{H^{-\beta}_{p}}+\|\partial^{\gamma/2}_{x}u\|_{L_{p}}).$
For $\ell_{2}$-valued $\overline{\mathcal{P}}$-measurable processes
$g=(g^{1},g^{2},\dots)$, we write $g\in\mathbb{H}^{\gamma}_{p}(T,\ell_{2})$ if
$\|g\|_{\mathbb{H}^{\gamma}_{p}(T,\ell_{2})}:=\left({\mathbb{E}}\int^{T}_{0}\|\,|(1-\Delta)^{\gamma/2}g(\omega,t)|_{\ell_{2}}\,\|^{p}_{p}\,dt\right)^{1/p}<\infty.$
(2.5)
Denote $\mathbb{L}_{p}(T):=\mathbb{H}^{0}_{p}(T)$ and
$\mathbb{L}_{p}(T,\ell_{2})=\mathbb{H}^{0}_{p}(T,\ell_{2})$. Finally, we say
$u_{0}\in U^{\gamma}_{p}$ if $u_{0}$ is $\mathcal{F}_{0}$-measurable function
$\Omega\to H^{\gamma}_{p}$ and
$\|u_{0}\|_{U^{\gamma}_{p}}:=\left({\mathbb{E}}\left[\|u_{0}\|^{p}_{H^{\gamma}_{p}}\right]\right)^{1/p}<\infty.$
###### Remark 2.4
It is easy to check (see Remark 3.2 in [14] for detailed proof) that for any
$\gamma\in(-\infty,\infty)$, $g\in\mathbb{H}^{\gamma}_{p}(T,\ell_{2})$ and
$\phi\in C^{\infty}_{0}(\mathbb{R}^{d})$ we have
$\sum_{k=1}^{\infty}\int^{T}_{0}(g^{k}(\omega,t),\phi)^{2}dt<\infty$ a.s., and
consequently the series of stochastic integral
$\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}(\omega,s),\phi)dW^{k}_{s}$ converges
uniformly in $t$ in probability on $[0,T]$.
###### Definition 2.5
Write $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ if
$u\in\mathbb{H}^{\gamma+\alpha}_{p}(T),u(0)\in
U^{\gamma+\alpha-\alpha/p}_{p}$, and for some $f\in\mathbb{H}^{\gamma}_{p}(T)$
and $g\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$
$du=fdt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad\hbox{for }t\in[0,T]$
in the sense of distributions, that is, for any $\phi\in
C^{\infty}_{0}(\mathbb{R}^{d})$,
$(u(t),\phi)=(u(0),\phi)+\int^{t}_{0}(f(s),\phi)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}(s),\phi)dW^{k}_{s}$
(2.6)
holds for all $t\leq T$ $a.s.$. In this case we write
$\mathbb{D}u:=f,\quad\mathbb{S}^{k}u:=g^{k},\quad\mathbb{S}u:=(\mathbb{S}^{1}u,\dots,\mathbb{S}^{k}u,\dots)$
and define the norm
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}:=\|u\|_{\mathbb{H}^{\gamma+\alpha}_{p}(T)}+\|\mathbb{D}u\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|\mathbb{S}u\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\|u(0)\|_{U^{\gamma+\alpha-\alpha/p}_{p}}.$
(2.7)
###### Theorem 2.6
The space $\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a Banach space, and for
every $0<t\leq T$
${\mathbb{E}}\Big{[}\sup_{s\leq
t}\|u(s,\cdot)\|^{p}_{H^{\gamma}_{p}}\Big{]}\leq
c(p,T,\alpha)\left(\|\mathbb{D}u\|^{p}_{\mathbb{H}^{\gamma}_{p}(t)}+\|\mathbb{S}u\|^{p}_{\mathbb{H}^{\gamma}_{p}(t,\ell_{2})}+\|u(0)\|^{p}_{U^{\gamma}_{p}}\right).$
(2.8)
In particular, for any $t\leq T$,
$\|u\|^{p}_{\mathbb{H}^{\gamma}_{p}(t)}\leq
c(p,T,\alpha)\int^{t}_{0}\|u\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(s)}\,ds.$
(2.9)
###### Remark 2.7
Note that $\alpha$ is not involved in (2.8).
Proof. See Theorem 3.7 in [14]. Actually in [14] the theorem is proved only
for $\alpha=2$, but the proof works for any $\alpha\in(0,2)$. We will give the
detailed proof of Theorem 3.4 below, which is the counterpart of Theorem 2.6
for pure-jump Lévy processes. $\Box$
###### Remark 2.8
It follows from (2.4) that for any $\mu,\gamma\in\mathbb{R}$, the operator
$(1-\Delta)^{\mu/2}:H^{\gamma}_{p}\to H^{\gamma-\mu}_{p}$ is an isometry.
Indeed,
$\|(1-\Delta)^{\mu/2}u\|_{H^{\gamma-\mu}_{p}}=\|(1-\Delta)^{(\gamma-\mu)/2}(1-\Delta)^{\mu/2}u\|_{p}=\|(1-\Delta)^{\gamma/2}u\|_{p}=\|u\|_{H^{\gamma}_{p}}.$
The same reason shows that
$(1-\Delta)^{\mu/2}:\mathcal{H}^{\gamma}_{p}(T)\to\mathcal{H}^{\gamma-\mu}_{p}(T)$
is an isometry.
###### Theorem 2.9
(i) For any deterministic functions $f=f(t,x)$ and $u_{0}=u_{0}(x)$ with
$\int^{T}_{0}\,\|f(t,\cdot)\|^{p}_{H^{\gamma}_{p}}\,dt<\infty,\quad\|u_{0}\|_{H^{\gamma+\alpha-\alpha/p}_{p}}<\infty,$
the (deterministic) equation
$u_{t}=\Delta^{\alpha/2}u+f,\quad u(0)=u_{0}$
has a unique solution $u$ with
$\int^{T}_{0}\,\|u(t,\cdot)\|^{p}_{H^{\gamma+\alpha}_{p}}\,dt<\infty$, and for
every $0<t\leq T$
$\int^{t}_{0}\,\|u(s,\cdot)\|^{p}_{H^{\gamma+\alpha}_{p}}ds\leq
c(p,T)\left(\int^{t}_{0}\,\|f(s,\cdot)\|^{p}_{H^{\gamma}_{p}}\,ds+\|u_{0}\|^{p}_{H^{\gamma+\alpha-\alpha/p}_{p}}\right).$
(2.10)
(ii) For any $f\in\mathbb{H}^{\gamma}_{p}(T)$ and $u_{0}\in
U^{\gamma+\alpha-\alpha/p}_{p}$, the equation
$u_{t}=\Delta^{\alpha/2}u+f,\quad u(0)=u_{0}$ (2.11)
has a unique solution $u\in\mathbb{H}^{\gamma+\alpha}_{p}(T)$ and for every
$0<t\leq T$
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}\leq
c(p,T)\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(t)}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right).$
(2.12)
Proof. (i). See, for instance, Theorem 2.1 in [18].
(ii). This result is also known. See, for instance, Lemma 3.2 and Lemma 3.4 of
[7]. Actually since equation (2.11) is deterministic for each fixed $\omega$,
the claim of (ii) can be obtained from (i). Indeed, the uniqueness and
estimate (2.12) are obvious by (i). For the existence of solution, assume that
$u_{0}$ and $f$ are sufficiently smooth in $x$, then using Fourier transform
one can easily check that
$u(t):=T_{t}u_{0}+\int^{t}_{0}T_{t-s}fds$
solves (2.11) and is in $\mathcal{H}^{\gamma+\alpha}_{p}(T)$. For general
$u_{0}$ and $f$ it is enough to use a standard approximation argument (see,
for instance, the proof Theorem 2.11).
$\Box$
Now we give our assumption on $a(\omega,t)$.
###### Assumption 2.10
The process $a(\omega,t)$ is predictable and there is a constant $\delta>0$ so
that
$\delta<a(\omega,t)<\delta^{-1},\quad\quad\forall\omega,t.$
Now we present an $L_{p}$-theory for linear stochastic parabolic equations
with random fractional Laplacian.
###### Theorem 2.11
Let $p\in[2,\infty)$ and $\gamma\in\mathbb{R}$. For any
$f\in\mathbb{H}^{\gamma}_{p}(T)$,
$g\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$ and $u_{0}\in
U^{\gamma+\alpha-\alpha/p}_{p}$, the linear equation
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f\right)\,dt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad
u(0)=u_{0},$ (2.13)
admits a unique solution $u$ in $\mathcal{H}^{\gamma+\alpha}_{p}(T)$, and for
this solution
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}\leq
c(p,T,\delta)\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|g\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right).$
(2.14)
###### Remark 2.12
(i) Recall that the unique solution $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$
is understood in the sense of distributions as in Definition 2.5, that is, for
any $\phi\in C^{\infty}_{0}(\mathbb{R}^{d})$,
$(u(t),\phi)=(u(0),\phi)+\int^{t}_{0}\left(a(\omega,s)(u,\Delta^{\alpha/2}\phi)+(f(s),\phi)\right)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}(s),\phi)dW^{k}_{s}$
holds for all $t\leq T$ $a.s.$.
(ii) A version of Theorem 2.11 is proved in [7] under stronger conditions on
$g$ and the processes. Precisely in [7] it is assumed that $a(\omega,t)=1$,
$g\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon}_{p}(T)$, $\varepsilon>0$, and
there are only finitely many Wiener processes in equation (2.13).
Proof. Step 1. Owing to Remark 2.8, we only need to show that the theorem
holds for a particular $\gamma=\gamma_{0}$. Indeed, suppose that the theorem
holds when $\gamma=\gamma_{0}$. Then it is enough to notice that
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a solution of the equation if and
only if
$\bar{u}:=(1-\Delta)^{(\gamma-\gamma_{0})/2}u\in\mathcal{H}^{\gamma_{0}+\alpha}_{p}$
is a solution of the equation with
$\bar{f}:=(1-\Delta)^{\frac{(\gamma-\gamma_{0})}{2}}f,\quad\bar{g}:=(1-\Delta)^{\frac{(\gamma-\gamma_{0})}{2}}g,\quad\bar{u}_{0}:=(1-\Delta)^{\frac{(\gamma-\gamma_{0})}{2}}u_{0},$
in place of $f,g$ and $u_{0}$, respectively. Furthermore,
$\displaystyle\|u\|_{\mathbb{H}^{\gamma+\alpha}_{p}(T)}=\|\bar{u}\|_{\mathbb{H}^{\gamma_{0}+\alpha}_{p}(T)}$
$\displaystyle\leq$ $\displaystyle
c\left(\|\bar{f}\|_{\mathbb{H}^{\gamma_{0}}_{p}(T)}+\|\bar{g}\|_{\mathbb{H}^{\gamma_{0}+\alpha/2+\varepsilon_{0}}_{p}(T,\ell_{2})}+\|\bar{u}_{0}\|_{U^{\gamma_{0}+\alpha-\alpha/p}_{p}}\right)$
$\displaystyle=$ $\displaystyle
c\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|g\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{0}}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right).$
Step 2. Next we assume $a(\omega,t)=1$ and prove the theorem for the equation:
$du=\Delta^{\alpha/2}udt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad u(0)=0.$
(2.15)
Remember that we may assume $\gamma=-\alpha/2$. Since the uniqueness of (2.15)
follows from results for the deterministic equations (Theorem 2.9), we only
need to show that there exists a solution $u\in\mathbb{H}^{\alpha/2}_{p}(T)$
of (2.15) and $u$ satisfies estimate (2.14) with $f=u_{0}=0$ and
$\gamma=-\alpha/2$.
For a moment, assume $N_{0}>0$ is a fixed non-random constant, $g^{k}=0$ for
all $k>N_{0}$ and
$g^{k}(t,x)=\sum_{i=0}^{m_{k}}I_{(\tau^{k}_{i},\tau^{k}_{i+1}]}(t)g^{k_{i}}(x)\qquad\text{for
}k\leq N_{0},$ (2.16)
where $\tau^{k}_{i}$ are bounded stopping times and $g^{k_{i}}(x)\in
C^{\infty}_{0}(\mathbb{R}^{d})$. Define
$v(t,x):=\sum_{k=1}^{N_{0}}\int^{t}_{0}g^{k}(s,x)dW^{k}_{s}=\sum_{k=1}^{N_{0}}\sum_{i=1}^{m_{k}}g^{k_{i}}(x)(W^{k}_{t\wedge\tau^{k}_{i+1}}-W^{k}_{t\wedge\tau^{k}_{i}})$
and
$u(t,x):=v(t,x)+\int^{t}_{0}\Delta^{\alpha/2}T_{t-s}v(s,x)\,ds=v(t,x)+\int^{t}_{0}T_{t-s}\Delta^{\alpha/2}v(s,x)\,ds.$
(2.17)
Using Fourier transform one can easily show (See, for instance, [7]) that if
functions $h_{1}=h_{1}(t,x)$ and $h_{2}=h_{2}(x)$ are sufficiently smooth in
$x$ then
$w_{1}(t,x):=\int^{t}_{0}T_{t-s}h_{1}(s,x)ds,\quad w_{2}(t,x)=T_{t}h_{2}(x)$
solve
$dw_{1}=(\Delta^{\alpha/2}w_{1}+h_{1})\,dt,\quad w_{1}(0)=0,$
$dw_{2}=\Delta^{\alpha/2}w_{2}\,dt,\quad w_{2}(0)=h_{2}.$
Therefore we have
$d(u-v)=(\Delta^{\alpha/2}(u-v)+\Delta^{\alpha/2}v)dt=\Delta^{\alpha/2}udt$,
and
$du=\Delta^{\alpha/2}udt+dv=\Delta^{\alpha/2}udt+\sum_{k=1}^{N_{0}}g^{k}dW^{k}_{t}.$
Also by (2.17) and stochastic Fubini theorem ([22, Theorem 64]), almost
surely,
$\displaystyle u(t,x)$ $\displaystyle=$ $\displaystyle
v(t,x)+\sum_{k=1}^{N_{0}}\int^{t}_{0}\int^{s}_{0}\Delta^{\alpha/2}T_{t-s}g^{k}(r,x)dW^{k}_{r}ds$
(2.18) $\displaystyle=$ $\displaystyle
v(t,x)-\sum_{k=1}^{N_{0}}\int^{t}_{0}\int^{t}_{r}\frac{\partial}{\partial
s}T_{t-s}g^{k}(r,x)dsdW^{k}_{r}$ $\displaystyle=$
$\displaystyle\sum_{k=1}^{N_{0}}\int^{t}_{0}T_{t-s}g^{k}(s,x)dW^{k}_{s}.$
Hence,
$\partial^{\alpha/2}_{x}u(t,x)=\sum_{k=1}^{N_{0}}\int^{t}_{0}\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,\cdot)(x)dW^{k}_{s},$
and by Burkholder-Davis-Gundy’s inequality, we have
${\mathbb{E}}\left[\big{|}\partial^{\alpha/2}_{x}u(t,x)\big{|}^{p}\right]\leq
c(p){\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{N_{0}}|\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,\cdot)(x)|^{2}ds\right)^{p/2}\right].$
Now we use a parabolic version of Littlewood-Paley inequality for fractional
Laplacian (Theorem 2.3 in [11])
$\int_{\mathbb{R}^{d}}\int^{T}_{0}\left[\int^{t}_{0}|\partial^{\alpha/2}_{x}T_{t-s}g(s,\cdot)(x)|^{2}_{\ell_{2}}ds\right]^{p/2}dtdx\leq
c(\alpha,p)\int_{\mathbb{R}^{d}}\int^{T}_{0}|g(t,x)|^{p}_{\ell_{2}}\,dtdx$
(2.19)
and get
${\mathbb{E}}\left[\int^{T}_{0}\|\partial^{\alpha/2}_{x}u(t,\cdot)\|^{p}_{p}\,dt\right]\leq
c(p){\mathbb{E}}\left[\int^{T}_{0}\||g(t,\cdot)|_{\ell_{2}}\|^{p}_{p}\,dt\right].$
(2.20)
Similarly, (2.18) and Burkholder-Davis-Gundy’s inequality yield
${\mathbb{E}}\left[|u(t,x)|^{p}\right]\leq
c(p)\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{N_{0}}|T_{t-s}g^{k}(s,x)|^{2}ds\right)^{p/2}\right].$
(2.21)
Since $(\sum_{k=1}^{N_{0}}|a_{n}|^{2})^{p/2}\leq
c(N_{0},p)\sum_{k=1}^{N_{0}}|a_{n}|^{p}$ and $\|T_{t}f\|_{p}\leq c\|f\|_{p}$,
we see that for every $t>0$
$\displaystyle\int_{\mathbb{R}^{d}}\left(\int^{t}_{0}\sum_{k=1}^{N_{0}}|T_{t-s}g^{k}(s,x)|^{2}ds\right)^{p/2}dx$
$\displaystyle\leq
t^{p/2-1}\int_{\mathbb{R}^{d}}\int^{t}_{0}\left(\sum_{k=1}^{N_{0}}|T_{t-s}g^{k}(s,x)|^{2}\right)^{p/2}dtdx$
$\displaystyle\leq
c(T,N_{0},p)\int^{t}_{0}\int_{\mathbb{R}^{d}}\sum_{k=1}^{N_{0}}|g^{k}(t,x)|^{p}dxdt.$
Consequently,
$\displaystyle{\mathbb{E}}\int_{0}^{T}\int_{\mathbb{R}^{d}}|u(t,x)|^{p}dxdt\leq
c(T,N_{0},p){\mathbb{E}}\int^{T}_{0}\int_{\mathbb{R}^{d}}|g(t,x)|^{p}_{\ell_{2}}dxds.$
(2.22)
Thus we proved $\partial^{\alpha/2}_{x}u,u\in\mathbb{L}_{p}(T)$, and hence by
Corollary 2.2 we have $u\in\mathcal{H}^{\alpha/2}_{p}(T)$. Note that by
Corollary 2.2(ii)
$\|\Delta^{\alpha/2}u\|_{H^{-\alpha/2}_{p}}=\|\Delta^{\alpha/4}(\partial^{\alpha/2}_{x}u)\|_{H^{-\alpha/2}_{p}}\leq
c\|\partial^{\alpha/2}_{x}u\|_{L_{p}}.$
By definition (2.7) and Remark 2.3, for any $t\leq T$,
$\displaystyle\|u\|^{p}_{\mathcal{H}^{\alpha/2}_{p}(t)}$ $\displaystyle\leq$
$\displaystyle
c(p)\left(\|u\|^{p}_{\mathbb{H}^{\alpha/2}_{p}(t)}+\|\Delta^{\alpha/2}u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|g\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}\right)$
(2.23) $\displaystyle\leq$ $\displaystyle
c\left(\|u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|\partial^{\alpha/2}_{x}u\|^{p}_{\mathbb{L}_{p}(t)}+\|g\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}\right).$
Combining this with (2.20) and (2.9) we have that for every $0<t\leq T$
$\displaystyle\|u\|^{p}_{\mathcal{H}^{\alpha/2}_{p}(t)}$ $\displaystyle\leq$
$\displaystyle
c(p,T,\alpha)\left(\|u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|g\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}\right)$
(2.24) $\displaystyle\leq$ $\displaystyle
c(p,T,\alpha)\int_{0}^{t}\|u\|^{p}_{\mathcal{H}^{\alpha/2}_{p}(s)}ds+c(p,T,\alpha)\|g\|^{p}_{\mathbb{L}_{p}(T,\ell_{2})}.$
Finally, Gronwall leads to (2.14).
Now we drop the additional assumptions on $g$ by using the following standard
approximation argument: By Theorem 3.10 in [14], for
$g\in\mathbb{L}_{p}(T,\ell_{2})$ we can take a sequence
$g_{n}\in\mathbb{L}_{p}(T,\ell_{2})$ so that $g_{n}\to g$ in
$\mathbb{L}_{p}(T,\ell_{2})$ and each $g_{n}=(g^{1}_{n},g^{2}_{n},\cdots)$
satisfies above assumed assumptions, that is, $g^{k}_{n}=0$ for all large $k$
and each $g^{k}_{n}$ is of type (2.16). By the above result, the equation
$du_{n}=\Delta^{\alpha/2}u_{n}dt+\sum_{k=1}^{\infty}g^{k}_{n}dW^{k}_{t},\quad
u_{n}(0)=0$
has a unique solution $u_{n}$. It also follows that $u_{n}-u_{m}$ is the
unique solution of
$d(u_{n}-u_{m})=\Delta^{\alpha/2}(u_{n}-u_{m})dt+(g^{k}_{n}-g^{k}_{m})dW^{k}_{t},\quad(u_{n}-u_{m})(0)=0$
and, by the previous argument
$\|u_{n}-u_{m}\|_{\mathcal{H}^{\alpha/2}_{p}(T)}\leq
c(p,T)\|g_{n}-g_{m}\|_{\mathbb{L}_{p}(T,\ell_{2})}.$
Consequently, there is $u\in\mathcal{H}^{\alpha/2}_{p}(T)$ so that $u_{n}\to
u$ in $\mathcal{H}^{\alpha/2}_{p}(T)$. We only need to prove $u$ is a solution
of (2.15). Equivalently, we need to prove that for any $\phi\in
C^{\infty}_{0}(\mathbb{R}^{d})$, the equality
$(u(t,\cdot),\phi)=\int^{t}_{0}(\Delta^{\alpha/2}u(s,\cdot),\phi)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}(s,\cdot),\phi)dW^{k}_{s}.$
(2.25)
holds for for all $t\leq T$ (a.s.), or equivalently
$((1-\Delta)^{-\alpha/2}u(t,\cdot),(1-\Delta)^{\alpha/2}\phi)=\int^{t}_{0}(\Delta^{\alpha/4}u(s,\cdot),\Delta^{\alpha/4}\phi)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}(s,\cdot),\phi)dW^{k}_{s}.$
(2.26)
By (2.8),
$\lim_{n\to\infty}{\mathbb{E}}\left[\sup_{t\leq
T}\|(1-\Delta)^{-\alpha/2}(u_{n}(t,\cdot)-u(t,\cdot))\|^{p}_{L_{p}}\right]=0,\quad\text{a.s.}$
which implies that one can take a subsequence $n_{j}$ so that
$(1-\Delta)^{-\alpha/2}u_{n_{j}}\to(1-\Delta)^{-\alpha/2}u$ in
$L_{p}(\mathbb{R}^{d})$ uniformly on $[0,T]$ (a.s) and consequently
$t\to((1-\Delta)^{-\alpha/2}u(t,\cdot),(1-\Delta)^{\alpha/2}\phi)$ is
continuous on $[0,T]$. By taking the limit from
$((1-\Delta)^{-\alpha/2}u_{n_{j}}(t,\cdot),(1-\Delta)^{\alpha/2}\phi)=\int^{t}_{0}(\Delta^{\alpha/4}u_{n_{j}}(s,\cdot),\Delta^{\alpha/4}\phi)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(g^{k}_{n_{j}}(s,\cdot),\phi)dW^{k}_{s}$
and remembering that both sides of (3.7) are continuous in $t$, one easily get
that equality (3.7) holds for all $t\leq T$ (a.s.).
Step 3. Next we prove the theorem for the equation
$du=(\Delta^{\alpha/2}u+f)dt+g^{k}dW^{k}_{t},\quad u(0)=u_{0}.$ (2.27)
Again we may assume $\gamma=-\alpha/2$, and due to Theorem 2.9 we only need to
show that there exists a solution $u$ and it satisfies estimate (2.14). By
Theorem 2.9, the equation
$dv=(\Delta^{\alpha/2}v+f)dt,\quad v(0)=u_{0}$
has a solution $v\in\mathcal{H}^{\alpha/2}_{p}(T)$ and
$\|v\|_{\mathcal{H}^{\alpha/2}_{p}(T)}\leq
c(p,T)\left(\|f\|_{\mathbb{H}^{-\alpha/2}_{p}(T)}+\|u_{0}\|_{U^{\alpha/2-\alpha/p}_{p}}\right).$
Also by the result of Step 2, the equation
$dw=\Delta^{\alpha/2}wdt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad w(0)=0$
has a unique solution and
$\|w\|_{\mathcal{H}^{\alpha/2}_{p}(T)}\leq
c(p,T)\|g\|_{\mathbb{L}_{p}(T,\ell_{2})}.$
Now it is enough to take $u=v+w$.
Step 4 (A priori estimate). We prove the a priori estimate (2.14) holds given
that a solution $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ of the following
equation already exists :
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f\right)\,dt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad
u(0)=u_{0}.$
This time we prove (2.14) only for $\gamma=0$. This is enough due to the
reason given in Step 1. By Step 3, the equation
$dv=(\Delta^{\alpha/2}v+f)dt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad
v(0)=u_{0}.$
has a solution $v\in\mathcal{H}^{\alpha}_{p}(T)$ and
$\|v\|_{\mathcal{H}^{\alpha}_{p}(T)}\leq
c(p,T)\left(\|f\|_{\mathbb{L}_{p}(T)}+\|g\|_{H^{\alpha/2}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\alpha-\alpha/p}_{p}}\right).$
Note that $\bar{u}:=u-v$ satisfies
$d\bar{u}=(a(\omega,t)\Delta^{\alpha/2}\bar{u}+\bar{f})dt,\quad\bar{u}(0)=0,$
where $\bar{f}:=(a(\omega,t)-1)\Delta^{\alpha/2}v$, and
$\|\bar{f}\|_{\mathbb{L}_{p}(T)}\leq
c\|\Delta^{\alpha/2}v\|_{\mathbb{L}_{p}(T)}\leq
c\left(\|f\|_{\mathbb{L}_{p}(T)}+\|g\|_{H^{\alpha/2}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\alpha-\alpha/p}_{p}}\right).$
Since $u=v+\bar{u}$,
$\|u\|_{\mathcal{H}^{\alpha}_{p}(T)}\leq\|v\|_{\mathcal{H}^{\alpha}_{p}(T)}+\|\bar{u}\|_{\mathcal{H}^{\alpha}_{p}(T)}$
and $\|\bar{u}\|_{\mathcal{H}^{\alpha}_{p}(T)}\leq
c\|\bar{u}\|_{\mathbb{H}^{\alpha}_{p}(T)}+c\|\bar{f}\|_{\mathbb{L}_{p}(T)}$,
to prove (2.14) we only need to show that for each $\omega\in\Omega$,
$\int^{T}_{0}\|\bar{u}(t,\cdot)\|^{p}_{H^{\alpha}_{p}}dt\leq
c(p,T,\alpha,\delta)\int^{T}_{0}\|\bar{f}(t,\cdot)\|^{p}_{L_{p}}dt.$ (2.28)
For fixed $\omega$, define a non-random functions
$\tilde{u}(t,x)=\bar{u}(\omega,\xi(\omega,t),x)\quad\tilde{f}(t,x)=a(\omega,t)^{-1}\bar{f}(\omega,\xi(\omega,t),x).$
(2.29)
where $\xi(\omega,t):=\int^{t}_{0}\frac{ds}{a(\omega,s)}$. Then clearly
$\tilde{u}$ satisfies
$\tilde{u}_{t}=\Delta^{\alpha/2}\tilde{u}+\tilde{f},\quad\tilde{u}(0)=0.$
Let $\tilde{T}(\omega,T)$ be such that
$T=\int^{\tilde{T}(\omega,T)}_{0}\frac{ds}{a(\omega,s)}$. Since $\delta
T<\tilde{T}(\omega,T)<\delta T$, applying (2.10), we get
$\int^{\tilde{T}(\omega,T)}_{0}\|\tilde{u}(t,\cdot)\|^{p}_{H^{\alpha}_{p}}dt\leq
c(p,T,\delta,\alpha)\int^{\tilde{T}(\omega,T)}_{0}\|\tilde{f}(t,\cdot)\|^{p}_{L_{p}}dt.$
This and relations in (2.29) easily lead to (2.28).
Step 5 (Method of continuity). The solvability of equation (2.27), the a
priori estimate (2.14) and the method of continuity obviously finish the proof
of the theorem. But below we show how the method of continuity works only for
reader’s convenience.
For $\lambda\in[0,1]$, denote $a_{\lambda}(\omega,t)=(1-\lambda)+\lambda
a(\omega,t)$. Then obviously $a_{\lambda}$ is predictable and
$a_{\lambda}\in(\delta,\delta^{-1})$. It follows from Step 4 that if
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a solution of the equation
$du=\left(a_{\lambda}(\omega,t)\Delta^{\alpha/2}u+f\right)\,dt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad
u(0)=u_{0},$ (2.30)
then the estimate (2.14) holds with the same constant $c=c(p,T,\delta)$. Now
let $J$ be the collection of $\lambda\in[0,1]$ so that for any
$f\in\mathbb{H}^{\gamma}_{p}(T),g\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$
and $u_{0}\in U^{\gamma+\alpha/2-\alpha/p}_{p}$, equation (2.30) has a
solution. By Step 3, $0\in J$. Note that to finish the proof of the theorem we
only need to show $1\in J$. Now let $\lambda_{0}\in J$. Obviously
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a solution of (2.30) if and only
if
$du=\left(a_{\lambda_{0}}(\omega,t)\Delta^{\alpha/2}u+[(a_{\lambda}(\omega,t)-a_{\lambda_{0}}(\omega,t))\Delta^{\alpha/2}u+f]\right)\,dt+\sum_{k=1}^{\infty}g^{k}dW^{k}_{t},\quad
u(0)=u_{0}.$ (2.31)
Now fix $u^{1}\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ with initial date $u_{0}$
(for instance take the solution of (2.27)), and define $u^{2},u^{3},\cdots$ so
that $u^{n+1}\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is the solution of
$du^{n+1}=\left(a_{\lambda_{0}}(\omega,t)\Delta^{\alpha/2}u^{n+1}+[(a_{\lambda}(\omega,t)-a_{\lambda_{0}}(\omega,t))\Delta^{\alpha/2}u^{n}+f]\right)dt\,dt+g^{k}dW^{k}_{t},\quad
u(0)=u_{0}.$ (2.32)
Then $v^{n+1}:=u^{n+1}-u^{n}$ satisfies
$dv^{n+1}=\left(a_{\lambda_{0}}(\omega,t)\Delta^{\alpha/2}v^{n+1}+(a_{\lambda}(\omega,t)-a_{\lambda_{0}}(\omega,t))\Delta^{\alpha/2}v^{n}\right)\,dt$
By the a priori estimate (2.14),
$\displaystyle\|v^{n+1}\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}$
$\displaystyle\leq$ $\displaystyle
c\|(a_{\lambda}-a_{\lambda_{0}})\Delta^{\alpha/2}v^{n}\|_{\mathbb{H}^{\gamma}_{p}(T)}$
$\displaystyle\leq$ $\displaystyle
N(p,T,\delta)|\lambda-\lambda_{0}|\|v^{n}\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}.$
Thus if $|\lambda-\lambda_{0}|<1/(2N(p,T,\delta))$, the map which send $u^{n}$
to $u^{n+1}$ is a contraction in $\mathcal{H}^{\gamma+\alpha}_{p}(T)$, and has
a unique fixed point $u$. Thus $u$ satisfies (2.30)–(2.31). Since the above
constant $N$ is independent of $\lambda$, it follows that $J=[0,1]$ and the
theorem is proved. $\Box$
Finally we consider the nonlinear equation
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f(u)\right)\,dt+\sum_{k=1}^{\infty}g^{k}(u)dW^{k}_{t},\quad
u(0)=u_{0},$ (2.33)
where $f(u)=f(\omega,t,x,u)$ and $g^{k}(u)=g^{k}(\omega,t,x,u)$.
###### Assumption 2.13
Assume $f(0)\in\mathbb{H}^{\gamma}_{p}(T)$ and
$g(0)\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$. Moreover, for any
$\varepsilon>0$, there exists a constant $K_{\varepsilon}$ so that for any
$u=u(x),v=v(x)\in H^{\gamma+\alpha}_{p}$ and $\omega,t$, we have
$\displaystyle\|f(t,\cdot,u(\cdot))-f(t,\cdot,v(\cdot))\|_{H^{\gamma}_{p}}+\|g(t,\cdot,u(\cdot))-g(t,\cdot,v(\cdot))\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}$
$\displaystyle\leq$
$\displaystyle\varepsilon\|u-v\|_{H^{\gamma+\alpha}_{p}}+K(\varepsilon)\|u-v\|_{H^{\gamma}_{p}}.$
(2.34)
To give an example of $f(u)$ and $g(u)$ satisfying Assumption 2.13, we
introduce the space of point-wise multipliers in $H^{\gamma}_{p}$. For each
$r\geq 0$, define
$B^{r}=\begin{cases}B(\mathbb{R}^{d})\qquad&\hbox{if }r=0,\\\
C^{r-1,1}(\mathbb{R}^{d})&\hbox{if }r=1,2,\cdots,\\\
C^{r}(\mathbb{R}^{d})&\hbox{otherwise},\end{cases}$ (2.35)
where $B(\mathbb{R}^{d})$ is the space of bounded Borel measurable functions
on $\mathbb{R}^{d}$, $C^{r-1,1}(\mathbb{R}^{d})$ is the space of $r-1$ times
continuously differentiable functions whose $(r-1)$st order derivatives are
Lipschitz continuous, and $C^{r}(\mathbb{R}^{d})$ is the usual Hölder space.
Also we use the space $B^{r}$ for $\ell_{2}$-valued functions. For instance,
if $g=(g^{1},g^{2},...)$, then $|g|_{B^{0}}=\sup_{x}|g(x)|_{\ell_{2}}$ and
$|g|_{C^{n-1,1}}=\sum_{|\alpha|\leq
n-1}|D^{\alpha}g|_{B^{0}}+\sum_{|\alpha|=n-1}\sup_{x\neq
y}\frac{|D^{\alpha}g(x)-D^{\alpha}g(y)|_{\ell_{2}}}{|x-y|}.$
Fix $\kappa_{0}=\kappa_{0}(\gamma)\geq 0$ so that $\kappa_{0}>0$ if $\gamma$
is not integer. It is known (see, for instance, Lemma 5.2 in [14]) that for
any $a\in B^{|\gamma|+\kappa_{0}}$ and $h\in H^{\gamma}_{p}$,
$\|ah\|_{H^{\gamma}_{p}}\leq
c(\gamma,\kappa_{0})|a|_{B^{|\gamma|+\kappa_{0}}}|h|_{H^{\gamma}_{p}}$ (2.36)
and the same inequality holds for $\ell_{2}$-valued functions $a$.
###### Example 2.14
Fix $\kappa_{0}=\kappa_{0}(\gamma)\geq 0$ so that $\kappa_{0}>0$ if $\gamma$
is not integer. Consider
$f(u)=b(\omega,t,x)\Delta^{\beta_{1}/2}u+\sum_{i=1}^{d}c^{i}(\omega,t,x)u_{x^{i}}I_{\alpha>1}+d(\omega,t,x)u+f_{0},$
$g^{k}(u)=\sigma^{k}(\omega,t,x)\Delta^{\beta_{2}/2}u+v^{k}(\omega,t,x)u+g^{k}_{0},$
where $\beta_{1}<\alpha$, $\beta_{2}<\alpha/2$,
$f_{0}\in\mathbb{H}^{\gamma}_{p}(T)$ and
$g_{0}\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$. Assume for each
$\omega,t$,
$|b|_{B^{|\gamma|+\kappa_{0}}}+\sum_{i=1}^{d}|c^{i}|_{B^{|\gamma|+\kappa_{0}}}+|d|_{B^{|\gamma|+\kappa_{0}}}+|\sigma|_{B^{|\gamma|+\alpha/2+\kappa_{0}}}+|\nu|_{B^{|\gamma|+\alpha/2+\kappa_{0}}}\leq
K.$
Then by (2.36), for each $t$
$\displaystyle\|f(t,\cdot,u(\cdot))-f(t,\cdot,v(\cdot))\|_{H^{\gamma}_{p}}+\|g(t,\cdot,u(\cdot))-g(t,\cdot,v(\cdot))\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}$
$\displaystyle\leq$ $\displaystyle
c\left(\|\Delta^{\beta_{1}/2}(u-v)\|_{H^{\gamma}_{p}}+I_{\alpha>1}\|D(u-v)\|_{H^{\gamma}_{p}}+\|u-v\|_{H^{\gamma+\alpha/2}_{p}}+\|\Delta^{\beta_{2}/2}(u-v)\|_{H^{\gamma+\alpha/2}_{p}}\right).$
Since for any $\alpha_{1}<\alpha$ and $\varepsilon>0$, by interpolation
theory,
$\|u\|_{H^{\gamma+\alpha_{1}}_{p}}\leq
c(\alpha,\alpha_{1})\|u\|_{H^{\gamma+\alpha}_{p}}^{\alpha_{1}/\alpha}\|u\|_{H^{\gamma}_{p}}^{1-\alpha_{1}/\alpha}\leq\varepsilon\|u\|_{H^{\gamma+\alpha}_{p}}+c(\varepsilon,\alpha_{1},\alpha)\|u\|_{H^{\gamma}_{p}},$
one easily gets (2.13).
Here is the main result of this section.
###### Theorem 2.15
Suppose Assumptions 2.10 and 2.13 hold. Then equation (2.33) has a unique
solution $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$, and for this solution
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}\leq
c\left(\|f(0)\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|g(0)\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha/2-\alpha/p}_{p}}\right),$
where $c=c(p,T,\delta)$.
Proof. Our proof is virtually identical to the that of Theorem 6.4 in [14],
where the theorem is proved when $\alpha=2$. The only difference is that one
has to use Theorem 2.11 in this article, in place the corresponding result in
[14]. We skip the proof here since we will give the proof for more general
case in next section. $\Box$
By Sobolev embedding theorem, we immediately get the following
###### Corollary 2.16
Suppose Assumptions 2.10 and2.13 hold. If $\gamma+\alpha>d/p$, then the unique
solution $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ of equation (2.33) is
$C^{\gamma+\alpha-d/p}$-valued process on $[0,T]\times\Omega$ a.s..
## 3 General case
Let $Z^{1}_{t},Z^{2}_{t},...$ be independent $m$-dimensional Lévy processes
relative to $\\{\mathcal{F}_{t},t\geq 0\\}$. For $t\geq 0$ and Borel set
$A\in\mathcal{B}(\mathbb{R}^{m}\setminus\\{0\\})$, define
$N_{k}(t,A):=\\#\left\\{0\leq s\leq t;\,Z^{k}_{s}-Z^{k}_{s-}\in
A\right\\},\quad\widetilde{N}_{k}(t,A):=N_{k}(t,A)-t\nu_{k}(A)$
where $\nu_{k}(A):={\mathbb{E}}[N_{k}(1,A)]$ is the Lévy measure of $Z^{k}$.
By Lévy-Itô decomposition, there exist a vector $\alpha^{k}$, a non-negative
definite matrix $\beta^{k}$ and $m$-dimensional Wiener process $B^{k}$ so that
$Z^{k}(t)=\alpha^{k}t+\beta^{k}B^{k}_{t}+\int_{|z|<1}z\widetilde{N}_{k}(t,dz)+\int_{|z|\geq
1}zN_{k}(t,dz).$ (3.1)
For any $q,k=1,2,\cdots$, denote
$\widehat{c}_{k,q}:=\left(\int_{\mathbb{R}^{m}}|z|^{q}\nu_{k}(dz)\right)^{1/q}.$
Now we fix $p\in[2,\infty)$ and denote
$\widehat{c}_{k}:=\left(\widehat{c}_{k,2}\vee\widehat{c}_{k,p}\right)$. In
this section we assume
$\widehat{c}:=\sup_{k\geq 1}\widehat{c}_{k}<\infty.$ (3.2)
((3.2) will be weaken in section 4). Then for any $2<q<p$, by Hölder’s
inequality,
$\widehat{c}_{k,q}\leq\left(\int_{\mathbb{R}^{m}}|z|^{2}\nu_{k}(dz)\right)^{(p-q)/(q(p-2))}\left(\int_{\mathbb{R}^{m}}|z|^{p}\nu_{k}(dz)\right)^{(q-2)/(q(p-2))}\leq\widehat{c}_{k}.$
By (3.2), $\int_{|z|\geq 1}|z|\nu_{k}(dz)\leq\int_{|z|\geq
1}|z|^{2}\nu_{k}(dz)<\infty$, and
$\int_{|z|\geq 1}zN_{k}(t,dz)=\int_{|z|\geq
1}z\widetilde{N}_{k}(t,dz)+t\int_{|z|\geq 1}z\nu_{k}(dz).$
Thus by absorbing $\widetilde{\alpha}_{k}:=\int_{|z|\geq 1}z\nu_{k}(dz)$ into
$\alpha_{k}$ we can rewrite (3.1) as
$Z^{k}_{t}=\tilde{\alpha}_{k}t+\beta_{k}B^{k}_{t}+\int_{\mathbb{R}^{m}}z\widetilde{N}_{k}(t,dz).$
We first consider the following linear equation:
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f\right)dt\,+\sum_{k=1}^{\infty}g^{k}\cdot
dZ^{k}_{t},\quad u(0)=u_{0}.$ (3.3)
Relocation of the term $\sum_{k=1}^{\infty}g^{k}\cdot\tilde{\alpha}_{k}dt$
into the deterministic part of (3.3) allow us to assume
$\tilde{\alpha}_{k}=(0,\dots,0)$. Moreover, since $B^{k,j}$’s are independent
1-dimensional Wiener processes where $B^{k}=(B^{k,1},\dots B^{k,m})$, (3.3)
can be written as
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f\right)\,dt+\sum_{i=1}^{\infty}h^{k}dW^{k}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}g^{k,j}dY^{k,j}_{t},\quad
u(0)=u_{0},$ (3.4)
for some $h=(h^{1},h^{2},\cdots)$ and independent one-dimensional Wiener
processes $W^{k}_{t}$ and
$Y^{k}_{t}:=\int_{\mathbb{R}^{m}}z\widetilde{N}_{k}(t,dz).$ Note that
$Y^{k}_{t}$ are independent $m$-dimensional pure jump Lévy processes with Lévy
measure of $\nu^{k}$.
Furthermore by considering $u-v$, where $v$ is the solution of
$dv=a(\omega,t)\Delta^{\alpha/2}v\,dt\,+\sum_{i=k}^{\infty}h^{k}dW^{k}_{t},\quad
u(0)=0$
from Theorem 2.11, we find that without loss of generality we may also assume
$h^{k}$’s are all zero.
By $<M,N>$ we denote the bracket of real-valued square integrable martingales
$M$ and $N$. Also let $[M]$ denote the quadratic variation of $M$.
###### Remark 3.1
(i) Note that, if $\widehat{c}_{k,2}<\infty$, then
$Y^{k,i}=\int_{\mathbb{R}^{m}}z^{i}\widetilde{N}_{k}(t,dz)$ is a square
integrable martingale for each $k\geq 1$ and $i=1,\cdots,m$. Also for any
$\overline{\mathcal{P}}$-measurable process $H=(H^{1},\dots,H^{m})\in
L_{2}(\Omega\times[0,T],\mathbb{R}^{m})$ which has a predictable version
$\bar{H}=(\bar{H}^{1},\dots,\bar{H}^{m})$,
$M^{k}_{t}:=\int_{0}^{t}H_{s}\cdot
dY^{k}_{s}=\sum_{i=1}^{m}\int_{0}^{t}\int_{\mathbb{R}^{m}}H^{i}_{s}z^{i}\widetilde{N}_{k}(ds,dz)=\sum_{i=1}^{m}\int_{0}^{t}\int_{\mathbb{R}^{m}}\bar{H}^{i}_{s}z^{i}\widetilde{N}_{k}(ds,dz)$
is a square integrable martingale with
$[M^{k}]_{t}=\sum_{i,j=1}^{m}\int^{t}_{0}\int_{\mathbb{R}^{m}}H^{i}H^{j}z^{i}z^{j}N(ds,dz),$
$E[M^{k}]_{t}=\sum_{i,j}(\int_{\mathbb{R}^{m}}z^{i}z^{j}\nu^{k}(dz)){\mathbb{E}}\int^{t}_{0}H^{i}(s)H^{j}(s)ds\leq\widehat{c}^{2}m^{2}{\mathbb{E}}\int_{0}^{t}|H_{s}|^{2}ds.$
(ii) Suppose that (3.2) holds. Then, for any $1\leq j\leq m$,
$g^{\cdot,j}\in\mathbb{H}^{\gamma}_{p}(T,\ell_{2})$ and $\phi\in
C^{\infty}_{0}(\mathbb{R}^{d})$, the series of stochastic integral
$\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)dY^{k,j}_{s}$
defines a square integrable martingale on $[0,T]$, which is right continuous
with left limits. Indeed, denote
$M_{n}:=\sum_{k=1}^{n}\sum_{j=1}^{m}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)dY^{k,j}_{s}$,
then the quadratic variation of $M_{n}$ is
$[M_{n}]_{t}=\sum_{k=1}^{n}\sum_{i,j=1}^{m}\int^{t}_{0}\int_{\mathbb{R}^{m}}(g^{k,i},\phi)(g^{k,j}(s,\cdot),\phi)z^{i}z^{k}N^{k}(ds,dz),$
and
$\displaystyle{\mathbb{E}}[M_{n}]_{t}=\sum_{k=1}^{n}\sum_{i,j=1}^{m}{\mathbb{E}}\int^{t}_{0}(g^{k,i},\phi)(g^{k,j}(s,\cdot),\phi)\int_{\mathbb{R}^{m}}z^{i}z^{j}\nu_{k}(dz)ds\leq
c(m,\widehat{c})\sum_{k=1}^{n}\sum_{j=1}^{m}{\mathbb{E}}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)^{2}ds.$
Also, with $q:=p/(p-2)$, for every $1\leq j\leq m$,
$\displaystyle\sum_{k=1}^{\infty}\,{\mathbb{E}}\left[\int^{T}_{0}(g^{k,j}(s,\cdot),\phi)^{2}ds\right]=\sum_{k=1}^{\infty}{\mathbb{E}}\left[\int_{0}^{T}((1-\Delta)^{\gamma/2}g^{k,j}(s,\cdot),(1-\Delta)^{-\gamma/2}\phi)^{2}\,ds\right]$
$\displaystyle\leq$
$\displaystyle\|(1-\Delta)^{-\gamma/2}\phi\|_{1}\,{\mathbb{E}}\left[\int^{T}_{0}\Big{(}\sum_{k=1}^{\infty}|(1-\Delta)^{\gamma/2}g^{k,j}(s,\cdot)|^{2},\,|(1-\Delta)^{-\gamma/2}\phi|\Big{)}\,ds\right]$
$\displaystyle\leq$
$\displaystyle\|(1-\Delta)^{-\gamma/2}\phi\|_{1}\,\|(1-\Delta)^{-\gamma/2}\phi\|_{q}\,{\mathbb{E}}\left[\int^{T}_{0}\Big{\|}\,\sum_{k=1}^{\infty}|(1-\Delta)^{\gamma/2}g^{k,j}(s,\cdot)|^{2}\Big{\|}_{p/2}\,ds\right]$
$\displaystyle\leq$
$\displaystyle\|(1-\Delta)^{-\gamma/2}\phi\|_{1}\,\|(1-\Delta)^{-\gamma/2}\phi\|_{q}\,T^{1-\frac{2}{p}}\,\|g^{\cdot,j}\|_{\mathbb{H}^{\gamma}_{p}(T,l^{2})}^{2}<\infty.$
It follows that $[M_{n}]_{t}$ converges in probability uniformly on $[0,T]$
and this certainly proves the claim.
###### Definition 3.2
Write $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ if
$u\in\mathbb{H}^{\gamma+\alpha}_{p}(T),u(0)\in
U^{\gamma+\alpha-\alpha/p}_{p}$, and for some $f\in\mathbb{H}^{\gamma}_{p}(T)$
$h\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$ and
$g^{\cdot,j}\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2}),1\leq j\leq m$
$du=f\,dt+\sum_{k=1}^{\infty}h^{k}dW^{k}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}g^{k,j}dY^{k,j}_{t},\quad
u(0)=u_{0},\quad\hbox{for }t\in[0,T]$
in the sense of distributions, that is, for any $\phi\in
C^{\infty}_{0}(\mathbb{R}^{d})$,
$(u(t,\cdot),\phi)=(u(0,\cdot),\phi)+\int^{t}_{0}(f(s,\cdot),\phi)ds+\sum_{k=1}^{\infty}\int^{t}_{0}(h^{k}(s,\cdot),\phi)ds+\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)dY^{k,j}_{s}$
(3.5)
holds for all $t\leq T$ $a.s.$. In this case we write
$\mathbb{D}u:=f,\quad\mathbb{S}_{c}u:=(h^{1},\dots
h^{k},\dots),\quad\mathbb{S}^{k,j}_{d}u:=g^{k,j},\quad\mathbb{S}^{\cdot,j}_{d}u:=(g^{1,j},\dots
g^{k,j},\dots)$
and define
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}:=\|u\|_{\mathbb{H}^{\gamma+\alpha}_{p}(T)}+\|\mathbb{D}u\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|\mathbb{S}_{c}u\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\sum_{j=1}^{m}\|\mathbb{S}^{\cdot,j}_{d}u\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\|u(0)\|_{U^{\gamma+\alpha-\alpha/p}_{p}}.$
To prove that $\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a Banach space we need
the following result, which is an infinite dimensional extension of Kunita’s
inequality (for example, see [2, Theorem 4.4.23]). In fact, if $m=1$ then the
proof is given in [5].
###### Lemma 3.3
Suppose $1\leq j\leq m$, $g^{\cdot,j}(\omega,t)=(g^{1,j},g^{2,j},\cdots)$’s
are $\ell_{2}$-valued predictable processes such that each
$g^{k}=(g^{k,1},\dots g^{k,m})$ is bounded. Then, under the assumption (3.2),
$\displaystyle{\mathbb{E}}\left[\left(\sum_{k=1}^{\infty}\int^{t}_{0}\int_{\mathbb{R}^{m}}|g^{k}(s)|^{2}\,|z|^{2}N_{k}(s,dz)ds\right)^{p/2}\right]$
(3.6) $\displaystyle\leq$ $\displaystyle
c(p)\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s)|^{2}ds\right)^{p/2}+\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s)|^{p}\,ds\right].$
Proof. Due to monotone convergence theorem we may assume $g^{k,j}=0$ for all
$i>M$ and $1\leq j\leq m$. By monotone convergence theorem,
$\displaystyle A$ $\displaystyle:=$
$\displaystyle{\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{\mathbb{R}^{m}}|g^{k}(s)|^{2}|z|^{2}N_{k}(s,dz)ds\right)^{p/2}\right]$
$\displaystyle=$
$\displaystyle\lim_{N\to\infty}{\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{2}|z|^{2}N_{k}(s,dz)ds\right)^{p/2}\right].$
Since $(a+b)^{p/2}\leq c(p)(|a|^{p/2}+|b|^{p/2})$ and
$\widetilde{N}_{k}(s,dz):=N_{k}(s,dz)-s\nu_{k}(dz)$,
$\displaystyle A\leq
c(p)\lim_{N\to\infty}{\mathbb{E}}\left[(J_{2,t})^{p/2}\right]+c(p){\mathbb{E}}\left[\left(\int^{t}_{0}\int_{\mathbb{R}^{m}}\sum_{k=1}^{M}|g^{k}(s)|^{2}\,|z|^{2}\nu_{k}(dz)ds\right)^{p/2}\right]$
where $J_{n,t}:=\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{n}|z|^{n}\widetilde{N}_{k}(s,dz)ds$, which is a square
integrable martingale becuase $g^{k}$ are bounded predictable processes. By
Burkholder-Davis-Gundy inequality (For example, see [22, Theorem 48].)
$\displaystyle{\mathbb{E}}\left[(J_{2,t})^{p/2}\right]\leq
c(p){\mathbb{E}}\left[[J_{2}]_{t}^{p/4}\right]=c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{4}|z|^{4}N_{k}(s,dz)ds\right)^{p/4}\right]$ (3.7)
$\displaystyle\leq c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\,\,\sum_{0\leq
s\leq t}|g^{k}(s)|^{4}\,|\Delta Y^{k}_{s}|^{4}\right)^{p/4}\right].$
Recall that for any $q>1$, $(\sum|a_{n}|^{q})^{1/q}\leq\sum|a_{n}|$. Thus if
$2<p\leq 4$, then
$\displaystyle{\mathbb{E}}\left[(J_{2,t})^{p/2}\right]\leq
c(p){\mathbb{E}}\left[\sum_{k=1}^{M}\,\,\sum_{0\leq s\leq
t}|g^{k}(s)|^{p}\,|\Delta Y^{k}_{s}|^{p}\right]\leq
c(p,\widehat{c}){\mathbb{E}}\left[\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{p}\,ds\right].$
If $4<p\leq 8$ then, by the relation
$\widetilde{N}_{k}(s,dz)=N_{k}(s,dz)-s\nu_{k}(dz)$ and Burkholder-Davis-Gundy
inequality,
$\displaystyle{\mathbb{E}}\left[(J_{2,t})^{p/2}\right]\leq
c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{4}\,|z|^{4}N_{k}(s,dz)ds\right)^{p/4}\right]$
$\displaystyle\leq$ $\displaystyle
c(p){\mathbb{E}}\left[(J_{8,t})^{p/8}\right]+c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g(s)|^{4}\,|z|^{4}\nu_{k}(dz)ds\right)^{p/4}\right]$ $\displaystyle\leq$
$\displaystyle
c(p,\widehat{c}){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{8}|z|^{8}N_{k}(s,dz)ds\right)^{p/8}+\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{4}ds\right)^{p/4}\right]$
$\displaystyle\leq$ $\displaystyle
c(p,\widehat{c}){\mathbb{E}}\left[\left(\sum_{k=1}^{M}\int^{t}_{0}\int_{|z|\leq
N}|g^{k}(s)|^{p}|z|^{p}N_{k}(s,dz)ds\right)+\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{4}ds\right)^{p/4}\right]$
$\displaystyle\leq$ $\displaystyle
c(p,\widehat{c}){\mathbb{E}}\left[\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{p}\,ds+\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{4}ds\right)^{p/4}\right].$
Similarly, in general, for $p\in(2^{n-1},2^{n}]$,
$\displaystyle A\leq
c(p,\widehat{c})\,\sum_{j=1}^{n}{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{2^{j}}ds\right)^{p2^{-j}}\right]+c(p,\widehat{c})\,{\mathbb{E}}\left[\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s,x)|^{p}ds\right].$
Also since for each $2\leq q\leq p$,
$\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{q}ds\right)^{1/q}\leq\left(\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{2}ds\right)^{1/2}+\left(\int^{t}_{0}\sum_{k=1}^{M}|g^{k}(s)|^{p}ds\right)^{1/p}\right),$
we get
$\displaystyle A\leq
c(p,\widehat{c})\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s)|^{2}ds\right)^{p/2}+\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s)|^{p}\,ds\right].$
(3.8)
Thus the lemma is proved. $\Box$
###### Theorem 3.4
Suppose that (3.2) holds. For any $p\in[2,\infty)$ and $\gamma\in\mathbb{R}$,
$\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is a Banach space with norm
$\|\cdot\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}$. Moreover, there is a
constant $c=c(d,p,T)>0$ such that for every
$u\in\mathcal{H}^{\gamma+2}_{p}(T)$ and $0<t\leq T$,
${\mathbb{E}}\left[\sup_{s\leq
t}\|u(s,\cdot)\|^{p}_{H^{\gamma}_{p}}\right]\leq
c(p,d,T)\left(\|\mathbb{D}u\|^{p}_{\mathbb{H}^{\gamma}_{p}(t)}+\|\mathbb{S}_{c}u\|_{\mathbb{H}^{\gamma}_{p}(t,\ell_{2})}+\sum_{j=1}^{m}\|\mathbb{S}^{\cdot,j}_{d}u\|^{p}_{\mathbb{H}^{\gamma}_{p}(t,\ell_{2})}+\|u_{0}\|^{p}_{U^{\gamma}_{p}}\right).$
(3.9)
Proof. By Theorem 2.6 and the reasons explained just before Remark 3.1,
without loss of generality we assume that
$Y^{k}_{t}=\int_{\mathbb{R}^{m}}z\widetilde{N}_{k}(t,dz)$. Moreover, due to
Remark 2.8 it suffices to prove the theorem only for $\gamma=0$. First we
prove (3.9). Let $du=fdt+\sum_{k=1}^{\infty}g^{k}\cdot dY^{k}_{t}$ with
$u(0)=u_{0}$.
For a moment, we assume that $g^{k,j}=0$ for all $k\geq N_{0},1\leq j\leq m$
and $g^{k,j}$ is of the type
$g^{k,j}(t,x)=\sum_{i=0}^{m_{k}}I_{(\tau^{k,j}_{i},\tau^{k,j}_{i+1}]}(t)g^{k_{i},j}(x),$
(3.10)
where $\tau^{k,j}_{i}$ are bounded stopping times and $g^{k_{i},j}\in
C^{\infty}_{0}(\mathbb{R}^{d})$. Define
$v(t,x)=\sum_{k=1}^{N_{0}}\int^{t}_{0}g^{k}(s,x)\cdot dY^{k}_{s}.$
Then by Burkholder-Davis-Gundy inequality and Lemma 3.3,
$\displaystyle{\mathbb{E}}\left[\sup_{s\leq
t}|v(s,x)|^{p}\right]={\mathbb{E}}\left[\sup_{s\leq
t}\left|\sum_{k=1}^{N_{0}}\sum_{j=1}^{m}\int^{t}_{0}\int_{\mathbb{R}^{m}}g^{k,j}(s,x)z^{j}\widetilde{N}_{k}(s,dz)ds\right|^{p}\right]$
$\displaystyle\leq$ $\displaystyle
c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{\infty}\sum_{i,j=1}^{m}\int^{t}_{0}\int_{\mathbb{R}^{m}}g^{k,i}(s,x)g^{k,j}(s,x)z^{i}z^{j}N_{k}(s,dz)ds\right)^{p/2}\right]$
$\displaystyle\leq$ $\displaystyle
c(p,\widehat{c}){\mathbb{E}}\left[\left(\sum_{k=1}^{\infty}\int^{t}_{0}\int_{\mathbb{R}^{m}}|g^{k}(s,x)|^{2}|z|^{2}N_{k}(s,dz)ds\right)^{p/2}\right]$
$\displaystyle\leq$ $\displaystyle
c(p,\widehat{c})\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s,x)|^{2}ds\right)^{p/2}+\int^{t}_{0}\sum_{k=1}^{\infty}|g^{k}(s,x)|^{p}ds\right].$
Since $\sum_{n}|a_{n}|^{p}\leq(\sum_{n}|a_{n}|^{2})^{p/2}$ and
$(\int^{t}_{0}|f|ds)^{p}\leq t^{p-1}\int^{t}_{0}|f|^{p}ds$, by integrating
over $\mathbb{R}^{d}$ we get that for every $t\leq T$
${\mathbb{E}}\left[\sup_{s\leq t}\|v\|^{p}_{p}\right]\leq
c(T,p)\sum_{j=1}^{m}\|g^{\cdot,j}\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}:=c(T,p)\,\sum_{j=1}^{m}{\mathbb{E}}\int^{t}_{0}\int_{\mathbb{R}^{d}}|g^{\cdot,j}|^{p}_{\ell_{2}}\,dxds.$
(3.11)
Next we prove (3.11) for general $g^{\cdot,j}\in\mathbb{L}_{p}(T,\ell_{2})$.
By Theorem 3.10 in [14], we can take a sequence
$g^{\cdot,j}_{n}\in\mathbb{L}_{p}(T,\ell_{2})$ so that for each fixed $n$,
$g^{k,j}_{n}=0$ for all large $k$ and each $g^{k,j}_{n}$ is of of the type
(3.10), and $g^{\cdot,j}_{n}\to g^{\cdot,j}$ in $\mathbb{L}_{p}(T,\ell_{2})$
as $n\to\infty$. Define
$v_{n}(t,x)=\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}g^{k,j}_{n}dY^{k,j}_{t}$,
then for every $t\leq T$
${\mathbb{E}}\left[\sup_{s\leq t}\|v_{n}\|^{p}_{p}\right]\leq
c(T,p)\sum_{j=1}^{m}\|g^{\cdot,j}_{n}\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})},\quad{\mathbb{E}}\left[\sup_{s\leq
t}\|v_{n_{1}}-v_{n_{2}}\|^{p}_{p}\right]\leq
c(T,p)\sum_{j=1}^{m}\|g^{\cdot,j}_{n_{1}}-g^{\cdot,j}_{n_{2}}\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}.$
Thus (3.11) follows by taking $n\to\infty$. Now note that
$d(u-v)=fdt\quad\hbox{ with}\quad(u-v)(0)=u_{0}.$
Thus it is easy to check that
${\mathbb{E}}\left[\sup_{s\leq t}\|u-v\|^{p}_{p}\right]\leq
N{\mathbb{E}}\left[\|u_{0}\|^{p}_{p}\right]+N{\mathbb{E}}\left[\int^{t}_{0}\|f(s,\cdot)\|^{p}_{p}\,ds\right].$
Consequently,
${\mathbb{E}}\left[\sup_{s\leq t}\|u\|^{p}_{p}\right]\leq
N\|f\|^{p}_{\mathbb{H}^{0}_{p}(t)}+c\sum_{j=1}^{m}\|g^{\cdot,j}\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}+N{\mathbb{E}}\|u_{0}\|^{p}_{p}.$
The completeness of the space $\mathcal{H}^{\alpha}_{p}(T)$ easily follows
from (3.9). Indeed, let $\\{u_{n}:n=1,2,\cdots\\}$ be a Cauchy sequence in
$\mathcal{H}^{\alpha}_{p}(T)$. Then $\\{u_{n}\\}$, $\\{\mathbb{D}u_{n}\\}$,
$\\{\mathbb{S}^{\cdot,j}_{d}u_{n}\\}$ and $\\{u_{n}(0)\\}$ are Cauchy
sequences in
$\mathbb{H}^{\alpha}_{p},\mathbb{L}_{p}(T),\mathbb{H}^{\alpha/2}_{p}(T,\ell_{2})$
and $U^{\alpha/2-\alpha/p}_{p}$ respectively. Thus there exist
$u\in\mathbb{H}^{\alpha}_{p}(T)$,
$f\in\mathbb{L}_{p}(T),g^{\cdot,j}\in\mathbb{H}^{\alpha/2}_{p}(T,\ell_{2})$
and $u_{0}\in U^{\alpha/2-\alpha/p}_{p}$ so that
$u_{n},\mathbb{D}u_{n},\mathbb{S}^{\cdot,j}_{d}u_{n},u_{n}(0)$ converge to
$u,f,g^{\cdot,j},u_{0}$ respectively, that is,
$\|u_{n}-u\|_{\mathbb{H}^{\alpha}_{p}(T)}+\|\mathbb{D}u_{n}-f\|_{\mathbb{L}_{p}(T)}+\|\mathbb{S}^{\cdot,j}u_{n}-g^{j}\|_{\mathbb{H}^{\alpha/2}_{p}(T,\ell_{2})}+\|u_{n}(0)-u_{0}\|_{U^{\alpha/2-\alpha/p}_{p}}\to
0$
as $n\to\infty$. Thus to prove $u\in\mathcal{H}^{\alpha}_{p}(T)$ and $u_{n}\to
u$ in $\mathcal{H}^{\alpha}_{p}(T)$, we only need to show that for any
$\phi\in C^{\infty}_{0}(\mathbb{R}^{d})$, the equality
$(u(t,\cdot),\phi)=(u_{0},\phi)+\int^{t}_{0}(f(s,\cdot),\phi)ds+\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)dY^{k,j}_{s}$
(3.12)
holds for all $t\leq T$ (a.s.). Taking the limit from
$(u_{n}(t,\cdot),\phi)=(u_{n}(0),\phi)+\int^{t}_{0}(\mathbb{D}u_{n}(s,\cdot),\phi)ds+\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}(\mathbb{S}^{k,j}_{d}u_{n}(s,\cdot),\phi)dY^{k,j}_{s}$
and using the argument used in Remark 3.1(ii) one can show that (3.12) holds
in $\Omega\times[0,T]$ (a.e.). Also using the inequality (see (3.9))
${\mathbb{E}}\left[\sup_{t\leq
T}\|u_{n}(\cdot,t)-u_{m}(\cdot,t)\|^{p}_{L_{p}}\right]\leq
N\|u_{n}-u_{m}\|_{\mathcal{H}^{\alpha}_{p}(T)}$
and taking $m\to\infty$, one finds that $(u(t,\cdot),\phi)$ is right
continuous with left limits, and consequently (3.12) holds for all $t\leq T$
(a.s.). The theorem is proved. $\Box$
###### Lemma 3.5
Let $p\in(2,\infty)$, $t>0$ and $f\in L_{p}([0,t\,]\times\mathbb{R}^{d})$.
Then for any $\varepsilon>\alpha(1/2-1/p)$,
$\int_{\mathbb{R}^{d}}\int^{t}_{0}\int^{s}_{0}|\partial^{\alpha/2}_{x}T_{s-r}f(r,x)|^{p}\,drds\,dx\leq
c\int^{t}_{0}\|f(s,\cdot)\|^{p}_{H^{\varepsilon}_{p}}\,ds,$ (3.13)
where $c=c(d,p,\alpha,\varepsilon)$ is independent of $t$.
Proof. Note that we may assume $\alpha(1/2-1/p)<\varepsilon<\alpha/2$. Let
$q>p$ be chosen so that
$\frac{1}{p}=(1-\frac{2\varepsilon}{\alpha})\times\frac{1}{2}+\frac{2\varepsilon}{\alpha}\times\frac{1}{q}.$
Such choice of $q$ is possible since
$1/p>(1-\frac{2\varepsilon}{\alpha})\times\frac{1}{2}$. We will use an
interpolation theorem. First, note that
$\varepsilon=(1-\frac{2\varepsilon}{\alpha})\times
0+\frac{2\varepsilon}{\alpha}\times\frac{\alpha}{2}.$
Define an operator $\mathcal{A}$ by
$\mathcal{A}f(s,r,x)=\begin{cases}\partial^{\alpha/2}T_{s-r}f\quad&\hbox{if
}r<s,\\\ 0&\hbox{otherwise}.\end{cases}$
Then, due to (2.19) and the inequality
$\|T_{s-r}\partial^{\alpha/2}f\|_{q}\leq\|\partial^{\alpha/2}f\|_{q}\leq\|f\|_{H^{\alpha/2}_{q}}$,
the linear mappings
$\mathcal{A}:L_{2}([0,t],L_{2}(\mathbb{R}^{d}))\to
L_{2}([0,t]\times[0,t]\times\mathbb{R}^{d})$
and
$\mathcal{A}:L_{q}([0,t],H^{\alpha/2}_{q})\to
L_{q}([0,t]\times[0,t]\times\mathbb{R}^{d})$
are bounded and their norms are independent of $t$. It follows from the
interpolation theory (see, for instance, [3, Theorem 5.1.2]) that the operator
$\mathcal{A}:L_{p}([0,t],H^{\varepsilon}_{p}(\mathbb{R}^{d}))\to
L_{p}([0,t]\times[0,t]\times\mathbb{R}^{d})$
is bounded and its norm is independent of $t$. The lemma is proved. $\Box$
###### Theorem 3.6
Fix a constant $\varepsilon_{1}$ so that $\varepsilon_{1}=0$ if $p=2$, and
$\varepsilon_{1}>\alpha(1/2-1/p)$ if $p>2$. suppose (3.2) holds. Then for any
$f\in\mathcal{H}^{\gamma}_{p}(T),h\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2}),g^{\cdot,j}\in\mathcal{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2}),1\leq
j\leq m$ and $u_{0}\in U^{\gamma+\alpha/2-\alpha/p}_{p}$, equation (3.4) has a
unique solution $u$ in $\mathcal{H}^{\gamma+\alpha}_{p}(T)$, and for this
solution
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}\leq
c(p,T,\delta)\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(t)}+\|h\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(t,\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(t,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right)$
(3.14)
for every $t\leq T$.
Proof. As explained before, without loss of generality we assume that
$h^{i}$’s are all zeros.
Step 1. As in the proof of Theorem 2.11, we only need to prove the theorem for
a particular $\gamma=\gamma_{0}$.
Step 2. We assume $a(\omega)=1$ and prove the theorem for the equation:
$du=\Delta^{\alpha/2}udt+\sum_{k=1}^{\infty}g^{k}dY^{k}_{t},\quad u(0)=0.$
(3.15)
By the result of Step 1, we may assume that $\gamma=-\alpha/2$. The uniqueness
is obvious and we only prove the existence and the estimate (3.14).
Considering approximation arguments, for a moment, we assume that $g^{k,j}=0$
for all $k>N_{0}$ and $1\leq j\leq m$ and that
$g^{k,j}(t,x)=\sum_{i=0}^{m_{k}}I_{(\tau^{k,j}_{i},\tau^{k,j}_{i+1}]}(t)g^{k_{i},j}(x),$
where $\tau^{k,j}_{i}$ are bounded stopping times and $g^{k_{i},j}(x)\in
C^{\infty}_{0}(\mathbb{R}^{d})$. Define
$v(t,x):=\sum_{k=1}^{N_{0}}\int^{t}_{0}g^{k}(s,x)\cdot
dY^{k}_{s}=\sum_{k=1}^{N_{0}}\sum_{i=1}^{m_{k}}\sum_{j=1}^{m}g^{k_{i},j}(x)(Y^{k,j}_{t\wedge\tau^{k}_{i+1}}-Y^{k,j}_{t\wedge\tau^{k}_{i}})$
and
$u(t,x):=v(t,x)+\int^{t}_{0}\Delta^{\alpha/2}T_{t-s}v\,ds=v(t,x)+\int^{t}_{0}T_{t-s}\Delta^{\alpha/2}v\,ds.$
(3.16)
Now we remember from the proof of Theorem 2.11 that if functions
$h_{1}=h_{1}(t,x)$ and $h_{2}=h_{2}(x)$ are sufficiently smooth, then
$w_{1}(t,x):=\int^{t}_{0}T_{t-s}h_{1}(s)ds,\quad w_{2}(t,x)=T_{t}h_{2}$
solve
$dw_{1}=(\Delta^{\alpha/2}w+h_{1})\,dt,\quad w_{1}(0)=0,$
$dw_{2}=\Delta^{\alpha/2}w_{2}\,dt,\quad w_{2}(0)=h_{2}.$
Therefore we have
$d(u-v)=(\Delta^{\alpha/2}(u-v)+\Delta^{\alpha/2}v)dt=\Delta^{\alpha/2}udt$,
and
$du=\Delta^{\alpha/2}udt+dv=\Delta^{\alpha/2}udt+\sum_{k=1}^{N_{0}}g^{k}\cdot
dY^{k}_{t}.$
Let $T_{t-s}g^{k}(r,x)=(T_{t-s}g^{k,1}(r,x),\dots T_{t-s}g^{k,m}(r,x))$. By
(3.16) and stochastic Fubini theorem ([22, Theorem 64]), almost surely,
$\displaystyle u(t,x)$ $\displaystyle=$ $\displaystyle
v(t,x)+\sum_{k=1}^{N_{0}}\int^{t}_{0}\int^{s}_{0}\Delta^{\alpha/2}T_{t-s}g^{k}(r,x)\cdot
dY^{k}_{r}ds$ (3.17) $\displaystyle=$ $\displaystyle
v(t,x)-\sum_{k=1}^{N_{0}}\sum_{j=1}^{m}\int^{t}_{0}\int^{t}_{r}\frac{\partial}{\partial
s}T_{t-s}g^{k,j}(r,x)dsdY^{k,j}_{r}$ $\displaystyle=$
$\displaystyle\sum_{k=1}^{N_{0}}\int^{t}_{0}T_{t-s}g^{k}(s,x)\cdot
dY^{k}_{s}.$
Hence,
$\partial^{\alpha/2}_{x}u(t,x)=\sum_{k=1}^{N_{0}}\int^{t}_{0}\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,x)\cdot
dY^{k}_{s}=\sum_{k=1}^{N_{0}}\sum_{j=1}^{m}\int^{t}_{0}\partial^{\alpha/2}_{x}T_{t-s}g^{k,j}(s,x)dY^{k,j}_{s}.$
By Burkholder-Davis-Gundy’s inequality and Lemma 3.3, we have for every
$0<t\leq T$
$\displaystyle{\mathbb{E}}\left[|\partial^{\alpha/2}_{x}u(t,x)|^{p}\right]\leq
c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{N_{0}}\sum_{i,j=1}^{m}\int^{t}_{0}\int_{\mathbb{R}^{m}}\partial^{\alpha/2}_{x}T_{t-s}g^{k,i}(s,x)\partial^{\alpha/2}_{x}T_{t-s}g^{k,j}(s,x)z^{i}z^{j}N^{k}(dz,ds)\right)^{p/2}\right]$
$\displaystyle\leq\,c(p){\mathbb{E}}\left[\left(\sum_{k=1}^{N_{0}}\int^{t}_{0}\int_{\mathbb{R}^{m}}|\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,x)|^{2}|z|^{2}N^{k}(dz,ds)\right)^{p/2}\right]$
$\displaystyle\leq
c(p)\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{\infty}|\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,x)|^{2}ds\right)^{p/2}+\int^{t}_{0}\sum_{k=1}^{\infty}|\partial^{\alpha/2}_{x}T_{t-s}g^{k}(s,x)|^{p}ds\right].$
By (2.19), Lemma 3.5 and the inequality
$\sum_{k=1}^{\infty}|a_{k}|^{p}\leq(\sum_{k=1}^{\infty}|a_{n}|^{2})^{p/2}$,
${\mathbb{E}}\left[\int^{t}_{0}\|\partial^{\alpha/2}_{x}u(s,\cdot)\|^{p}_{p}\,ds\right]\leq
c(p,\alpha)\sum_{j=1}^{m}{\mathbb{E}}\left[\int^{t}_{0}\|g^{\cdot,j}(s,\cdot)\|^{p}_{H^{\varepsilon_{1}}_{p}(\ell_{2})}\,dt\right].$
(3.18)
Similarly from (3.17) we also get, for every $0<t\leq T$,
${\mathbb{E}}\left[|u(t,x)|^{p}\right]\leq
c(p)\,{\mathbb{E}}\left[\left(\int^{t}_{0}\sum_{k=1}^{N_{0}}|T_{t-s}g^{k}(s,x)|^{2}ds\right)^{p/2}+\int^{t}_{0}\sum_{k=1}^{N_{0}}|T_{t-s}g^{k}(s,x)|^{p}\,ds\right].$
(3.19)
By the same argument which leads to (2.22), we see that the right side of
(3.19) is finite. Thus we proved
$\partial^{\alpha/2}_{x}u,u\in\mathbb{L}_{p}(T)$, and hence
$u\in\mathcal{H}^{\alpha/2}_{p}(T)$. As in (2.23) and (2.24),
$\displaystyle\|u\|^{p}_{\mathcal{H}^{\alpha/2}_{p}(t)}$ $\displaystyle\leq$
$\displaystyle
c(p)\left(\|u\|^{p}_{\mathbb{H}^{\alpha/2}_{p}(t)}+\|\Delta^{\alpha/2}u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|g\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}\right)$
$\displaystyle\leq$ $\displaystyle
c\left(\|u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|\partial^{\alpha/2}_{x}u\|^{p}_{\mathbb{L}_{p}(t)}+\|g\|^{p}_{\mathbb{L}_{p}(t,\ell_{2})}\right)$
$\displaystyle\leq$ $\displaystyle
c(p,T,\alpha)\left(\|u\|^{p}_{\mathbb{H}^{-\alpha/2}_{p}(t)}+\|g\|^{p}_{\mathbb{H}^{\varepsilon_{1}}_{p}(t,\ell_{2})}\right)$
$\displaystyle\leq$ $\displaystyle
c(p,T,\alpha)\int_{0}^{t}\|u\|^{p}_{\mathcal{H}^{\alpha/2}_{p}(s)}ds+c(p,T,\alpha)\|g\|^{p}_{\mathbb{H}^{\varepsilon_{1}}_{p}(T,\ell_{2})}.$
Finally, Gronwall leads to (3.14). Once one has a unique solvability of
equation (3.15) and estimate (3.14) for sufficiently smooth $g$, we repeat the
same approximation argument used in the Step 2 of the proof of Theorem 2.11.
Step 3. Now, we follow Step 3–Step 5 of the proof of Theorem 2.11 word for
word except obvious changes from $W^{k}_{t}$ to $Y^{k}_{t}$. The theorem is
proved. $\Box$
Finally we consider the nonlinear equation
$du=\left(a(\omega,t)\Delta^{\alpha/2}u+f(u)\right)\,dt+\sum_{k=1}^{\infty}h^{k}(u)dW^{i}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}g^{k,j}(u)dY^{k,j}_{t},\quad
u(0)=u_{0},$ (3.20)
where $f(u)=f(\omega,t,x,u)$, $h^{k}(u)=h^{k}(\omega,t,x,u)$,
$g^{k}(u)=(g^{k,1}(\omega,t,x,u),\dots,g^{k,m}(\omega,t,x,u))$, $W_{t}$ are
independent $1$-dimensional Wiener processes and
$Y^{k}_{t}:=\int_{\mathbb{R}^{m}}z\widetilde{N}_{k}(t,dz)$ are independent
$m$-dimensional pure jump Lévy processes with Lévy measure $\nu_{k}$.
###### Assumption 3.7
Fix a constant $\varepsilon_{1}$ so that $\varepsilon_{1}=0$ if $p=2$, and
$\varepsilon_{1}>\alpha(1/2-1/p)$ if $p>2$. Assume that
$f(0)\in\mathbb{H}^{\gamma}_{p}(T)$,
$g^{\cdot,j}(0)\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2})$
and $h(0)\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2}).$ Moreover, for any
$\varepsilon>0$, there exists a constant $K_{\varepsilon}$ so that for any
$u=u(x),v=v(x)\in H^{\gamma+\alpha}_{p}$ and $t,\omega$ we have
$\displaystyle\|f(t,\cdot,u(\cdot))-f(t,\cdot,v(\cdot))\|_{H^{\gamma}_{p}}+\sum_{j=1}^{m}\|g^{\cdot,j}(t,\cdot,u(\cdot))-g^{\cdot,j}(t,\cdot,v(\cdot))\|_{H^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\ell_{2})}$
$\displaystyle+\|h(t,\cdot,u(\cdot))-h(t,\cdot,v(\cdot))\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}\leq\varepsilon\|u-v\|_{H^{\gamma+\alpha}_{p}}+K(\varepsilon)\|u-v\|_{H^{\gamma}_{p}}.$
(3.21)
###### Example 3.8
Recall that the space $B^{r}$ is defined in (2.35). Fix
$\kappa_{0}=\kappa_{0}(\gamma)\geq 0$ so that $\kappa_{0}>0$ if $\gamma$ is
not integer. Consider
$\displaystyle f(u)$
$\displaystyle=b(\omega,t,x)\Delta^{\beta_{1}/2}u+\sum_{i=1}^{d}c^{i}(\omega,t,x)u_{x^{i}}I_{\alpha>1}+d(\omega,t,x)u+f_{0},$
$\displaystyle h^{k}(u)$
$\displaystyle=\eta^{k}(\omega,t,x)\Delta^{\beta_{2}/2}u+l^{k}(\omega,t,x)u+h^{k}_{0},$
$\displaystyle g^{k,j}(u)$
$\displaystyle=\sigma^{k,j}(\omega,t,x)\Delta^{\beta^{j}_{3}/2}u+v^{k,j}(\omega,t,x)u+g^{k,j}_{0},\quad
j=1,\dots m.$
Here $\beta_{1}<\alpha$, $\beta_{2}<\alpha/2$,
$\beta^{j}_{3}<\alpha/2-\varepsilon_{1}$ and
$f_{0}\in\mathbb{H}^{\gamma}_{p}(T)$,
$h_{0}\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})$,
$g^{\cdot,j}_{0}\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2})$.
Assume for each $\omega,t,i,j$,
$\displaystyle|b|_{B^{|\gamma|+\kappa_{0}}}+|c^{i}|_{B^{|\gamma|+\kappa_{0}}}+|d|_{B^{|\gamma|+\kappa_{0}}}+|\eta|_{B^{|\gamma|+\alpha/2+\kappa_{0}}}+|l|_{B^{|\gamma|+\alpha/2+\kappa_{0}}}$
$\displaystyle+|\sigma^{\cdot,j}|_{B^{|\gamma|+\alpha/2+\varepsilon_{1}+\kappa_{0}}}+|v^{\cdot,j}|_{B^{|\gamma|+\alpha/2+\varepsilon_{1}+\kappa_{0}}}\leq
K<\infty.$
Then the calculus in Example 2.14 shows that (3.7) holds.
Here is the main result of this section.
###### Theorem 3.9
Suppose (3.2) and Assumptions 2.10 and 3.7 hold. Then the equation (3.20) has
a unique solution $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$, and for this
solution we have
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}\leq
c\left(\|f(0)\|_{\mathbb{H}^{\gamma}_{p}(t)}+\|h(0)\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(t,\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}(0)\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(t,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha/2-\alpha/p}_{p}}\right),$
(3.22)
for every $t\leq T$, where $c=c(p,T,\delta)$.
Proof. As we mentioned in the previous section, our proof is a repetition of
that of Theorem 6.4 in [14]. By Theorem 3.6, for any
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ with initial data $u_{0}$ we can
define $v=\mathcal{R}u$ as the solution of
$\displaystyle
dv=\left(a(\omega,t)\Delta^{\alpha/2}v(t,x)+f(t,x,u)\right)dt+\sum_{i=1}^{\infty}h^{i}(t,x,u)dW^{i}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}g^{k,j}(t,x,u)dY^{k,j}_{t},\quad
v(0)=u_{0}.$
Then for any $u,v$ initial data $u_{0}$, we have
$(\mathcal{R}u-\mathcal{R}v)(0,x)=0$ and
$\displaystyle d(\mathcal{R}u-\mathcal{R}v)=$
$\displaystyle\left(a(\omega,t)\Delta^{\alpha/2}(\mathcal{R}u-\mathcal{R}v)+(f(t,x,u)-f(t,x,v))\right)dt$
$\displaystyle+\sum_{i=1}^{\infty}\int_{0}^{t}(h^{i}(t,x,u)-h^{i}(t,x,u))dW^{i}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int_{0}^{t}(g^{k,j}(t,x,u)-g^{k,j}(t,x,v))dY^{k,j}_{t}.$
By Theorems 2.11 and Assumption 3.7, for every $t\in(0,T]$,
$\displaystyle\|\mathcal{R}u-\mathcal{R}v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}$
$\displaystyle\leq
c(p,T,\delta)\Big{(}\|f(u)-f(v)\|^{p}_{\mathbb{H}^{\gamma}_{p}(t)}+\|h(u)-h(v)\|^{p}_{\mathbb{H}^{\gamma+\alpha/2}_{p}(t,\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}(u)-g^{\cdot,j}(v)\|^{p}_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(t,\ell_{2})}\Big{)}$
$\displaystyle\leq\varepsilon^{p}c(p,T,\delta)\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}+K(\varepsilon)c(p,T,\delta)\int_{0}^{t}{\mathbb{E}}\|u(s,\cdot)-v(s,\cdot)\|^{p}_{H^{\gamma}_{p}}ds$
$\displaystyle\leq\theta\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}+N\int_{0}^{t}\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(s)}ds$
where $\theta:=\varepsilon^{p}c(p,T,\delta)$ and
$N=c(p,T,\delta,\varepsilon)$. Denote
$\mathcal{R}^{n+1}u:=\mathcal{R}(\mathcal{R}^{n}u)$. Then by induction, for
every $t\in(0,T]$
$\displaystyle\|\mathcal{R}^{n}u-\mathcal{R}^{n}v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}$
$\displaystyle\leq$
$\displaystyle\theta^{n}\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(t)}+\sum_{k=1}^{n}{{n}\choose{k}}\theta^{n-k}N^{k}\int_{0}^{t}\frac{(t-s)^{k-1}}{(k-1)!}\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(s)}ds$
Therefore,
$\displaystyle\|\mathcal{R}^{n}u-\mathcal{R}^{n}v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}$
$\displaystyle\leq$
$\displaystyle\theta^{n}\sum_{k=0}^{n}{{n}\choose{k}}\frac{(NT/\theta)^{k}}{k!}\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}$
$\displaystyle\leq$ $\displaystyle(2\theta)^{n}\left(\sup_{k\geq
0}\frac{(NT/\theta)^{k}}{k!}\right)\|u-v\|^{p}_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}.$
Choose $\varepsilon>0$ so that $2\theta<1/2$, and then fix $n$ large enough so
that $(2\theta)^{n}\left(\sup_{k\geq 0}\frac{(NT/\theta)^{k}}{k!}\right)<1/2$.
Then $\bar{\mathcal{R}}:=\mathcal{R}^{n}$ is a contraction in
$\mathcal{H}^{\gamma+\alpha}_{p}(T)$ and obviously the unique fixed point $u$
under this map becomes the unique solution of (3.20). Moreover, the estimate
(3.22) also easily from Assumption 3.7, Theorems 3.6 and 3.4. We leave the
details to the readers as an exercise. $\Box$
## 4 Application and Extension
First, we consider equations with the random fractional Laplacian driven by
(Lévy) space-time white noise; Let $d=1$ and consider the equation
$du=(a(\omega,t)\Delta^{\alpha/2}u(t,x)+f(\omega,t,x,u(t,x)))dt+\xi(\omega,t,x)h(\omega,t,x,u(t,x))d\mathcal{Z}_{t}$
(4.23)
where $\mathcal{Z}_{t}$ is a cylindrical Lévy process on $L_{2}(\mathbb{R})$,
that is $\mathcal{Z}_{t}$ has an expansion of the form
$\mathcal{Z}_{t}=\sum_{k=1}^{\infty}\eta^{k}(x)Z^{k}_{t}$
where $\\{\eta^{k}:k=1,2,\dots\\}$ is an orthonormal basis in $L_{2}$ and
$Z^{k}_{t}$ are i.i.d. one-dimensional $\mathcal{F}_{t}$-adapted Lévy
processes (see [10] for the details). Using this expansion we can rewrite
(4.23) as follows :
$du=(a(\omega,t)\Delta^{\alpha/2}u+f(u))dt+\sum_{k=1}^{\infty}g^{k}(u)dZ^{k}_{t},$
(4.24)
where $g^{k}(u):=\xi(\omega,t,x)h(\omega,t,x,u(t,x))\eta^{k}(x)$.
Let $\gamma,p,s,r$ be constants satisfying
$0>\gamma+\alpha/2>-1,\quad p\geq 2r\geq 2,\quad 1\leq
r<(2\gamma+\alpha+2)^{-1},\quad s^{-1}+r^{-1}=1\,\,(1\leq s\leq\infty).$
(4.25)
Define
$R_{\gamma}(x):=|x|^{-(\gamma+\alpha/2+1)}\int^{\infty}_{0}t^{-(\gamma+\alpha/2+3)/2}e^{-tx^{2}-1/(4t)}dt.$
It is known that there exists a constant $c>0$ so that $cR_{\gamma}(x)$ is the
kernel of the operator $(1-\Delta)^{(\gamma+\alpha/2)/2}$, that is
$(1-\Delta)^{(\gamma+\alpha/2)/2}f=(cR_{\gamma}*f)(x)$.
###### Assumption 4.1
(i) For each $x$, $\xi=\xi(\omega,t,x)$ is predictable, and
$\|\xi(\omega,t,\cdot)\|_{L_{2s}}\leq K$ for each $\omega,t$.
(ii) For each $x,u$, the processes $f(\omega,t,x,u),h(\omega,t,x,u)$ are
predictable, and
$|f(\omega,t,x,u)-f(\omega,t,x,v)|\leq
K|u-v|,\quad|h(\omega,t,x,u)-h(\omega,t,x,v)|\leq K|u-v|.$
By following the arguments in the proof of [14, Lemma 8.4], we get the
following
###### Lemma 4.2
Let (4.25) hold. Take some functions $h_{0}=h_{0}(x)\in L_{p}(\mathbb{R})$,
$\xi_{0}=\xi_{0}(x)\in L_{2s}(\mathbb{R})$, and set
$g^{k}_{0}=\xi_{0}h_{0}\eta^{k}$. Then $g_{0}=\\{g^{k}_{0}\\}\in
H^{\gamma+\alpha/2}_{p}(\ell_{2})$ and
$\|g_{0}\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}=\|\overline{h}_{0,\gamma}\|_{p}\leq
N\|\xi_{0}\|_{2s}\|h_{0}\|_{p},$
where $N=\|R_{\gamma}\|_{2r}<\infty$ and
$\bar{h}_{0,\gamma}(x):=\left(\int_{\mathbb{R}}R_{\gamma}^{2}(x-y)\xi^{2}_{0}(y)h^{2}_{0}(y)dy\right)^{1/2}.$
We first discuss the case when $Z^{k}_{t}$ are independent one-dimensional
Wiener processes.
###### Theorem 4.3
Let $Z^{k}_{t}$ be independent one-dimensional Wiener processes. Suppose
(4.25) and Assumption 4.1 hold. Also assume
$\gamma\in(-\alpha,\frac{-1-\alpha}{2})$, $u_{0}\in
U^{\gamma+\alpha-\alpha/p}_{p}$ and
$I(p,T):=\left({\mathbb{E}}\int^{T}_{0}\left(\|f(t,\cdot,0)\|^{p}_{H^{\gamma}_{p}}+\|\bar{h}(t,\cdot,0)\|^{p}_{p}\right)ds\right)^{1/p}<\infty,$
(4.26)
where
$\bar{h}(t,x,0):=\left(\int_{\mathbb{R}}R_{\gamma}^{2}(x-y)\xi^{2}(y)h^{2}(t,y,0)dy\right)^{1/2}.$
Then equation (4.23) with initial data $u_{0}$ has a unique solution
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ and for this solution,
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}\leq
c\left(I(p,T)+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right).$
Proof. We check whether $f(u)$ and $g(u)$ satisfy condition (2.13). Since
$\gamma<0$ and $\gamma+\alpha>0$,
$\|f(u)-f(v)\|_{H^{\gamma}_{p}}\leq\|f(u)-f(v)\|_{L_{p}}\leq
K\|u-v\|_{L_{p}}\leq\varepsilon\|u-v\|_{H^{\gamma+\alpha}_{p}}+K(\varepsilon)\|u-v\|_{H^{\gamma}_{p}}.$
Also for $g(u)=\\{g^{k}(u)\\}$, by Lemma 4.2,
$\|g(0)\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}\leq\|R_{\gamma}\|_{2r}\|\xi\|_{2s}\|h(0)\|_{p}\leq
c\|h(0)\|_{p},$
$\|g(u)-g(v)\|_{H^{\gamma+\alpha/2}_{p}(\ell_{2})}\leq\|R_{\gamma}\|_{2r}\xi\|_{2s}\|h(u)-h(v)\|_{p}\leq
c\|u-v\|_{L_{p}}\leq\varepsilon\|u-v\|_{H^{\gamma+\alpha}_{p}}+K(\varepsilon)\|u-v\|_{H^{\gamma}_{p}}.$
Therefore condition (2.13) is satisfied and the theorem is proved. $\Box$
Now we consider space-time white noise with jump Lévy processes. Unlike
Theorem 4.3, in the case space-time white noise with jump Lévy processes,
$L_{p}$-theory is not satisfactory due to the condition
$\varepsilon_{1}>\alpha(1/2-1/p)$ if $p>2$. Thus we only give an
$L_{2}$-theory.
###### Theorem 4.4
Suppose $Z^{k}_{t}$ are independent one-dimensional jump Lévy processes with
Lévy measure $\nu$. Suppose (3.2), (4.25) and Assumption 4.1 hold with $p=2$.
Also assume $\gamma\in(-\alpha,\frac{-1-\alpha}{2})$, $u_{0}\in
U^{\gamma+\alpha-\alpha/2}_{2}$ and $I(2,T)<\infty$, where $I(2,T)$ is taken
from (4.26). Then equation (4.23) with initial data $u_{0}$ has a unique
solution $u\in\mathcal{H}^{\gamma+\alpha}_{2}(T)$ and for this solution,
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{2}(T)}\leq
c\left(I(2,T)+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/2}_{2}}\right).$
Proof. There is nothing to prove since conditions on $f$ and $g$ were already
checked in the proof of Theorem 4.3. $\Box$
For a stopping time $\tau$ relative to $\\{\mathcal{F}_{t}\\}$, denote
$(\\![0,\tau]\\!]:=\\{(\omega,t):0<t\leq\tau(\omega)\\}.$
Then obviously the process ${\bf 1}_{(\\![0,\tau]\\!]}(\omega,t)$ is left-
continuous and predictable. For an $H^{\gamma}_{p}$-valued
$\mathcal{P}^{dP\times dt}$-measurable process $u$, write
$u\in\mathbb{H}^{\gamma}_{p}(\tau)$ if
$\|u\|^{2}_{\mathbb{H}^{\gamma}_{p}(\tau)}:={\mathbb{E}}\left[\int^{\tau}_{0}\|u\|^{2}_{H^{\gamma}_{p}}ds\right]<\infty.$
We define the Banach spaces $\mathbb{L}_{p}(\tau)$,
$\mathbb{L}_{p}(\tau,\ell_{2})$ and $\mathcal{H}^{\gamma}_{p}(\tau)$
similarly. The following theorem plays the key role when we weaken condition
(3.2) later in the next section.
###### Theorem 4.5
Let $\tau\leq T$ be a stopping time. Fix a constant $\varepsilon_{1}$ so that
$\varepsilon_{1}=0$ if $p=2$, and $\varepsilon_{1}>\alpha(1/2-1/p)$ if $p>2$.
Then, under Assumption 2.10 and (3.2), for any
$f\in\mathbb{H}^{\gamma}_{p}(\tau)$,
$h\in\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau,\ell_{2})$,
$g^{\cdot,j}\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau,\ell_{2}),1\leq
j\leq m$ and $u_{0}\in U^{\gamma+\alpha/2-\alpha/p}_{p}$, equation (3.4) has a
unique solution $u$ in $\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$, and for this
solution
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(\tau)}\leq
c\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(\tau)}+\|h\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau,\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right),$
(4.27)
where $c=c(p,T,\delta)$ independent of $\tau$.
Proof. First we prove the existence and (4.27). Obviously we have
$\bar{f}:={\bf
1}_{(\\![0,\tau]\\!]}\,f\in\mathbb{H}^{\gamma}_{p}(T),\quad\bar{h}:={\bf
1}_{(\\![0,\tau]\\!]}\,h\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2}),\quad\bar{g^{\cdot,j}}:={\bf
1}_{(\\![0,\tau]\\!]}\,g^{\cdot,j}\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2}).$
Let $u\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ be the solution of (2.14) with
$\bar{f},\bar{h}$ and $\bar{g}$ instead of $f,h$ and $g$ respectively. Then,
since $\tau\leq T$, we have
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(\tau)}\leq\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(T)}$,
and by Theorem 3.6,
$\displaystyle\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(\tau)}$
$\displaystyle\leq$ $\displaystyle
c\left(\|\bar{f}\|_{\mathbb{H}^{\gamma}_{p}(T)}+\|\bar{h}\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2})}+\sum_{j=1}^{m}\|\bar{g}^{\cdot,j}\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha/2-\alpha/p}_{p}}\right)$
$\displaystyle=$ $\displaystyle
c\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(\tau)}+\|h\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau,\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau,\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha/2-\alpha/p}_{p}}\right).$
Now we prove the uniqueness. Let $u\in\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$
be a solution of equation (2.14). Then obviously,
${\bf
1}_{(\\![0,\tau]\\!]}\cdot(\mathbb{D}u-a(\omega,t)\Delta^{\alpha/2}u)\in\mathbb{H}^{\gamma}_{p}(T),\quad{\bf
1}_{(\\![0,\tau]\\!]}\cdot\mathbb{S}_{c}u\in\mathbb{H}^{\gamma+\alpha/2}_{p}(T,\ell_{2}),\quad{\bf
1}_{(\\![0,\tau]\\!]}\cdot\mathbb{S}^{\cdot,j}_{d}u\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(T,\ell_{2}).$
According to Theorem 3.6 we can define
$v\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ as the solution of
$\displaystyle dv$ $\displaystyle=$
$\displaystyle(a(\omega,t)\Delta^{\alpha/2}v+{\bf
1}_{(\\![0,\tau]\\!]}\,(\mathbb{D}u-a(\omega,t)\Delta^{\alpha/2}u))dt+\sum_{k=1}^{\infty}1_{(\\![0,\tau]\\!]}\,\mathbb{S}^{k}_{c}u\,dW^{k}_{t}$
(4.28)
$\displaystyle\quad+\sum_{k=1}^{\infty}\sum_{j=1}^{m}1_{(\\![0,\tau]\\!]}\,\mathbb{S}^{k,j}_{d}u\,dZ^{k}_{t},\qquad
v(0)=u(0).$
Then for $t\leq\tau$, $d(u-v)=\Delta^{\alpha/2}(u-v)dt$ and $(u-v)(0)=0$.
Therefore by Theorem 2.9, we conclude that $u(t)=v(t)$ for all $t\leq\tau$
a.s.. By replacing $u$ by $v$ for $t\leq\tau$, from (4.28) we find that $v$
satisfies
$dv=\left(a\Delta^{\alpha/2}v+f{\bf
1}_{(\\![0,\tau]\\!]}\right)dt+\sum_{k=1}^{\infty}1_{(\\![0,\tau]\\!]}\,h^{k}\,dW^{k}_{t}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}1_{(\\![0,\tau]\\!]}\,g^{k,j}\,dZ^{k}_{t},\quad
v(0)=u_{0}.$ (4.29)
We proved that if $u\in\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$ is a solution of
equation (2.14) then $u(t)=v(t)$ for all $t\leq\tau$ a.s.. This proves the
uniqueness of solution of equation (2.14) in the class
$\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$ because by Theorem 3.6
$v\in\mathcal{H}^{\gamma+\alpha}_{p}(T)$ is the unique solution of equation
(4.29). The theorem is proved. $\Box$
For a stopping time $\tau\leq T$ and $\gamma\in\mathbb{R}$, write
$u\in\mathbb{H}^{\gamma}_{p,{\rm loc}}(\tau)$ if there exists a sequence of
stopping times $\tau_{n}\uparrow\infty$ so that
$u\in\mathbb{H}^{\gamma}_{p}(\tau\wedge\tau_{n})$ for each $n$.
The following is a weakened version of (3.2).
###### Assumption 4.6
There exists an integer $N_{0}\geq 1$ so that $\widehat{c}_{k}<\infty$ for all
integer $k>N_{0}$.
###### Definition 4.7
Let $u_{0}\in U^{\gamma+\alpha-\alpha/p}_{p}$,
$f(0)\in\mathbb{H}^{\gamma}_{p}(\tau)$,
$h(0)\in\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau,\ell_{2})$ and
$g^{\cdot,j}(0)\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau,\ell_{2})$,
$1\leq j\leq m$. We say that $u\in\mathcal{H}^{\gamma+\alpha}_{p,{\rm
loc}}(\tau)$ is a path-wise solution to (3.4) if the followings hold;
(i) $u\in\mathbb{H}^{\gamma+\alpha}_{p,{\rm loc}}(\tau)$ and $u(t)$ is right
continuous with left limits in $H^{\gamma}_{p}$ for $t<\tau$ ($a.s.$),
(ii) for any $\phi\in C^{\infty}_{0}(\mathbb{R}^{d})$, the equality
$\displaystyle(u(t,\cdot),\phi)=$
$\displaystyle(u_{0},\phi)+\int^{t}_{0}a(\omega,s)(u(s,\cdot),\Delta^{\alpha/2}\phi)ds+\int^{t}_{0}(f(s,\cdot),\phi)ds$
$\displaystyle+\sum_{k=1}^{\infty}\int^{t}_{0}(h^{k}(s,\cdot),\phi)dW^{k}_{s}+\sum_{k=1}^{\infty}\sum_{j=1}^{m}\int^{t}_{0}(g^{k,j}(s,\cdot),\phi)dY^{k,j}_{s}$
(4.30)
holds for all $t<\tau$ $a.s.$.
###### Theorem 4.8
Let $\tau\leq T$. Suppose that Assumptions 2.10 and 4.6 hold. Then for any
$u_{0}\in U^{\gamma+\alpha-\alpha/p}_{p}$,
$f\in\mathbb{H}^{\gamma}_{p}(\tau)$,
$h\in\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau,\ell_{2})$,
$g^{\cdot,j}\in\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau,\ell_{2}),1\leq
j\leq m$, there exists a unique path-wise solution
$u\in\mathcal{H}^{\gamma+\alpha}_{p,{\rm loc}}(\tau)$ to (3.4). In particular,
if $\gamma+\alpha>d/p$, then the unique path-wise solution $u$ is
$C^{\gamma+\alpha-d/p}$-valued process (for $t\leq\tau$) a.s..
Proof. Step 1. First, additionally assume that (3.2) holds. Then the existence
of path-wise solution under (3.2) in $\mathcal{H}_{p}^{\gamma+\alpha}(\tau)$
(hence in $\mathcal{H}^{\gamma+\alpha}_{p,{\rm loc}}(\tau))$) follows from
Theorem 4.5. Now we show that the pathwise solution is unique in
$\mathcal{H}^{\gamma+\alpha}_{p,{\rm loc}}(\tau)$. Let
$u\in\mathcal{H}^{\gamma+\alpha}_{p,{\rm loc}}(\tau)$ be a path-wise solution.
Define
$\tau_{n}=\tau\wedge\inf\\{t:\int^{t}_{0}\|u\|^{2}_{H^{\gamma+\alpha}_{p}}ds>n\\}$.
Then $u\in\mathbb{H}^{\gamma+\alpha}_{p}(\tau_{n})$ and $\tau_{n}\uparrow\tau$
since $\int^{t}_{0}\|u\|^{2}_{H^{\gamma+\alpha}_{p}}ds<\infty$ for all
$t<\tau$, a.s. By Theorem 4.5,
$\|u\|_{\mathcal{H}^{\gamma+\alpha}_{p}(\tau_{n})}\leq
c(T,d,\alpha)\left(\|f\|_{\mathbb{H}^{\gamma}_{p}(\tau_{n})}+\|h\|_{\mathbb{H}^{\gamma+\alpha/2}_{p}(\tau_{n},\ell_{2})}+\sum_{j=1}^{m}\|g^{\cdot,j}\|_{\mathbb{H}^{\gamma+\alpha/2+\varepsilon_{1}}_{p}(\tau_{n},\ell_{2})}+\|u_{0}\|_{U^{\gamma+\alpha-\alpha/p}_{p}}\right).$
By letting $n\to\infty$ we find that
$u\in\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$, and the uniqueness of the
pathwise solution under (3.2) follows from the uniqueness result of Theorem
4.5.
Step 2. For the general case, note that for each $n>0$ and $k\leq N_{0}$,
$\widehat{c}_{k,n}:=\left(\int_{\\{z\in\mathbb{R}^{m}:|z|\leq
n\\}}|z|^{2}\nu_{k}(dz)\right)^{1/2}\vee\left(\int_{\\{z\in\mathbb{R}^{m}:|z|\leq
n\\}}|z|^{p}\nu_{k}(dz)\right)^{1/p}<\infty.$
Consider Lévy processes $(Z^{1}_{n},\cdots,Z^{N_{0}}_{n},Z^{N_{0}+1},\cdots)$
in place of $(Z^{1},Z^{2}\cdots)$, where $Z^{k}_{n}(k\leq N_{0})$ is obtained
from $Z^{k}$ by removing all the jumps that has absolute size strictly large
than $n$. Note that condition (3.2) is valid with $\widehat{c}_{k}$ replaced
by $\widehat{c}_{k,n}$. By Step 1, there is a unique path-wise solution
$v_{n}\in\mathcal{H}^{\gamma+\alpha}_{p}(\tau)$ with $Z^{k}_{n}$ in place of
$Z^{k}$ for $k=1,2,\cdots,N_{0}$. Let $T_{n}$ be the first time that one of
the Lévy processes $\\{Z^{k},1\leq k\leq N_{0}\\}$ has a jump of (absolute)
size in $(n,\infty)$. Define $u(t)=v_{n}(t)$ for $t<T_{n}\wedge\tau$. Note
that for $n<m$, by Step 1, we have $v_{n}(t)=v_{m}(t)$ for
$t<T_{n}\wedge\tau$. This is because, for $t<T_{n}\wedge\tau$, both $v_{n}$
and $v_{m}$ satisfy (4.30) with each term inside the stochastic integral
multiplied by $1_{s<T_{n}}$ (and with $Z^{k}_{n}$, $k\leq N_{0}$, in place of
$Z^{k}$). Thus $u$ is well defined. By letting $n\to\infty$, one constructs a
unique pathwise solution $u$ in
$\mathcal{H}^{\gamma+\alpha}_{p,\text{loc}}(\tau)$. The last claim follows
from Sobolev embedding theorem. The theorem is proved. $\Box$
## References
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* [2] D. Applebaum, Lévy processes and stochastic calculus. Second edition. Cambridge University Press, Cambridge, 2009.
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* [4] K. Bogdan, A. Stós and P. Sztonyk, Harnack inequality for stable processes on $d$-sets, Studia Math. 158 (2003), no. 2, 163-198.
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* [6] Z.-Q. Chen and K.H. Kim, An $L^{p}$-theory of non-divergence form SPDE driven by Lévy processes. Preprint, 2011.
* [7] T. Chang and K. Lee, On a stochastic partial differential equation with a fractional Laplacian operator, preprint.
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|
arxiv-papers
| 2011-11-21T01:50:37 |
2024-09-04T02:49:24.514708
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kyeong-Hun Kim and Panki Kim",
"submitter": "Kyeong-Hun Kim",
"url": "https://arxiv.org/abs/1111.4712"
}
|
1111.4719
|
# The stellar metallicity distribution in intermediate latitude fields with
BATC and SDSS data
Xiyan Peng1, Cuihua Du1 , Zhenyu Wu2
1College of Physical Sciences, Graduate university of the Chinese Academy of
Sciences, Beijing 100049, P. R. China
2National Astronomical Observatories, Chinese Academy of Sciences, Beijing
100012, P. R. China E-mail:ducuihua@gucas.ac.cn
(Received)
###### Abstract
Based on the Beijing-Arizona-Taiwan-Connecticut (BATC) and Sloan Digital Sky
Survey (SDSS) photometric data, we adopt SEDs fitting method to evaluate the
metallicity distribution for $\sim$40, 000 main-sequence stars in the Galaxy.
According to the derived photometric metallicities of these sample stars, we
find that the metallicity distribution shift from metal-rich to metal-poor
with the increase of distance from the Galactic center. The mean metallicity
is about of $-1.5\pm 0.2$ dex in the outer halo and $-1.3\pm 0.1$ dex in the
inner halo. The mean metallicity smoothly decreases from $-0.4$ to $-0.8$ in
interval $0<r\leqslant 5$ kpc. The fluctuation in the mean metallicity with
Galactic longitude can be found in interval $4<r\leqslant 8$ kpc. There is a
vertical abundance gradients d[Fe/H]/dz $\sim-0.21\pm 0.05$ dex kpc-1 for the
thin disk ($z\leqslant 2$ kpc). At distance $2<z\leqslant 5$ kpc, where the
thick disk stars are dominated, the gradients are about of $-0.16\pm 0.06$ dex
kpc-1, it can be interpreted as a mixture of stellar population with different
mean metallicities at all $z$ levels. The vertical metallicity gradient is
$-0.05\pm 0.04$ dex kpc-1 for the halo ($z$ $>$ 5 kpc). So there is little or
no metallicity gradient in the halo.
###### keywords:
Galaxy: structure-Galaxy: metallicity-Galaxy: formation.
††pagerange: The stellar metallicity distribution in intermediate latitude
fields with BATC and SDSS data–References††pubyear: 2002
## 1 Introduction
The structure, formation and evolution of the Galaxy are very important issues
in contemporary astrophysics. The basic components of Galaxy are the thin
disk, the thick disk, the halo and central bulge, albeit that the inter-
relationships and distinction among different components remain subject to
some debate (Gilmore 1983; 1984; Lemon et al. 2004). Recent studies, based on
accurate large-area surveys, have revealed that the Galaxy is marked by
numerous irregular substructure such as the Sagittarius dwarf tidal stream in
the halo (Ivezic et al. 2000; Yanny et al. 2000; Vivas et al. 2001; Majewski
et al. 2003) and the Monoceros stream closer to the Galactic plane (Newberg et
al. 2002; Rocha-Pinto et al. 2003). Carollo et al. (2007) shown that the the
halo is clearly divisible into two broadly overlapping structural components -
an inner and outer halo from a local kinematic analysis. It is now apparent
that our Galaxy is much more complex system than we thought before. The
formation of galaxies was long thought to be a steady process resulting in a
smooth distribution of stars (Bahcall & Soneira 1981; Gilmore et al. 1989;
Majewski 1993). But the view of the formation of the Galaxy has changed
dramatically since the discoveries of complex substructures (Newberg et al.
2002; Belokurov et al. 2007). The presence of these lumpy and complex
substructure are in qualitatively agreement with models for the formation of
the stellar halo through the accretion and merging of nearby dwarf galaxies.
Numerical simulations also suggest that this merger process plays a crucial
role in setting the structure and motions of stars within galaxies (Bullock &
Johnston 2005).
The abundance distribution is particular importance to understanding the
formation and chemical evolution of the Galaxy (Freeman & Bland-Hawthorn
2002). Researchers have long sought to constrain models for the Galactic
formation and evolution on the basis of observation of the stellar and
clusters populations that it contains. Specific models of galaxy formation
make specific predictions about the stellar abundance distribution. For
example, stars on more radial orbits are more metal-poor than stars on planar
orbits. This may indicate that the Milky Way formation began with a relatively
rapid collapse of the initial proto-galactic cloud, which means the halo stars
formed during the initial collapse, the disk stars formed after the gas had
settled into the galactic plane. But the global collapse theory was unable to
account for the lack of an abundance gradient in the Galactic halo (Searle &
Zinn 1978). The current view is that the Galactic halo formed at least partly
through the accretion of small satellite galaxies or merger of larger systems
(Freeman et al., 2002), which is well-supported by studies of stellar
kinematics and spatial distribution (Yannny et al. 2003; Juric et al. 2008).
For the thick disk, it may be one of the most significant components for
studying signatures of galaxy formation because it presents a snap frozen
relic of the state of the early disk (Freeman 2002). An intrinsic abundance
gradient in the thick disk would favor a scenarios which the thick disk was
formed either in the slow late stages of the early Galactic collapse or the
gradual kinematical diffusion of disk stars. On the contrary, a irregular
metallicity distribution or absence of gradient would favor the thick disk
having formed via the kinematical heating of thin disk or from merger debris
(Siegel et al. 2009).
The metallicity distribution of the Galaxy is best probed directly through
spectroscopic surveys (Yoss et al. 1987; Allende-Prieto et al. 2006). However,
it has the advantage of using the photometric metallicity of many more stars
out to limiting magnitude of photometric survey. Accurate determination of the
properties of the Galactic components requires surveys with sufficient sky
coverage to assess the overall geometry, sufficient depth for mapping stars to
larger distance and sufficient information to obtain reasonable distance
estimates for these stars (De Jong et al. 2010). Over the past few years,
numerous surveys have been used to investigate the existence and size of the
Galactic abundance gradient in the disk and halo. The existence of a radial
gradient in the Galaxy is now well established. An average gradient of about
$-0.06$ dex kpc-1 is observed in the Galactic disk for most of the elements
(Chen et al, 2003). De Jong et al. (2010) provided evidence for a radical
metallicity gradients in the Galactic stellar halo. However, there is
considerable disagreement about whether there is a vertical metallicity
gradient among field and/or open cluster stars of the Galaxy.
The BATC multicolor photometric survey accumulated a large data base which is
very useful for studying the Galactic structure and formation. Du et al.
(2003) provided some information on the density distribution of the main
components of the Galaxy, which can present constraints on the parameters of
models of the Galactic structure. Later, they use F and G dwarfs from the BATC
data to study the metal-abundance information (Du et al., 2004). With the new
improved observation and improved knowledge regarding galaxy formation, it
becomes possible to further discuss the metallicity gradient from different
observation direction of the Galaxy.
In this paper, we attempt to study the metallicity gradient of the Milky Way
galaxy using the 21 BATC photometric survey fields combined with the SDSS
photometric data. The outlines of this paper is as follows: The BATC
photometric system and data reduction are introduced briefly in Section 2. In
Sect. 3 we describe the theoretical model atmospheric spectra and synthetic
photometry. The metallicity distribution is discussed in Sect. 4. Finally in
Sect. 5 we summarize our main conclusions in this study.
## 2 Observations and data
### 2.1 BATC photometric system and SDSS photometric system
The BATC survey performs photometric observations with a large field multi-
colour system. There are 15 intermediate-band filters in the BATC filter
system, which covers an optical wavelength range from 3000 to 10000 Å (Fan et
al. 1996; Zhou et al. 2001). The 60/90 cm f/3 Schmidt Telescope of National
Astronomical Observatories (NAOC) was used in the BATC program, with a Ford
Aerospace 2048 $\times$ 2048 CCD camera at its main focus. The field of view
of the CCD is $58^{\prime}$ $\times$ $58^{\prime}$ with a pixel scale of
$1\arcsec{\mbox{}.}7$. The BATC magnitudes adopt the monochromatic AB
magnitudes as defined by Oke & Gunn (1983). The PIPELINE II reduction
procedure was performed on each single CCD frame to get the point spread
function (PSF) magnitude of each point source (Zhou et al. 2003). The detailed
description of the BATC photometric system and flux calibration of the
standard stars can be found in Fan et al. (1996) and Zhou et al. (2001, 2003).
In order to apply more color information to accurately estimate the
photometric stellar metallicity, we combine the BATC colors with the SDSS
colors message for the sample stars.
The SDSS used a dedicated 2.5 m telescope which has an imaging camera and a
pair of spectrographs. The imaging camera (Gunn et al. 1998) contained 30 2048
$\times$ 2048 CCDs in the focal plane of the telescope. The flux densities of
observed objects were measured almost simultaneously in five broad bands [
$u$, $g$, $r$, $i$, $z$] (Fukugita et al. 1996; Gunn et al. 1998; Hongg et al.
2001). For distinguishing explicitly between BATC and SDSS filter names, we
refer to the SDSS filters and magnitudes as $u^{\prime}$, $g^{\prime}$,
$r^{\prime}$, $i^{\prime}$, $z^{\prime}$. The photometric pipeline (Luption et
al. 2001) detected the objects, matched the data from the five filters and
measured instrumental fluxes, positions and shape parameters. The shape
parameters allowed the classification of objects as point source or extended.
The magnitudes derived from fitting a PSF are currently accurate to about 2
percents in $g$, $r$ and $i$, and 3-5 percents in $u$ and $z$ for bright
($\leq$ 20 mag) point source. In Table 1, we list the parameters of the BATC
and SDSS filters. Col. (1) and Col. (2) represent the ID of the BATC and SDSS
filters, Col. (3) and Col. (4) the central wavelengths and FWHM of the 20
filters, respectively. The reddening extinction for each star are determined
from the SDSS catalog.
Table 1: Parameters of the BATC and SDSS filters No. | Filter | Wavelength | FWHM
---|---|---|---
| | (Å) | (Å)
1 | $a$ | 3371.5 | 359
2 | $b$ | 3906.9 | 291
3 | $c$ | 4193.5 | 309
4 | $d$ | 4540.0 | 332
5 | $e$ | 4925.0 | 374
6 | $f$ | 5266.8 | 344
7 | $g$ | 5789.9 | 289
8 | $h$ | 6073.9 | 308
9 | $i$ | 6655.9 | 491
10 | $j$ | 7057.4 | 238
11 | $k$ | 7546.3 | 192
12 | $m$ | 8023.2 | 255
13 | $n$ | 8484.3 | 167
14 | $o$ | 9182.2 | 247
15 | $p$ | 9738.5 | 275
1 | $u^{\prime}$ | 3543 | 569
2 | $g^{\prime}$ | 4770 | 1387
3 | $r^{\prime}$ | 6231 | 1373
4 | $i^{\prime}$ | 7625 | 1526
5 | $z^{\prime}$ | 9134 | 9500
### 2.2 Direction and data reduction
Accurate determination of the properties of components of the Milky Way
requires surveys with sufficient sky coverage to assess the overall geometry
(De Jong et al. 2010). Most previous investigations about the Galactic
metallicity distribution use only one or a few selected lines-of-sight
directions (Du et al. 2004; Siegel et al. 2009; Karatas et al. 2009). In this
paper, the BATC photometry survey presented 21 intermediate-latitude fields in
the multiple directions. These fields used in this paper are towards the
Galactic center, the anticentre, the antirotation direction at median and high
latitudes, $|b|>35^{\circ}$. Since metallicity distribution at high Galactic
latitudes are not strongly related to the radial distribution, they are well
suited to study the vertical distribution of the Galaxy. Table 2 lists the
locations of the observed fields and their general characteristics. In Table
2, column 1 represents the BATC field name, columns 2 and 3 represent the
Galactic longitude and latitude, columns 4 and 5 represent the limit magnitude
and the number of sample star used in this study. As shown in the Table 2, the
most photometric depth of our data is 21.0 mag in $i$ band. In total, there
are about 40, 000 sample stars in our study.
Table 2: Relative information for the BATC observation fields Observed field | $l$ (deg) | $b$ (deg) | $i$ (Comp) | star number
---|---|---|---|---
T485 | 175.7 | 37.8 | 21.0 | 1550
T518 | 238.9 | 39.8 | 19.5 | 1584
T288 | 189.0 | 37.5 | 20.0 | 2115
T477 | 175.7 | 39.2 | 20.0 | 2001
T328 | 160.3 | 41.9 | 19.5 | 1666
T349 | 224.1 | 35.3 | 20.5 | 2436
TA26 | 191.1 | 44.4 | 20.0 | 1285
T291 | 167.8 | 46.4 | 20.0 | 1670
T362 | 245.7 | 53.4 | 20.0 | 1237
T330 | 147.2 | 68.3 | 20.5 | 1020
U085 | 121.6 | 60.2 | 21.0 | 1679
T521 | 56.1 | -36.8 | 20.5 | 4430
T491 | 62.9 | -44.0 | 20.0 | 2824
T359 | 79.7 | -37.8 | 20.5 | 2932
T350 | 251.3 | 67.3 | 19.5 | 1353
T534 | 91.6 | 51.1 | 21.0 | 2044
T193 | 59.8 | -39.7 | 20.0 | 2830
T516 | 125.0 | -62.0 | 20.0 | 1113
T329 | 169.9 | 50.4 | 21.0 | 1704
TA01 | 135.7 | -62.1 | 20.5 | 1342
T517 | 188.6 | -38.2 | 20.0 | 1296
Figure 1: Two-color diagrams for the BATC T329 field stars. The panel (a) is
the distribution of sample stars in the ($j-m$)/($g-j$) two-colour diagram,
and panel (b) is the ($j-m$)/($g-j$) two-colour diagram after removing those
objects that lay significantly off the stellar locus. The line denotes the
stellar locus from equation (1).
Here, because our fields in this work have also been observed by the Sloan
Digital Space Survey (SDSS-DR7) and each object type (stars-galaxies-QSO) has
been given. Thus, we can obtain a relative reliable star catalogue. In these
star sample, it is complete to 20.5 mag with an error of less than 0.1 mag in
the BATC $i$ band. Owing to the BATC observing strategy, data of stars
brighter than $m$ = 14 mag, are saturated, and star counts are not completed
for visual magnitudes fainter than $m$ = 21.0 mag. So our work is restricted
to the magnitudes range 15 $\leqslant$ $i$ $<$ 21\. From two-color diagrams
for all objects in any field as an example, we can find that most stars are
plotted on a diagonal band. In Fig. 1 show the BATC ($g-j$) versus ($j-m$)
two-color diagram for T329 field stars. The line which can be described by
equation (1) donates the stellar locus of the sample in the T329 field .
$\displaystyle(j-m)=-0.10245+0.81172(g-j)$ (1)
The sample shows a main-sequence (MS) stellar locus but has significant
contamination from giants, manifested in the broader distribution overlying
the narrow stellar locus. In order to determine the MS star sample, we use
multi-color selection criteria outlined in Karaali et al. (2003) and Juric et
al. (2008), which remove objects based on their location relative to the
dominant stellar locus. For example, Juric et al. (2008) applied an extra
procedure which consists of rejecting objects at distance larger than 0.3 mag
from the stellar locus in order to remove hot dwarfs, low-red shift quasars,
and white /red dwarf unresolved binaries from their sample. The procedure does
well for high latitude field data from SDSS (Karaali et al. 2003; Yaz et al.
2010). Fig. 1b gives the cleaning sample stars after rejecting those objects
which is lay significantly off the stellar locus.
## 3 THEORETICAL MODEL AND CALIBRATION FOR METALLICITY
### 3.1 theoretical stellar library and synthetic photometric
A homogeneous and complete stellar library can match any ambitious goals
imposed on a standard library. Lejeune et al. (1997) presented a hybrid
library of synthetic stellar spectra. The library covers a wide range of
stellar parameters $T_{\rm eff}$= 50, 000 K to 2, 000 K in intervals of 250 K,
log $g$= $-1.02$ to 5.50 in main increments of 0.5, and [M/H]= $-5.0$ to +1.0.
For each model in the library, a flux spectrum is given for the same set of
1221 wavelength points covering the range 9.1 to 160, 000 nm, with a mean
resolution of 20Å in the visible. The spectra are thus in a format which has
proved to be adequate for synthetic photometry of wide- and intermediate-band
systems (Du et al. 2004).
On the basis of the theoretical library, we calculate synthetic colors of the
BATC and SDSS photometric system. Here, we synthesize colors for simulated
stellar spectra with $T_{\rm eff}$ and log $g$ characteristic of main
sequences (log $g$ = 3.5, 4.0, 4.5 for dwarfs) and 19 values of metallicity
([M/H]= $-5.0,-4.5,-4.0,-3.5$,
$-3.0,-2.5,-2.0,-1.5,-1.0,-0.5$,$-0.3,-0.2,-0.1$, 0.0, +0.1, +0.2, +0.3, +0.5
and +1.0), where [M/H] denotes metallicity relative to hydrogen. The synthetic
$i$th filter magnitude can be calculated with equation (2).
$m=-2.5~{}{\rm log}\frac{\int{F_{\lambda}\phi_{i}({\lambda}){\rm
d}\lambda}}{\int{\phi_{i}({\lambda}){\rm d}\lambda}}-48.60,$ (2)
Where ${F_{\lambda}}$ is the flux per unit wavelength, $\phi_{i}$ is the
transmission curve of the $i$th filter of the BATC or SDSS filter system (Du
et al. 2004) .
The bluer colors are sensitive to metallicity down to the lowest observed
metallicities because most of the line-blanketing from heavy elements occurs
in the shorter wavelength regions. In contrast, the redder colors are
primarily sensitive to temperature index. The BATC $a$, $b$ bands contain the
Balmier jump, a stellar spectral feature which is sensitive to surface
gravity. Since our sample includes only main sequences, it conveys little
gravity information. It should be mentioned that, although the metallicity or
temperature derived from synthetic photometry is not very accurate for a
single star, perhaps which can be distorted by a poor point, it is meaningful
for the statistic analysis of sample stars.
### 3.2 METALLICITY AND PHOTOMETRIC PARALLAX
The most accurate measurements of stellar metallicity are based on
spectroscopic observation. Despite the recent progress in the availability of
stellar spectra (e.g. SDSS-III and RAVE), the stellar number detected in
photometric surveys is much more than spectroscopic observation. So
photometric methods have also often been used to give the stellar metallicity.
For example, Sandage (1969) detailed a technique using UBV photometry indices
to measure approximate abundance. Karaali (2003) evaluated the metal abundance
by ultraviolet-excess photometric parameter using CCD UBVI data. Karaali
(2005) extended this method to the SDSS photometry. Ivezic et al. (2008)
obtained the mean metallicity of stars as a function of $u-g$ and $g-r$ colors
of SDSS data. For the BATC multicolor photometric system, there are 15
intermediate-band filters covering an optical wavelength range from 3000 to
10000 Å. There are 5 filters for the SDSS photometric system. So the SEDs of
20 filters for every object are equivalent to a low resolution spectrum.
The sample SEDs simulation with template SEDs can be used to derive the
parameter of sample stars (Du et al. 2004). The standard minimization,
computing and minimizing the deviations between the photometric SEDs of the
star and the templates SEDs obtained with the same photometric system, is used
in the fitting process. The minimum indicated the best fit to the observed SED
by the set of spectra (Du et al. 2004).
${\chi^{2}=\sum\limits_{l=1}^{N_{filt}=20}\left[\frac{m_{obs,l}-m_{temp,l}}{\sigma_{l}}\right]^{2}},$
(3)
where ${m_{obs,l}}$, $m_{temp,l}$ and $\sigma_{l}$ are the observed magnitude,
template magnitude and their uncertainty in filter $l$, respectively, and
$N_{filt}$ is the total number of filters in the photometry.
According to the results of SEDs fitting, the metallicity and temperature of
about 40, 000 sample stars are obtained in 21 fields. In addition, we extract
the spectroscopic metallicities from sdss DR7 database for our studied 21
fields, and there are about 870 stars for which they also have photometric
metallicities from our method. Using these stars, we present a calibration of
our SED fitting method. After applying calibration, it is reliable for the
derived photometric metallicities from SEDs fitting method. In Figure 2., we
present the difference of photometric and spectroscopic metallicity as a
function of $(g-r)$ color. The uncertainties of metallicity obtained from
comparing SEDs between photometry and theoretical models are due to the
observational error and the finite grid of the models. For the metal-poor
stars ([Fe/H]$<-1.0$), the metallicity uncertainty is about 0.5 dex, and 0.2
dex for the stars [Fe/H]$>-0.5$ (Du et al. 2004). The metallicity distribution
diagram for all sample stars was given in Fig. 3. One local maximum appear at
[Fe/H] from -0.5 to 0 dex, and a tail down to -3.0 dex.
Figure 2: The difference of photometric and spectroscopic metallicity for 870
stars as a function of $(g-r)$ color is shown. Figure 3: The metallicity
distribution for all the sample stars selected in our study is shown.
The stellar type can be derived from effective temperature of dwarfs, then the
stellar distance relative to the sun can be obtained by equation (4).
$\displaystyle m_{v}-M_{v}=5lgr-5+A_{v}$ (4)
where $m_{v}$ is the visual magnitude, absolute magnitude $M_{v}$ can be
obtained according to the stellar type. The reddening extinction $A_{v}$ is
small for most fields. We adopted the absolute magnitude versus stellar type
relation for main-sequence stars from Lang (1992). $r$ is the stellar
distances. The vertical distance of the star to the galactic plane can be
evaluated by equation (5):
$\displaystyle z=rsinb$ (5)
A variety of errors affect the determination of stellar distances. The first
source of errors is from photometric uncertainty less than 0.1 mag in the BATC
$i$ band; the second from the misclassification, which should be small due to
the multicolor photometry. For luminosity class V, types F/G, the absolute
magnitude uncertainty is about 0.3 mag. In addition, there may exist an error
from the contamination of binary stars in our sample. We neglect the effect of
binary contamination on distance derivation due to the unknown but small
influence from mass dsitribution in binary components (Kroupa et al. 1993;
Ojha et al. 1996).
Figure 4: The metallicity distribution for the T291 field in different
distance range is shown. In the short distance, r $<$ 1kpc, most stars are in
the range [0, $-0.5$], in the larger distance, 15 $<$ r $<$ 20 kpc, most stars
are poorer than $-1$.
## 4 METALLICITY DISTRIBUTION
It is well known that the chemical abundance of stellar population contains
much information about the population,s early evolution. The stellar
metallicity distribution in the Galaxy has been the subject of photometric and
spectroscopic surveys (Gilmore et al. 1985; Ratnatunga et al. 1989; Friel
1988). In this study, we want to explore possible stellar metallicity
distribution variation with the observation direction. A method of SED
combination for the SDSS and BATC photometries has been adopted to give the
steller metallicity distribution. At first, the metallicity for the sample
stars can be derived by comparing SEDs between photometry and the theoretical
models. The SED fitting method is described in Section 3\. 2. The mean
metallicity distribution for each field is determined in the following
distance intervals (in kpc): $15<r\leqslant 20$, $10<r\leqslant 15$,
$8<r\leqslant 10$, $6<r\leqslant 8$, $5<r\leqslant 6$, $4<r\leqslant 5$,
$3<r\leqslant 4$, $2.5<r\leqslant 3$, $2<r\leqslant 2.5$, $1.5<r\leqslant 2$,
$1<r\leqslant 1.5$, $0.25<r\leqslant 1$. As an example the metallicity
distributions as a function of vertical distance for the field T291 is
presented in Fig. 4. From the figures (Fig. 4), it is clear that there is a
number-shift from metal-rich stars to metal-poor ones with the increasing of
distance. In star counts the younger metal rich stars are confined to regions
close to the Galactic mid-plane, while the older, metal-poorer stars with a
larger scale height dominated at larger vertical distances from the Galactic
plane.
### 4.1 Metallicity variation with Galactic longitude
Mean metallicity distribution as a function of Galactic longitude for
different distance intervals are presented in Fig. 5. The mean metallicity
shift from metal-rich to metal-poor with the increase of distance from the
Galactic center can be found in Fig. 5. The solid points represent the south
galactic latitude fields, The open square points represent the north galactic
latitude fields. As shown in Fig. 5, the mean metallicity in interval
$10<r\leqslant 20$ kpc is around of $-1.5$ dex. In intervals $8<r\leqslant 10$
kpc our result indicates that the mean metallicity is $\sim$ $-1.3$. The mean
metallicity in intervals $5<r\leqslant 8$ kpc is about $-1.0$ and the mean
metallicity smoothly decreases from $-0.6$ to $-0.8$ in intervals
$2.5<r\leqslant 5$ kpc. Our results are consistent with the results of Siegel
et al. (2009) and Karatas et al. (2009). Siegel et al. find a monometallic
thick disk and halo with abundances of [Fe/H] = $-0.8$ and $-1.4$
respectively. Karatas et al. derive mean abundance values of [Fe/H]= $-0.77\pm
0.36$ dex for the thick disk, and [Fe/H] = $-1.42\pm 0.98$. The mean
metallicity decreases from $-0.4$ to $-0.6$ in intervals $0<r\leqslant 2.5$
kpc. The mean metallicity in interval $0.25<r\leqslant 1$ kpc is [Fe/H] $\sim$
$-0.3$, which is consistent with the result of Yaz et al. (2010).
As shown in Fig. 5, at larger distance, r $>$ 10 kpc, compared to the typical
error bars, the mean metallicity distributions variation with Galactic
longitude is almost flat. For $4<r\leqslant 8$ kpc, there is a fluctuation in
the mean metallicity with Galactic longitude. The overall distribution of mean
metallicity has a maxminum at $l$ $\sim$ 200∘. For $2<r\leqslant 8$ kpc, the
T517 filed (Galactic coordinates : $l=188.6^{\circ},b=-38.2^{\circ}$;
Equatorial coordinates : $\alpha=58.59^{\circ}$, $\delta=-0.35^{\circ}$) and
the TA26 field (Galactic coordinates : $l=191.1^{\circ},b=44.4^{\circ}$;
Equatorial coordinates : $\alpha=139.956^{\circ}$, $\delta=33.745^{\circ}$)
show metal rich character related to other fields. This feature may reflect a
fluctuation from streams (such as Monocers stream) which are accreted from
nearby galaxies. Juric et al.(2008) detect two overdensities in the thick disk
region. Klement et al.(2009) also find individual stream from the SSPP in the
direction with central coordinates (Equatorial coordinates): $\alpha$ =
$58.58^{\circ}$ and $\delta$ = $-4.99^{\circ}$. Maybe the deviant behaviors of
the two fields result from systematical error in the observation. In the work
of AK et al. (2007), they find that the metallicity distributions for both
(relatively) short and large vertical distances show systematic fluctuations.
The scaleheight of thick disk varies with the observed direction were found in
the works of Du et al. (2006) and Bilir et al. (2008).
Figure 5: The mean metallicity distribution as a function of galactic
longitude in different distance intervals are shown.
### 4.2 THE VERTICAL METALLICITY GRADIENT
Detailed information about the vertical metallicity gradient can provide
important clue about the formation scenario of stellar population. Here, we
used the mean metallicity to described the metallicity distribution function.
As an example, The distribution trend of mean metallicity [Fe/H] with height
above the galactic plane [$z$] for the T291 field is shown in Fig. 6. The
metallicity gradients for all the fields in different $z$ intervals $z<2$ kpc,
$2<z\leqslant 5$ kpc and $5<z\leqslant 15$ kpc are given in Fig. 7 and
detailed in Table 3. In Table 3, Column 1 represents the BATC field name,
Columns 2 - 7 represent gradient and error of gradient in different $z$
distance: $z\leqslant 2$ kpc, $2<z\leqslant 5$ kpc and $5<z\leqslant 15$ kpc,
respectively.
From the Fig. 7 we can find that the variation of the gradient for the halo
with galactic longitude is flat and the mean gradient of halo is about
$-0.05\pm 0.04$ dex kpc-1 ($5<z\leqslant 15$ kpc), which is essentially in
agreement with the conclusion of Yaz et al. (2010) and Du et al. (2004). Du et
al.(2004) find the small or zero gradient d[Fe/H]/dz = $-0.06\pm 0.09$ in the
halo. Yaz et al. (2010) find $d[M/H]/dz=-0.01$ dex kpc-1 for the inner
spheroid. The result of Karaali et al. (2003) is slightly steeper than the
value of our result. Karaali et al. (2003) find that there is a metallicity
gradient d[Fe/H]/dz $\sim-0.1$ dex kpc-1 in the inner halo ($5<z\leqslant 8$
kpc ) and zero in the outer part ($8<z\leqslant 10$ kpc). From Fig. 6 we can
find that the incompleteness of the star sample causes significant statistical
uncertainties at large distance. Probably, there is little or no metallicity
gradient in the halo. It is consistent with the merger or accretion origin of
the outer halo.
As shown in Fig. 7, at distance $0<z<2$ kpc, the mean vertical abundance
gradient is about d[Fe/H]/dz $\sim-0.21\pm 0.05$ dex kpc-1. The value for the
vertical metallicity at distance $0<z<2$ kpc is in agreement with the
canonical metallicity gradients with the same $z$ distances. For example, Yaz
et al. (2009) find the metallicity gradient is d[Fe/H]/dz $\sim-0.3$ dex kpc-1
for short distance. The metallicity gradient is found to be d[Fe/H]/dz
$\sim-0.37$ dex kpc-1 for $z$ $<$ 4 kpc in the work of Du et al. (2004). The
result of Karaali et al. (2003) can be described as d[Fe/H]/dz $\sim-0.2$ dex
kpc-1 for the thin and thick disk.
At distance $2<z\leqslant 5$ kpc, where the thick disk stars dominated, the
gradient is about $-0.16\pm 0.06$ dex kpc-1 in our work which is consistent
with the work of karaali et al. (2003) and less than the value of Du et al.
(2004). Du et al. (2004) point out that the metallicity gradient is d[Fe/H]/dz
$\sim$ $-0.37$ dex kpc-1. In our study, the thick disk gradient is interpreted
as different contribution from three components of the Galaxy at different $z$
distance. The existence of a clear vertical metallicity of the thick disk
would be an important clue about the origin of the thick disk. However, it is
an open question for the formation of the thick disk component. A number of
models have been put forward since the confirmation if its existence. Chen et
al. (2001) support that the thick disk formed through the heating of a
preexisting thin disk, with the heating mechanism being the merging of a
satellite galaxy. Here, we also favor the thick disk having formed via the
kinematical heating of thin disk and from merger debris. Thus, there is a
irregular metallicity distribution or absence of intrinsic gradient.
Figure 6: The mean metallicity as a function of vertical distance $z$ for the T291 field. The metallicity gradients of the thin disk, thick disk and halo are $-0.23\pm 0.03$, $-0.18\pm 0.07$, $-0.05\pm 0.01$ dex kpc-1, respectively. Figure 7: The metallicity gradients distribution for all fields in this study are shown for the intervals $z$ $<$ 2 kpc , 2 kpc $<$ $z$ $<$ 5kpc and 5 kpc $<$ $z$ $<$ 15kpc. Table 3: The gradient distribution in different distance interval for the selected fields | 0-2 (kpc) | | 2-5 (kpc) | | 5-15 (kpc) |
---|---|---|---|---|---|---
Observed field | gradient | error | gradient | error | gradient | error
T193 | -0.136 | 0.017 | -0.080 | 0.020 | -0.020 | 0.022
T288 | -0.258 | 0.027 | -0.267 | 0.032 | -0.007 | 0.037
T291 | -0.228 | 0.028 | -0.185 | 0.071 | -0.045 | 0.014
T328 | -0.110 | 0.027 | -0.126 | 0.032 | -0.137 | 0.043
T329 | -0.170 | 0.028 | -0.100 | 0.042 | -0.037 | 0.015
T330 | -0.171 | 0.042 | -0.217 | 0.055 | -0.080 | 0.010
T349 | -0.180 | 0.017 | -0.225 | 0.050 | -0.029 | 0.031
T350 | -0.267 | 0.033 | -0.178 | 0.074 | -0.069 | 0.035
T359 | -0.203 | 0.017 | -0.126 | 0.023 | -0.054 | 0.023
T362 | -0.237 | 0.037 | -0.126 | 0.060 | -0.037 | 0.022
T477 | -0.181 | 0.025 | -0.275 | 0.030 | -0.053 | 0.022
T485 | -0.247 | 0.033 | -0.119 | 0.044 | 0.028 | 0.040
T491 | -0.194 | 0.017 | -0.179 | 0.020 | -0.070 | 0.027
T516 | -0.291 | 0.039 | -0.204 | 0.061 | -0.005 | 0.026
T517 | 0.054 | 0.041 | -0.173 | 0.059 | -0.101 | 0.067
T518 | -0.163 | 0.035 | -0.026 | 0.045 | -0.051 | 0.052
T521 | -0.197 | 0.012 | -0.175 | 0.013 | -0.045 | 0.011
T534 | -0.254 | 0.028 | -0.110 | 0.035 | -0.035 | 0.012
TA01 | -0.199 | 0.037 | -0.192 | 0.046 | -0.047 | 0.015
TA26 | 0.307 | 0.030 | -0.187 | 0.038 | -0.085 | 0.047
U085 | -0.146 | 0.024 | -0.148 | 0.037 | -0.035 | 0.013
## 5 CONCLUSIONS AND SUMMARY
In this work, based on the BATC and SDSS photometric data, we evaluated the
stellar metallicity distribution for 40, 000 main-sequence stars in the Galaxy
by adopting SEDs fitting method. These selected fields are towards the
Galactic center, the anticentre, the antirotation direction at median and high
latitudes . The metallicity distribution could be obtained up to distances
$r=20$ kpc, which covers the thin disk, thick disk and halo. We determined the
mean stellar metallicity as a function of vertical distance in different
direction. It can be clearly seen that the metallicity distribution shift from
metal-rich to metal-poor with the increase of distance from the Galactic
center. The mean metallicity is about $-1.5\pm 0.2$ dex in intervals
$10<r\leqslant 20$ kpc and $-1.3\pm 0.1$ dex in interval $8<r\leqslant 10$
kpc. The mean metallicity smoothly decreases from $-0.6$ to $-0.8$ in interval
$2.5<r\leqslant 5$ kpc, while the mean metallicity decreases from $-0.4$ to
$-0.6$ in interval $0<r\leqslant 2.5$ kpc. In addition, a fluctuation in the
mean metallicity with Galactic longitude can be found and the overall
distribution has a maximum at about $l$ $\sim$ 200∘ in interval $4<r\leqslant
8$ kpc. Maybe this feature can be related with the substructure or streams
(such as Monoceros stream) which are accreted from nearby galaxies. At the
same time, we find the vertical abundance gradients for the thin disk ($0<z<2$
kpc) is d[Fe/H]/dz $\sim-0.21\pm 0.05$ dex kpc-1, and a vertical gradient
$-0.16\pm 0.06$ dex kpc-1 at distances $2<z\leqslant 5$ kpc where the thick
disk stars are dominated. Here, we consider the thick disk gradient may be the
result from the different contributions from three components of the Galaxy at
different $z$ distance. The vertical gradient d[Fe/H]/dz $\sim-0.05\pm 0.04$
dex kpc-1 is found in distance $5<z\leqslant 15$ kpc. So, there is little or
no gradient in the halo. These results are in agreement with the values in the
literature (Yoss et al. 1987; Trefzger et al. 1995; Karaali et al. 2003; AK et
al. 2007; Yaz et al. 2010). It is possible that additional observational
investigations (some projects aimed at spectroscopic sky surveys such as
SEGUE, LAMOST, GAIA) will give more evidence for the metallicity gradient of
the Galaxy and therefore provide a powerful clue to the disk and halo
formation.
## Acknowledgments
We especially thank the anonymous referee for numerous helpful comments and
suggestions which have significantly improved this manuscript. This work was
supported by the joint fund of Astronomy of the National Nature Science
Foundation of China and the Chinese Academy of Science, under Grants 10778720
and 10803007. This work was also supported by the GUCAS president fund.
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|
arxiv-papers
| 2011-11-21T02:56:43 |
2024-09-04T02:49:24.527261
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Xiyan Peng, Cuihua Du, Zhenyu Wu",
"submitter": "Cuihua Du",
"url": "https://arxiv.org/abs/1111.4719"
}
|
1111.4723
|
Self-dual interval orders and row-Fishburn matrices
Sherry H. F. Yan, Yuexiao Xu
Department of Mathematics, Zhejiang Normal University, Jinhua 321004, P.R.
China
huifangyan@hotmail.com
Abstract. Recently, Jelínek derived that the number of self-dual interval
orders of reduced size $n$ is twice the number of row-Fishburn matrices of
size $n$ by using generating functions. In this paper, we present a bijective
proof of this relation by establishing a bijection between two variations of
upper-triangular matrices of nonnegative integers. Using the bijection, we
provide a combinatorial proof of the refined relations between self-dual
Fishburn matrices and row-Fishburn matrices in answer to a problem proposed by
Jelínek.
Key words: self-dual interval order, self-dual Fishburn matrix, row-Fishburn
matrix.
AMS Mathematical Subject Classifications: 05A05, 05C30.
## 1 Introduction
A poset is said to be an interval order ( also known as $(2+2)$-free poset) if
it does not contain an induced subposet that is isomorphic to $2+2$, the union
of two disjoint $2$-element chains. Let $P$ be a poset with a strict order
relation $\prec$. A strict down-set of an element $x\in P$ is the set $D(x)$
of all the elements of P that are smaller than $y$, i.e., $D(y)=\\{y\in
P:y\prec x\\}$. Similarly, the strict up-set of $x$, denoted by $U(x)$, is the
set $\\{y\in P:y\succ x\\}$. A poset $P$ is $(2+2)$-free if and only if its
sets of strict down-sets, $D(P)=\\{D(x):x\in P\\}$ can be written as
$D(P)=\\{D_{1},D_{2},\ldots,D_{m}\\}$
where $\emptyset=D_{1}\subset D_{2}\subset\ldots\subset D_{m}$, see [1, 2]. In
such context, we say that $x\in P$ has level $i$ if $D(x)=D_{i}$. An element
$x$ is said to be a minimal element if $x$ has level $1$. Following Fishburn
[7], we call the number $m$ of distinct strict down-sets the magnitude of P.
It turns out that $m$ is also equal to the number of distinct strict up-sets,
and we can order the strict up-sets of $P$ into a decreasing chain
$D(P)=\\{U_{1},U_{2},\ldots,U_{m}\\}$
where $U_{1}\supset U_{2}\supset\ldots\supset U_{m}=\emptyset$, see [7, 8]. We
say that $x$ has up-level $i$ if $U(x)=U_{i}$. An element $x$ is said to be a
maximal element if $x$ has up-level $m$.
The dual of a poset $P$ is the poset $\overline{P}$ with the same elements as
$P$ and an order relation $\overline{\prec}$ defined by $x\overline{\prec}y$
$\Longleftrightarrow$ $y\prec x$. A poset is self-dual if it is isomorphic to
its dual.
Fishburn [7, 9] did pioneering work on interval orders; for instance, he
showed the basic theorem that a poset is an interval order if and only if it
is $(2+2)$-free and established a bijection between interval orders and a
certain kind of integer matrices, called Fishburn matrices. Recently,
Bousquet-Mélou et al. [2] constructed bijections between interval orders and
ascent sequences, between ascent sequences and permutations avoiding a certain
pattern, between interval orders and regular linearized chord diagrams by
Stoimenow [12]. Several other papers have focused on bijections between
interval orders and other objects. For instance, Dukes and Parviainen [4] have
described a direct bijection between Fishburn matrices and ascent sequences,
while the papers of Claesson et al. [3] and Dukes et al. [6] extend the
bijection between interval orders and Fishburn matrices to more general
combinatorial structures.
A Fishburn matrix of size $n$ is an upper-triangular matrix with nonnegative
integers which sum to $n$ and each row and each column contains a nonzero
entry. Throughout this paper that each matrix has its rows numbered from top
to bottom, and columns numbered left-to-right, starting with row and column
number one. We let $M_{i,j}$ denote the entry of M in row $i$ and column $j$.
The size of a matrix $M$ is the sum of all its entries. Moreover, the
dimension of an upper triangular matrix is defined to the number of rows.
The dual matrix of $M$, denoted by $\overline{M}$, is obtained from $M$ by
transposition along the diagonal running from bottom-left to top-right. More
precisely, for $1\leq i,j\leq m$, we have $\overline{M}_{i,j}=M_{m+1-j,m+1-i}$
where $m$ is the dimension of $M$. If a matrix M is equal to $\overline{M}$,
we call it self-dual.
Fishburn [7, 9] showed that an interval order $P$ of magnitude m corresponds
to an $m\times m$ Fishburn matrix $M$ with $M_{i,j}$ being equal to the number
of elements of $P$ that have level $i$ and up-level $j$. Jelínek [10] showed
that the Fishburn’s bijection turns out to be a bijection between self-dual
interval orders of size $n$ and self-dual Fishburn matrices of size $n$.
Following the terminologies given in [10], we distinguish three types of cells
in a Fishburn matrix $M$ of dimension $k$ : a cell $(i,j)$ is a diagonal cell
if $i+j=k+1$, i.e., $(i,j)$ belongs to the north-east diagonal of the matrix.
If $i+j<k+1$ (i.e., $(i,j)$ is above and to the left of the diagonal) then
$(i,j)$ is a North-West cell, or NW-cell, while if $i+j>k+1$, then $(i,j)$ is
an SE-cell. Clearly, NW-cells and diagonal cells together determine a self-
dual Fishburn matrix. The reduced size of a self-dual fishburn matrix $M$ is
the sum of all diagonal cells and NW-cells. The reduced size of a self-dual
interval order $P$ is the reduced size of its corresponding self-dual Fishburn
matrix under Fishburn’s bijection.
A row-Fishburn matrix of size $n$ is defined to be an upper-triangular matrix
with nonnegative integers which sum to $n$ and each row contains a nonzero
entry. In a matrix $A$, the sum of a column (resp. row) is defined to the sum
of all the entries in this column (resp. row). A column or a row is said to be
zero if it contains no nonzero entries. The set of self-dual Fishburn matrices
of reduced size $n$ is denoted by $\mathcal{M}(n)$. Denote by
$\mathcal{M}(n,k,p)$ be the set of self-dual Fishburn matrices of reduced size
$n$ whose first row has sum $k$ and diagonal cells have sum $p$. Let
$\mathcal{RM}(n)$ be the set of row-Fishburn matrices of size $n$. The set of
row-Fishburn matrices in $\mathcal{RM}(n)$ whose last column has sum $k$ is
denote by $\mathcal{RM}(n,k)$. Denote by $\mathcal{RM}(n,k,p)$ be the set of
row-Fishburn matrices in $\mathcal{RM}(n,k)$ whose first row has sum $p$.
Moreover, the set of self-dual interval orders of reduced size $n$ is denoted
by $\mathcal{I}(n)$.
Based on the bijection between interval orders and Fishburn matrices, Jelínek
[10] presented a new method to derive formulas for the generating functions of
interval orders, counted with respect to their size, magnitude, and number of
minimal and maximal elements, which generalize previous results on refined
enumeration of interval orders obtained by Bousquet-Mélou et al. [2], Kitaev
and Remmel [11], and Dukes et al. [5]. Applying the new method, Jelínek [10]
obtained formulas for the generating functions of self-dual interval orders
with respect to analogous statistics. From the obtained generating functions,
relations between self-dual Fishburn matrices and row-Fishburn matrices were
derived, that is,
$|\mathcal{M}(n,k,0)|=|\mathcal{RM}(n,k)|,$ (1.1)
and for $p\geq 1$
$|\mathcal{M}(n,k,p)|=|\mathcal{RM}(n,k,p)|.$ (1.2)
Combining the bijection between self-dual interval orders and self-dual
Fishburn martices, formulas (1.1) and (1.2), Jelínek derived that for $n\geq
1$,
$|\mathcal{I}(n)|=|\mathcal{M}(n)|=2|\mathcal{RM}(n)|,$ (1.3)
and asked for bijective proofs of (1.1) and (1.2). The main objective of this
paper is to present bijective proofs of these formulas by establishing a one-
to-one correspondence between two variations of upper-triangular matrices of
nonnegative integers.
Let $\mathcal{M}(n,k)$ be the set of self-dual Fishburn matrices of reduced
size $n$ whose first row has sum $k$. Denote by $\mathcal{EM}(n,k)$ (resp.
$\mathcal{OM}(n,k)$ ) be the set of self-dual Fishburn matrices in
$\mathcal{M}(n,k)$ whose dimension are even (resp. odd). Using the bijection
between two variations of upper-triangular matrices of nonnegative integers,
we derive that
$|\mathcal{EM}(n,k)|=|\mathcal{OM}(n,k)|=|\mathcal{RM}(n,k)|.$ (1.4)
## 2 The bijective proofs
Recall that a self-dual Fishburn matrix is determined by its NW-cells and
diagonal cells. Given a self-dual Fishburn matrix $M$, the reduced matrix of
$M$, denoted by $R(M)$, is a matrix obtained from $M$ by filling all the SE-
cells with zeros. An upper-triangular matrix is said to a super triangular
matrix if all its SE-cells are zero.
###### Lemma 2.1
Let $M^{\prime}$ be a super triangular matrix of dimension $m$. Then
$M^{\prime}$ is a reduced matrix of a self-dual Fishburn matrix if and only if
it satisfies the following two conditions:
* (i)
for $1\leq i\leq\lceil{m\over 2}\rceil$, each column $i$ contains a nonzero
entry;
* (ii)
for $1\leq i\leq\lceil{m\over 2}\rceil$, either row $i$ or column $m+1-i$
contains a nonzero entry.
Proof. Let $M$ be a self-dual Fishburn matrix with $R(M)=M^{\prime}$. Clearly,
$M^{\prime}$ is a super triangular matrix. Since the first $\lceil{m\over
2}\rceil$ columns of $M^{\prime}$ are the same as those in $M$, the condition
$(i)$ follows immediately. It remains to show that $M^{\prime}$ satisfies
condition $(ii)$. Since $M$ is self-dual Fishburn matrix, for all $1\leq i\leq
m$, row $i$ must contains a nonzero entry, that is,
$\sum_{j=1}^{m}M_{i,j}=\sum_{j=1}^{m-i}M_{i,j}+\sum_{j=m+1-i}^{m}M_{i,j}=\sum_{j=1}^{m-i}M_{i,j}+\sum_{j=1}^{i}M_{j,m+1-i}>0.$
Hence, for $1\leq i\leq\lceil{m\over 2}\rceil$, either row $i$ or column
$m+1-i$ of $R(M)$ contains a nonzero entry. Therefore, the condition $(ii)$
holds for $R(M)$.
Conversely, given a super triangular matrix $M^{\prime}$ satisfying conditions
$(i)$ and $(ii)$, We can recover a self-dual matrix $M$ from $M^{\prime}$ by
filling the SE-cell $(m+1-j,m+1-i)$ with $M^{\prime}_{i,j}$. If $1\leq
i\leq\lceil{m\over 2}\rceil$, the sum of row $i$ of $M$ is given by
$\sum_{j=i}^{m}M_{i,j}=\sum_{j=i}^{m-i}M_{i,j}+\sum_{j=m+1-i}^{m}M_{i,j}=\sum_{j=1}^{m-i}M_{i,j}+\sum_{j=1}^{i}M_{j,m+1-i}=\sum_{j=1}^{m-i}M^{\prime}_{i,j}+\sum_{j=1}^{i}M^{\prime}_{j,m+1-i}.$
By the condition $(ii)$, we have
$\sum_{j=i}^{m}M_{i,j}=\sum_{j=1}^{m-i}M^{\prime}_{i,j}+\sum_{j=1}^{i}M^{\prime}_{j,m+1-i}>0$,
which implies that row $i$ contains a nonzero entry. If $\lceil{m\over
2}\rceil+1\leq i\leq m$, the sum of row $i$ of $M$ is given by
$\sum_{j=i}^{m}M_{i,j}=\sum_{j=i}^{m}M^{\prime}_{m+1-j,m+1-i}=\sum_{j=1}^{m+1-i}M^{\prime}_{j,m+1-i},$
which implies that the sum of row $i$ of $M$ is the same as that of column
$m+1-i$ of $M^{\prime}$. By condition $(i)$, row $i$ contains a nonzero entry.
Hence $M$ is a self-dual Fishburn matrix with $R(M)=M^{\prime}$. This
completes the proof.
Denote by $\mathcal{SM}_{k}(n)$ the set of all super triangular matrices of
size $n$ and dimension $2k+1$ having the following two properties:
1. $(a)$
for $1\leq i\leq k$, each column $i$ contains a nonzero entry;
2. $(b)$
for $1\leq i\leq k$, either row $k+1-i$ or column $k+1+i$ contains a nonzero
entry.
Let $\mathcal{SM}(n)=\bigcup_{k\geq 0}\mathcal{SM}_{k}(n)$.
Now we proceed to present a map $\alpha$ from $\mathcal{M}(n)$ to
$\mathcal{SM}(n)$. Given a nonempty self-dual matrix $M$ of dimension $m$, let
$\alpha(M)$ be the matrix obtained from $M$ by the following procedure.
* •
If $m=2k+1$ for some integer $k\geq 0$, then let $\alpha(M)$ be the matrix
obtained from the reduced matrix $R(M)$ of $M$ by interchanging the cell
$(i,k+1)$ and the diagonal cell $(i,m+1-i)$ for $1\leq i\leq k$.
* •
If $m=2k$ for some integer $k\geq 1$, then let $A$ be the matrix obtained from
$R(M)$ by adding one zero row and one zero column immediately after column $k$
and row $k$. Define $\alpha(M)$ to be the matrix obtained from $A$ by
interchanging the cell $(i,k+1)$ and the diagonal cell $(i,m+1-i)$ of the
resulting matrix $A$.
Obviously, $\alpha(M)$ is a super triangular matrix of dimension $2k+1$ and
size $n$. It easy to check that the map $\alpha$ preserves the first $k$
columns and the total sum of row $i$ and column $m+1-i$ of the reduced matrix
$R(M)$. By Lemma 2.1, the matrix $\alpha(M)$ has properties $(a)$ and $(b)$.
Hence $\alpha(M)$ is a super triangular matrix in $\mathcal{SM}(n)$.
Conversely, given a super triangular matrix $M^{\prime}$ in $\mathcal{SM}(n)$
of dimension $2k+1$, we can recover a matrix $M\in\mathcal{M}(n)$ with
$\alpha(M)=M^{\prime}$. First we interchange the cell $(i,k+1)$ with the
diagonal cell $(i,m-i)$ for $1\leq i\leq k$. Then we obtain a matrix $A$ by
deleting column $k+1$ and row $k+1$ if they are zero. It is easy to check that
properties $(a)$ and $(b)$ ensure that the obtained matrix $A$ is the reduced
matrix of a self-dual Fishburn matrix. Let $M$ be a self dual Fishburn matrix
with $R(M)=A$. Hence $\alpha$ is a bijection between $\mathcal{M}(n)$ and
$\mathcal{SM}(n)$.
Let $M$ be a super triangular matrix of dimension $2k+1$, then column $k+1$ is
called a center column. From the construction of the bijection $\alpha$, we
see that the map $\alpha$ transforms the sum of the diagonal cells of a self-
dual matrix to the sum of the center column of a super triangular matrix.
Hence, we have the following result.
###### Theorem 2.2
The map $\alpha$ is a bijection between $\mathcal{M}(n)$ and
$\mathcal{SM}(n)$. Moreover, the bijection $\alpha$ preserves the sum of the
first row, and transforms the sum of the diagonal cells of a self-dual matrix
to the sum of the center column of a super triangular matrix.
###### Example 2.3
Consider a matrix $A\in\mathcal{M}(5)$,
$A=\begin{bmatrix}1&0&1&0&0\\\ 0&1&1&1&0\\\ 0&0&0&1&1\\\ 0&0&0&1&0\\\
0&0&0&0&1\\\ \end{bmatrix}.$
The reduced matrix of $A$ is given by
$R(A)=\begin{bmatrix}1&0&1&0&0\\\ 0&1&1&1&0\\\ 0&0&0&0&0\\\ 0&0&0&0&0\\\
0&0&0&0&0\\\ \end{bmatrix},$
and we have
$\alpha(A)=\begin{bmatrix}1&0&0&0&1\\\ 0&1&1&1&0\\\ 0&0&0&0&0\\\ 0&0&0&0&0\\\
0&0&0&0&0\\\ \end{bmatrix}.$
Let $\mathcal{B}(n)$ the set of upper-triangular matrices of size $n$ in which
each row contains a nonzero entry except for the first row. Given a nonempty
matrix $A\in\mathcal{RM}(n)$, we can get two distinct matrices in
$\mathcal{B}(n)$ from $A$ by either doing nothing or adding a zero row and a
zero column before the first row and the first column. Thus for $n\geq 1$ we
have the following relation
$|\mathcal{B}(n)|=2|\mathcal{RM}(n)|.$ (2.1)
Now we proceed to construct a bijection between the set $\mathcal{SM}(n)$ and
the set $\mathcal{B}(n)$. Before constructing the bijection, we need some
definitions. In a matrix $A$ with $m$ rows, the operation of adding column $i$
to column $j$ is defined by increasing $A_{k,j}$ by $A_{k,i}$ for each
$k=1,2,\ldots,m$.
Let $\mathcal{B}(n,k,p)$ be the set of matrices in $\mathcal{B}(n)$ whose
whose first row has sum $p$ and last column has sum $k$. Similarly, let
$\mathcal{SM}(n,k,p)$ be the set of matrices in $\mathcal{SM}(n)$ whose first
row has sum $k$ and center column has sum $p$.
###### Theorem 2.4
There is a bijection $\beta$ between $\mathcal{SM}(n)$ and $\mathcal{B}(n)$.
Moreover, the map $\beta$ is essentially a bijection between
$\mathcal{SM}(n,k,p)$ and $\mathcal{B}(n,k,p)$.
Proof. Given a nonempty triangular matrix $A\in\mathcal{SM}(n)$ of dimension
$2k+1$, we recursively construct a sequence of super triangular matrices
$A^{(0)},A^{(1)},\ldots,A^{(l)}$. Let $A^{(0)}=A$ and assume that we have
obtained the matrix $A^{(j)}$. Let $A^{(j)}$ be a super triangular matrix of
dimension $2r+1$ for some integer $r$. For $1\leq i\leq r$, if each column
$r+1+i$ is zero, then let $A^{(l)}=A^{(j)}$. Otherwise, we proceed to generate
the matrix $A^{(j+1)}$ by the following insertion algorithm.
* •
Find the largest value $i$ such that column $r+1+i$ contains a nonzero entry.
Then fill the entries of column $r+1+i$ with zeros.
* •
Insert one column immediately after column $r+1-i$, one zero row immediately
after row $r+1-i$, one zero column immediately before column $r+1+i$ and one
zero row immediately before row $r+1+i$. Let the entry in row $j$ of the new
inserted column after column $r+1-i$ be filled with the entry in row $j$ of
column $r+1+i$ of $A^{(j)}$ for $1\leq j\leq 2r+1$.
Suppose that $A^{(l)}$ is of dimension $2q+1$. Then the last $q$ rows and $q$
columns of $A^{(l)}$ are zero rows and columns. Let $B$ be an upper-triangular
matrix obtained from $A^{(l)}$ by deleting the last $q$ columns and $q$ rows.
From the above insertion procedure to generate $A^{(j+1)}$ form $A^{(j)}$ , we
see that the inserted column after column $r+1-i$ contains a nonzero entry.
This ensures that each matrix $A^{(j)}$ has property $(a)$ with $0\leq j\leq
l$. Hence each column $i$ of $A^{(l)}$ contains a nonzero entry with $1\leq
i\leq q$. Hence, $B$ is an upper-triangular matrix in which each column
contains a nonzero entry except for the last column. Moreover, the insertion
algorithm preserves the sum of each nonzero row of $A$, which implies that $B$
is of size $n$. Let $\beta(A)$ be the dual matrix of $B$. Hence we have
$\beta(A)\in\mathcal{B}(n)$.
Conversely, we can construct a matrix $A=\beta^{\prime}(A^{\prime})$ in
$\mathcal{SM}(n)$ from a matrix $A^{\prime}$ of dimension $k+1$ in
$\mathcal{B}(n)$. Let $B$ be the dual matrix of $A^{\prime}$. Define $M$ to be
a matrix of dimension $2k+1$ obtained from $B$ by adding $k$ consecutive zero
rows and $k$ consecutive zero columns immediately after column $k+1$ and row
$k+1$. Clearly, the obtained matrix is a super triangular matrix having
property $(a)$. If for all $1\leq i\leq k$, either row $k+1-i$ or column
$k+1+i$ contains a nonzero entry, then we do nothing for $M$ and let $A=M$.
Otherwise, we can construct a new super triangular matrix $A$ by the following
removal algorithm.
* •
Find the least value $i$ such that neither row $k+1-i$ nor column $k+1+i$
contains a nonzero entry. Then we obtain a super triangular matrix by adding
column $k+1-i$ to column $k+2+i$ and removing columns $k+1+i$, $k+1-i$ and
rows $k+1-i$, $k+1+i$.
* •
Repeat the above procedure for the resulting matrix until the obtained matrix
has property $(b)$.
Obviously, the obtained matrix $A$ is a super triangular matrix having
properties $(a)$ and $(b)$. Since the algorithm preserves the sums of entries
in each non-zero row of $B$, the matrix $A$ is of size $n$ and the sum of the
first row of $A$ is the same as that of $B$. The property $(b)$ ensures that
the inserted columns in the insertion algorithm are the removed columns in the
removal algorithm. Thus the map $\beta^{\prime}$ is the inverse of the map
$\beta$. From the construction of the removal algorithm, the sum of the center
column of $A$ is equal to the sum of the last column of $B$ as well as the the
sum of the first row of $A^{\prime}$. This completes the proof.
###### Example 2.5
Consider a matrix $A\in\mathcal{SM}(6)$,
$A=\begin{bmatrix}1&0&0&1&1\\\ 0&1&1&1&0\\\ 0&0&0&0&0\\\ 0&0&0&0&0\\\
0&0&0&0&0\\\ \end{bmatrix}.$
Let $A^{(0)}=A$. By applying the insertion algorithm, we get
$A^{(1)}=\begin{bmatrix}1&{\textbf{1}}&0&0&1&\textbf{0}&0\\\
\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}\\\
0&\textbf{0}&1&1&1&\textbf{0}&0\\\ 0&\textbf{0}&0&0&0&\textbf{0}&0\\\
0&\textbf{0}&0&0&0&\textbf{0}&0\\\
\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}\\\
0&\textbf{0}&0&0&0&\textbf{0}&0\\\ \end{bmatrix},$
$A^{(2)}=\begin{bmatrix}1&1&0&\textbf{1}&0&\textbf{0}&0&0&0\\\
0&0&0&\textbf{0}&0&\textbf{0}&0&0&0\\\ 0&0&1&\textbf{1}&1&\textbf{0}&0&0&0\\\
\bf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}\\\
0&0&0&\textbf{0}&0&\textbf{0}&0&0&0\\\ \textbf{0}&\textbf{0}&\textbf{0
}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}&\textbf{0}\\\
0&0&0&\textbf{0}&0&\textbf{0}&0&0&0\\\ 0&0&0&\textbf{0}&0&\textbf{0}&0&0&0\\\
0&0&0&\textbf{0}&0&\textbf{0 }&0&0&0\\\ \end{bmatrix},$
where the inserted rows and columns are illustrated in bold at each step of
the insertion algorithm. Removing the last $4$ zero rows and $4$ zero columns,
we get
$B=\begin{bmatrix}1&1&0&1&0\\\ 0&0&0&0&0\\\ 0&0&1&1&1\\\ 0&0&0&0&0\\\
0&0&0&0&0\\\ \end{bmatrix}.$
Finally, we obtain
$A^{\prime}=\beta(A)=\begin{bmatrix}0&0&1&0&0\\\ 0&0&1&0&1\\\ 0&0&1&0&0\\\
0&0&0&0&1\\\ 0&0&0&0&1\\\ \end{bmatrix}$
Conversely, given $A^{\prime}\in\mathcal{B}(6)$, by applying removal
algorithm, we can get $A\in\mathcal{SM}(6)$, where the removed rows and
columns are illustrated in bold at each step of the removal algorithm.
Combining the bijection between self-dual interval orders and self-dual
Fishburn matrices and Theorems 2.2 and 2.4, we get a bijective proof of
$(\ref{eq3})$.
From Theorems 2.2 and 2.4, we have
$|\mathcal{M}(n,k,p)|=|\mathcal{SM}(n,k,p)|=|\mathcal{B}(n,k,p)|.$
Given a matrix $M\in\mathcal{B}(n,k,0)$, we can get a matrix
$A\in\mathcal{RM}(n,k)$ by deleting the first row and the first column.
Conversely, given a matrix $A^{\prime}\in\mathcal{RM}(n,k)$, we can obtain a
matrix $M^{\prime}\in\mathcal{B}(n,k,0)$ by inserting a zero row and a zero
column before the first row and the first column. This yields that
$|\mathcal{M}(n,k,0)|=|\mathcal{B}(n,k,0)|=|\mathcal{RM}(n,k)|.$ (2.2)
If $p>0$, then $\mathcal{B}(n,k,p)$ is the same as $\mathcal{RM}(n,k,p)$.
Hence, if $p>0$ then we have
$|\mathcal{M}(n,k,p)|=|\mathcal{B}(n,k,p)|=|\mathcal{RM}(n,k,p)|.$ (2.3)
Therefore, we get combinatorial proofs of (1.1) and (1.2), in answer to the
problem posed by Jelínek [10].
Now we proceed to prove (1.4). Given a matrix $A\in\mathcal{EM}(n,k)$ of
dimension $2m$ for some integers $m\geq 1$, let $R(A)$ be its reduced matrix.
We obtain a super triangular matrix $A^{\prime}$ from $A$ by inserting a zero
column and a zero row immediately after column $m$ and row $m$. By Lemma 2.1,
we have $A^{\prime}\in\mathcal{SM}(n,k,0)$.
Conversely, given a matrix $A^{\prime}\in\mathcal{SM}(n,k,0)$ of dimension
$2m+1$ for some integer $m\geq 1$, we can recover a self-dual matrix
$A\in\mathcal{EM}(n,k)$ as follows. First, we get a super triangular matrix
$B$ from $A^{\prime}$ by deleting column $m+1$ and row $m+1$. Let $A$ be a
matrix with $B=R(A)$. Obviously, we have the matrix $A\in\mathcal{EM}(n,k)$.
Hence, we get $|\mathcal{EM}(n,k)|=|\mathcal{SM}(n,k,0)|.$ By (2.2), we deduce
that
$|\mathcal{EM}(n,k)|=|\mathcal{SM}(n,k,0)|=|\mathcal{RM}(n,k)|.$ (2.4)
From (2.2) and (2.3), we have
$\begin{array}[]{lll}|\mathcal{M}(n,k)|&=&|\mathcal{M}(n,k,0)|+\sum_{p\geq
1}|\mathcal{M}(n,k,p)|\\\ &=&|\mathcal{RM}(n,k)|+\sum_{p\geq
1}|\mathcal{RM}(n,k,p)|\\\ &=&2|\mathcal{RM}(n,k)|.\end{array}$
Meanwhile, we have
$|\mathcal{M}(n,k)|=|\mathcal{EM}(n,k)|+|\mathcal{OM}(n,k)|$. Hence, (1.4)
follows from (2.4).
Acknowledgments. This work was supported by the National Natural Science
Foundation of China (10901141).
## References
* [1] K.P. Bogart, An obvious proof of Fishburn’s interval order theorem, Discrete Math. 118 (1993), 239–242.
* [2] M. Bousquet-Mélou, A. Claesson, M. Dukes, S. Kitaev, $(2+2)$-free posets, ascent sequences and pattern avoiding permutations, J. Combin. Theory Ser. A 117 (2010), 884–909.
* [3] A. Claesson, M. Dukes, and M. Kubitzke, Partition and composition matrices, J. Combin. Theory, Ser. A, 118 (2011), 1624–1637.
* [4] M. Dukes, R. Parviainen, Ascent sequences and upper triangular matrices containing non-negative integers, Electronic J. combin. 17 (2010), R53.
* [5] M. Dukes, S. Kitaev, J. Remmel, and E. Steingrímsson, Enumerating (2+2)-free posets by indistinguishable elements, arXiv:1006.2696, 2010\.
* [6] M. Dukes, V. Jelínek, and M. Kubitzke, Composition matrices, (2+2)-free posets and their specializations, Electronic J. Combin., 18 (2011), P44.
* [7] P. C. Fishburn, Interval lengths for interval orders: A minimization problem, Discrete Mathematics, 47 (1983), 63–82.
* [8] P. C. Fishburn, Interval graphs and interval orders, Discrete Mathematics, 55 (1985), 135–149.
* [9] P. C. Fishburn, Interval orders and interval graphs: A study of partially ordered sets, John Wiley & Sons, 1985.
* [10] V. Jelínek, Counting self-dual interval orders, arXiv:1106.2261, 2011.
* [11] S. Kitaev, J. Remmel, Enumerating $(2+2)$-free posets by the number of minimal elements and other statistics, Disctere Appl. Math., 159 (2011), 2098–2108.
* [12] A. Stoimenow, Enumumeration of chord diagrams and an upper bound for Vassiliev invariants, J. Knot Theory Ramifications 7 (1998), 93–114.
* [13] D. Zagier, Vassiliev invariants and a stange identity related to the Dedeking eta-function, Topology 40 (2001), 945–960.
|
arxiv-papers
| 2011-11-21T03:33:39 |
2024-09-04T02:49:24.535265
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Sherry H. F. Yan, Yuexiao Xu",
"submitter": "Sherry H.F. Yan",
"url": "https://arxiv.org/abs/1111.4723"
}
|
1111.4831
|
# Analytical calculation of optimal POVM for unambiguous discrimination of
quantum states using KKT method
N. Karimi a
aDepartment of Physics, Azarbaijan University of Tarbiat Moallem, 53714-161
Tabriz, Iran. E-mail:${na}_{-}karimi@yahoo.com$
###### Abstract
In the present paper, an exact analytic solution for the optimal unambiguous
state discrimination(OPUSD) problem involving an arbitrary number of pure
linearly independent quantum states with real and complex inner product is
presented. Using semidefinite programming and Karush-Kuhn-Tucker convex
optimization method, we derive an analytical formula which shows the relation
between optimal solution of unambiguous state discrimination problem and an
arbitrary number of pure linearly independent quantum states.
Keywords: Unambiguous discrimination, Semidefinite programming, linearly
independent States.
PACs Index: 03.67.Hk, 03.65.Ta
## 1 Introduction
Many applications in quantum communication and quantum cryptography are based
on transmitting quantum systems that, with given prior probabilities, are
prepared in one from a set of known mutually nonorthogonal states[1]. A
fundamental aspect of quantum information theory is that nonorthogonal quantum
states cannot be perfectly distinguished. Therefore, a central problem in
quantum mechanics is to design measurements optimized to distinguish between a
collection of nonorthogonal quantum states. The topic of quantum state
discrimination was firmly established in the 1970s by the pioneering work of
Helstrom [2], who considered a minimum error discrimination of two known
quantum states. In this case, the state identification is probabilistic.
Another possible discrimination strategy is the so-called unambiguous state
discrimination (USD) where the states are successfully identified with nonunit
probability, but without error.
USD was originally formulated and analyzed by Ivanovic, Dieks, and Peres [3,
4, 5] in 1987. The solution for unambiguous discrimination of two known pure
states appearing with arbitrary prior probabilities was obtained by Jaeger and
Shimony [6]. Although the two-state problem is well developed, the problem of
unambiguous discrimination between multiple quantum states has received
considerably less attention. The problem of discrimination among three
nonorthogonal states was first considered by Peres and Terno [5]. They
developed a geometric approach and applied it numerically on several examples.
A different method was considered by Duan and Guo [7] and Sun et al.[8] .
Chefles [9] showed that a necessary and sufficient condition for the existence
of unambiguous measurements for distinguishing between N quantum states is
that the states are linearly independent. He also proposed a simple suboptimal
measurement for unambiguous discrimination for which the probability of an
inconclusive result is the same regardless of the state of the system.
Equivalently, the measurement yields an equal probability of correct detection
of each one of the ensemble states.
Over the past years, semidefinite programming (SDP) has been recognized as a
valuable numerical tool for control system analysis and design. In SDP, one
minimizes a linear function subject to the constraint that an affine
combination of symmetric matrices is positive semidefinite. SDP has been
studied (under various names) as far back as the 1940s. Subsequent research in
semidefinite programming during the 1990s was driven by applications in
combinatorial optimization [10], communications and signal processing [11, 12,
13] and other areas of engineering [14]. Although semidefinite programming is
designed to be applied in numerical methods, it can be used for analytic
computations, too. In the context of quantum computation and quantum
information, Barnum, Saks, and Szegedy have reformulated quantum query
complexity in terms of a semidefinite program [15]. The problem of finding the
optimal measurement to distinguish between a set of quantum states was first
formulated as a semidefinite program in 1974 by Holevo, who gave optimality
conditions equivalent to the complementary slackness conditions[2]. Recently,
Eldar, Megretski, and Verghese showed that the optimal measurements can be
found efficiently by solving the dual followed by the use of linear
programming[16]. Also in Ref. [17], SDP has been used to show that the
standard algorithm implements the optimal set of measurements. All of the
above mentioned applications indicate that the method of SDP is very useful.
The reason why the area has shown relatively slow progress until recently
within the rapidly evolving field of quantum information is that it poses
quite formidable mathematical challenges. Except for a handful of very special
cases, no general exact solution has been available involving more than two
arbitrary states and mostly numerical algorithms are proposed for finding
optimal measurements for quantum-state discrimination, where the theory of the
semidefinite programming provides a simple check of the optimality of the
numerically obtained results. In Ref. [18] we obtained the feasible region in
terms of the inner product of the reciprocal states and the feasible region in
terms of the inner product of the states which enables us to solve the problem
without using reciprocal states. Moreover, for the real inner product of
states, we obtained an exact analytic solution for OPUSD problem involving an
arbitrary number of pure linearly independent quantum states by using KKT
convex optimization method. In this paper, an exact analytic solution for the
optimal unambiguous state discrimination(OPUSD) problem involving an arbitrary
number of pure linearly independent quantum states with real and complex inner
product of states,using the Karush- Kuhn-Tucker convex optimization method, is
given.
The organization of the paper is as follows. First, the definition of the
unambiguous quantum state discrimination is given. Then, Semidefinite
programming, Karush-Kuhn-Tucker (KKT) theorem and SDP formulation of
unambiguous discrimination is studied. Finally for the real and complex inner
product of reciprocal states, an exact analytic solution for OPUSD problem
involving an arbitrary number of pure linearly independent quantum states is
presented by using KKT convex optimization method. The paper ends with a brief
conclusion.
## 2 Unambiguous quantum state discrimination
In quantum theory, measurements are represented by positive operator valued
measures (POVMs). A measurement is described by a collection ${M_{k}}$ of
measurement operators. These operators are acting on the state space of the
system being measured. The index k refers to the measurement outcomes that may
occur in the experiment. In quantum information theory the measurement
operators ${M_{k}}$ are often called Kraus operators [19]. If we define the
operator
$\Pi_{k}=M^{\dagger}_{k}M_{k},$ (2.1)
the probability of obtaining the outcome $k$ for a given state $\rho_{i}$ is
given by $p(k|i)=Tr(\Pi_{k}\rho_{i})$. Thus, the set of operators $\Pi_{k}$ is
sufficient to determine the measurement statistics.
### 2.1 Definition of the POVM
A set of operators $\\{\Pi_{k}\\}$ is named a positive operator valued measure
if and only if the following two conditions are met: $(1)$ each operator
$\Pi_{k}$ is positive positive
$\Leftrightarrow\langle\psi|\Pi_{k}|\psi\rangle\geq 0$, $\forall|\psi\rangle$
and $(2)$ the completeness relation is satisfied, i.e.,
$\sum_{k}\Pi_{k}=1.$ (2.2)
The elements of $\\{\Pi_{k}\\}$ are called effects or POVM elements. On its
own, a given POVM $\\{\Pi_{k}\\}$ is enough to give complete knowledge about
the probabilities of all possible outcomes; measurement statistics is the only
item of interest. Consider a set of known states $\rho_{i},i=1,...,N$ with
their prior probabilities $\eta_{i}$. We are looking for a measurement that
either identifies a state uniquely (conclusive result) or fails to identify it
(inconclusive result). The goal is to minimize the probability of inconclusive
results. The measurements involved are typically generalized measurements. A
measurement described by a POVM $\\{\Pi_{k}\\}_{k=1}^{N}$ is called an
unambiguous state discrimination measurement (USDM) on the set of states
$\\{\rho_{i}\\}_{i=1}^{N}$ if and only if the following conditions are
satisfied. (1) The POVM contains the elements $\\{\Pi_{k}\\}_{k=0}^{N}$ where
$N$ is the number of different signals in the set of states. The element
$\Pi_{0}$ is related to an inconclusive result, while the other elements
correspond to an identification of one of the states $\rho_{i},i=1,...,N$. (2)
No states are wrongly identified, that is, $Tr(\rho_{i}\Pi_{k})=0\ \forall
i\neq k,k=1,...,N$. Each USD measurement gives rise to a failure probability,
that is, the rate of inconclusive result. This can be calculated as
$Q=\sum_{i}\eta_{i}Tr(\eta_{i}\Pi_{0}).$ (2.3)
The success probability can be calculated as
$P=1-Q=\sum_{i}\eta_{i}Tr(\eta_{i}\Pi_{i}).$ (2.4)
A measurement described by a POVM $\\{\Pi_{k}^{opt}\\}$ is called an optimal
unambiguous state discrimination measurement (OUSDM) on a set of states
$\\{\rho_{i}\\}$ with the corresponding prior probabilities $\\{\eta_{i}\\}$
if and only if the following conditions are satisfied. (1) The POVM
$\\{\Pi_{k}^{opt}\\}$ is a USD measurement on $\\{\rho_{i}\\}$. (2) The
probability of inconclusive result is minimal, that is,
$Q(\\{\Pi_{k}^{opt}\\})=min\ Q(\\{\Pi_{k}\\})$, where the minimum is taken
over all USDM. Unambiguous state discrimination is an error-free
discrimination. This implies a strong constraint on the measurement. Suppose
that a quantum system is prepared in a pure quantum state drawn from a
collection of given states $\\{|\psi_{i}\rangle\\},1\leq i\leq N$ in
d-dimensional complex Hilbert space $\mathcal{H}$ with $d\geq N$. These states
span a subspace $\mathcal{U}$ of $\mathcal{H}$. In order to detect the state
of the system, a measurement is constructed comprising $N+1$ measurement
operators $\\{\Pi_{i},0\leq i\leq N\\}$. Given that the state of the system is
$|\psi_{i}\rangle$, the probability of obtaining outcome k is
$\langle\psi_{i}|\Pi_{k}|\psi_{i}\rangle$. Therefore, in order to ensure that
each state is either correctly detected or an inconclusive result is obtained,
we must have
$\langle\psi_{i}|\Pi_{k}|\psi_{i}\rangle=p_{i}\delta_{ij},\ 1\leq i,k\leq N$
(2.5)
for some $0\leq p_{i}\leq 1$. Since $\Pi_{0}=I_{d}-\sum_{i=1}^{N}\Pi_{i}$, we
have $\langle\psi_{i}|\Pi_{0}|\psi_{i}\rangle=1-p_{i}$. So a system with given
state $|\psi_{i}\rangle$ , the state of the system is correctly detected with
probability $p_{i}$ and an inconclusive result is obtained with probability
$1-p_{i}$. It was shown in Ref. [20] that Eq.(2.5) is satisfied if and only if
the vectors $|\psi_{i}\rangle$ are linearly independent, or equivalently,
$dim\ \mathcal{U}=N$. Therefore, we will take this assumption throughout the
paper. In this case, we may choose [21]
$\Pi_{i}=p_{i}|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|,\ 1\leq
i\leq N,$ (2.6)
where the vectors $|\widetilde{\psi_{i}}\rangle\in\mathcal{U}$ are the
reciprocal states associated with the states $|\psi_{i}\rangle$ , i.e., there
are unique vectors in $\mathcal{U}$ such that
$\langle\widetilde{\psi}_{k}|\psi_{i}\rangle=\delta_{ij},\ 1\leq i,k\leq N.$
(2.7)
With $\Phi$ and $\widetilde{\Phi}$ we denote the matrices such that their
columns are $|\psi_{i}\rangle$ and $|\widetilde{\psi_{i}}\rangle$,
respectively. Then, one can show that $\widetilde{\Phi}$ is
$\widetilde{\Phi}=\Phi(\Phi^{\ast}\Phi)^{-1}.$ (2.8)
Since the vectors $|\psi_{i}\rangle,i=1,...,N$ are linearly independent,
$\Phi^{\ast}\Phi$ is always invertible. Alternatively,
$\widetilde{\Phi}=(\Phi\Phi^{\ast})^{{\ddagger}}\Phi,$ (2.9)
so that
$|\widetilde{\psi_{i}}\rangle=(\Phi\Phi^{\ast})^{{\ddagger}}|\psi_{i}\rangle,$
(2.10)
where $(...)^{{\ddagger}}$ denotes the Moore-Penrose pseudoinverse [22, 23].
The inverse is taken on the subspace spanned by the columns of the matrix. If
the state $|\psi_{i}\rangle$ is prepared with prior probability $\eta_{i}$,
then the total probability of correctly detecting the state is
$P=\sum_{i=1}^{N}\eta_{i}\langle\psi_{i}|\Pi_{i}|\psi_{i}\rangle=\sum_{i=1}^{N}\eta_{i}p_{i}$
(2.11)
and the probability of the inconclusive result is given by
$Q=1-P=\sum_{i=1}^{N}\eta_{i}\langle\psi_{i}|\Pi_{0}|\psi_{i}\rangle=1-\sum_{i=1}^{N}\eta_{i}p_{i}$
(2.12)
In general, an optimal measurement for a given strategy depends on the quantum
states and the prior probabilities of their appearance. In the unambiguous
discrimination for a given strategy and a given ensemble of states, the goal
is to find a measurement which minimizes the inconclusive result. In fact, it
is known that USD (of both pure and mixed states) is a convex optimization
problem. Mathematically, this means that the quantity which is to be optimized
as well as the constraints on the unknowns, are convex functions. Practically,
this implies that the optimal solution can be computed in an extremely
efficient way. This is therefore a very useful tool. Nevertheless, our aim is
to understand the structure of USD in order to relate it with neat and
relevant quantities and to find feasible region for numerical and analytic
solutions.
## 3 Semidefinite programming
A SDP problem requires minimizing a linear function subject to a linear matrix
inequality (LMI) constraint
$\mathbf{minimize}\ p=c^{T}x,such\ that\ F(x)\geq 0,$ (3.13)
where $c^{T}$ is a given vector, $x=(x_{1},...,x_{n})$, and
$F(x)=F_{0}+\sum_{i}x_{i}F_{i}$, for some fixed Hermitian matrices $F_{i}.$
The inequality sign in $F(x)\geq 0$ means that $F(x)$ is positive
semidefinite. This problem is called the primal problem. Vectors $x$ whose
components are the variables of the problem and satisfy the constraint
$F(x)\geq 0$ are called primal feasible points, and if they satisfy $F(x)>$
they are called strictly feasible points. The minimal objective value $c^{T}x$
is by convention denoted by $P^{\ast}$ and is called the primal optimal value.
Due to the convexity of set of feasible points, SDP has a nice duality
structure with the associated dual program being
$\mathbf{maximize}\ -Tr[F_{0}Z],\ Z\geq,\ Tr[F_{i}Z]=c_{i}.$ (3.14)
Here the variable is the real symmetric (or Hermitian) matrix $Z$, and the
data $c$, $F_{i}$ are the same as in the primal problem. Correspondingly,
matrix $Z$ satisfying the constraints are called dual feasible (or strictly
dual feasible if $Z>$). The maximal objective value of $-Tr[F_{0}Z]$, i.e.,
the dual optimal value, is denoted by $d^{\ast}$. The objective value of a
primal (dual) feasible point is an upper (lower) bound on
$P^{\ast}(d^{\ast})$. The main reason why one is interested in the dual
problem is that one can prove that $d^{\ast}\leq P^{\ast}$, and under
relatively mild assumptions, we can have $P^{\ast}=d^{\ast}$. If the equality
holds, one can prove the following optimality condition on $x$. A primal
feasible $x$ and a dual feasible $Z$ are optimal which is denoted by $\hat{x}$
and $\hat{Z}$ if and only if
$F(\hat{x})\hat{Z}=\hat{Z}F(\hat{x}).$ (3.15)
This latter condition is called the complementary slackness condition. In one
way or another, numerical methods for solving SDP problems always exploit the
inequality $d\leq d^{\ast}\leq P^{\ast}\leq P,$ where $d$ and $P$ are the
objective values for any dual feasible point and primal feasible point,
respectively. The difference
$P^{\ast}-d^{\ast}=c^{T}x+Tr[F_{0}Z]=Tr[F_{x}Z]\geq 0$ (3.16)
is called the duality gap. If the equality holds $d^{\ast}=P^{\ast}$, i.e.,
the optimal duality gap is zero, then we say that strong duality holds.
### 3.1 Karush-Kuhn-Tucker (KKT) theorem
Assuming that functions $g_{i}$, $h_{i}$ are differentiable and that strong
duality holds, there exists vectors $\xi\in R^{k}$ and $y\in R^{m}$ such that
the gradient of dual Lagrangian
$L(x^{\ast}.\xi^{\ast},y^{\ast})=f(x^{\ast})+\sum_{i}\xi_{i}^{\ast}h_{i}(x^{\ast})+\sum_{i}y_{i}^{\ast}g_{i}(x^{\ast})$
over $x$ vanishes at $x^{\ast}$:
$h_{i}(x^{\ast})=0\ (primal\ feasible),$ $g_{i}(x^{\ast})\leq 0\ (primal\
feasible),$ $y_{i}^{\ast}\geq 0\ (dual\ feasible),\
y_{i}^{\ast}g_{i}(x^{\ast})=0,$ $\bigtriangledown
f(x^{\ast})+\sum_{i}\xi_{i}^{\ast}\bigtriangledown
h_{i}(x^{\ast})+\sum_{i}y_{i}^{\ast}\bigtriangledown g_{i}(x^{\ast})=0.$
(3.17)
Then $x^{\ast}$ and $(\xi_{i}^{\ast},y_{i}^{\ast})$ are primal and dual
optimal with zero duality gap. In summary, for any convex optimization problem
with differentiable objective and constraint functions, the points which
satisfy the KKT conditions are primal and dual optimal, and have zero duality
gap. Necessary KKT conditions satisfied by any primal and dual optimal pair
and for convex problems, KKT conditions are also sufficient. If a convex
optimization problem with differentiable objective and constraint functions
satisfies Slater s condition, then the KKT conditions provide necessary and
sufficient conditions for optimality: Slater s condition implies that the
optimal duality gap is zero and the dual optimum is attained, so x is optimal
if and only if there are $(\xi_{i}^{\ast},y_{i}^{\ast})$ such that they,
together with $x$, satisfy the KKT conditions.
### 3.2 Slater s condition
Suppose $x^{\ast}$ solves
$\mathbf{minimize}\ f(x)g_{i}(x)\geq b_{i},\ i=1,...,m$ (3.18)
and the feasible set is nonempty. Then there is a nonnegative vector $\xi$
such that for all $x$
$L(x,\xi)=f(x)+\xi^{T}[b-g(x)]\leq f(x^{\ast})=L(x^{\ast},\xi).$ (3.19)
In addition, if $f(...)$, $g_{i}(...)$, $i=1,...,m,$ are continuously
differentiable, then
$\frac{\partial f(x^{\ast})}{\partial(x_{j})}-\xi\frac{\partial
g(x^{\ast})}{\partial(x)}=0.$ (3.20)
In the spatial case the vector x is a solution of the linear program
$\mathbf{minimize}\ c^{T}x,such\ that\ Ax=bx\geq 0,$ (3.21)
if and only if there exist vectors $\xi\in R^{k}$ and $y\in R^{m}$ for which
the following conditions hold for $(x,\xi,y)=(x^{\ast},\xi^{\ast},y^{\ast})$:
$A^{T}\xi+y=c,Ax=b,x_{i}\geq 0;\ y\geq 0;$ $x_{i}y_{i}=0,\ i=1,...,m.$ (3.22)
A solution $(x^{\ast},\xi^{\ast},y^{\ast})$ is called strictly complementary,
if $x^{\ast}+y^{\ast}>0$ , i.e., if there exists no index $i\in 1,...m$ such
that $x_{i}^{\ast}=y_{i}^{\ast}$.
## 4 SDP Formulation of unambiguous discrimination
Eldar, Megretski, and Verghese in Ref.[24] have showed that the unambiguous
discrimination problem can be reduced to SDP method and the KKT conditions can
be defined as
$F(p)=\sum_{i=1}^{N}p_{i}F_{i}+F_{0}\geq 0,\ Tr(\Pi_{i}X)=z_{i}+\eta_{i},$
$z_{i}\geq 0,1\leq i\leq N,\ AF(p)=0,\ X(I_{d}-\sum_{i=1}^{N}p_{i}\Pi_{i})=0,$
$z_{i}p_{i}=0,\ 1\leq i\leq N,\exists p_{i}:\sum_{i=1}^{N}p_{i}F_{i}+F_{0}\geq
0,$
$\exists X,z:X\geq 0,z\geq 0$ (4.23)
such that
$|p\rangle=:\left(\begin{array}[]{c}p_{1}\\\ \vdots\\\ p_{N}\\\
\end{array}\right),|c\rangle=:\left(\begin{array}[]{c}\eta_{1}\\\ \vdots\\\
\eta_{2}\\\ \end{array}\right),$
$F_{0}=:\left(\begin{array}[]{cccc}I_{d}&0&\cdots&0\\\ 0&0&\cdots&0\\\
\vdots&\vdots&\ddots&\vdots\\\ 0&0&\cdots&0\\\
\end{array}\right),F_{1}=:\left(\begin{array}[]{ccccc}-\Pi_{1}&0&0&\cdots&0\\\
0&1&0&\cdots&0\\\ 0&\cdots&0&0&0\\\ \vdots&\vdots&\vdots&\ddots&\vdots\\\
0&0&\cdots&0&0\\\
\end{array}\right),F_{2}=:\left(\begin{array}[]{ccccc}-\Pi_{2}&0&0&\cdots&0\\\
0&0&0&\cdots&0\\\ 0&\cdots&1&0&0\\\ \vdots&\vdots&\vdots&\ddots&\vdots\\\
0&0&\cdots&0&0\\\ \end{array}\right),...,$
$F_{N}=:\left(\begin{array}[]{ccccc}-\Pi_{N}&0&0&\cdots&0\\\ 0&0&0&\cdots&0\\\
0&\cdots&0&0&0\\\ \vdots&\vdots&\vdots&\ddots&\vdots\\\ 0&0&\cdots&0&1\\\
\end{array}\right),A=:\left(\begin{array}[]{cc}X_{d}&Y\\\
Y^{T}&\left(\begin{array}[]{ccc}z_{1}&\cdots&0\\\ 0&\ddots&0\\\
0&\cdots&z_{N}\\\ \end{array}\right)\\\ \end{array}\right),$ (4.24)
where $X_{d}$ is the $d\times d$ matrix and $Y$ is the $N\times d$ matrix.
## 5 Analytical solution of unambiguous state discrimination problem
When the set of states $\\{\psi_{i}\\}_{i=1}^{N}$ and the prior probabilities
$\\{\eta_{i}\\}_{i=1}^{N}$ are given, the POVM elements for the measurement
$(\Pi_{i}=\sum_{i=1}^{N}p_{i}|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|,\
\Pi_{0}=I-\sum_{i=1}^{N}\prod_{i})$ depend on only the variables $p_{i}$ and
these variables lie in the feasible region. Feasible region that gives the
domain of acceptable values of $p_{i}$ is determined by the following equation
[25, 26, 27, 28]:
$1-\sum_{i}D_{i}p_{i}+\sum_{i<j}D_{ij}p_{i}p_{j}-\sum_{i<j<k}D_{ijk}p_{i}p_{j}p_{k}+...+(-1)^{N}\sum_{i_{1}<i_{2}...<i_{N}}p_{i_{1}}p_{i_{2}}...p_{i_{N}}D_{i_{1}i_{2}...i_{N}}=0,$
(5.25)
where the set of $\\{D_{i_{1}i_{2}...i_{N}}\\}$ are the subdeterminants
(minor) of matrix D defined by
$D=\left(\begin{array}[]{cccc}\tilde{a}_{11}&\tilde{a}_{12}&\cdots&\tilde{a}_{1N}\\\
\tilde{a}_{21}&\tilde{a}_{22}&\cdots&\tilde{a}_{2N}\\\
\vdots&\vdots&\ddots&\vdots\\\
\tilde{a}_{N1}&\tilde{a}_{N2}&\cdots&\tilde{a}_{NN}\\\ \end{array}\right)$
(5.26)
with $\tilde{a}_{ij}=\langle\tilde{\psi}_{i}|\tilde{\psi}_{j}\rangle$.
In this section, for the real and complex inner product of reciprocal states,
an exact analytic solution for OPUSD problem involving an arbitrary number of
pure linearly independent quantum states is presented by using KKT convex
optimization method. The KKT conditions for unambiguous discrimination of N
linearly independent states are given by
$I-\sum_{i=1}^{N}p_{i}|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|\geq
0,\ p_{i}\geq 0,$
$(I-\sum_{i=1}^{N}p_{i}|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|)X=X(I-\sum_{i=1}^{N}p_{i}|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|)=0,$
$z_{i}p_{i}=0,$
$Tr(X|\widetilde{\psi_{i}}\rangle\langle\widetilde{\psi_{i}}|)=z_{i}+\eta_{i},\eta_{i}\geq
0,\ i=1,...,N$ (5.27)
Then using KKT conditions one can show that
$\left(\begin{array}[]{cccc}1-p_{1}\tilde{a}_{11}&-p_{2}\tilde{a}_{12}&\cdots&-p_{N}\tilde{a}_{1N}\\\
-p_{1}\tilde{a}_{21}&1-p_{2}\tilde{a}_{22}&\cdots&-p_{N}\tilde{a}_{2N}\\\
\vdots&\vdots&\ddots&\vdots\\\
-p_{1}\tilde{a}_{N1}&-p_{2}\tilde{a}_{N2}&\cdots&1-p_{N}\tilde{a}_{NN}\\\
\end{array}\right)\times\left(\begin{array}[]{cccc}x_{11}&x_{12}&\cdots&x_{1N}\\\
x_{21}&x_{22}&\cdots&x_{2N}\\\ \vdots&\vdots&\ddots&\vdots\\\
x_{N1}&x_{N2}&\cdots&x_{NN}\\\ \end{array}\right)=0$ (5.28)
If $p_{i}\geq 0,(i=1,...,N)$ usin Eq.(5.27) we can conclude that $X$ is the
rank one matrix and we have
$x_{ij}=e^{i\varphi_{i}}e^{i\varphi_{j}}\sqrt{\eta}_{i}\sqrt{\eta_{j}}$ (5.29)
Consequently, the Eq. (5.28) can be written as
$\left(\begin{array}[]{cccc}1-p_{1}\tilde{a}_{11}&-p_{2}\tilde{a}_{12}&\cdots&-p_{N}\tilde{a}_{1N}\\\
-p_{1}\tilde{a_{21}}&1-p_{2}\tilde{a_{22}}&\cdots&-p_{N}\tilde{a_{2N}}\\\
\vdots&\vdots&\ddots&\vdots\\\
-p_{1}\tilde{a}_{N1}&-p_{2}\tilde{a}_{N2}&\cdots&1-p_{N}\tilde{a}_{NN}\\\
\end{array}\right)\times\left(\begin{array}[]{c}x_{11}\\\ x_{12}\\\ \vdots\\\
x_{1N}\\\ \end{array}\right)=0$ (5.30)
Or, equivalently,
$\left(\begin{array}[]{cccc}\tilde{a}_{11}&\tilde{a}_{12}&\cdots&\tilde{a}_{1N}\\\
\tilde{a}_{21}&\tilde{a}_{22}&\cdots&\tilde{a}_{2N}\\\
\vdots&\vdots&\ddots&\vdots\\\
\tilde{a}_{N1}&\tilde{a}_{N2}&\cdots&\tilde{a}_{NN}\\\
\end{array}\right)\times\left(\begin{array}[]{cccc}x_{11}&0&\cdots&0\\\
0&x_{22}&\cdots&0\\\ \vdots&\vdots&\ddots&\vdots\\\ 0&0&\cdots&x_{NN}\\\
\end{array}\right)\times\left(\begin{array}[]{c}p_{1}\\\ p_{2}\\\ \vdots\\\
p_{N1}\\\ \end{array}\right)=\left(\begin{array}[]{c}x_{11}\\\ x_{12}\\\
\vdots\\\ x_{1N}\\\ \end{array}\right)$ (5.31)
If the first condition in Eq. (5.27) is satisfied, using above equation one
can show that the optimal solution for $p_{i}$ is given by
$p_{i}=\frac{1}{e^{i\varphi_{i}}\sqrt{\eta_{i}}det(D)}\times\det\left(\begin{array}[]{ccccccc}\tilde{a}_{1,1}&\cdots&\tilde{a}_{1,i-1}&e^{i\varphi_{1}}\sqrt{\eta_{1}}&\tilde{a}_{1,i+1}&\cdots&\tilde{a}_{1,N}\\\
\tilde{a}_{2,1}&\cdots&\tilde{a}_{2,i-1}&e^{i\varphi_{2}}\sqrt{\eta_{N}}&\tilde{a}_{2,i+1}&\cdots&\tilde{a_{2,N}}\\\
\vdots&\ddots&\vdots&\vdots&\vdots&\ddots&\vdots\\\
\tilde{a}_{N,1}&\cdots&\tilde{a}_{N,i-1}&e^{i\varphi_{i}}\sqrt{\eta_{i}}&\tilde{a}_{N,i+1}&\cdots&\tilde{a}_{N,N}\\\
\end{array}\right)$ (5.32)
Using the fact that the matrix whose components are $\tilde{a}_{ij}$ is
inverse matrix of the one whose components are
$a_{ij}=\langle\psi_{i}|\psi_{j}\rangle$, we have
$\left(\begin{array}[]{c}x_{11}p_{1}\\\ x_{12}p_{2}\\\ \vdots\\\
x_{1N}p_{N}\\\
\end{array}\right)=\left(\begin{array}[]{cccc}a_{11}&a_{12}&\cdots&a_{1N}\\\
a_{21}&a_{22}&\cdots&a_{2N}\\\ \vdots&\vdots&\ddots&\vdots\\\
a_{N1}&a_{N2}&\cdots&a_{NN}\\\
\end{array}\right)\left(\begin{array}[]{c}x_{11}\\\ x_{12}\\\ \vdots\\\
x_{1N}\\\ \end{array}\right)\Rightarrow
x_{ii}p_{i}=\sum_{i=1}^{N}a_{ij}x_{ij}$
Using Eq. (5.29) we have
$p_{i}=\sum_{i=1}^{N}\frac{e^{i(\varphi_{i}+\varphi_{j})}}{e^{2i\varphi_{i}}}\sqrt{\frac{\eta_{i}}{\eta_{j}}}a_{ij}$
Then we can rewrite the probability $p_{i}$ in a simpler form
$p_{i}=e^{-i\varphi_{i}}\sum_{i=1}^{N}e^{i\varphi_{j}}\sqrt{\frac{\eta_{i}}{\eta_{j}}}a_{ij}$
(5.33)
where $e^{-i\varphi_{j}}a_{ij}$ is determined such that $p_{i}$ satisfy the
condition $0\leq p_{i}\leq 1$. Then the success probability with respect to
$\tilde{a}_{ij}$ components is given by
$P=-\frac{1}{\det(D)}\times\det\left(\begin{array}[]{cccc}0&e^{i\varphi_{1}}\sqrt{\eta_{1}}&\cdots&e^{i\varphi_{N}}\sqrt{\eta_{N}}\\\
e^{i\varphi_{1}}\sqrt{\eta_{1}}&&&\\\ \vdots&&D&\\\
e^{i\varphi_{N}}\sqrt{\eta_{N}}&&&\\\ \end{array}\right)$ (5.34)
The success probability with respect to $a_{ij}$ components is as follows
$P=\sum_{i=1}^{N}\eta_{i}p_{i}=\sum_{i,j=1}^{N}e^{i(\varphi_{i}-\varphi_{j})}\sqrt{\eta_{i}\eta_{j}}a_{ij}=\|\sum_{j=1}^{N}\sqrt{\eta_{j}}e^{i\varphi_{j}}|\psi_{j}\rangle\|^{2}$
(5.35)
If the first condition of KKT is not satisfied, for a specified i, $p_{i}$
does not lie in the feasible region $(p_{i}\leq 0\ or\ p_{i}\geq 1)$. In this
case, one can omit the jth row and jth column in the square $N\times N$
matrices and jth row in the row matrices in Eqs. (5.32) and (5.34) , and find
the optimal solutions with $p_{i}=0$. After this reduction, if any other
$p_{i}$ does not lie in the feasible region, the same procedure will be
repeated. Noted that a similar result was derived in Ref. [29] with another
method. The equations (5.33) and (5.35) gives an analytical relation between
the maximum average success probability and the $N$ pure linearly independent
quantum states to be discriminated. In general, finding an exact analytic
solution for the OPUSD problem involving an arbitrary number of pure linearly
independent quantum states is hard, since the explicit expressions of the
phases $e^{i\varphi_{j}}(j=1,...,N)$ are not given in equations (5.33) and
(5.35). However using (5.33) and (5.35) one can simplify the calculation of
the optimal solution in special cases and it may also drive some bounds for
the average success probability [18, 29]. Then, the approximated methods are
useful for unambiguous discrimination of N linearly independent quantum
states. Since we have presented an analytic relation for the feasible region
of N linearly independent quantum states, and this region is convex, then one
can easily obtain the optimal POVM by some well-known numerical methods such
as constrained linear or nonlinear least-squares, interior points, and simplex
and quadratic programming methods [30].
## 6 Conclusion
In conclusion, for the real and complex inner product of states, we have been
able to obtain an exact analytic solution for OPUSD problem involving an
arbitrary number of pure linearly independent quantum states by using KKT
convex optimization method. Moreover, Using semidefinite programming and
Karush-Kuhn-Tucker convex optimization method, we have been able to obtain an
analytical formula which shows the relation between optimal solution of
unambiguous discrimination problem and an arbitrary number of pure linearly
independent quantum states to be identified. Using this analytical formula one
can simplify the calculation of the optimal solution in special cases and it
may also drive some bounds for the average success probability.
## References
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* [7] L. M. Duan and G. C. Guo, Phys. Rev. Lett. 80, 4999 (1998).
* [8] Y. Sun, M. Hillery, and J. A. Bergou, Phys. Rev. A 64, 022311 (2001).
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* [10] M. X. Goemans and D. P. Williamson, J. Assoc. Comput. Mach. 42, 1115 (1995).
* [11] Z. Q. Luo, Math. Program. 97, 587 (2003).
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* [13] Wing-Kin Ma, T. N. Davidson, K. M. Wong, Z.-Q. Luo, and P. C. Ching, IEEE Trans. Signal Process. 50, 912 (2002).
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* [15] H. Barnum, M. Saks, and M. Szegedy, in Proceedings of the 18th IEEE Annual Conference on Computational Complexity (IEEE Computer Society, New York, 2003), pp. 179 193.
* [16] Y. C. Eldar, A. Megretski, and G. C. Verghese, IEEE Trans. Inf. Theory 49, 1007 (2003).
* [17] L. Ip, http://www.qcaustralia.org
* [18] M. A. Jafarizadeh, M. Rezaei, N. Karimi and A. R. Amiri, Phys. Rev. A 77, 042314(2008).
* [19] K. Kraus, States, Effects, and Operations, Lecture Notes in Physics No. 190 (Springer, Berlin, 1983).
* [20] A. Chefles, Phys. Lett. A 239, 339 (1998).
* [21] Y. C. Eldar, IEEE Trans. Inf. Theory 49, 446 (2003).
* [22] G. H. Golub and C. F. Van Loan, Matrix Computations, 3rd ed.(Johns Hopkins University Press, Baltimore, MD, 1996).
* [23] L. R. Welch, IEEE Trans. Inf. Theory 20, 397 (1974).
* [24] Y. C. Eldar, A. Megretski, and G. C. Verghese, IEEE Trans. Inf. Theory 49, 1007 (2003).
* [25] M. Lewenstein and A. Sanpera, Phys. Rev. Lett. 80, 2261(1998).
* [26] S. Karnas and M. Lewenstein, J. Phys. A 34, 6919 (2001).
* [27] M. A. Jafarizadeh, M. Mirzaee, M. Rezaee, Int. J. Quantum Inf. 2, 541 (2004).
* [28] M. A. Jafarizadeh, M. Mirzaee, and M. Rezaee, Quantum Inf. Process. 4, 199 (2005).
* [29] Shengshi Pang, Shengjun Wu, Phys. Rev. A 80, 052320(2009).
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|
arxiv-papers
| 2011-11-21T11:12:53 |
2024-09-04T02:49:24.543791
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "N. Karimi",
"submitter": "Nasser Karimi",
"url": "https://arxiv.org/abs/1111.4831"
}
|
1111.5037
|
# The TW Hya Disk at 870 $\mu$m: Comparison of CO and Dust Radial Structures
Sean M. Andrews11affiliation: Harvard-Smithsonian Center for Astrophysics, 60
Garden Street, Cambridge, MA 02138 , David J. Wilner11affiliation: Harvard-
Smithsonian Center for Astrophysics, 60 Garden Street, Cambridge, MA 02138 ,
A. M. Hughes22affiliation: University of California at Berkeley, Department of
Astronomy, 601 Campbell Hall, Berkeley, CA 94720 , Chunhua Qi11affiliation:
Harvard-Smithsonian Center for Astrophysics, 60 Garden Street, Cambridge, MA
02138 , Katherine A. Rosenfeld11affiliation: Harvard-Smithsonian Center for
Astrophysics, 60 Garden Street, Cambridge, MA 02138 , Karin I.
Öberg11affiliation: Harvard-Smithsonian Center for Astrophysics, 60 Garden
Street, Cambridge, MA 02138 33affiliation: Hubble Fellow , T.
Birnstiel44affiliation: Ludwig-Maximilians-Universität, University Observatory
Munich, Scheinerstrasse 1, D-81679 Munich, Germany , Catherine
Espaillat11affiliation: Harvard-Smithsonian Center for Astrophysics, 60 Garden
Street, Cambridge, MA 02138 55affiliation: NSF Astronomy & Astrophysics
Postdoctoral Fellow , Lucas A. Cieza66affiliation: University of Hawaii
Institute for Astronomy, 2680 Woodlawn Drive, Honolulu, HI 96822 , Jonathan P.
Williams66affiliation: University of Hawaii Institute for Astronomy, 2680
Woodlawn Drive, Honolulu, HI 96822 , Shin-Yi Lin77affiliation: Center for
Astrophysics & Space Science, University of California San Diego, La Jolla, CA
92093 , and Paul T. P. Ho11affiliation: Harvard-Smithsonian Center for
Astrophysics, 60 Garden Street, Cambridge, MA 02138 88affiliation: Academia
Sinica Institute of Astronomy and Astrophysics, P.O. Box 23-141, Taipei 106,
Taiwan
###### Abstract
We present high resolution ($0\farcs 3=16$ AU), high signal-to-noise ratio
Submillimeter Array observations of the 870 $\mu$m (345 GHz) continuum and CO
$J$=3$-$2 line emission from the protoplanetary disk around TW Hya. Using
continuum and line radiative transfer calculations, those data and the
multiwavelength spectral energy distribution are analyzed together in the
context of simple two-dimensional parametric disk structure models. Under the
assumptions of a radially invariant dust population and (vertically
integrated) gas-to-dust mass ratio, we are unable to simultaneously reproduce
the CO and dust observations with model structures that employ either a
single, distinct outer boundary or a smooth (exponential) taper at large
radii. Instead, we find that the distribution of millimeter-sized dust grains
in the TW Hya disk has a relatively sharp edge near 60 AU, contrary to the CO
emission (and optical/infrared scattered light) that extends to a much larger
radius of at least 215 AU. We discuss some possible explanations for the
observed radial distribution of millimeter-sized dust grains and the apparent
CO-dust size discrepancy, and suggest that they may be hallmarks of
substructure in the dust disk or natural signatures of the growth and radial
drift of solids that might be expected for disks around older pre-main
sequence stars like TW Hya.
circumstellar matter — protoplanetary disks — planetary systems: formation —
stars: individual (TW Hya)
## 1 Introduction
The physical conditions of the gas and dust in young circumstellar disks shape
the formation and early evolution of planetary systems. The spatial
distribution of disk material – the density structure – plays many fundamental
roles, and is especially relevant for dictating when and where planets can
form and migrate in the disk. Naturally, empirical constraints on
protoplanetary disk structures are of significant value in efforts to develop
realistic models of the complex processes involved in planet formation. Some
recent studies have made progress in extracting the radial (surface) density
profiles in these disks from high angular resolution measurements of their
millimeter-wave continuum emission (Andrews et al., 2009, 2010b; Isella et
al., 2009; Guilloteau et al., 2011). However, such observations are only
sensitive to the density structure of the (presumably) trace population of
disk solids, and not the gas that wields far more influence over the evolution
of disk structure and the formation of planets. The major challenge is that
the gas in these disks is primarily cold H2, which is not directly observable.
And while there have been many investigations of less abundant molecular
species (see Dutrey et al., 2007), limited sensitivity to low optical depth
lines and an incomplete understanding of the complicated chemistry in these
disks have so far made it difficult to convert observations of the gas into
density constraints (see Williams & Cieza, 2011).
Given those obstacles, there have not been many attempts to directly compare
the gas and dust structures of protoplanetary disks with spatially resolved
data. A few studies used simple models to fit the millimeter continuum (dust)
and CO line (gas) emission independently, finding inconsistent structures
where the dust is much more compact than the gas (Piétu et al., 2005; Isella
et al., 2007). Hughes et al. (2008) suggested that this discrepancy was likely
an optical depth illusion, an artifact of the assumed sharp outer edge in the
density model. They showed that a surface density profile with a smooth taper
at large radii could better reproduce the dust and gas observations. However,
Panić et al. (2009) argued that similar modifications to the model structure
of the IM Lup disk (see Pinte et al., 2008) were insufficient to account for
the much different sizes they measure for its gas and dust emission. Recently,
Qi et al. (2011) successfully matched their observations of multiple CO
transitions and continuum emission from the HD 163296 disk with a single
density model, although refinements to that model based on the detailed
morphology of the continuum emission were not a priority in that study. In
general, there is still substantial uncertainty that the gas and dust trace
the same structures in protoplanetary disks. From a practical standpoint, that
uncertainty is disconcerting. Our knowledge of the mass contents of these
disks is entirely based on their (easy to measure) dust emission, but we are
not sure if those dust structures also describe the gas reservoirs that
effectively control the key disk evolution and planet formation processes.
Although direct measurements of H2 densities are not possible, in principle an
indirect comparison of the dust and gas structures in a protoplanetary disk
can be made with sufficiently sensitive and high resolution observations of
the thermal continuum and emission lines from trace gas species. To meet those
requirements, a relatively massive and nearby disk would be an ideal candidate
target. TW Hya is an isolated, $\sim$0.8 M⊙ T Tauri star located only $54\pm
6$ pc from the Sun (Rucinski & Krautter, 1983; Wichmann et al., 1998; van
Leeuwen, 2007). Despite its advanced age ($\sim$10 Myr; Kastner et al., 1997;
Webb et al., 1999), it hosts a massive, gaseous accretion disk with a rich
molecular spectrum and strong continuum emission out to centimeter wavelengths
(Wilner et al., 2000, 2003, 2005; Qi et al., 2004, 2006, 2008). Those
observations and a suite of resolved optical/infrared scattered light images
confirm that the disk is well-resolved and viewed nearly face-on (Krist et
al., 2000; Trilling et al., 2001; Weinberger et al., 2002; Apai et al., 2004;
Roberge et al., 2005). The inner edge of this “transition” disk is truncated
at a radius of 4 AU, likely by an unseen (and maybe planetary) companion
(Calvet et al., 2002; Hughes et al., 2007). Given its proximity, favorable
orientation, and rich gas and dust content, the TW Hya disk is uniquely well-
suited for a comparative investigation of the radial behavior of its gas and
dust structures.
In this article, we present new sub-arcsecond resolution observations of the
870 $\mu$m continuum and CO $J$=3$-$2 emission from the TW Hya disk. Using
state-of-the-art radiative transfer modeling tools, we use those measurements
to compare the radial distributions of its CO gas and dust. In §2 we describe
the observations and data calibration. The basic characteristics of the data
are reviewed in §3. A detailed description of the modeling and results are
provided in §4. The implications of that analysis are discussed in §5, and our
principle conclusions are summarized in §6.
## 2 Observations and Data Reduction
TW Hya was observed at 345 GHz (870 $\mu$m) with the Submillimeter Array (SMA;
Ho et al., 2004) on 8 occasions in 2006-2010 in the compact (C), extended (E),
and very extended (V) array configurations, providing baseline lengths of
16-70 m, 28-226 m, and 68-509 m, respectively. To accomodate various studies
of the CO $J$=3$-$2 emission line (at 345.796 GHz), the SMA double sideband
receivers and correlator were configured with several different local
oscillator (LO) settings and spectral resolutions throughout this campaign.
The “default” setup used an LO frequency of 340.755 GHz (880 $\mu$m) with the
CO line in the center of the upper sideband, with a resolution of 0.70 km s-1
in one spectral chunk. The other 23 partially-overlapping 104 MHz spectral
chunks in each sideband were sampled coarsely with 32 channels each. Some of
the 2008 observations utilized a “high” resolution correlator mode, with a
large portion of the bandwidth devoted to probing the CO line with 0.044 km
s-1 channels (see Hughes et al., 2011). While those data sample the continuum
with the default spectral resolution, they have a reduced continuum bandwidth
(2.6 GHz; $\sim$70% of the available bandwidth at the time) and a higher LO
frequency at 349.935 GHz (857 $\mu$m). The 2010 observations were conducted in
a “medium” resolution mode that employed the new expanded bandwidth
capabilities at the SMA. In that case, two different 2 GHz IF bands were
centered $\pm$5 GHz (the normal setup) and $\pm$7 GHz from the LO frequency,
339.853 GHz (883 $\mu$m). The CO line was observed at a moderate resolution of
0.18 km s-1 in the upper sideband of the first IF band, and the continuum was
coarsely sampled across the remaining expanded bandwidth. Finally, a “hybrid”
mode was utilized at the end of 2006 to cover the $J$=4$-$3 transition of
H13CO+ (see Qi et al., 2008). A summary of the observational setups is
provided in Table 1.
Table 1: Summary of SMA Observations
UT Date | array | spectral | RMS noise | beam size | beam PA
---|---|---|---|---|---
| config. | setup | [mJy beam-1] | [″] | [°]
(1) | (2) | (3) | (4) | (5) | (6)
2006 Dec 28 | C | hybrid | 3.7 | $4.2\times 1.7$ | 4
2008 Jan 23 | E | high | 3.2 | $0.9\times 0.7$ | 21
2008 Feb 20 | E | high | 3.1 | $1.0\times 0.7$ | $-$20
2008 Feb 21 | E | default | 4.1 | $1.2\times 0.5$ | 14
2008 Mar 2 | C | high | 3.6 | $3.9\times 1.9$ | $-$17
2008 Mar 9 | C | default | 3.2 | $3.7\times 1.8$ | $-$13
2008 Apr 3 | V | default | 2.2 | $0.6\times 0.3$ | 19
2010 Feb 9 | V | medium | 1.8 | $0.5\times 0.3$ | $-$3
combined | all | $\cdots$ | 2.0 | $0.8\times 0.6$ | $-$3
Note. — Col. (1): UT date of observations. Col. (2): SMA array configuration,
C = compact, E = extended, and V = very extended. Col. (3): Adopted correlator
mode (see §2). Col. (4): Continuum RMS noise in a naturally-weighted map. Col.
(5): Naturally-weighted synthesized beam dimensions. Col. (6): Naturally-
weighted synthesized beam orientation (major axis position angle, measured
east of north).
TW Hya was observed in an alternating sequence with the nearby quasar
J1037-295 (a projected separation of 7.3°), with a total cycle time of 15
minutes in the C and E configurations and 8-10 minutes in the V configuration.
The quasars 3C 279 or J1146-289 were observed every other cycle. For the
tracks on 2008 February 21, March 9, and April 3, the target portion of the
cycle was shared with the nearby sources HD 98800 and Hen 3-600 (see Andrews
et al., 2010a). Planets, satellites, and bright quasars were observed as
bandpass and absolute flux calibrators when TW Hya was at low elevations
($<$18°), depending on their availability and the array configuration. Most of
these data were obtained in the best atmospheric conditions available at Mauna
Kea, with zenith opacities of 0.03-0.05 at 225 GHz (0.6-1.0 mm of precipitable
water vapor) and well-behaved phase variations on timescales longer than the
calibration cycle. The conditions on 2008 March 2 and 9 were worse, but
typical for Mauna Kea (with 1.4-1.8 mm of precipitable water vapor).
The data from each individual observation were reduced independently with the
IDL-based MIR software package. The bandpass response was calibrated with
observations of bright planets and quasars, and broadband continuum channels
in each sideband (and IF band, where applicable) were generated by averaging
and then combining the central 82 MHz in each line-free spectral chunk. The
visibility amplitude scale was derived from observations of Uranus, Titan,
Callisto, or Vesta, with a typical systematic uncertainty of 10-15%. The
antenna-based complex gain response of the system was determined with
reference to J1037-295. The observations of 3C 279 or J1146-289 were used to
assess the quality of phase transfer in the gain calibration process. Based on
those data, we estimate that the “seeing” generated by atmospheric phase noise
and any small baseline errors is small, 0.10-0$\farcs$15\. This result is a
testament to the high quality of the observing conditions, especially given
the low target elevation and wide projected separations between the
calibrators (for reference, 3C 279 is located 39.7° from TW Hya and 38.5° from
J1037-295, while the corresponding separations for J1146-289 are 11.0° and
15.1°, respectively).
Before the observations could be combined, we had to account for the proper
motion of TW Hya over the $\sim$3 year observing baseline
($\mu_{\alpha}\cos{\delta}=-0\farcs 066$ yr-1, $\mu_{\delta}=-0\farcs 014$
yr-1; van Leeuwen, 2007). Fortunately, each individual dataset exhibited
bright, symmetric, centrally-peaked continuum emission that we associate with
dust in the TW Hya disk. Using elliptical Gaussian fits to those continuum
visibilities, we determined centroid positions for each individual set of
observations and associated them with the stellar position at that epoch. The
measured emission centroids are consistent with the expected stellar
positions, well within the $\sim$0$\farcs$1 absolute astrometric accuracy of
the SMA (set primarily by small baseline uncertainties). Based on those
centroid measurements, each individual dataset was aligned to a common
coordinate system. After confirming that the aligned visibility sets showed
excellent agreement on all overlapping baselines, all datasets were combined
to produce composite 345 GHz continuum and CO $J$=3$-$2 visibilities.
The composite visibilities were Fourier inverted, deconvolved with the CLEAN
algorithm, and restored with a synthesized beam using the MIRIAD software
package. The natural weighting of the continuum data produces an image with a
$0\farcs 80\times 0\farcs 58$ beam and an RMS noise level of 2.0 mJy beam-1
(which is dynamic-range limited). Composite CO $J$=3$-$2 channel maps were
synthesized with a velocity resolution of 0.2 km s-1 and a circular beam with
a FWHM of 1$\farcs$0 (the naturally-weighted resolution is $0\farcs 96\times
0\farcs 73$). The RMS noise level is 0.13 Jy beam-1 in each channel.
Ancillary information in the literature was used to construct the TW Hya
spectral energy distribution (SED) that will be used with the SMA data. We
adopted the optical monitoring results of Mekkaden (1998) and near-infrared
photometry from Weintraub et al. (2000) and the 2MASS point source catalog
(Cutri et al., 2003). In the thermal infrared, Spitzer flux densities measured
by Hartmann et al. (2005) and Low et al. (2005) were supplemented with IRAS
and Herschel data (Weaver & Jones, 1992; Thi et al., 2010). We also include a
Spitzer IRS spectrum, kindly provided by E. Furlan (see Uchida et al., 2004).
At submillimeter wavelengths, we relied on the calibrator flux densities from
the SHARC-II111http://www.submm.caltech.edu/$∼$sharc/analysis/calibration.htm
and
SCUBA222http://www.jach.hawaii.edu/JCMT/continuum/calibration/sens/potentialcalibrators.html
instruments at the Caltech Submillimeter Observatory and James Clerk Maxwell
Telescope, respectively ($F_{\nu}=6.13\pm 0.68$, $3.9\pm 0.7$, and $1.37\pm
0.01$ Jy at 350, 443, and 869 $\mu$m, respectively). Additional millimeter-
wave measurements were collected from Weintraub et al. (1989), Qi et al.
(2004), and Wilner et al. (2003).
## 3 Results
The astrometrically aligned and combined SMA data are shown together in Figure
1, along with the broadband SED. The synthesized continuum image in Figure
1$a$ has an effective frequency of 344.4 GHz (870 $\mu$m), with an integrated
flux density of $1.34\pm 0.13$ Jy and a peak flux density of $0.337\pm 0.034$
Jy beam-1 (a peak signal-to-noise ratio of $\sim$170), including the
systematic calibration uncertainties. The continuum emission is regular and
symmetric on the angular scales probed here, with a synthesized beam size of
$43\times 31$ AU projected on the sky. Essentially no emission is detected
outside a $\sim$1″ radius. Given that the integrated flux density determined
from the SMA data is in good agreement with single-dish measurements (see §2;
Weintraub et al., 1989; Di Francesco et al., 2008), it is clear that this is
no optical depth effect: all of the millimeter-wave dust emission from the TW
Hya disk is concentrated inside a projected radius of $\sim$60 AU.
Figure 1: ($a$) Naturally weighted composite image of the 870 $\mu$m continuum
emission from the TW Hya disk. Contours start at 10 mJy beam-1 (5 $\sigma$)
and increase in 30 mJy beam-1 (15 $\sigma$) increments. The synthesized beam
dimensions are shown in the lower left. ($b$) Azimuthally averaged 870 $\mu$m
continuum visibility profile as a function of the deprojected baseline length
(real part only; the imaginary terms are effectively zero on all baselines).
The uncertainties are typically smaller than the symbol sizes. Note the low-
amplitude oscillations beyond $\sim$150 k$\lambda$. ($c$) Complete SED for the
TW Hya star+disk system (references are in the text; see §2). The Spitzer IRS
spectrum is shown as a thick gray curve. Our adopted stellar photosphere model
is overlaid as a thin gray curve (see §4.1). ($d$) CO $J$=3$-$2 channel maps
from the TW Hya disk, on the same angular scale as the continuum map in panel
$a$. Contours are drawn at 0.4 Jy beam-1 intervals ($\sim$3 $\sigma$) in each
0.2 km s-1-wide channel. The 1″ synthesized beam is shown in the lower left.
The central channel represents the TW Hya systemic velocity, at $V_{\rm
LSR}=+2.86$ km s-1.
While the continuum image appears rather plain, some interesting emission
features are apparent in a direct examination of the visibilities. Figure 1$b$
displays the azimuthally averaged 870 $\mu$m continuum real visibilities as a
function of the deprojected baseline length (see Andrews et al., 2009),
assuming the inclination ($i=6\arcdeg$) and major axis position angle (PA =
335°) derived from high spectral resolution CO emission line data by Hughes et
al. (2011). The imaginary visibilities are effectively zero within the noise
on all sampled baselines, confirming the accuracy of the astrometric alignment
and reinforcing that there are no obvious departures from axisymmetry in the
TW Hya disk. The visibility profile in Figure 1$b$ shows a smooth decrease out
to $\sim$150 k$\lambda$, followed by low-amplitude oscillations on longer
baselines and an apparent null near 500 k$\lambda$. These features are
distinct in independent datasets (particularly the “dip” near 180 k$\lambda$,
where 4 different observations span that range of baselines), are present
regardless of the bin sizes used for profile averaging, and are not noted in
the visibility profiles for the test calibrators (J1146-289 or 3C 279).
Despite the challenges of calibrating SMA data for low-elevation targets, the
persistence of these visibility modulations make us confident that they are
real features. Moreover, we will demonstrate in §4 that they can be reproduced
with models that incorporate a sharp outer edge in their emission profiles.
The null at 500 k$\lambda$ is consistent with the 4 AU-radius inner disk
cavity inferred from the SED (see Figure 1$c$; Calvet et al., 2002) and a VLA
7 mm continuum image (Hughes et al., 2007).
The panels in Figure 1$d$ show the composite CO $J$=3$-$2 emission line
channel maps for the TW Hya disk, resampled to a velocity resolution of 0.2 km
s-1 with a circular 1$\farcs$0 (54 AU) synthesized beam. The emission is
firmly detected ($>$3 $\sigma$) out to $\pm$1.2 km s-1 from the systemic
velocity ($V_{\rm LSR}=2.86$ km s-1), with an integrated intensity of $34.8\pm
3.5$ Jy km s-1 and a peak flux of $4.2\pm 0.4$ Jy beam-1 ($43\pm 4$ K),
including the calibration uncertainties. Those values are in good agreement
with previous single-dish and SMA measurements (van Zadelhoff et al., 2001; Qi
et al., 2004; Hughes et al., 2011). The channel maps in Figure 1$d$ show a
clear rotation pattern, from northwest (blueshifted) to southeast
(redshifted), with a narrow line-width due to a face-on viewing geometry. Near
the systemic velocity, the CO emission subtends $\sim$4″ (215 AU) in radius.
## 4 Modeling Analysis
These SMA observations offer some new insights into the TW Hya disk structure.
Naturally, as one of the nearest pre-main sequence stars, TW Hya and its
associated disk have been the subject of intense observational scrutiny. The
global structure of the TW Hya disk has been investigated previously, using
the SED (Calvet et al., 2002), optical/infrared scattered light observations
(Krist et al., 2000; Trilling et al., 2001; Weinberger et al., 2002; Apai et
al., 2004; Roberge et al., 2005), millimeter/radio-wave continuum images
(Wilner et al., 2000, 2003, 2005; Hughes et al., 2007), and resolved molecular
line maps (Qi et al., 2004, 2006, 2008; Hughes et al., 2008, 2011). However,
none of those previous studies had the combination of angular resolution and
sensitivity for the optically thin dust and high-quality gas tracers that are
available from the SMA datasets presented here.
In the following, a technique is described for extracting the structure of the
TW Hya disk from these data using radiative transfer models. We focus
specifically on enabling a comparison between the radial distributions of the
dust and CO tracers. The approach we have adopted to make that comparison
consists of three key steps. First, we construct a model of the radial density
structure of the dust that is able to reproduce our resolved observations of
870 $\mu$m continuum emission and the broadband SED. Next, we make the
assumption of a radially constant (vertically-integrated) dust-to-gas mass
ratio and use the model structure we derive from the dust to predict the
emission morphology of the CO $J$=3$-$2 line. Then, we show that this
assumption implies an inconsistency with the observations, highlighting a
clear difference in the radial distributions of millimeter-sized dust grains
and CO gas in the TW Hya disk. Some of the potential implications of this
inconsistency are discussed further in §5.
### 4.1 Dust Structure
The dust disk structure is determined following the technique outlined by
Andrews et al. (2011), with some modifications for generality. We assume the
dust is spatially distributed with a parametric two-dimensional density
structure in cylindrical-polar coordinates {$r$, $z$},
$\rho_{d}(r,z)=\frac{\Sigma_{d}}{\sqrt{2\pi}z_{d}}\exp{\left[-\frac{1}{2}\left(\frac{z}{z_{d}}\right)^{2}\right]},$
(1)
where $\Sigma_{d}$ and $z_{d}$ are surface densities and characteric heights,
which both vary radially (see below). As will be explained further in §4.3, we
investigated two different models for the radial surface density profile.
First, we employed the similarity solution for simple viscous accretion disk
structures (Lynden-Bell & Pringle, 1974) that we have used successfully to
characterize both normal and transition disks in the past (Andrews et al.,
2009, 2010a, 2010b, 2011; Hughes et al., 2010). In that case,
$\Sigma_{d}(r)=\Sigma_{c}\left(\frac{r}{r_{c}}\right)^{-\gamma}\exp{\left[-\left(\frac{r}{r_{c}}\right)^{2-\gamma}\right]},$
(2)
where $\Sigma_{c}$ is a normalization, $r_{c}$ is a characteristic scaling
radius, and $\gamma$ is a gradient parameter. As an alternative, we considered
a less physically motivated (but perhaps more commonly used) model that
incorporates a power-law density profile with a sharp cut-off (see Andrews et
al., 2008),
$\Sigma_{d}(r)=\Sigma_{0}\left(\frac{r}{r_{0}}\right)^{-p}\,\,\,\,\,\,\,({\rm
if}\,\,\,r\leq r_{0};\,\,{\rm else}\,\,\,\Sigma_{d}=0),$ (3)
where $\Sigma_{0}$ is a normalization, $r_{0}$ is the outer edge of the disk,
and $p$ is a gradient parameter. In either case, the surface densities at
small radii are modified to account for the TW Hya disk cavity (Calvet et al.,
2002; Hughes et al., 2007). To simplify the inner disk model of Andrews et al.
(2011), we set the surface densities to a constant value $\Sigma_{\rm in}$
between the sublimation radius ($r_{\rm sub}$) and a “gap” radius ($r_{\rm
gap}$). No dust is present between that gap radius and the cavity edge,
$r_{\rm cav}$. In the vertical dimension, the dust is distributed like a
Gaussian with a variance $z_{d}^{2}$. The characteristic height varies with
radius like $z_{d}=z_{0}(r/r_{0})^{1+\psi}$. Following Andrews et al. (2011),
we employ a cavity “wall” to reproduce the infrared spectrum of TW Hya (no
such feature was required at the sublimation radius). The local value of
$z_{d}$ is scaled up to $z_{\rm wall}$ at $r_{\rm cav}$, and then
exponentially joined to the global $z_{d}$ distribution over a small radial
width, $\Delta r_{\rm wall}$.
This structure model has 11 parameters: three describe the base surface
density profile, {$\Sigma_{c}$, $r_{c}$, $\gamma$} or {$\Sigma_{0}$, $r_{0}$,
$p$}, five determine the cavity and inner disk properties, {$\Sigma_{\rm in}$,
$r_{\rm sub}$, $r_{\rm gap}$, $r_{\rm cav}$, $\Delta r_{\rm wall}$}, and three
others characterize the vertical distribution of dust, {$z_{0}$, $z_{\rm
wall}$, $\psi$}. To simplify the modeling, we fixed some of the parameters
that are of less direct interest here. The sublimation radius was set to
$r_{\rm sub}=0.05$ AU, the location where dust temperatures reach 1400 K (see
also Eisner et al., 2006). The gap radius was set to $r_{\rm gap}=0.3$ AU and
the (constant) inner disk density to $\Sigma_{\rm in}=5\times 10^{-4}$ g cm-2.
The cavity edge was fixed at $r_{\rm cav}=4$ AU (see Hughes et al., 2007), the
wall height was set to $z_{\rm wall}=0.25$ AU, and the wall width to $\Delta
r_{\rm wall}=1$ AU. Since the details of this gap are not the focus, no
attempt was made to reconcile the models with infrared interferometric data
(but see Eisner et al., 2006; Ratzka et al., 2007; Akeson et al., 2011). After
extensive experimentation with modeling the SED, we also fixed the scale
height gradient to $\psi=0.25$. The interplay and degeneracies between these
free parameters were discussed in detail by Andrews et al. (2011). For our
purposes here, it is worth emphasizing that the parameters we have fixed have
little quantitative impact on the derived radial structures (i.e., sizes and
density gradients).
We used the dust composition advocated by Pollack et al. (1994), consisting of
a mixture of astronomical silicates, water ice, troilite, and organics with
the abundances, optical properties, and sublimation temperatures discussed by
D’Alessio et al. (2001). Based on the efforts of Uchida et al. (2004) to
faithfully reproduce the details of the Spitzer IRS spectrum, we let 25% of
the total silicate abundance inside the disk cavity ($r\leq r_{\rm cav}$) be
composed of crystalline forsterite (using optical constants from Jäger et al.,
2003). Two grain populations were employed, with a power-law size ($s$)
distribution, $n(s)\propto s^{-3.5}$, between $s_{\rm min}=0.005$ $\mu$m and a
given $s_{\rm max}$. Outside the cavity wall, 95% of the dust (by mass) has
$s_{\rm max}=1$ mm and the remaining 5% has $s_{\rm max}=1$ $\mu$m. The dust
in the wall itself and the tenuous inner disk was assumed to have $s_{\rm
max}=1$ $\mu$m. No effort was made to distinguish the vertical distributions
of these dust populations. Opacity spectra for each population were determined
from Mie calculations, assuming segregated spherical grains. For these dust
assumptions, the 870 $\mu$m dust opacity in the outer disk is $\kappa_{\rm
mm}=3.4$ cm2 g-1.
We assumed the central star has a K7 spectral type with $T_{\rm eff}=4110$ K,
$R_{\ast}=1.04$ R⊙, and $M_{\ast}=0.8$ M⊙ ($\log{g}=4.3$), based on an effort
to match Lejeune et al. (1997) spectral synthesis models to the broadband SED
and the detailed optical/infrared spectral analysis work of Yang et al.
(2005). The best-fit stellar spectrum template is overlaid on the SED in
Figure 1$d$ as a light gray curve. Recently, Vacca & Sandell (2011) have
argued instead for a M2.5 spectral type in the near-infrared, and a
correspondingly cooler stellar photosphere (3400 K), larger radius (1.29 R⊙),
and lower mass (0.4 M⊙). While that stellar model provides a good match to the
broadband infrared photometry for TW Hya, it underpredicts the observed
optical fluxes by a factor of $\sim$3 (in the $BVR$ bandpasses, and its known
variability does not bridge that gap; see Mekkaden, 1998). We prefer the
parameters for the warmer photosphere because they produce a template spectrum
that better matches the SED across the complete set of optical and infrared
bandpasses.
For a given set of parameters, we simulated the stellar irradiation and
emission output of a model dust structure using the two-dimensional,
axisymmetric Monte Carlo radiative transfer code RADMC (see Dullemond &
Dominik, 2004a). Assuming the fixed viewing geometry determined by Hughes et
al. (2011), a raytracing algorithm was then used to compute a synthetic model
SED and set of 870 $\mu$m continuum visibilities sampled at the same spatial
frequencies observed with the SMA. For each surface density model, we found
the best simultaneous fit to the observed SED and SMA visibilities over a
coarse grid of the gradient parameter $\gamma$ or $p$, by varying the
parameters {$\Sigma_{c}$ or $\Sigma_{0}$, $r_{c}$ or $r_{0}$, $z_{0}$}. Based
on those results, we refined our search and permitted the gradients to vary
freely to find the best-fit parameter sets for each model type (see §4.3 for
results).
### 4.2 CO Gas Structure
Unlike for the dust, the radial density profile of the gas disk cannot be
inferred directly from models of the optically thick $J$=3$-$2 transition of
CO. Therefore, we make a fundamental assumption that the gas traces the dust
in the radial dimension. For any given $\Sigma_{d}$, we define the gas surface
density profile as $\Sigma_{g}=\Sigma_{d}/\zeta$, where $\zeta$ is a
(radially) constant dust-to-gas mass ratio.
However, we have elected to permit some freedom in the vertical distribution
of the gas to facilitate a more faithful reproduction of the CO channel maps.
Using a multi-transition CO dataset, Qi et al. (2006) noted that models of the
TW Hya disk structure had a difficult time reproducing an appropriate vertical
temperature gradient of the gas. The intensity of the high-excitation
$J$=6$-$5 line indicated that the gas in the disk atmosphere was significantly
hotter than the dust, presumably due to substantial X-ray heating from the
central star. To be able to reproduce the observed CO spectral images, we have
characterized the vertical temperature profile of the gas in parametric form,
based on the modeling analysis of Dartois et al. (2003). We assume that
$T_{g}(r,z)=T_{a}+(T_{m}-T_{a})\cos{\left(\frac{\pi
z}{2z_{q}}\right)}^{2\delta},$ (4)
where $T_{a}=T_{1}(r/{\rm 1\,AU})^{-q}$ is a parametric radial temperature
profile in the disk atmosphere, $T_{m}$ is the midplane temperature determined
from the RADMC simulations of the dust, $\delta$ describes the shape of the
vertical profile, and $z_{q}$ defines the height of the atmosphere layer such
that $T_{g}(z\geq z_{q})=T_{a}$. In our modeling, we fix $\delta=2$ and
$z_{q}=4H_{p}$, where $H_{p}$ is the pressure scale height assuming the
midplane temperature at each radius ($H_{p}=c_{s}/\Omega$, the ratio of the
midplane sound speed to the Keplerian angular velocity). In practice, Eq. (4)
is a reasonable parametric approximation of the vertical temperature profile
in an irradiated disk (e.g., D’Alessio et al., 1999). We have grounded the
models by forcing $T_{g}=T_{d}$ at the midplane ($z=0$), but allowed the gas
temperatures to increase faster with height than the dust to simulate any
additional external heating sources. For a given $T_{g}(r,z)$ specified by
{$T_{1}$, $q$}, we then calculate the vertical density structure of the gas by
numerically integrating the equation of vertical hydrostatic equilibrium,
$\frac{\partial\ln{\rho_{g}}}{\partial
z}=-\left[\left(\frac{GM_{\ast}z}{(r^{2}+z^{2})^{3/2}}\right)\left(\frac{\mu
m_{H}}{kT_{g}}\right)+\frac{\partial\ln{T_{g}}}{\partial z}\right]$ (5)
using $\Sigma_{g}$ as a boundary condition, where $G$ is the gravitational
constant, $\mu=2.37$ is the mean molecular weight of the gas, $m_{H}$ is the
mass of a hydrogen atom, and $k$ is the Boltzmann constant. For reference, a
vertical slice of a representative model structure is shown in Figure 2.
Figure 2: Schematic demonstration of our model vertical structure, shown as a
vertical slice at a fixed radius. (top) The density profile of the gas (red)
and dust (blue) as a function of height above the midplane. The latter has a
parametric Gaussian distribution with a variance $z_{d}^{2}$, while the former
is computed assuming it is in hydrostatic pressure balance for its specified
temperature structure. The relative normalizations of each are represented
accurately, such that $\Sigma_{d}=\zeta\Sigma_{g}$, where $\zeta$ is a fixed
dust-to-gas mass ratio. The hatched region near the midplane ($z\leq z_{m}$)
marks where CO is depleted from the gas phase because it is frozen onto dust
grain mantles (where $T_{g}\leq T_{\rm frz}$). The comparable surface CO
depletion zone due to photodissociation ($z\geq z_{s}$) is well off the right
of the plot. (bottom) The corresponding temperature profiles, where $T_{d}$ is
computed from the RADMC radiative transfer calculations and $T_{g}$ is
determined parametrically, as described in the text. The gas and dust
temperatures are equivalent ($T_{g}=T_{d}=T_{m}$) in the midplane, but $T_{g}$
rises more rapidly than the dust before it saturates to a value $T_{a}$ at a
height $z_{q}$.
To quantify the density of CO molecules from any given gas model structure, we
adopt the layered approach of Qi et al. (2008, 2011). Based on the detailed
chemical calculations of Aikawa & Nomura (2006), we define two vertical
boundaries {$z_{m}$, $z_{s}$} at any given radius such that the CO mass
fraction (relative to H2) is $X_{\rm co}$ if $z_{m}(r)\leq z\leq z_{s}(r)$ and
$10^{-4}X_{\rm co}$ elsewhere. The “midplane” boundary, $z_{m}$, marks the
maximum height where CO molecules are expected to be frozen out of the gas
phase and affixed to dust grain mantles. In practice, $z_{m}$ is defined as
the minimum height where $T_{g}\geq T_{\rm frz}$, where the freeze-out
temperature $T_{\rm frz}$ is a parameter considered to be constant with
radius. The “surface” boundary, $z_{s}$, is meant to represent the height
where CO molecules can be photodissociated by X-rays or cosmic rays. Following
Qi et al. (2008), we define $z_{s}$ such that
$N_{\rm pd}=\frac{1}{\mu m_{H}}\int_{\infty}^{z_{s}}\rho_{g}\,\,dz,$ (6)
where $N_{\rm pd}$ is a vertically-integrated H2 column density that
effectively represents the penetration depth of the photodissociating
radiation field: $N_{\rm pd}$ is treated as a radially constant parameter.
So, for any dust model a corresponding CO model can be characterized with a
set of six additional parameters: two describe the abundances of dust and CO
relative to the total gas mass, {$\zeta$, $X_{\rm co}$}, two others
characterize the gas temperatures in the disk atmosphere, {$T_{1}$, $q$}, and
the last two define the spatial distribution of CO in the gas phase, {$T_{\rm
frz}$, $N_{\rm pd}$}. For a given set of these parameters, we generate a two-
dimensional grid of $n_{\rm co}(r,z)$ and $T_{g}(r,z)$ values and define a
velocity field based on Keplerian rotation around a point mass $M_{\ast}$,
assuming a minimal turbulent velocity line width of 10 m s-1 based on the
analysis of Hughes et al. (2011). We then feed that information into the
radiative transfer modeling code LIME (Brinch & Hogerheijde, 2010) to solve
the non-LTE molecular excitation conditions of the model and generate a
synthetic high-resolution model cube for the CO $J$=3$-$2 transition. That
model cube was then re-sampled at the velocity resolution of the data, and its
Fourier transform was sampled at the same spatial frequencies observed by the
SMA. In practice, we fixed the dust-to-gas ratio based on the assumed dust
composition, where $\zeta=0.014$ (a gas-to-dust mass ratio of 71, see Pollack
et al., 1994; D’Alessio et al., 2001). For each dust model, we varied
{$T_{1}$, $q$} for a coarse grid of CO abundance layer parameters {$X_{\rm
co}$, $T_{\rm frz}$, $N_{\rm pd}$} to find the best available match to the SMA
spectral visibilities. Since we are using only a single CO transition in this
investigation, the abundance layer parameters do not have a strong,
independent effect on the synthetic CO visibilities. For simplicity, we adopt
the best-fit models where $X_{\rm co}=2\times 10^{-6}$, $T_{\rm frz}=20$ K,
and $N_{\rm pd}=10^{21}$ cm-2 as representative.
Table 2: Estimated Model Parameters and Fit Results
Model | $\Sigma_{d}$(10) | $\gamma$, $p$ | $r_{c}$, $r_{0}$ | $z_{d}$(10) | $T_{a}$(10) | $q$ | $\chi_{\rm sed}^{2}$ | $\chi_{\rm cont}^{2}$ | $\chi_{\rm co}^{2}$
---|---|---|---|---|---|---|---|---|---
| [g cm-2] | | [AU] | [AU] | [K] | | | |
(1) | (2) | (3) | (4) | (5) | (6) | (7) | (8) | (9) | (10)
sA | 0.29 | 1.0 | 35 | 0.62 | 104 | 0.55 | 201 | 313,914 | 319,974
sB | 0.51 | 0.5 | 43 | 0.58 | 101 | 0.54 | 204 | 304,303 | 319,989
sC | 0.28 | 0.0 | 45 | 0.57 | 98 | 0.54 | 210 | 302,132 | 320,013
sD | 0.20 | -0.5 | 42 | 0.53 | 103 | 0.53 | 216 | 304,329 | 320,024
sE | 0.14 | -1.0 | 40 | 0.50 | 102 | 0.53 | 231 | 307,785 | 320,040
pA | 0.79 | 1.5 | 77 | 0.62 | 104 | 0.54 | 205 | 309,370 | 320,056
pB | 0.46 | 1.0 | 67 | 0.60 | 99 | 0.54 | 211 | 302,747 | 320,047
pC | 0.39 | 0.75 | 60 | 0.58 | 99 | 0.54 | 215 | 301,785 | 320,044
pD | 0.31 | 0.5 | 58 | 0.58 | 100 | 0.54 | 227 | 302,570 | 320,035
pE | 0.23 | 0.0 | 51 | 0.55 | 104 | 0.54 | 249 | 306,980 | 320,068
Note. — Col. (1): Model designation, where ‘s’ = similarity solution and ‘p’ =
power-law with sharp edge as defined in Eq. 2 and 3, respectively. Col. (2):
Dust surface density at $r=10$ AU. Col. (3): Surface density gradient. Col.
(4): Characteristic scaling radius (‘s’ models) or outer edge radius (‘p’
models). Col. (5): Characteristic dust height at $r=10$ AU. Col. (6): Gas
temperature in the disk atmosphere at $r=10$ AU (see Eq. 4). Col. (7): Radial
gradient of the atmosphere gas temperature profile. Col. (8): $\chi^{2}$
statistic for the SED (including Spitzer IRS spectrum). Col. (9): $\chi^{2}$
statistic for the 870 $\mu$m visibilities. Col. (10): $\chi^{2}$ statistic for
the CO $J$=3$-$2 spectral visibilities. The quoted parameter values and fit
results are valid for a set of additional fixed parameters, including: the
gradient of the dust height profile $\psi=0.25$, the sublimation radius
$r_{\rm sub}=0.05$ AU, the inner radius of the gap $r_{\rm gap}=0.3$ AU, the
cavity edge $r_{\rm cav}=4$ AU, the cavity wall height $z_{\rm wall}=0.25$ AU,
the gas temperature profile parameters $\delta=2$ and $z_{q}=4H_{p}$, the
dust-to-gas ratio $\zeta=0.014$, the CO/H2 abundance ratio $X_{\rm co}=2\times
10^{-6}$, the CO freezeout temperature $T_{\rm frz}=20$ K, and the CO
photodissociation column $N_{\rm pd}=10^{21}$ cm-2. Furthermore, we assume a
disk inclination of 6°, major axis position angle of 335°, and stellar mass of
0.8 M⊙.
Figure 3: (top) Best-fit gas surface densities (where
$\Sigma_{g}=\Sigma_{d}/\zeta$). Similarity solution models are shown on the
left and power-law + sharp edge models are shown on the right. (bottom)
Temperature profiles at the disk midplane ($T_{m}$) are shown as solid curves,
and were determined from RADMC radiative transfer calculations. The parametric
atmosphere temperatures ($T_{a}$; at a height $z_{q}$, see §4.2) are overlaid
as dashed curves. All models have a gap at $r=4$ AU, marked with a dotted gray
line. The CO freezeout temperature, $T_{\rm frz}=20$ K, is marked as a
horizontal gray line. The temperature profiles shown here are similar for all
models: the midplane values do not change much because the total irradiated
dust mass is roughly the same in each case, and the atmosphere temperatures
are determined from the same CO emission data. However, each model does have
slight variation in the vertical temperature profile, which is not displayed
here for the sake of clarity. Figure 4: A comparison of the model structures
in Table 2 with observations of the dust continuum emission from the TW Hya
disk, including the 870 $\mu$m visibility profile (top, with a zoomed-in view
of the emission on longer baselines in the middle) and SED (bottom). The
similarity solution surface density models (Eq. 2) are shown on the left, and
the power-law models with sharp outer edges (Eq. 3) are shown on the right.
The model SEDs are essentially indistinguishable, but there are clear
differences in the 870 $\mu$m radial emission profiles. The overall best match
to the continuum emission is Model pC, which has a density gradient $p=0.75$
and a sharp outer edge at $r_{0}=60$ AU (for the similarity solution models,
the best match is Model sC).
### 4.3 Modeling Results
The best-fit parameters for a range of representative surface density
gradients of each model type are compiled in Table 2. For clear notation, each
model is labeled with a letter designation corresponding to its $\gamma$ or
$p$ value, preceded by a lower-case ‘s’ for the similarity solution models
(based on Eq. 2) or ‘p’ for the power-law models (based on Eq. 3). The surface
densities, characteristic dust heights, and gas atmosphere temperatures are
listed for a radius of 10 AU, where the two surface density model types have
comparable behavior. The four parameters describing the dust densities,
{$\Sigma_{d}(10)$, $\gamma$ or $p$, $r_{c}$ or $r_{0}$, $z_{d}(10)$}, were
determined from joint fits to the SED and 870 $\mu$m continuum visibilities.
With those parameters fixed, the gas temperature parameters {$T_{a}(10)$, $q$}
were estimated from the CO $J$=3$-$2 visibilities only. The individual
$\chi^{2}$ values for the SED, continuum visibilities, and CO visibilities are
listed in Cols. (8)-(10). There are 77 independent datapoints in the SED
(including 45 equally-spaced points across the Spitzer IRS spectrum), 166,796
continuum visibilities, and 83,740 CO visibilities in each of 15 spectral
channels (a total of 1,256,100). The $\chi^{2}$ values were calculated with
weights that incorporate the quadrature sum of the formal uncertainties and
absolute calibration uncertainties on each datapoint (e.g., a 10% systematic
uncertainty on the amplitudes). Figure 3 shows the radial structures for each
of the disk models in Table 2, including their surface densities (top) and
temperature profiles (bottom).
Figure 4 directly compares the 870 $\mu$m visibilities and SEDs for these
model structures with the observations. All of the model structures provide
excellent fits to the broadband SED and Spitzer IRS spectrum. Individual SED
model behaviors are indistinguishable, aside from small deviations near the
far-infrared turnover where the measurement uncertainties are largest.
However, the different structures and model types exhibit distinctive
signatures in their resolved 870 $\mu$m emission profiles. The similarity
solution models with positive density gradients ($\gamma>0$; Models sA and sB)
substantially over-predict the emission on $\sim$100-200 k$\lambda$ scales,
while those with negative gradients ($\gamma<0$; Models sD and sE) have
visibility nulls at $\sim$150 k$\lambda$ that are clearly not commensurate
with the data. With an intermediate $\gamma=0$, Model sC provides the best
match to the continuum data for this model type. However, it and all other
similarity solution models fail to reproduce the visibility oscillations
beyond 200 k$\lambda$. This “ringing” is a classic sign of a sharp edge in the
radial emission profile, which is naturally produced with the power-law models
described by Eq. (3). For those structures, steep density gradients
($p=1.0$-1.5, Models pA and pB) produce too much emission on 100-200
k$\lambda$ baselines and shallow gradients ($p=0.0$-0.5, Models pD and pE)
generate visibility nulls that are not observed. However, an intermediate case
with $p=0.75$ (Model pC) is an excellent match to the data, with only a small
(albeit significant) departure on 280-380 k$\lambda$ scales. That mismatch is
likely related to the shape of the outer edge, although we have not pursued
that speculation further. Of all the dust structures explored here, Model pC
is clearly favored.
Figure 5: Moment maps of the CO $J$=3$-$2 emission from the TW Hya disk and
the various disk structure models compiled in Table 2. The leftmost panels
show the SMA observations. The top panels make a direct comparison with the
similarity solution models, and the bottom panels do the same for the power-
law models with sharp edges. In all plots, contours mark the velocity-
integrated CO intensities (0th moment) at 0.4 Jy km s-1 ($\sim$3 $\sigma$)
intervals and the color scale corresponds to the intensity-weighted line
velocities (1st moment). Only Model sA provides a good match to the observed
CO emission; all others predict gas distributions that are too small relative
to observations.
Figure 6: A direct comparison of the observed CO $J$=3$-$2 channel maps with
synthetic data from the similarity solution models in Table 2. The top block
of panels show the model spectral visibilities synthesized in the same way as
the data, and the bottom block displays the imaged residual visibilities. All
maps have the same contour intervals as in Figure 1$d$. The viewing geometry
of the TW Hya disk is marked with a gray cross in each channel map.
Figure 7: Same as Figure 6, but for the power-law models with sharp outer
edges.
A consistent feature of all the dust-based model structures is their
compactness. The similarity solution models require characteristic radii of
$r_{c}=35$-45 AU and the power-law models call for sharp outer edges at
$r_{0}=51$-77 AU. That compact dust distribution was noted in §3, where we
pointed out that all of the 870 $\mu$m dust continuum emission is concentrated
inside a radius of roughly 1″ ($\sim$60 AU; see Figure 1$a$). Given the much
larger radial extent of the CO $J$=3$-$2 emission – out to radii of at least
4″ (215 AU; see Figure 1$d$) – the best-fit models face major problems
reconciling the CO and dust observations. The moment maps in Figure 5 confirm
this tension between the dust-based structure models and CO data, highlighting
a clear CO-dust size discrepancy in nearly all cases. More direct comparisons
of the CO channel maps with each model structure can be made in Figures 6 and
7, which show both the models (top panels) and the imaged residual
visibilities (bottom panels) synthesized in the same way as the data. Thanks
to the freedom afforded by the parametric treatment of the gas temperatures,
all models successfully reproduce the central CO emission core that is
generated in the line wings. However, only Model sA has a sufficient amount of
mass at large disk radii to account for the observations near the systemic
velocity. Unfortunately, Model sA provides a poor match to the resolved 870
$\mu$m continuum emission profile.
Before investigating potential physical explanations for the apparent CO-dust
size discrepancy, it is important to evaluate the possibility that this is an
artifact of some assumption in the modeling (aside from the underlying
parametric models themselves; see §5). Admittedly, we adopted a relaxed
approach for treating the relative vertical distributions of gas and dust in
these models. Only two basic requirements were enforced for consistency: the
gas and dust are in thermal equilibrium at the disk midplane, and the gas is
distributed vertically according to hydrostatic pressure equilibrium. There
are more sophisticated treatments of the gas in disk atmospheres that
accurately account for the effects that X-ray heating, chemical
differentiation, and gas-dust temperature departures have on the vertical
structure of the disk (e.g., van Zadelhoff et al., 2001; Jonkheid et al.,
2004; Kamp & Dullemond, 2004; Kamp et al., 2010; Woitke et al., 2009, 2010;
Aresu et al., 2011). A nearby disk like TW Hya would benefit from a more
detailed analysis with such models, particularly using multiple resolved
molecular lines (see Thi et al., 2010, for a start). With a more physically
motivated treatment of the vertical structure, we might infer a quantitatively
different set of atmosphere properties that would affect the normalization and
shape of the CO emission profile (currently controlled by {$T_{1}$, $q$} and
other fixed parameters; see §4.2). However, the added complexity of those
models would not change the fact that a compact density distribution is
required to explain the dust emission. And since we assumed that the CO gas
traces the dust in the radial direction, the size discrepancy would remain:
these dust-based density profiles do not have enough material to be able to
excite CO emission lines at large disk radii (see Figure 3).
To summarize, the fundamental conclusion from the radiative transfer modeling
analysis of the SMA data is that the spatial distributions of the CO and
millimeter-sized dust in the TW Hya disk are different: the dust is
systematically inferred to be more compact than the CO gas.
## 5 Discussion
We used the SMA to observe the 870 $\mu$m continuum and CO $J$=3$-$2 line
emission from the disk around the nearby young star TW Hya. These data
represent the most sensitive high spatial resolution (down to scales of 16 AU)
probes of the millimeter-wave dust and molecular gas content in any
circumstellar disk to date. Along with the SED, these observations were used
in concert with continuum and line radiative transfer calculations in an
effort to extract the disk density structure. We were unable to identify a
consistent model structure that simultaneously accounts for the observed
radial distributions of CO and dust. Assuming a radially constant grain size
distribution and (vertically integrated) gas-to-dust mass ratio, the
millimeter-sized dust structure is significantly more compact than the CO. The
resolved continuum emission profile demonstrates that the radial distribution
of the millimeter-sized solids in the TW Hya disk has a relatively sharp outer
edge near 60 AU, which is considerably smaller than the observed extent of the
CO emission (out to at least 215 AU).
The TW Hya disk structure has been studied extensively with millimeter-wave
observations at lower angular resolution. In a series of investigations
focused almost exclusively on spectral line emission, Qi et al. (2004, 2006,
2008) and Hughes et al. (2011) relied on a slightly modified version of the
physical structure model that was developed by Calvet et al. (2002) to match
the TW Hya SED. While that model has been used quite successfully to explain
the molecular gas structure, there has always been some tension between
observations and its predicted millimeter/radio continuum emission profile
(see Qi et al., 2004; Wilner et al., 2003, 2005). Given the modest quality of
the previous continuum data and the fact that this structure model was not
designed (or fitted) with access to resolved observations of any kind, the
disagreement was understandably dismissed. As might be expected from its
assumption of a “steady” accretion disk structure (i.e., with a large,
positive surface density gradient), the Calvet et al. (2002) model exhibits
the same type of behavior as Models sA/B or pA/B: it agrees with the data on
large scales ($<$100 k$\lambda$), but significantly over-predicts the amount
of emission on smaller scales. The same is true for the parametric similarity
solution models (where $\gamma\approx 1$) explored by Hughes et al. (2008,
2011), and would certainly apply to the analogous models developed by Thi et
al. (2010) and Gorti et al. (2011). All of this previous modeling work
admirably reproduces the extended molecular line emission that is observed,
and in many cases the SED and millimeter-wave continuum emission on large
angular scales as well. It is only with the sensitive, high angular resolution
data presented here that we recognize a problem: the radial distribution of
the millimeter-sized dust grains is much more compact than for the CO.
Unlike the thermal radiation at millimeter wavelengths, the optical and near-
infrared light that scatters off small ($\leq$1 $\mu$m) grains in the TW Hya
disk surface is detected out to large distances from the central star – at
least $\sim$4″, comparable to what is inferred from CO spectral images (Krist
et al., 2000; Trilling et al., 2001; Weinberger et al., 2002; Apai et al.,
2004; Roberge et al., 2005). So some dust traces the molecular gas, even if it
is only a limited mass of small grains up in the disk atmosphere. However,
these exquisitely detailed scattered light images exhibit subtle structural
complexities. Krist et al. (2000) identified four distinct radial zones in the
optical scattered light disk, with a prominent steepening of the brightness
distribution just outside a radius of 50 AU ($\sim$1″; their Zone 1/2
boundary). Similar infrared behavior is noted in the studies by Weinberger et
al. (2002) and Apai et al. (2004), which both suggested a break in the
emission profile in the 50-80 AU ($\sim$1.0-1.5″) range. Those results were
confirmed in a comprehensive analysis of new data by Roberge et al. (2005),
who also called attention to a color change at a similar radius as well as an
azimuthal asymmetry out to a slightly larger distance from the star ($\sim$135
AU). The physical origin of these scattered light features has been a mystery,
although speculation centered around variations in the dust height (shadowing)
and gradients in the dust scattering properties (either mineralogical or size-
related). But in light of our discovery of an abrupt drop in the millimeter-
wave continuum emission at the same location as these features, it is only
natural to suspect that a more fundamental change occurs in the physical
structure of the TW Hya dust disk near 60 AU.
Perhaps the most straightforward explanation of the apparent CO-dust size
discrepancy inferred from the SMA data is that we have used an incomplete
description of the disk structure. As an example, consider a modification of
either model type that incorporates an abrupt decrease in the surface
densities (or millimeter-wave dust opacities) – not the dust-to-gas ratio –
outside $r\approx 60$ AU. If that drop in $\Sigma$ (or $\kappa_{\rm mm}$) was
not too large (perhaps a factor of $\sim$100), there would still be enough
disk material to produce bright emission from the optically thick CO lines and
scattered light while also accounting for the sharp edge feature noted in the
optically thin 870 $\mu$m emission profile. That “substructure” in the outer
dust disk might actually enhance the local gas-phase CO abundance, as
ultraviolet radiation can penetrate deeper into the disk interior and
photodesorb CO from the (small) reservoir of cold dust grains that remains at
large radii (e.g., Hersant et al., 2009).
Nevertheless, a physical origin for such a dramatic drop in the dust densities
and/or the millimeter-wave dust opacities is not obvious. One possibility is
that the disk has been perturbed by a long-period (as yet unseen) companion.
If a faint object is embedded in the disk near the apparent edge of the 870
$\mu$m emission distribution, it might open a gap that splits the disk into
two distinct reservoirs and generate the warp asymmetry suggested by Roberge
et al. (2005). But, a narrow gap alone would not account for the SMA continuum
observations. The millimeter-wave luminosity exterior to the gap would still
need to be decreased, perhaps because the particles at those larger radii were
preferentially unable to grow to millimeter sizes. Weinberger et al. (2002)
quote deep limits on $H$-band point sources in the TW Hya disk that suggest
there are no companions more massive than $\sim$6 MJup near 60 AU, according
to the Baraffe et al. (2003) models (the corresponding mass limit is higher
for the models of Marley et al., 2007). Certainly this kind of truncation or
other forms of substructure could be invoked to explain the sharp radial edge
in the SMA dust observations. But rather than engage in further speculation on
the details, it should suffice to point out that the potential for
substructure or other anomalies in the TW Hya disk can be tested with a
substantial increase in resolution and sensitivity. Moreover, spectral imaging
of optically thinner gas tracers (e.g., the CO isotopologues) would make for
an ideal test of the origins of the apparent CO-dust size discrepancy.
Fortunately, such observations will shortly be available as the Atacama Large
Millimeter Array (ALMA) begins routine science operations.
There is a compelling alternative explanation that has a more concrete
physical motivation. In any protoplanetary disk, the thermal pressure of the
gas is thought to cause it to orbit the star at slightly sub-Keplerian rates,
generating a small velocity difference relative to the particles embedded in
it (Weidenschilling, 1977). Depending on their size, particles can experience
a head-wind from this gas drag that decays their orbits and sends them
spiraling in toward the central star. This radial drift of particles modifies
the radial dust-to-gas mass ratio profile and introduces a pronounced spatial
gradient in the particle size distribution – in essence, it causes $\zeta$ and
$\kappa_{\rm mm}$ to decrease with radius (Brauer et al., 2007, 2008;
Birnstiel et al., 2009). In the outer disk, the drift rates are expected to be
largest for millimeter-sized particles (e.g., Takeuchi & Lin, 2002). Since
thermal emission peaks at a wavelength comparable to the particle size, the
drift process should then naturally produce a millimeter-wave emission profile
that is considerably more compact than would be inferred from tracers of the
gas (Takeuchi & Lin, 2005). Moreover, the dust particles that reflect light in
the disk atmosphere are small enough to be dynamically coupled to the gas –
therefore, scattered light images should be extended like the probes of the
gas phase.
From a qualitative perspective, our analysis of the CO and dust structures in
the TW Hya disk are certainly consistent with a scenario where growth and
radial drift have had an observable impact. However, it is unclear if
realistic models of the growth and migration of solids in a disk like this can
quantitatively account for the details. One particular challenge worth
highlighting is related to timescales. The Takeuchi & Lin (2005) models
suggest that millimeter-wave dust emission should be strongly attenuated on a
timescale of $\sim$1 Myr without constant replenishment (presumably from
growth and/or fragmentation). Given the advanced age of TW Hya ($\sim$8-20
Myr), this model would require either that replenishment shuts off after
several Myr or that drift is inefficient at early times before becoming more
important later in the disk evolution process. If this is indeed the mechanism
responsible for our results, additional observations of disks at a range of
ages could be used to help calibrate models of the long-term evolution of
drift rates. One potential way to test this hypothesis relies on the particle
size dependence of the radial drift rates. The theory implies that larger
particles will end up with more centrally concentrated density distributions.
At long and optically thin wavelengths, we expect that the size of the
continuum emission region should be anti-correlated with wavelength – longer
wavelengths imply more compact emission. In principle, this could be tested in
the near future by combining high resolution ALMA and Expanded Very Large
Array (EVLA) observations of the TW Hya dust disk that span the
millimeter/radio spectrum.
In many ways, TW Hya and its disk are unique and may not be representative of
the bulk population of pre-main sequence stars and their circumstellar
material. Nevertheless, it is worth considering the broader implications of
our findings for this specific example – they may prove to be more generally
applicable. It is possible that the CO-dust size discrepancy found here is
present in most other millimeter-wave disk observations, but it would likely
be difficult to identify with current sensitivity and resolution limitations.
If that feature is common and its underlying cause is a drop in the dust-to-
gas ratio in the outer disk, there would be serious consequences for disk mass
estimates based on dust continuum measurements. There could be large and
hidden mass reservoirs of molecular gas in the outer reaches of protoplanetary
disks, implying that disk masses might be substantially under-estimated. If
true, gas densities at large disk radii may be higher than typically assumed,
with profound implications for facilitating giant planet formation by
gravitational instability (e.g., Boley, 2009; Kratter et al., 2010; Boss,
2011). However, if radial drift is the key process responsible for the
apparent CO-dust size discrepancy, it is also possible that any “primordial”
dust-to-gas ratio integrated over the entire disk is preserved. This would be
the case if the millimeter-sized dust originally present at large radii had
its inward migration halted before it was accreted onto the central star. In
that scenario, the total disk mass estimates from millimeter-wave luminosities
would still be relatively accurate (assuming a proper model for the dust
opacities that considers the simultaneous particle size evolution is
available), although the densities in the outer disk would remain uncertain
without measurements of optically thin gas emission lines.
Independent of the apparent CO-dust size discrepancy, our finding that the
millimeter-wave continuum emission from the TW Hya disk is so sharply
truncated comes as a surprise. Modern interferometric datasets generally do
not have sufficient sensitivity to differentiate between dust density models
with sharp edges or smooth tapers (e.g., see Isella et al., 2010; Guilloteau
et al., 2011). Given that ambiguity, it is possible that the edge feature we
have identified in the 870 $\mu$m emission profile of the TW Hya disk could be
relatively common. Moreover, it is tempting to associate this $r\approx 60$ AU
edge in the distribution of larger solid particles in the TW Hya disk with the
similarly abrupt truncation of the classical Kuiper Belt at $r\approx 40$-50
AU in our solar system (Trujillo et al., 2001; Gladman et al., 2001). If this
behavior ends up being a generic feature in protoplanetary disks, it may
signify an important diagnostic of the radial migration of disk solids and
provide new insights into the structural origins and evolution of the outer
solar system (e.g., see Kenyon & Luu, 1999; Levison & Morbidelli, 2003).
Further speculation on the generality of the features identified in the TW Hya
disk is unnecessary. With the recent start of ALMA science operations, the
quality of the data presented here will be matched (and exceeded) routinely
for large samples, and the basic trends of disk properties like those probed
here will be clarified. If the TW Hya disk is not anomalous, it is clear that
the general methods used to interpret observations of dust in disks will need
to be modified to focus less on their bulk density structures and more on the
dynamical evolution of their solid contents.
## 6 Summary
We have presented sensitive, high resolution ($0\farcs 3=16$ AU) SMA
observations of the 870 $\mu$m continuum and CO $J$=3$-$2 line emission from
the disk around the nearby young star TW Hya. Based on two different
parametric formulations for the disk densities, we used radiative transfer
calculations to compare the predicted radial structures of the dust and CO gas
in the TW Hya disk with these SMA observations and ancillary measurements of
the spectral energy distribution. The key conclusions from this modeling
analysis of these high-quality data include:
1. 1.
Under the assumption that the dust-to-gas surface density ratio is constant
with radius, we were not able to find any model structure that can
simultaneously reproduce the resolved brightness profiles of the 870 $\mu$m
continuum and CO line emission. We have identified a clear CO-dust size
discrepancy that is present regardless of whether the assumed surface density
profile has a sharp outer edge or a smooth (exponential) taper at large radii.
2. 2.
The radial distribution of millimeter-sized dust grains in the TW Hya disk is
substantially more compact than its CO gas reservoir. The 870 $\mu$m dust
emission has a sharp outer edge near 60 AU, while the CO emission (and
optical/infrared scattered light from small grains that are dynamically
coupled to the gas) extends to a radius of at least 215 AU.
3. 3.
The observationally inferred CO-dust size discrepancy could potentially be
explained with a more complex dust density profile that exhibits a sudden
decrease by a large factor near a radius of 60 AU. That “break” in the density
profile might be consistent with a tidal perturbation by a long-period
(unseen) companion, although the dust at larger radii would then also have to
be preferentially depleted of millimeter-sized grains.
4. 4.
Alternatively, the observations might have uncovered some preliminary evidence
for a key evolutionary mechanism related to the planet formation process: the
growth and inward migration of disk solids. The radial drift of millimeter-
sized particles is expected to naturally concentrate long-wavelength thermal
emission near the star relative to tracers of the molecular gas reservoir.
However, a more detailed exploration of disk evolution models is needed to
verify if the observations of the TW Hya disk are quantitatively consistent
with this scenario.
We thank Mark Gurwell for his assistance with the multi-epoch calibration of
the SMA data, Elise Furlan and the Spitzer IRS Disks team for providing a
reduced TW Hya spectrum, and an anonymous referee for thoughtful suggestions
that contributed significant value to the article. We are especially grateful
to Christian Brinch for his generous help with the LIME code, Kees Dullemond
for technical support with the RADMC package, and Paola D’Alessio for her
valuable advice on dust populations and optical constants. The Submillimeter
Array (SMA) is a joint project between the Smithsonian Astrophysical
Observatory and the Academia Sinica Institute of Astronomy and Astrophysics
and is funded by the Smithsonian Institution and the Academia Sinica.
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|
arxiv-papers
| 2011-11-21T21:39:13 |
2024-09-04T02:49:24.555643
|
{
"license": "Public Domain",
"authors": "Sean M. Andrews, David J. Wilner, A. M. Hughes, Chunhua Qi, Katherine\n A. Rosenfeld, Karin I. Oberg, T. Birnstiel, Catherine Espaillat, Lucas A.\n Cieza, Jonathan P. Williams, Shin-Yi Lin, and Paul T. P. Ho",
"submitter": "Sean Andrews",
"url": "https://arxiv.org/abs/1111.5037"
}
|
1111.5042
|
# Evidence for a Universal Scaling of Length, Time and Energy in the Cuprate
High Temperature Superconductors
J.D. Rameau Z.-H. Pan H.-B. Yang G.D. Gu P.D. Johnson Brookhaven National
Laboratory, Upton, NY 11973, USA
###### Abstract
A microscopic scaling relation linking the normal and superconducting states
of the cuprates in the presence of a pseudogap is presented using Angle
Resolved Photoemission Spectroscopy. This scaling relation, complementary to
the bulk universal scaling relation embodied by Homes’ law, explicitly
connects the momentum dependent amplitude of the d-wave superconducting order
parameter at T$\sim 0$ to quasiparticle scattering mechanisms operative at
T$\gtrsim T_{c}$. The form of the scaling is proposed to be a consequence of
the Marginal Fermi Liquid phenomenology and the inherently strong dissipation
of the normal pseudogap state of the cuprates.
###### pacs:
Not long after the discovery of high temperature superconductivity in the
cuprates it was hypothesized that the transition temperature $T_{c}$ of these
materials might be governed by the onset of phase coherence amongst
“preformed” Cooper pairsEmery and Kivelson (1995); Emergy and Kivelson (1994).
This scenario, essentially postulating a form of Bose condensation of such
pairs, gives rise to a situation in which $T_{c}$ is lower than $T_{pair}$,
the temperature at which the pairing amplitude of the superconducting order
parameter develops. This point of view was bolstered early on by the
observation of Uemura et al.Uemura _et al._ (1989) that underdoped cuprates
obey a seemingly universal scaling law, $\rho_{s0}\propto T_{c}$, where
$\rho_{s0}$ is the superfluid density, or phase stiffness, at $T=0$, implying
that the mechanism for high $T_{c}$ superconductivity does indeed entail a
Bose condensation of well defined, preformed pairs rather than the traditional
BCS mechanism in which the pairing amplitude of the order parameter and global
phase coherence arise simultaneously. Recently however the Uemura relation was
shown to be accompanied by another universal scaling law, “Homes’ law”Homes
_et al._ (2004, 2005), stating that $\rho_{s0}\propto\sigma_{DC}(T_{c})T_{c}$
where $\sigma_{DC}(T_{c})$ is the DC optical conductivity at $T\gtrsim T_{c}$.
While Homes’ law is valid over a much wider swath of the cuprate phase diagram
than the Uemura relation, having been shown to apply to optimally and
overdoped materials as well as the underdoped variety and even the new Fe base
high $T_{c}$ superconductorsHomes _et al._ (2010), a transparent picture of
what it portends for the mechanism of high $T_{c}$ superconductivity in these
materials has yet to emerge.
In this Report it is shown that angle resolved photoemission spectroscopy
(ARPES) provides evidence for a complementary scaling relation between the
momentum dependent single particle scattering rates of carriers at $T\gtrsim
T_{c}$, at the Fermi energy $E_{F}$ on the Fermi surface (FS) ’arcs’, and the
magnitude of the superconducting gap at $T\sim 0$ K, respectively. This
finding, deriving from an examination of Homes’ lawHomes _et al._ (2004,
2005), extends and clarifies the microscopic origins of that relationship,
which was derived originally in the context of optical conductivity. As such,
the present work represents a long sought after correlation between
_microscopic_ spectral properties of the normal and superconducting states of
high $T_{c}$ materials.
Homes’ law has previously been interpreted as arising from a universal clean
limit superconductivity ($\xi_{0}\ll\ell_{TC}$ where $\xi_{0}$ is Pippard’s
coherence length at $T=0$ and $\ell_{TC}$ is the electronic mean free path at
$T\sim T_{c}$), universal dirty limit superconductivity Homes _et al._ (2005)
($\xi_{0}\geq\ell_{TC}$), “hard core” boson scattering Lindner and Auerbach
(2010) and as indicative of normal state cuprates obeying a quantum critical-
like relation of the form $k_{B}T_{c}\approx\hbar/\tau_{TC}$ where $\tau_{TC}$
is the mean free time (here in the sense of transport) of a normal state
electron at $T\sim T_{c}$Zaanen (2004)Abdel-Jawad _et al._ (2006). This
“Planckian” dissipation, viewable as a limit of the Marginal Fermi Liquid
(MFL) phenomenologyVarma _et al._ (1989), signifies that the observed
electronic scattering is as rapid as is causally allowed. Separately, it has
been suggested that Homes’ law implies $\ell_{TC}\approx 2\xi_{0}$Tallon _et
al._ (2006). Altogether this implies the central issue in distinguishing
various interpretations of Homes’ law rests upon understanding the ratio
$R=\ell_{TC}/\xi_{0}$, or equivalently, $R=\Delta_{0}\tau_{TC}$, a quantity
often used to quantify the strength of scattering in a superconductor relative
to the robustness of its pairing. To accomplish this in an ARPES experiment we
must take full account of the d-wave nature of the superconducting order
parameter and generalize from the coherence length and mean free path measured
in transport to momentum dependent quantities $\xi_{0}(\phi)$ and
$\ell_{TC}(\phi)$, $\phi$ being the FS angle as measured from the node
(defined in the inset of Fig.1e). While such a generalization must be treated
with care, especially near the node where $\xi_{0}(\phi)$ diverges, the result
is nonetheless phenomenologically simple.
$\xi_{0}(\phi)$ can be measured in ARPES assuming a generalization of the
coherence length, $\xi_{0}(\phi)=\frac{\hbar
v_{F}^{0}(\phi)}{\pi\Delta_{0}(\phi)}$, where $v_{F}^{0}(\phi)$ and
$\Delta_{0}(\phi)$ are the momentum dependent bare Fermi velocity at $T\sim
T_{c}$ and the anisotropic superconducting gap at $T\sim 0$, respectively.
Similarly $\ell_{TC}(\phi)=1/\Delta k_{TC}(\phi)$ is the momentum dependent
mean free path measured at $T\sim T_{c}$ with $\Delta k_{TC}(\phi)$ being the
Lorentzian full width at half maximum of the momentum distribution curve (MDC)
at $E=E_{F}$Valla _et al._ (1999). Noting that $\hbar v_{F}^{0}(\phi)\Delta
k(\phi)=\hbar/\tau_{TC}(\phi)=\Gamma_{TC}(\phi)=2\Im\Sigma_{TC}$, where
$\Im\Sigma$ is the imaginary part of the electron self energy, we find that
the quantity of interest from the point of view of ARPES is
$R(\phi)=\pi\Delta_{0}(\phi)/\Gamma_{TC}(\phi)$ where we recall that the
inverse quasiparticle (QP) lifetime $\Gamma_{TC}(\phi)$ is the full width at
half maximum of the peak in the ARPES energy distribution curve (EDC) line
shape. Expressing $R(\phi)$ in terms of energy rather than length scales via
the MDC equation has the advantage of obviating the need to infer the bare
Fermi velocity from measurement.
It is evident from our definition of $R(\phi)$ that it can only have meaning
when measured over the Fermi arc, understood to be the visible side of a nodal
hole pocketYang _et al._ (2011) as it exists at $T_{c}$, because
$\Gamma_{TC}(\phi)$ is formally undefined for $\phi>\phi_{c}$ (where
$\phi_{c}$ is the FS angle of the arc tip) due to the presence of the
pseudogap (PG) at $E_{F}$ in optimal, underdoped and some lightly overdoped
materials. Taken altogether the program for examining the quantity $R(\phi)$
using ARPES is to measure the QP lifetime at $E_{F}$ from the node to the tip
of the Fermi arc, $\phi_{c}$, in the normal state at $T\gtrsim T_{c}$ (which
are the states probed in transport experiments) and to compare these to
measurements of the momentum dependent superconducting gap at the same points
in k-space, for the same samples, at low temperature.
Figure 1: (Color online) Raw ARPES spectra of OP91 a) and UD70 b) at the tip
of the Fermi arc, $\phi_{c}$, in the normal state just above $T_{c}$. Low
temperature spectra ($T=15$ K, after deconvolution, are shown for the OP91
sample in c) and the UD70 sample in d). e) and f) show the $T_{c}$ and $T=15$
K symmetrized EDC’s for OP91 and UD70, respectively, as well as Lorentzian
fits thereof. The energy scale of the low temparture EDC’s in panels e) and f)
have been scaled up by a factor of $\pi$ to better illustrate Eq. 1 evaluated
at $\phi=\phi_{c}$ and the intensity scaled to half that of the normal state
EDC’s. The inset of panel e) illustrates the Fermi surface angle $\phi$ in
relation to the normal state Fermi surface and the cuts represented by the
current experiment (solid lines).
Figure 2: (Color online) a) Superconducting gap versus Fermi surface angle at
T=15 K for UD70 (black squares) and OP91 (red circles). The gap is measured as
the distance of the coherence peak in the superconducting state to $E_{F}$. b)
Inverse lifetime at $T\sim T_{c}$ versus Fermi surface angle for UD70 (black
squares) and OP91 (red circles). Fits described in the text are shown as
dotted lines extrapolated towards the antinodal regionPushp _et al._ (2009).
All error bars are derived from the fits.
The experiment described above was carried out at beamline U13UB of the NSLS
with a Scienta-2002 electron spectrometer. The end station was equipped with
an open flow Helium cryostat. Two samples, high quality single crystals of
Bi2Sr2CaCu2O8+x (Bi2212) grown by the floating zone method, were used for
measurements of $R(\phi)$: an optimally doped sample with $T_{c}$=91 K (OP91)
and an underdoped sample with $T_{c}$=70 K (UD70). The lower $T_{c}$ of the
UD70 sample was achieved by annealing in vacuum at 500 C for two days.
Transition temperatures for both samples were ascertained prior to the ARPES
measurement by SQUID magnetometry. The OP91 and UD70 samples were measured in
their normal states at $T=95$ K and $T=70+$ K, respectively, and in their
superconducting states at T=15 K. The overall resolution of the experiment was
set to 12.5 meV in energy and $0.1^{o}$ in angle. The photon energies used
were 18 eV (OP91) and 17.46 eV (UD70). All measurements were acquired within
48 hours of cleaving the samples at a chamber pressure of $1\times 10^{-10}$
Torr.
Low temperature spectra used for acquiring $\Delta_{0}(\phi)$ were resolution
corrected using the Lucy-Richardson methodRameau _et al._ (2010)Yang _et
al._ (2008)Plumb _et al._ (2010) of deconvolution, yielding an effective
energy resolution of 4 meV. Values for the superconducting gap were determined
by the binding energy of the coherence peak. Raw normal state and deconvolved
low temperature data for $\phi=\phi_{c}$ are shown in Fig. 1a)-d). Values of
$\Delta_{0}(\phi)$ for UD70 and OP91 are presented in Fig. 2a along with pure
d-wave ($\Delta(\phi)=\Delta_{0}^{AN}\sin(2\phi)$) fits, $\Delta_{0}^{AN}$
being determined by extrapolation from the nodal regionPushp _et al._ (2009).
Normal state values for $\Gamma_{TC}(\phi)$ were acquired by fitting
Lorentzians on a linear background to spectra symmetrized about $E_{F}$.
Strictly speaking this procedure is only valid for states residing at $E_{F}$
and $k_{F}$ on the Fermi arc in the normal PG state of the copper oxides due
to the presence of the previously observed particle-hole asymmetry in the
nodal regionYang _et al._ (2008). This is, by design, where the measurement
is carried out. It has similarly been affirmed that the superconducting state
is particle-hole symmetric well below $T_{c}$Yang _et al._ (2008)Lee _et
al._ (2007) so that symmetrization for our purposes is allowed at low
temperature. The angular dependence of the inverse lifetimes, presented in
Fig. 2b, were fit with the “offset” d-wave Abdel-Jawad _et al._ (2006)
$\Gamma_{TC}(\phi)=\Gamma^{N}_{TC}+\delta\Gamma_{TC}\sin(2\phi)$ where
$\delta\Gamma_{TC}=\Gamma_{TC}^{AN}-\Gamma_{TC}^{N}$, $\Gamma_{TC}^{N}$ and
$\Gamma_{TC}^{AN}$ being the nodal and antinodal inverse lifetimes at $T_{c}$,
respectively. $R(\phi)$, extracted from the data in Fig. 2, is plotted in Fig.
3 along with analytical fits. We emphasize that $R(\phi)$ is obtained using
total scattering rate (at each point in k-space) rather than just the
anistropic component, as is appropriate to a comparison with the DC optical
conductivity used in obtaining Homes’ law.
Figure 3: (Color online) a) Plot of $R(\phi)$ for OP91 (red circles) and UD70
Bi2212 (black squares). Error bars are formally propagated from those in Fig.
2. $R(\phi)$ obtained from the fits is superimposed on the data.
Figure 4: (Color online) Summary of available ARPES data from the current
study as well as Ref. Kondo _et al._ (2009) for $R(\psi)$. $\phi_{c}$ is
taken in every case to simply be the last point to which a lifetime at the
Fermi level could be reasonably ascertained. Error bars in the abscissa have
been suppressed.
To augment the present experimental results we reanalyzed data from a previous
ARPES experiment performed on the single layer Bi2201 system under similar
experimental conditionsKondo _et al._ (2009). In Fig. 4, which constitutes
our main finding, the quantity $R(\psi)=\pi\Delta_{0}(\psi)/\Gamma_{TC}(\psi)$
versus the reduced Fermi surface angle $\psi\equiv\phi/\phi_{c}$ is plotted.
Plotting $R$ in terms of $\psi$, which ranges between 0 at the node and 1 at
the arc tip, rather than $\phi$, serves the purpose of collapsing data from
samples with varying $T_{c}$ onto an equal footing. Remarkably, though the
$T_{c}$’s of the materials thus investigated range between 25 K and 91 K,
including single and bilayer systems, all materials are found to be very well
approximated by a simple expression:
$\Delta_{0}(\psi)\cong\frac{\psi\Gamma_{TC}(\psi)}{2\pi}.$ (1)
While Eq. 1 might be taken as a purely phenomenological expression it can be
related to the important length scales of the system, allowing our previous
derivation, such that
$\xi_{0}(\psi)\cong 2\psi^{-1}\ell_{TC}(\psi).$ (2)
Eq. 2 must be treated carefully in order to avoid divergence at the node.
Superconductors well described by the BCS theory have long been characterized
as being in a “clean” or “dirty” limit based on comparisons of the type
represented by Eqs. 1 and 2DeGennes (1999). What such distinctions mean for
superconductors possessing anisotropic order parameters is far from clear. In
the present case those terms should evidently be eschewed because while the
ratio of $\xi_{0}$ to $\ell_{TC}$, for example, is clearly a useful metric for
parameterizing BCS superconductors there is no evidence of, or prescription
for, a universal relationship between these quantities for a generic system as
there is, say, between $T_{c}$ and $\Delta_{0}$. The present findings thus
constitute evidence of a fundamental physical process in the cuprate
superconductors that is not an obvious consequence of the BCS theory.
We postulate that a simple explanation for this behavior can be found by
invoking the MFL phenomenology at $T_{c}$, $\Gamma_{TC}\propto T_{c}$, and the
BCS superfluid, $\Delta_{0}\propto T_{c}$. If both of these properties hold
across the Fermi arc then it is natural to conclude that
$\Delta_{0}\propto\Gamma_{TC}$ will also hold across the fermi arc. Indeed,
there is mounting evidence from ARPES and Scanning Tunneling Microscopy (STM)
that $\Delta_{0}(\phi_{c})\propto T_{c}$ Pushp _et al._ (2009)Kurosawa _et
al._ (2010)Lee _et al._ (2007) and it was shown long ago that the MFL
phenomenology is maintained in the ARPES spectrum of Bi2212Valla _et al._
(1999)Valla _et al._ (2000). Transport studies have also repeatedly reported
observations of a correlation between the anisotropic dissipation of the
normal state and $T_{c}$ Abdel-Jawad _et al._ (2006)Zaanen (2004)Homes _et
al._ (2005) in the cuprates. It has also been shown recently that spin
fluctuations, a leading candidate for the pairing mechanism of the curpates,
leads (at least in some cases) to a T-linear scattering rate in the normal
pseudogap stateJin _et al._ (2011). Regardless, the physical content of Eq. 1
is to imply that the interaction responsible for the anomalous non-Fermi
liquid normal state scattering rate is intimately related to the interaction
that gives rise to the pairing strength observed as a single particle gap on
the Fermi arcs below $T_{c}$. It is hard to escape this conclusion given the
ultimate proportionality between the superconducting gap at low $T$ and the
imaginary part of the self energy at $T_{c}$ reported here.
The link between the pairing interaction and the electronic dissipation at
$T_{c}$ has been remarked upon previously Basov and Chubukov (2011), and
indeed is explicitly predicted within the MFL phenomenologyAbrahams and Varma
(2000), though the existence of such a direct experimental relationship
between the two, on the microscopic level, has not to our knowledge been
previously reported. Eq. 2 offers a more intuitive, real space picture of what
the maximal dissipation of the normal state as a function of $T$ implies for
superconductivity. Evidently, even if carriers were to experience a strong
pairing interaction well above $T_{c}$, true pairs could not arise on the
Fermi arc because the constituents rescatter before that information can be
coherently propagated to a mate. $T_{c}$ appears to occur when all the
carriers on the Fermi arc have a pairing amplitude _and_ can propagate that
information, implying the relevance of an “intra-pair” phase coherence to the
magnitude of $T_{c}$. This length scale, $\xi_{0}$, set by the pairing
strength plays a role fundamentally different than in BCS materials. If the
scattering length of a single particle is too short relative to the size of
the pairing potential it’s in, it won’t sense that potential. Such a
dependence of the phase transition temperature $T_{c}$ purely on the relevant
length and time scales of the system, rather than the details of the
interaction, is the essence of quantum critical phenomenaSachdev (1999). Eq. 2
shows that in the cuprates, the new longer length scale is introduced a priori
by the pairing interaction. Additionally, all carriers on the Fermi arc must
be able to pair coherently before a gap can open, otherwise the symmetry of
the d-wave order parameter would be violated.
Finally, we note that states at the Fermi arc tip appear to play a unique role
in the phenomenology of the cuprates, on par with the high symmetry points of
the nodal and antinodal states. There, Eqs. 1 and 2 reduce to
$\Delta_{0}(\phi_{c})=\Gamma_{TC}(\phi_{c})/2\pi$ and
$\xi_{0}(\phi_{c})=2\ell_{TC}(\phi_{c})$, respectively. This scaling, in
relation to Eq. 1, is illustrated graphically in Fig. 1e)-f). That PG states
at higher momenta cannot satisfy this conditional relationship highlights a
fundamental, if subtle, difference between the nodal and antinodal region of
the Brillouin zone and suggests an inability of carriers tied up in the PG
state above $T_{c}$ to ever condense into a true superfluid state.
Our findings evidence a universal, microscopic scaling relation between two
fundamental properties of the cuprate superconductors: the normal state
lifetime of carriers on the Fermi surface and the superconducting gap that
arises from those states well below $T_{c}$. This relationship represents a
clear departure from BCS theory by itself, yet suggests several key concepts
of the BCS superfluid survive the quantum critical nature of the cuprates’
anomalous normal state. The transition temperature is shown to be governed by
a competition between length and time scales - pairing and single particle -
both of which appear to be modulated by the same interaction. Lastly we note
that the successful application of Homes’ law to the pnictide and chalcogenide
superconductors raises the intriguing possibility of performing ARPES
experiments similar to those presented here in those systems.
###### Acknowledgements.
We would like to acknowledge illuminating conversations with C.C. Homes, A.M.
Tsvelik, M. Khodas and T.M. Rice. We would also like to acknowledge that the
original inspiration for this work arose from conversations with Myron
Strongin. This work is supported by the US DOE under Contract No. DE-
AC02-98CH10886 and by the Center for Emergent Superconductivity, an Energy
Frontier Research Consortium supported by the Office of Basic Energy Science
of the Department of Energy.
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|
arxiv-papers
| 2011-11-21T21:54:11 |
2024-09-04T02:49:24.568199
|
{
"license": "Public Domain",
"authors": "J. D. Rameau, Z.-H. Pan, H.-B. Yang, G. D. Gu and P. D. Johnson",
"submitter": "Jonathan Rameau",
"url": "https://arxiv.org/abs/1111.5042"
}
|
1111.5254
|
# Markov Chains application to the financial-economic time series prediction
Vladimir Soloviev vnsoloviev@rambler.ru Cherkasy National University named
after B. Khmelnitsky, Cherkassy, Ukraine Vladimir Saptsin
saptsin@sat.poltava.ua Kremenchug National University named after M.
Ostrogradskii, Kremenchuk, Ukraine Dmitry Chabanenko chdn6026@mail.ru
Cherkasy National University named after B. Khmelnitsky, Ukraine
###### Abstract
In this research the technology of complex Markov chains is applied to predict
financial time series. The main distinction of complex or high-order Markov
Chains and simple first-order ones is the existing of aftereffect or memory.
The technology proposes prediction with the hierarchy of time discretization
intervals and splicing procedure for the prediction results at the different
frequency levels to the single prediction output time series. The hierarchy of
time discretizations gives a possibility to use fractal properties of the
given time series to make prediction on the different frequencies of the
series. The prediction results for world’s stock market indices is presented.
Prediction, time series, complex Markov chains, discrete time, fractal
properties, discrete Fourier prediction.
###### Contents
1. 1 Introduction
2. 2 Analysis of prominent publications relevant to the subject
3. 3 Aims of the paper, problem statement
4. 4 Classical modeling problems of ESE systems dynamics
5. 5 Modern concepts in ESE systems modeling
6. 6 Markov chains prediction technology
7. 7 Prediction construction algorithm
8. 8 States in complex Markov chains and approaches for defining them
9. 9 Step-by-step prediction procedure. Defining the most probable state on the next step, prediction scenarios
10. 10 Time increments hierarchy and splicing procedure
11. 11 Results of stock indices prediction
12. 12 Conclusions and further work
## 1 Introduction
Successful modeling and prediction of processes peculiar to complex systems,
such as ecological, social, and economical (ESE) ones, remain one of the most
relevant problems as applied to the whole complex of natural, human and social
sciences (r001SamarskyMihaylov01 ; IvakhnenkoMGUA68 ;
r002BogoboyashyKurbanov04 ; AndersenGluzmanSornette2000 ;
PincakCurrencyStringTheory2011 ; PincakStringPrediction2011 ).
The diversity of possible approaches to modeling such systems and, usually,
more than modest success in the dynamics prediction, compel us to look for the
reasons of failure, finding them not only in details, but also in the
axiomatics, which relates to problem statement, chosen modeling methods,
results interpretation, connections with other scientific directions.
With the appearance of quantum mechanics and relativity theory in early
twentieth century new philosophical ideas on physical values, measuring
procedures and system state have been established, the ones that are
completely different from Newtonian notions r003ElutinKrivchenko76 ;
r004LandauQuantNerelyativistic .
For more than 70 years basic concepts of classical and neoclassical economic
theories have been discussed by leading scientists, generating new approaches
r005SapirEconTheoryNeodnorSyst . The general systems theory has acquired
recognition in the middle of the 20th century giving way to development of the
new, systemic, emergent, and quantum in essence approach to investigation of
complex objects, which postulates the limited nature of any kind of modeling
and is based upon fixed and closed system of axioms r006BertalanfyGenSyst62 .
However, the development of this new philosophical basis of ESE systems
modeling is still accompanied with numerous difficulties, and new principles
are often merely declared.
Current research is devoted to investigation and application of the new
modeling and prediction technology, suggested in r007KurbanovSaptsin07 ;
r008SaptsinMarkovPaper09_ , based on concepts of determined chaos, complex
Markov chains and hierarchic (in terms of time scale) organization of
calculating procedures.
## 2 Analysis of prominent publications relevant to the subject
Prediction of financial-economic time series is an extremely urgent task.
Modern approaches to the problem can be characterized by the following
directions: 1) approximation of a time series using an analytical function and
extrapolation of the derived function towards future – so-called trend models
r009LukashinAdapt03 ; 2) investigation of the possible influence various
factors might have on the index, which is being predicted, as well as
development of econometric or more complicated models using the Group Method
of Data Handling (GMDH) IvakhnenkoMGUA68 ; r010ZaychenkoMonogr08 ; 3) modeling
future prices as the decisions-making results using neuronal networks, genetic
algorithms, fuzzy sets r010ZaychenkoMonogr08 ; r011EzhovShumsky98 ;
r012Zaencev .
Unfortunately, these techniques don’t produce stable forecasts, what can be
explained by complexity of the investigated systems, constant changes in their
structure. Although we are trying to join these directions in one algorithm,
it is the latter option that we prefer, with it consisting in creating a model
adequate to the process generating a price time series r013ChabM10_ . This
very approach gives a chance to approach the complexity of the system, which
generates the observed series, develop the model and use its properties as the
prognosis.
## 3 Aims of the paper, problem statement
Assume the time series is set by a sequence of discrete levels with constant
step of time sampling $\Delta t$. We need to generate variants of the time
series continuation (prognosis scenarios) according to the relations between
the sequences of absolute and relative changes discovered with the help of
complex Markov chains.
## 4 Classical modeling problems of ESE systems dynamics
Another peculiar feature of ESE systems, apart from complexity, is a memory,
including the long-term one, as well as nonlinear and unstable nature of
interactions and components, which makes it harder to predict their future
behavior.
Unfortunately, mathematical models based on differential equations have no
memory (there is no aftereffect), while for models with memory, where integral
interrelations are used, it is not always possible to take into account
nonlinearity (the integration procedure is linear by definition).
In reality, in the Cauchy problem future systems behavior is defined by its
initial state and doesn’t depend on the way the system reached its current
state. However, it is hardly true that future behavior of a real socio-
economic or socio-ecological system can be predicted by giving an immediate
time “slice” of a variables set that describe its state.
Let us consider possible ways to take into account past events while modeling
ESE systems’ dynamics, which goes beyond the boundaries of classical
differential and integral equations.
Functional differential lagging equation can serve as a simple example of the
dynamic model with memory, where present time is defined by the state variable
$x(t)$ and depends on the past state $x(t-\tau)$ with constant time lag
$\tau=const$:
$x(t)=f\left(x(t-\tau)\right);t\geq t_{0},$ (1)
where $f(x)$ is the known function, with initial conditions being set for the
half-interval $t_{0}-\tau\leq t<t_{0}$ by the function $\phi(t)$:
$x(t)=\phi(t);t_{0}-\tau\leq t<t_{0}.$ (2)
Given the 2 equation 1 has the only solution, defined by recurrent ratios:
$x(t)=\begin{cases}f\left(\phi(t-\tau)\right);\text{ if }t_{0}\leq
t<t+\tau;\\\ f\left(f\left(\phi(t-\tau)\right)\right);\text{ if
}t_{0}+\tau\leq t<t+2\tau;\\\
f\left(\left(f\left(\phi(t-\tau)\right)\right)\right);\text{ if
}t_{0}+2\tau\leq t<t+3\tau;\\\ \cdots\end{cases}$ (3)
Using Dirac delta function, as defined by ratios:
$\delta(t)=0,\text{ if }x\neq
0;\int\limits_{-\infty}^{+\infty}{\delta(t)dt}=1,$ (4)
we can formally rewrite equation 1 in the integral form:
$x(t)=\int\limits_{-\infty}^{t}{dt_{1}f\left(x(t_{1})\right)H(t_{1},t)};H(t_{1},t)\equiv\delta\left(t_{1}-(t-\tau)\right);t\geq
t_{0}.$ (5)
Delta function is not a function in the conventional interpretation and is
related to the class of generalized functions that were mathematically
described only in the middle of the last century r014SobloevFunctProstr
(physics started using this function much earlier). Its classical form is
considered to be a limit of the “peak” sequence, with its centre set in the
point of origin. The afore-mentioned “peaks” indefinitely converge widthway,
indefinitely increase throughout the height and have a unit area.
An approximate classic integral analogue of the equation 5 can be derived by
substituting $\delta(t)$ with an ordinary function - some specific narrow
enough “peak” of a unit area, a certain finite width $\backsim\Delta t$ as
well as a finite height . The derivative of the Fermi function is one of the
possible examples:
$\Phi(t)=\frac{1}{1+exp\left(\frac{-t}{\theta}\right)};\delta(t)\approx\frac{1}{\theta\left(2+exp\left(\frac{-t}{\theta}\right)+exp\left(\frac{t}{\theta}\right)\right)}.$
(6)
If the system’s state in the moment $t$, $x(t)$, is defined not by one, as in
1, but $k$ ($k=2,3,4,\cdots$) of her past states
$x(t-\tau_{1}),x(t-\tau_{2}),\cdots x(t-\tau_{k})$ in the following moments of
time $(t-\tau_{1}),(t-\tau_{2}),\cdots,(t-\tau_{k})$ respectively
$(\tau_{1}=const,\tau_{2}=const,\cdots,\tau_{k}=const,\tau_{1}>\tau_{2}>\cdots>\tau_{k}>0)$,
then instead of (1), (2), (5) we get:
$x(t)=f\left(x(t-\tau_{1});x(t-\tau_{2});\cdots;x(t-\tau_{k})\right);t\geq
t_{0};$ (7) $x(t)=\phi(t);t_{0}-\tau_{1}\leq t<t_{0};$ (8)
$\begin{array}[]{c}x(t)=\int\limits_{-\infty}^{t}{dt_{1}}\int\limits_{-\infty}^{t}{dt_{2}}\cdots\int\limits_{-\infty}^{t}{dt_{k}f\left(x(t_{1}),x(t_{2}),\cdots,x(t_{k})\right)}.\\\
\delta\left((t_{1}-(t-\tau_{1})\right)\delta\left((t_{2}-(t-\tau_{2})\right)\cdots\delta\left((t_{k}-(t-\tau_{k})\right);t\geq
t_{0}.\end{array}$ (9)
Therefore if the system’s state in the time moment $t$ depends on the infinite
sequence of its past states, the integral analogue of the functional
differential lagging equation will, generally speaking, contain an integral of
the infinite multiplicity. At the same time the infinite amount of past states
can relate to both finite $(t-\tau_{1};t)$ (short-term memory) and infinite
$(-\infty;t)$ (long-term memory) time span.
Pay attention that the classic integral lagging equation is one of the
Volterra type r015HeliFunctDiffury :
$x(t)=\int\limits_{-\infty}^{t}{F\left(x(\tilde{t});t;\tilde{t}\right)d\tilde{t}},$
(10)
where $F\left(x(\tilde{t});t;\tilde{t}\right)$ \- is an arbitrary (generally
nonlinear) function of variables $x(\tilde{t});t;\tilde{t}$, which allows to
take into account system’s memory of its past states only in the additive
approximation, which becomes evident, if the right section 10 is rewritten in
the following way:
$\begin{array}[]{c}\int\limits_{-\infty}^{t}{F\left(x(\tilde{t});t;\tilde{t}\right)d\tilde{t}}\equiv\int\limits_{t_{1}}^{t}{F\left(x(\tilde{t});t;\tilde{t}\right)d\tilde{t}}+\int\limits_{t_{2}}^{t_{1}}{F\left(x(\tilde{t});t;\tilde{t}\right)d\tilde{t}}+\cdots=\\\
F\left(x(\tilde{t_{1}});t;\tilde{t_{1}}\right)\cdot(t-t_{1})+F\left(x(\tilde{t_{2}});t;\tilde{t_{2}}\right)\cdot(t_{1}-t_{2})+\cdots;\\\
t>t_{1}>t_{2}>\cdots;\tilde{t_{1}}\in\left[t_{1},t\right];\tilde{t_{2}}\in\left[t_{2},t_{1}\right];\cdots\end{array}$
(11)
In connection with it note that the equation 9 in case of an additive
dependency of contemporaneity on the past, i.e. in case:
$f\left(x(t-\tau_{1});x(t-\tau_{2});\cdots\right)\equiv
f_{1}\left(x(t-\tau_{1})\right);f_{2}\left(x(t-\tau_{2}))\right);\cdots$ (12)
becomes a particular case of the equation 10 with the following integrand:
$F\left(x(\tilde{t});t;\tilde{t}\right)\equiv
f_{1}\left(x(\tilde{t})\right)\delta\left(\tilde{t}-(t-\tau_{1})\right)+f_{2}\left(x(\tilde{t})\right)\delta\left(\tilde{t}-(t-\tau_{2})\right)+\cdots$
(13)
Meaningful analysis of nonlinear models dynamics with memory, in which the
future is defined by the infinite amount of states in the past is generally
possible only in case of a discrete representation. The results of such
analysis will be approximated, i.e. will contain uncertainty, which has to be
considered endogenous, i.e. internal, and peculiar to this very system.
With a certain level of time sampling, models with memory both 7 and 10
becomes:
$x(n+1)=f\left(x(n);x(n-1);x(n-2)\cdots\right).$ (14)
To take into account and quantify the uncertainties, observed in ESE as well
as other complex systems probability models are normally used. However their
application is based on doubtful hypotheses, while the statistical
interpretation of the results is not always informative enough and results
might not correspond with the real process occurring within the system. In
particular, the well-known problem of $1/f$–noise (look for example
r016BukingemShumy ), closely connected to the presence of long-term memory in
complex systems, implies the absence of the mean temporary value (as a limit
of a certain time span converging to infinity, which serves as the basis for
averaging) for any process occurring in such kind of system. Therefore such
processes can’t have a rigorous statistical substantiation.
## 5 Modern concepts in ESE systems modeling
New approaches to modeling and prediction of complex nonlinear systems
dynamics with memory are based on the use of determined chaos and neural
networks technologies (cf. e.g. r011EzhovShumsky98 ; r017LorenzNonlinear89 ;
r018PetersChaosOrder ). Both investigation and realization of such techniques
has become possible only with the appearance of quick-operating computers. Use
of the recurrent computational process has become the general feature for all
these technologies:
$x_{n+1}=f_{n}\left(f_{n-1}\left(\cdots\left(f_{1}(x_{1})\cdots\right)\right)\right),n=1,2,\cdots,$
(15)
where $f_{i}(x_{i})$ is a certain nonlinear mapping of a multi-dimensional
vector $x_{i}$, $i$ \- discrete, real or fictitious, time. Identification of
the model 15 is reduced to the determination of functions $f_{i}(x_{i})$,
while the differences between the models of determined chaos and neural
networks are connected with the function type and methods of its definition
(neural network models normally use rather narrow class of $f_{i}(x_{i})$
mappings r012Zaencev ). Generally speaking, stability or convergence of the
process 15 is not required, whereas a single-step set of vector $x_{i}$
components as well as their time dynamics can be of great interest.
For the particular case of the model 15, introducing corresponding lagged
variables, a model 14 can be transformed.
Both determined stable processes, described by integro-differential equations,
and random processes, which also include complex Markov chains (CMC), can be
formally considered as separate extreme cases of determined chaos models
realization 15. Given the sampling scale, which tends to zero, if such a
tendency makes sense and corresponding limits exist, we derive classical
differential and integral problem statement. Finite $\Delta t$ allows to get
models with discrete time, which in the general case in the corresponding
phase space (which also includes lagged variables) can produce both measurable
sets (discrete or continuous) that allow probabilistic interpretation and
those of the special structure – fractals r019FederFractals , that can’t be
always interpreted in that way.
Various digital generators of so-called random sequences used in imitational
modeling can be an example of determined chaos models that allow probabilistic
interpretation.
Let us note that in reality there are no accurate procedures that would give
an opportunity to distinguish a “real” random sequence from the pseudorandom
one.
## 6 Markov chains prediction technology
Suppose there is a sequence of a certain system discrete states. From this
sequence we can determine transitions probabilities between the two states.
Simple Markov chain is a random process, in which the next state probability
depends solely on the previous state and is independent from the rest of them.
Complex Markov chain, unlike the simple one, stands for the random process, in
which the next state probability depends not only on the current, but also on
the sequence of several previous states (history). The amount of states in
history is the order of the Markov chain.
Theory of simple Markov chains is widely presented in literature, for example
r020TihonovMironovMarkovProcesses . As for the high order Markov chains,
modern literature r021KornVMSpravochnik73 can offer us a mere definition.
Developing complex or high order Markov chain’s properties is not widely
presented in modern scientific publications. It’s necessary to mention the
papers RafteryHighOrderMarkovChain ; RafteryTavare94 where properties of
complex Markov Chains are developed, but no prediction algorithm is proposed
there. The development of prediction method, based on complex Markov chains,
is proposed in this paper.
Markov chain of the higher order can be brought to a simple Markov chain by
introducing the notion of a “generalized state” and including a series of
consequent system’s states into it. In this case, tools of simple Markov
chains can be applied to the complex ones.
Investigated dynamic series is a result of a certain process. It is assumed
that this process is determined, which implies the existence of a causal
dependence of further states on history. It is impossible to fix and analyze
the infinite history, which puts obstacles in the way of an accurate detection
of this influence and making precise predictions.
The problem consists in the maximal use of information, which is contained in
the known segment of the time series, and subsequent modeling of the most
probable future dynamics scenario.
The observed process is described as a time series of prices $p_{t}$ with the
given sampling time span $\Delta t$
$p_{ti}=p(t_{0}+i\Delta t).$ (16)
Discrete presentation of the time series is in fact a way of existence of this
very system. New prices are formed on the basis of contracts or deals, made on
the market in certain discrete moments of time, while the price time series is
a series of the averaged price levels during the chosen time intervals. While
making a decision each trader, who is an active part of the pricing system,
works solely with discrete series of the chosen time interval (e.g. minute,
5-minute, hourly, daily etc.). For $\Delta t\to 0$ the accuracy of data
presentations reaches a certain limit, since for relatively small $\Delta t$
the price leaps in the moment of deal, while staying unchanged and equal to
the last deal during the time between the two deals. Hence, the discreteness
of time series has to be understood not only as a limited presentation of
activity of the complex financial system, but also as one of the principles of
its operation r007KurbanovSaptsin07 ; r008SaptsinMarkovPaper09_ ;
r022SaptsinSoloviev_ ; r022SapSolArxiv_ ; r023SolDerbMonogr2010_ .
The time series of initial conditions has to be turned into a sequence of
discrete states. Let us denote the amount of chosen states as $s$, each of
them being connected to the change in the quantity of the initial signal
(returns). For example, consider the classification with two states, first of
which corresponds to positive returns as the price increases, while the second
one – to negative as it descends. Generally all possible increments of the
initial time series are divided into $s$ groups. Ways of division will be
discussed further.
Next we develop predictions for the time series of sampled states. For the
given order of the Markov chain and the last generalized state the most
probable state is chosen to be the next one. In case if ambiguity occurs while
the state of maximum probability is being evaluated, an algorithm is used that
allows reducing the amount of possible prediction scenarios. Therefore we get
the series of predicted states that can be turned into a sampled sequence of
prognostic values.
Evaluation of increments, prediction, and subsequent restoration are conducted
for the given hierarchy of time increments $t$. To use the given information
as effectively as possible, the prediction is conducted for time increments
$t=1,2,4,8,...$, or a more complex hierarchy of increments and subsequent
“splicing” of the results derived from different prediction samplings.
The procedure of prediction and splicing is iterative and conducted starting
from smaller increments, adding a prediction with the bigger time increment on
every step.
As the sampling time step $t$ increases, the statistics for the investigation
of Markov chains decreases, whereas the biggest sampling step, which takes
part in the prognostication, limits itself. To supplement the prediction with
the low-frequency component the approximation of zero order is being used in
the form of a linear trend or a combination of a linear trend and harmonic
oscillations r024SapChFourierKharkov ; r025ChabS10_ .
## 7 Prediction construction algorithm
Let us consider the consequence of operations, required for the prognostic
time series construction. To do this we need to set the following parameters:
1) The type of time increments hierarchy (simple – powers of two, complex –
product of powers of the first simple numbers).
2) Values of $s$ – the amount of states and $r$ – the order of the Markov
chain. These parameters can be individual for every sampling level; finding of
optimal parameters is done experimentally.
3) Threshold values $\delta$, and minimal number of transitions $N_{min}$.
Prediction construction algorithm includes the following steps:
1) Generating hierarchy of time increments - $t$ sequence. The maximal of them
has to correspond to the length of a prognostic interval $N_{max}$.
2) For every time increment $\Delta t$, as the increments increase, a
prediction of states and restoration of the time series along the prognostic
states is conducted. Current stage includes following actions:
2.1. Evaluating increments (returns) of the series with $\Delta t$ sampling.
2.2. Transforming the time series of increments into the series of state
numbers ($1..s$).
2.3. Calculating transition probabilities for generalized states.
2.4. Constructing the series of prognostic states using the procedure of
defining the most probable next state.
2.5. Restoring the value series from the state series with $\Delta t$
sampling.
2.6. Splicing the prediction of $\Delta t$ sampling with the time series
derived from splicing of the previous layers (with the lesser step $\Delta
t$). In case if the current time series is the first one, the unchanged time
series will come as a result of splicing.
3) To splice the last spliced time series with the continuation of the linear
trend, created along all previously known points.
The time series, spliced with the linear trend, is the result of prediction.
Let us consider the stages of the given algorithm in detail.
## 8 States in complex Markov chains and approaches for defining them
In everything that concerns current technology, states are connected to the
measuring of a prognostic value. There is a number of ways to classify returns
in states, from which the following are suggested. One of them is the
classification based on the homogeneity principle as concerning the amount of
representatives in classes; based on the homogeneity principle of deviation,
as well their combinations for different deviation modules.
Increment or returns of the time series serves as the basis for states
classification r025ChabS10_ ; r026solovievmatheconomics_ . Absolute $r_{a}$
and relative $r_{t}$ increments of the time series are considered:
$r_{a}=p_{t}-p_{t-\Delta t},$ (17) $r_{t}=\frac{p_{t}-p_{t-\Delta t}}{p_{t}},$
(18)
where $p_{t}$ – is the input time series of price dynamics, $\Delta t$ –
sampling interval, which is chosen for subsequent analysis. It is known that
mathematical expectation of the returns time series equals zero, whereas
variation comes as the measure of time series volatility. Based on returns
values $r_{t}$ classification and transformation of values to the time series
of discrete states are conducted. One of the classification principles is
homogeneity according to the amount of class representatives. This
classification divides the set of all increments into $s$ groups equal in
number. Calculated with the given sampling, time series increments are then
systematized in growth and divided into equal parts. Thus we define limit
values $\\{r_{lim,i}\\}$, which are used afterwards during transformation of
the returns into class numbers. Large number of identical states can cause
certain problems, such as identical bounds of several neighbouring states. It
creates a number of states with no representatives, which makes correction of
the division a necessary action. In that way we will reach the largest
possible homogeneity in state division. Classification is conducted along the
following algorithm r027SolSapChDrezden09 ; r028ChabRiga10 :
$s_{t}=(i\mid r_{lim,i-1}>r_{t}>r_{lim,i})$ (19)
where $s_{t}$ is the number of state, which corresponds to the moment of time
$t$, for which the returns level was computed $r_{t}$; i is the number of
state $[1\dots s]$, which is characterized by the interval
$[r_{lim,i-1},r_{lim,i}]$ corresponding to the calculated returns level
$r_{t}$.
Apart from the returns interval, given by the aforementioned values
$[r_{lim,i-1},r_{lim,i}]$, a mean returns value is chosen for every state
$r_{avg,i}$, which will be used in time series values transformation according
to the prognostic discrete states.
Another way of dividing the time series into states implies dividing the
interval of returns values into equal parts, from minimal to maximal
deviation. In this case homogeneity according to the amount of representatives
in states does not occur. In fact this method differs from the previous one in
terms of defining limit values $\\{r_{lim,i}\\}$. Possible combined ways of
division, in case of which the limit value, dependent on standard deviation,
is used instead of maximal and minimal value, and division is conducted
homogenously according to the deviation.
Since the real causal dependence is unknown during the process, to find it
adequate state classification, which would allow to reveal vital dependencies
of the time series, is required. We suggest a couple of ways to divide the
time series into states, which in the first place allow to preserve adequate
transition probabilities between states, as well as prevent averaged
deviations inside the states from affecting the accuracy of the derived
prediction.
To check the efficiency of division we conduct the sampling procedure and
classify the increments according to each hierarchy. Having completed that, we
restore the time series using known states for each hierarchy and finish the
splicing procedure. Since the state series correspond to the initial time
series, we get the curve, with deviation, caused exclusively by the state
averaging mistake (quantum mistake). Thus, having set a certain value of state
numbers $s$ and carried out sampling, restoring, and sampling procedures
(excluding prediction), we get absolute sampling (quantum) mistake.
Increasing the number of states, we improve the accuracy of restoration,
however one should remember, that the choice of the quantum levels is limited
by the fact that the transition probabilities definition with sufficient
accuracy is required, which is confirmed by artificial test time series
prediction experiments.
## 9 Step-by-step prediction procedure. Defining the most probable state on
the next step, prediction scenarios
Predicting procedure uses the most probable state as the next one under
current circumstances. Probability matrix of state transitions is used for the
afore-mentioned purpose. In this case, you have to take into account that
probabilities are calculated with a certain mistake. We cannot precisely
compute the probabilities, since it is impossible to derive an infinite time
series, and only a part of the time series is known – the known part serves as
the basis for probabilities. The second important aspect implies the case of
several states with maximal probability.
To prevent the omission of the states, for which the probabilities are
computed with a mistake, one should add a state with maximal probability to
the states, which are located in the distance of $\delta$ from the maximal
one. The value of parameter $\delta$ depends on the probability evaluation
mistake and requires experimental refinement.
If $\delta>0$, the number of states with maximal probability increases in
comparison to the value $\delta=0$. Let us call a couple of neighbouring
states with maximal probability a cluster. Cluster states with average
deviation values are supposed to have the largest probability. To predict the
dynamics, let us confine ourselves to one or two most probable states. To
define them a following algorithm is suggested:
1) If levels (discretized increments) create several clusters (cluster is a
group of several neighbouring levels – cluster elements, minimal cluster is a
single isolated level) with maximal probability, we choose the largest
cluster.
2) If the number of cluster elements is odd, as $k_{max}$ we choose a central
element.
3) If the number of cluster elements is even, we consider two central cluster
elements and choose as $k_{max}$ the one, which is closer to the centre of
distribution.
4) If two central cluster elements are equidistant to the centre of
distribution, we consider both cases as possible variants of $k_{max}$ values
(bifurcation point).
5) If there are several clusters of maximal size, we consider them as new
elements, which can also form clusters that will undergo the same steps 1)-4).
This principle is based on the following ideas:
1) If there are two neighbouring states of maximal probability, it is better
to take the one, which is closer to the centre of distribution, in order to
minimize the risk of occurrence of false linear trends in the prediction.
2) If levels of maximal probability are not the neighbouring ones, at least
two variants have to be considered, as it can be connected to the bifurcations
that should not be omitted.
3) If the prediction is carried out according to 1) (on all stages of the
hierarchy), we receive a certain approximation of the lower limit of the
prediction, whereas in case of 2) – we get an approximation of the upper
limit.
Hence this algorithm can adequately restore the case of possible bimodal
probability distribution, it is proposed to consider 2 prediction scenarios.
In case of the complex Markov chains, probability of the next state depends
not only on the previous state, but also on the sequence of $r$ states, which
have occurred before given. In this case, it is necessary to calculate
transition probabilities from the sequence of $r$ states into the $r+1$ state.
Formally, these probabilities can be written into the rectangular table of
$(r^{s},s)$ size.
Having generalized the notion of “present state” and included a sequence of
$r$ preceding states into it, we can reduce Markov chains of $r$ order to the
chain of the first order. Thus transition probabilities can be written into
rectangular matrices of $(r^{s},r^{s})$, that come as transition probability
matrices for generalized states.
The process of prediction implies the following: the last state is chosen (in
case of Markov chains of an order $r>1$ a sequence of $r$ latest states is
taken). The probability of transition from current state to all possible
states is defined. From all possible states a state with maximal probability
is chosen. It is possible that several states with maximal probability occur,
which can be explained by the bimodal probability distribution. The process of
decision-making in this case is described later.
The chosen most probable state is taken as the next prognostic state and the
procedure is repeated for the next (last added) state. Thus we receive a time
series of prognostic states for the given sampling time $\Delta t$.
Further according to the received state sequence and known initial value the
time series is being restored for the given time sampling $\Delta t$. In this
case every state implies $\Delta t$ points of the time series. On the stage of
state classification every state was connected to the average increment
$r_{avg,i}$, which is added to the value of the last point in the time series,
and the next discrete point is computed. Intermediate points are filled as
linear interpolation of two known neighbouring points. Algorithm of $y_{t}$
time series values restoration according to the initial price $p_{t}$ and a
series of average increments $r_{avg,ik}$, corresponding to the prognostic
states $s_{k}$, can be given by a sequence of calculations:
$\begin{array}[]{c}y_{t}=p_{t},\\\ y_{t+1}=y_{t}+r_{avg,i1}/\Delta
t=p_{t}+r_{avg,i1}/\Delta t,\\\ y_{t+2}=y_{t+1}+r_{avg,i1}/\Delta
t=p_{t}+2r_{avg,i1}/\Delta t,\\\ \dots\\\ y_{t+\Delta t-1}=y_{t+\Delta
t-2}+r_{avg,i1}/\Delta t=p_{t}+(\Delta t-1)r_{avg,i}/t,\\\ y_{t+\Delta
t}=y_{t+\Delta t-1}+r_{avg,i1}/\Delta t=p_{t}+\Delta tr_{avg,i}/\Delta
t=p_{t}+r_{avg,i1},\\\ y_{t+\Delta t+1}=y_{t+\Delta t}+r_{avg,i2}/\Delta
t=p_{t}+r_{avg,i1}+r_{avg,i2}/\Delta t,\\\ \dots\\\ y_{t+n\Delta
t-1}=y_{t+n\Delta t-2}+r_{avg,in}/\Delta
t=p_{t}+\sum_{k=1}^{n-1}r_{avg,ik}+\frac{\left(\Delta t-1\right)}{\Delta
t}r_{avg,ik},\\\ y_{t+n\Delta t}=y_{t+n\Delta t-1}+r_{avg,in}/\Delta
t=p_{t}+\sum_{k=1}^{n}r_{avg,ik}.\\\ \end{array}$ (20)
## 10 Time increments hierarchy and splicing procedure
Time series increments will be computed with different steps. For example,
analogous to the discrete Fourier transform, time increments are equal the
powers of 2 are considered. First, we calculate increments as a remainder of
two nearest neighbouring time series values, then next nearest values are
considered with the step of 2, 4, 8, 16 etc. Let us mark this difference in
time as $\Delta t$.
For every $\Delta t$ we conduct an increment time series transformation
leading to a time series of states. Further we predict the future sequence of
states and restore the time series with the given sampling rate according to
the prognostic series of states.
Time series, received as a result of restoring for different $\Delta t$,
undergo the splicing procedure, which gives out an actual prognostic time
series.
Thus an increment hierarchy is chosen, where each one is responsible for its
own sampling rate, which serves as a basis for predicting, restoring and
splicing.
The splicing process implies the following. The procedure is iterative. With
every next (along with the increasing step) sampling time the series corrects
itself, driving the prediction, formed under lower $\Delta t$, to its actual
point. Transformations that are conducted during splicing can be written down
in the form of the following calculations.
Suppose the splicing procedure has been finished for all time increments
$\Delta t<Deltat_{i}$, the prediction has been done under the $\Delta t_{i}$
sampling according to formulae 20, and as a result a time series $y_{i}$ has
been derived. Let us consider the iterative splicing procedure of the received
series $y_{i}$ with the series, acquired during all preceding splicing
procedures $g_{i}$.
Since the series $y_{i}$ contains system points only in moments aliquot to
$\Delta t_{i}$, and other points of the series are interpolated, the process
of splicing implies the substitution of these interpolated points with the
values of system points from previous $\Delta t<\Delta t_{i}$, which are
contained in the series of results of previous splicing procedures $g_{i}$.
Splicing algorithm can be written in the sequence of computations:
$\begin{array}[]{c}z_{t}=g_{t}=p_{t},\\\ z_{t+1}=g_{t+1}+\left(y_{t+\Delta
t_{i}}-g_{t+\Delta ti}\right)/\Delta t_{i},\\\
z_{t+2}=g_{t+2}+2\left(y_{t+\Delta t_{i}}-g_{t+\Delta ti}\right)/\Delta
t_{i},\\\ \dots\\\ z_{t+\Delta t-1}=g_{t+\Delta t-1}+\left(\Delta
t_{i}-1\right)\left(y_{t+\Delta t_{i}}-g_{t+\Delta ti}\right)/\Delta t_{i},\\\
z_{t+\Delta t}=g_{t+\Delta t}+\left(\Delta t_{i}\right)\left(y_{t+\Delta
t_{i}}-g_{t+\Delta ti}\right)/\Delta t_{i}=y_{t+\Delta_{ti}},\\\ z_{t+\Delta
t+1}=g_{t+\Delta t+1}+\left(\left(y_{t+2\Delta ti}-g_{t+2\Delta
ti}\right)-\left(y_{t+\Delta ti}-g_{t+\Delta ti}\right)\right)/\Delta
t_{i},\\\ z_{t+\Delta t+2}=g_{t+\Delta t+2}+2\left(\left(y_{t+2\Delta
ti}-g_{t+2\Delta ti}\right)-\left(y_{t+\Delta ti}-g_{t+\Delta
ti}\right)\right)/\Delta t_{i},\\\ \dots\\\ z_{t+n\Delta t-1}=g_{t+n\Delta
t-1}+\frac{\left(\Delta t-1\right)}{\Delta t}\left(\left(y_{t-n\Delta
t}-g_{t-n\Delta t}\right)-\left(y_{t-(n-1)\Delta t}-g_{t-(n-1)\Delta
t}\right)\right),\\\ z_{t+n\Delta t}=g_{t+n\Delta t}+\left(\left(y_{t-n\Delta
t}-g_{t+n\Delta t}\right)-\left(y_{t+(n-1)\Delta t}-g_{t+(n-1)\Delta
t}\right)\right)=\\\ =g_{t+(n-1)\Delta t}-y_{t+(n-1)\Delta t}-y_{t-n\Delta
t}.\end{array}$ (21)
## 11 Results of stock indices prediction
In this section we offer the results of stock indices prediction. The stock’s
indices databases are available from Finance_yahoo . Point 2000 indicates the
starting moment of the prognosis: March 24, 2011. The green line on the next
figures indicates real indice’s or price’s values. Our software for time
series forecasting by the proposed methods is available from our website:
http://kafek.at.ua/MarkovChains1_2_20100505.rar.
Figure 1: Stock indices prediction. a) Dow Jones Industrial Average - DJI
(USA). b) FTSE 100 (Great Britain)
Figure 2: Financial companie’s share prices forecasting. a) Morgan Stanley
(USA). b) BNP Paribas (France)
Prediction time series with different input learning set’s length are shown at
the fig.3 and 4. Prediction series for DJI at the figure 3 are more
correlated, than FTSE index at the figure 4. At the subplot b) of the above
mentioned plots the mean value and standard deviations of the prediction’s
series are presented. The time of prediction series beginning on the next
figures is the point 1000 and correspond to October 14, 2011.
Figure 3: Dow Jones Industrial Average - DJI (USA). a) Prediction series,
calculated with different learning set’s length. b) Mean value and standard
deviation for prediction series.
Figure 4: FTSE 100 index prediction. a) Prediction series, calculated with
different learning set’s length. b) Mean value and standard deviation for
prediction series.
The normalization procedure is proposed in order to compare indices and it’s
prediction series with different absolute values. The normalized values
calculated with the following formula:
$y_{n}(t)=\frac{y(t)-min\left(y(t)\right)}{max\left(y(t)\right)-min\left(y(t)\right)}.$
(22)
Normalized prediction time series are shown at the fig.5 (America), fig.7
(Europe, developed countries), fig.7 (Europe, PIIGS), fig.9 (Asian markets).
All the figures contain mean time series, which are weighted average of
countrie’s stock indices predictions, weihted with GDP values EconomyWatch_com
for the corresponding countries.
Figure 5: Normalized mean values for the prediction series of America’s stock
indices. Brazil (BVSP), Mexico (MXX), Canada (GSPTSE), Argentina (MERV), USA
(S&P 500)
Figure 6: Normalized mean values for the prediction series of European stock
indices. Developed countries: FTSE (Great Britain), DAX (Germany) FCHI
(France), Netherlands (AEX)
Figure 7: Normalized mean values for the prediction series of European stock
indices. Portugal (PSI20), Italy (FTSEMIB), Ireland (ISEQ), Greece (GD) and
Spain (IBEX).
Figure 8: Normalized mean values for the prediction series of Asian stock
indices. China (SSEC, HSI), Korea (KS11), Japan (NIKKEI), India (BSESN), New
Zealand (NZ50).
Figure 9: Mean values of normalized World’s powerful economies indices
prediction series.
## 12 Conclusions and further work
Current paper suggests an algorithm of time series prediction based on complex
Markov chains. Hierarchy of time increments principle allows to use the
information, which is contained in the time series during the prognosis
construction, to its fullest. Experimental work on stock market indices time
series prediction shows the efficiency of the algorithm and confirms the
relevance of further research of the offered method.
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|
arxiv-papers
| 2011-11-22T17:10:09 |
2024-09-04T02:49:24.578985
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Vladimir Soloviev and Vladimir Saptsin and Dmitry Chabanenko",
"submitter": "Vladimir Saptsin",
"url": "https://arxiv.org/abs/1111.5254"
}
|
1111.5289
|
# HEISENBERG UNCERTAINTY PRINCIPLE AND ECONOMIC ANALOGUES OF BASIC PHYSICAL
QUANTITIES
Soloviev V., Prof. Dr. Sc vnsoloviev@rambler.ru Cherkasy National University
named after B. Khmelnitsky, Cherkassy, Ukraine Saptsin V., PhD
saptsin@sat.poltava.ua Kremenchug National University named after M.
Ostrogradskii, Kremenchuk, Ukraine
###### Abstract
From positions, attained by modern theoretical physics in understanding of the
universe bases, the methodological and philosophical analysis of fundamental
physical concepts and their formal and informal connections with the real
economic measurings is carried out. Procedures for heterogeneous economic time
determination, normalized economic coordinates and economic mass are offered,
based on the analysis of time series, the concept of economic Plank’s constant
has been proposed. The theory has been approved on the real economic dynamic’s
time series, including stock indices, Forex and spot prices, the achieved
results are open for discussion.
quantum econophysics, uncertainty principle, economic dynamics time series,
economic time.
###### Contents
1. 1 Introduction
2. 2 About nature and interrelations of basic physical notions
3. 3 Dynamical peculiarities of economic measurements, economical analog of Heisenberg’s uncertainty ratio
4. 4 Experimental results and their discussion
5. 5 Conclusions
## 1 Introduction
The instability of global financial systems depending on ordinary and natural
disturbances in modern markets and highly undesirable financial crises are the
evidence of methodologial crisis in modelling, predicting and interpretation
of current socio-economic conditions.
In papers r001SolSapSimferopolThesis ; r002SaptsinSoloviev ; r003SapSolArxiv
we have suggestsed a new paradigm of complex systems modelling based on the
ideas of quantum as well as relativistic mechanics. It has been revealed that
the use of quantum-mechanical analogies (such as the uncertainty principle,
notion of the operator, and quantum measurement interpretation) can be applied
to describing socio-economic processes. In papers r001SolSapSimferopolThesis ;
r002SaptsinSoloviev ; r003SapSolArxiv we have suggestsed a new paradigm of
complex systems modelling based on the ideas of quantum as well as
relativistic mechanics. It has been revealed that the use of quantum-
mechanical analogies (such as the uncertainty principle, notion of the
operator, and quantum measurement interpretation) can be applied to describing
socio-economic processes.
It is worth noting that quantum analogies in economy need to be considered as
the subject of new inter-disciplinary direction – quantum econophysics (e.g.
r004BqaqueQuantumFin2004 ; r005MaslovQuantumEconomics06 ;
r006GuevaraQuantumEconophysics06 ; r007 ; r008Goncalez2011 ), which, despite
being relatively young, has already become a part of classical econophysics
r009MantegnaStanley00 ; r010Romanovsky07 ; r011solovievmatheconomics ;
r012SolDerbMonogr2010 .
Ideas r001SolSapSimferopolThesis ; r002SaptsinSoloviev ; r003SapSolArxiv were
anticipated and further developed in our works on modelling, predicting, and
identification of socio-economic systems r013SaptsinGenComplMarkovChains08 ;
r014SaptsinMarkovPaper09 ; r015SolSapChDrezden09 ; r016ChabM09 ; r017ChabM10 ;
r018ChabRiga10 ; r019SolSapChPsepGlava ; r020KuznetzKoleb2011 (complex Markov
chains), r019SolSapChPsepGlava ; r020KuznetzKoleb2011 ; r021Fourier2009 ;
r022SapChFourierKharkov ; r023ChabS10 (discrete Fourier Transorm),
r024SapDvufaznye05 ; r025OlhovayaKonkur2010 ; r026SapChabAGranUstoych11 ;
r027SapSolBatyrNelnConcur2agents10 (multi-agent modelling),
r028SolSapDynNetwMaths ; r029SolSapChabAPedZbornik (dynamic network
mathematics), r030SaptsinSetPrir ; r031SaptsinMerezhPrirody (network
measurements), r032SolSapShokotko ; r033HeisenbergRiga (uncertainty principle
in economics) etc.
However significant differences between physical and socio-economical
phenomena, diversity and complexity of mathematical toolset (which, on account
of historical circumstances, has been developed as the language of sciences),
as well as lack of deep understanding of quantum ideology among the
scientists, working at the joint of different fields require a special
approach and attention while using quantum econophysical analogies.
Aim of work. Our aim is to conduct detailed methodological and phylosophical
analysis of fundamental physical notions and constants, such as time, space
and spatial coordinates, mass, Planck’s constant, light velocity from the
point of view of modern theoretical physics, and search of adequate and useful
analogues in socio-economic phenomena and processes.
## 2 About nature and interrelations of basic physical notions
Time, distance and mass are normally considered to be initial, main or basic
physical notions, that are not strictly defined. It is thought that they can
be matched with certain numerical values. In this case other physical values,
e.g. speed, acceleration, pulse, force, energy, electrical charge, current
etc. can be conveyed and defined with the help of the three above-listed ones
via appropriate physical laws.
Let us emphasize that none of the modern physical theories, including
relativistic and quantum physics, can exist without basic notions.
Nevertheless, we would like to draw attention to the following aspects.
As Einstein has shown in his relativity theory, presence of heterogeneous
masses leads to the distortion of 4-dimensional time-space in which our world
exists. As a result Cartesian coordinates of the 4-dimensional Minkowski space
$(x,y,z,ict)$, including three ordinary Cartesian coordinates $(x,y,z)$ and
the forth formally introduced time-coordinate $ict$ ($i=\sqrt{-1}$ \-
imaginary unit, $c$ \- speed of light in vacuum, $t$ \- time), become
curvilinear r034landau1975classical .
It is also possible to approach the interpretation of Einstein’s theory from
other point of view, considering that the observed heterogeneous mass
distribution is the consequence of really existing curvilinear coordinates
$(x,y,z,ict)$. Then the existence of masses in our world becomes the
consequence of geometrical factors (presence of time-space and its curvature)
and can be described in geomatrical terms.
If we step away from global macro-phenomena that are described by the general
relativity theory, and move to micro-world, where laws of quantum physics
operate, we come to the same conclusion about the priority of time-space
coordinates in the definition of all other physical values, mass included.
To demonstrate it, let us use the known Heisenberg’s uncertainty ratio which
is the fundamental consequence of non-relativistic quantum mechanics axioms
and appears to be (e.g. r002SaptsinSoloviev ):
$\Delta x\cdot\Delta v\geq\frac{\hbar}{2m_{0}},$ (1)
where $\Delta x$ and $\Delta v$ are mean square deviations of $x$ coordinate
and velocity $v$ corresponding to the particle with (rest) mass $m_{0}$,
$\hbar$ \- Planck’s constant. Considering values $\Delta x$ $\Delta v$ to be
measurable when their product reaches its minimum, we derive (from (1)):
$m_{0}=\frac{\hbar}{2\cdot\Delta x\cdot\Delta v},$ (2)
i.e. mass of the particle is conveyed via uncertainties of its coordinate and
velocity – time derivative of the same coordinate.
Nowadays, scientists from different fields occupy themselves with the
investigation of structure and other fundamental properties of spacetime from
physical, methodological, psychological, philosophical and other points of
view. However, theoretical physics r035CramerTransactional86 ;
r036kaku1999introduction , including its most advanced and developing spheres
(e.g. string theory r036kaku1999introduction ;
r037BalasubramanianWhatWeDontKnow ) is expected to show the most significant
progress in understanding the subject, though there is no single concept so
far r035CramerTransactional86 ; r036kaku1999introduction ;
r037BalasubramanianWhatWeDontKnow ; r038Vladimirov_p1 ; r039Vladimirov_p2 ;
r040kaku2006parallel ; r041kaku2008physics .
Within fundamental physical science we can mark out two investigational
directions: 1) receipt of quantitive patterns, possible to verify
experimentally or empirically and 2) interpretation of existing theories or
development of new theories, that allow accurate and laconic (involving as
little as possible mathematical notions and formalisms) interpretation of
basic physical facts. The second direction is especially important when
speaking of transferring physical notions and mathematical formalisms into
other spheres, e.g. economics.
Not claiming to be exhaustive, aiming to make the audience (professional
economists included) as wide as possible, we will confine ourselves to the
examination of the most typical and clear examples.
According to the concept r038Vladimirov_p1 ; r039Vladimirov_p2 , having been
developed for the last couple of decades by the Moscow school of theoretical
physicists (headed by Y. Vladimirov), space, time, and four fundamental
physical interactions (gravitational, electromagnetic, strong and weak) are
secondary notions. They share common origins and are generated by the so-
called world matrix which has special structure and peculiar symmetrical
properties. Its elements are complex numbers which have double transitions in
some abstract pre-space.
At the same time, physical properties of spacetime in this very point are
defined by the nonlocal (“immediate”) interaction of this point with its close
and distant neighbourhood, and acquire statistical nature. In other words,
according to Vladimirov’s concept, the observed space coordinates and time
have statistical nature.
It is worth noting that similar ideas as of interpreting quantum mechanics,
different from those of the Copenhagen school were proclaimed by John Cramer
r035CramerTransactional86 (Transactional interpretation of quantum
mechanics).
In our opinion the afore-metioned conception of nonlocal statistical origin of
time and space coordinates can be qualitatively illustrated on the assuptions
of quantum-mechanical uncertainty principle using known ratios (e.g.
r002SaptsinSoloviev :)
$\Delta p\cdot\Delta x\sim\hbar;$ (3)
$\Delta E\cdot\Delta t\sim\hbar;$ (4)
$\Delta p\cdot\Delta t\sim\frac{\hbar}{c}.$ (5)
Interpreting values $\Delta E,\Delta p,\Delta x,\Delta t$ as uncertainties of
particle’s energy $E$, its pulse $p$, coordinate $x$ and time localization $t$
(the latter ratio relates to the relativistic case $E=pc$, and is formally
derived from the ratio (4), if $\Delta E=\Delta p\cdot c$, and takes into
account maximum speed $c$ limitations in an explicit form), let us conduct the
following reasoning.
While $\Delta x\to 0$ uncertainty of pulse, and thus particle energy,
uncertainty, formally becomes as big as possible, which can be provided only
by its significant and nonlocal energetical interaction with the rest of the
neighbourhood 3. On the other side, while $\Delta p\to 0$ the particle gets
smeared along the whole space (according to (3) $\Delta x\to\infty$), i.e.
becomes delocalized. It might be supposed that the fact of “delocalized” state
of the particle takes place in any other, not necessarily marginal $\Delta x$
and $\Delta p$ value ratios.
Similar results can be acquired while analyzing ratios (4)-(5), and for
temporary localization $\Delta t$.
Vladimirov’s concept probably becomes more graphic (at least for those, who
are familiar with the basics of the band theory), if one remembers that so-
called “electrones” and “holes” are considered to be really existing charge
bearers in semiconductors. These “particles” have negative and positive charge
respectively, accurate to the decimal place, which corresponds to the charge
of a free electron, and are characterised by effective masses $m_{e}$ and
$m_{h}$, different from the mass of a free electron (generally $m_{e}$ and
$m_{h}$ can also be tensor values). However, in reality, these particles are
virtual results of the whole semicondoctor crystal – so-called quasi-particles
– and don’t exist beyond its bounds.
Drawing the analogy with crystal it can be supposed that all structural
formations of our Universes are such “quasi-particles”, caused by nonlocal
interaction and non-existent beyond its spacetime bounds.
In conclusion we would like to note that the conept of nonlocal interaction is
quite capable of giving the logical explanation to the empirical fact of
indistinguishability and identity of all microparticles of this kind, which
always takes place during the observation (identification) regardless of
spacetime localization of this very observation.
## 3 Dynamical peculiarities of economic measurements, economical analog of
Heisenberg’s uncertainty ratio
Main physical laws are normally distinguished with the presence of constants,
that have been staying unchanged for the past $\sim 10^{11}$ years (the age of
our Universe since so-called “big bang”- the most widespread hypothesis of its
origin). Gravitational constant, speed of light in vaccuum, Planck’s constant
are among the above-listed.
Speaking of economic laws, based on the results of both physical (e.g.
quantities of material resources) and economical (e.g. their value) dynamic
measurements, the situation will appear to be somewhat different. Adequacy of
the formalisms used for mathematical descriptions has to be constantly checked
and corrected if necessary. The reason is that measurements always imply a
comparison with something, considered to be a model, while there are no
constant standards in economics (they change not only quantitavely, but also
qualitatively – new standards and models appear). Thus, economic measurements
are fundamentally relative, are local in time, space and other socio-economic
coordinates, and can be carried out via consequent and/or parallel comparisons
“here and now”, “here and there”, “yesterday and today”, “a year ago and now”
etc. (see r030SaptsinSetPrir ; r031SaptsinMerezhPrirody for further
information on the subject).
Due to these reasons constant monitoring, analysis, and time series prediction
(time series imply data derived from the dynamics of stock indices, exchange
rates, spot prices and other socio-economic indicators) becomes relevant for
evaluation of the state, tendencies, and perspectives of global, regional, and
national economies.
Let us proceed to the description of structural elements of our work and
building of the model.
Suppose there is a set of $M$ time series, each of $N$ samples, that
correspond to the single distance $T$, with an equal minimal time step $\Delta
t_{\min}$:
$X_{i}(t_{n}),\begin{array}[]{c}{}\hfil\end{array}t_{n}=\Delta
t_{\min}n;\begin{array}[]{c}{}\hfil\end{array}n=0,1,2,...N-1;\begin{array}[]{c}{}\hfil\end{array}i=1,2,...M.$
(6)
To bring all series to the unified and non-dimentional representation,
accurate to the additive constant, we normalize them, having taken a natural
logarithm of each term of the series:
$x_{i}(t_{n})=\ln
X_{i}(t_{n}),\begin{array}[]{c}{}\hfil\end{array}t_{n}=\Delta
t_{\min}n;\begin{array}[]{c}{}\hfil\end{array}n=0,1,2,...N-1;\begin{array}[]{c}{}\hfil\end{array}i=1,2,...M.$
(7)
Let us consider that every new series $x_{i}(t_{n})$ is a one-dimensional
trajectory of a certain fictitious or abstract particle numbered $i$, while
its coordinate is registered after every time span $\Delta t_{\min}$, and
evaluate mean square deviations of its coordinate and speed in some time
window $\Delta T$:
$\Delta T=\Delta N\cdot\Delta t_{\min}=\Delta
N,\begin{array}[]{c}{}\hfil\end{array}1<<\Delta N<<N.$ (8)
The “immediate” speed of $i$ particle at the moment $t_{n}$ is defined by the
ratio:
$v_{i}\left(t_{n}\right)=\frac{x_{i}(t_{n+1})-x_{i}(t_{n})}{\Delta
t_{\min}}=\frac{1}{\Delta t_{\min}}\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})},$
(9)
its variance $D_{v_{i}}$:
$D_{v_{i}}=\frac{1}{(\Delta
t_{\min})^{2}}\left(<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right),$ (10)
and mean square deviation $\Delta v_{i}$:
$\Delta v_{i}=\sqrt{D_{v_{i}}}=\frac{1}{(\Delta
t_{\min})}\left(<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)^{\frac{1}{2}},$ (11)
where $<...>_{n,\Delta N}$ means averaging on the time window of $\Delta
T=\Delta N\cdot\Delta t_{\min}$ length. Calculated according to (11) value of
$\Delta v_{i}$ has to be ascribed to the time, corresponding with the middle
of the avaraging interval $\Delta T$.
To evaluate dispersion $D_{x_{i}}$ coordinates of the $i$ particle are used in
an approximated ratio:
$2D_{x_{i}}\approx D_{\Delta x_{i}},$ (12)
where
$D_{\Delta
x_{i}}=\begin{array}[]{c}{}\hfil\end{array}<\left(x_{i}(t_{n+1})-x_{i}(t_{n})\right)^{2}>_{n,\Delta
N}-\left(<x_{i}(t_{n+1})-x_{i}(t_{n})>_{n,\Delta N}\right)^{2}=$
$=\begin{array}[]{c}{}\hfil\end{array}<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta N}\right)^{2},$
(13)
which is derived from the supposition that $x$ coordinates neighbouring
subject to the time of deviation from the average value $\bar{x}$ are weakly
correlated:
$<\left(x_{i}\left(t_{n}\right)-\bar{x}\right)\left(x_{i+1}\left(t_{n}\right)-\bar{x}\right)>_{n,\Delta
N}\approx 0.$ (14)
Thus, taking into account 12 and 13 we get:
$\Delta x_{i}=\sqrt{\frac{D_{\Delta
x_{i}}}{2}}=\frac{1}{\sqrt{2}}\left(<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)^{\frac{1}{2}}.$ (15)
Pay attention that it was not necessary for us to prove the connection 12, as
it was possible to postulate statement (15) as the definition of $\Delta
x_{i}$.
It is also worth noting that the value
$\left|v_{i}\left(t_{n}\right)\right|\cdot\Delta
t_{\min}=\left|\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}\right|,$
which, accurate to multiplier $\Delta t_{\min}$ coincides with
$\left|v_{i}\left(t_{n}\right)\right|$ (see (9)), is commonly named absolute
returns, while dispersion of a random value
$\ln\left({X_{i}(t_{n+1})\mathord{\left/{\vphantom{X_{i}(t_{n+1})X_{i}(t_{n})}}\right.\kern-1.2pt}X_{i}(t_{n})}\right)$,
which differs from $D_{v_{i}}$ by $(\Delta t_{\min})^{2}$ (see (13)) –
volatility.
The chaotic nature of real time series allows to $x_{i}(t_{n})$ as the
trajectory of a certain abstract quantum particle (observed at $\Delta
t_{\min}$ time spans). Analogous to (1) we can write an uncertainty ratio for
this trajectory:
$\Delta x_{i}\cdot\Delta v_{i}\sim\frac{h}{m_{i}},$ (16)
or, taking into account (11) and (15):
$\frac{1}{\Delta
t_{\min}}\left(<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)\sim\frac{h}{m_{i}},$ (17)
where $m_{i}$ \- economic “mass” of an $i$ series, $h$ \- value which comes as
an economic Planck’s constant.
Having rewritten the ration 17:
$\Delta t_{\min}\cdot\frac{m_{i}}{(\Delta
t_{\min})^{2}}\left(<\ln^{2}\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+1})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)\sim h$ (18)
and interpreting the multiplier by $\Delta t_{\min}$ in the left part as the
uncertainty of an “economical” energy (accurate to the constant multiplier),
we get an economic analog of the ratio (4).
Since the analogy with physical particle trajectory is merely formal, $h$
value, unlike the physical Planck’s constant $\hbar$, can, generally speaking,
depend on the historical period of time, for which the series are taken, and
the length of the averaging interval (e.g. economical processes are different
in the time of crisis and recession), on the series number $i$ etc. Whether
this analogy is correct or not depends on particular series’ roperties.
Let us generalize the ratios (17), (18) for the case, when economic
measurements on the time span $T$, used to derive the series (6), are
conducted with the time step $\Delta t=k\cdot\Delta t_{\min}$, where $k\geq 1$
\- is a certain given integer positive number. From the formal point of view
it would mean that all terms, apart from those numbered $n=0,k,2k,3k,...$. are
discarded from the initial series (6). As a result the ratios would be the
following:
$\frac{1}{k\Delta
t_{\min}}\left(<\ln^{2}\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)\sim\frac{h}{m_{i}},$ (19) $k\Delta
t_{\min}\cdot\frac{1}{(k\Delta
t_{\min})^{2}}\left(<\ln^{2}\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}>_{n,\Delta
N}-\left(<\ln\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}>_{n,\Delta
N}\right)^{2}\right)\sim\frac{h}{m_{i}}$ (20)
and would be dependent on $k$.
Let us proceed to the analysis of the acquired results, that have to be
considered as intermediate.
In case of $h=const$, the formal analogy with the physical particle would be
complete, and in this case, as appears from (19), variance of a random
$i$-numbered value:
$\ln\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}\approx\frac{X_{i}(t_{n+k})-X_{i}(t_{n})}{X_{i}(t_{n})}$
\- practically coinciding with the relative increment of terms of the $i$
initial series – would keep increasing in a linear way with $k\Delta t_{\min}$
(interval between the observations) growing. Such dynamics is peculiar to the
series with statistically independent increments.
However, both in cases of a real physical particle and its formal economic
analogue any kind of change influences on the result. Therefore statistic
properties of the “thinned” series, used to create the ratio (19), have to
depend on real measurements in the intermediate points if there were any.
Besides, presence of “long” and “heavy” “tails” increasing along the amplitude
with decreasing $\Delta t$ on distributions of corresponding returns ${\Delta
X\mathord{\left/{\vphantom{\Delta XX}}\right.\kern-1.2pt}X}$, are in our
opinion the evidence of this thesis (see for example r031SaptsinMerezhPrirody
).
Thus, generalizing everything said above,
${h\mathord{\left/{\vphantom{hm_{i}}}\right.\kern-1.2pt}m_{i}}$ratio on the
right side of (19) (or (20)) has to be considered a certain unknown function
of the series number $i$, size of the averaging window$\Delta N$, time
$\bar{n}$ (centre of the averaging window), and time step of the observation
(registration) $k$.
To get at least an approximate, yet obvious, formula of this function and
track the nature of dependencies, we postulate the following model
presentation of the right side (19):
$\frac{h}{m_{i}}\simeq\frac{\tau\left(\bar{n},\Delta N_{\tau}\right)\cdot
H_{i}\left(k,\bar{n},\Delta N_{H}\right)}{\Delta t_{\min}\cdot m_{i}},$ (21)
where
$\frac{1}{m_{i}}=\begin{array}[]{c}{}\hfil\end{array}<\varphi_{i}\left(n,1\right)>_{(0\leq
n\leq N-2)},$ (22)
$m_{i}$ is a non-dimentional economic mass of an $i$-numbered series,
$\tau\left(\bar{n}\right)=\frac{<\varphi_{i}\left(n,1,\Delta
N_{\tau}\right)>_{(\bar{n}-{\Delta N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2}}\right.\kern-1.2pt}2}}}\right.\kern-1.2pt}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2}}\right.\kern-1.2pt}2}}),\begin{array}[]{c}{}\hfil\end{array}(1\leq
i\leq M)}}{<\left(<\varphi_{i}\left(n,1,\Delta
N_{\tau}\right)>_{(\bar{n}-{\Delta N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2}}\right.\kern-1.2pt}2}}}\right.\kern-1.2pt}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{\tau}\mathord{\left/{\vphantom{\Delta
N_{\tau}2}}\right.\kern-1.2pt}2}}),\begin{array}[]{c}{}\hfil\end{array}(1\leq
i\leq M)}\right)>_{\bar{n}}}$ (23)
\- local physical time compression ($\tau\left(\bar{n}\right)<1$) or
magnification ($\tau\left(\bar{n}\right)>1$) ratio, which allows to introduce
the notion of heterogenous economic time (for a homogenous
$\tau\left(\bar{n}\right)=1$),
$H_{i}\left(k,\bar{n}\right)=\frac{<\varphi_{i}\left(n,k,\Delta
N_{H}\right)>_{\bar{n}-{\Delta N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2}}\right.\kern-1.2pt}2}}}\right.\kern-1.2pt}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2}}\right.\kern-1.2pt}2}}}}{<\varphi_{i}\left(n,1,\Delta
N_{H}\right)>_{\bar{n}-{\Delta N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2}}\right.\kern-1.2pt}2}}}\right.\kern-1.2pt}2\begin{array}[]{c}{}\hfil\end{array}<n<\begin{array}[]{c}{}\hfil\end{array}\bar{n}+{\Delta
N_{H}\mathord{\left/{\vphantom{\Delta
N_{H}2}}\right.\kern-1.2pt}2}}}};\begin{array}[]{c}{}\hfil\end{array}k=1,2,...k_{\max}$
(24)
\- non-dimentional coefficient of the order of unit, which indicates
differences in the dependence of variance $D_{\Delta x_{i}}$ (see (13) taking
into account the case of $k\geq 1$) on the law $D_{\Delta x_{i}}\sim k$ for
the given $i$ and $\bar{n}$.
$\varphi_{i}\left(n,k,\tilde{N}\right)=\frac{1}{k}\left(\ln^{2}\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}-\left(<\ln\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})}>_{n,\tilde{N}}\right)^{2}\right)$
(25)
(index $\tilde{N}=N,\Delta N_{\tau},\Delta N_{H}$ in the last formula
indicates the averaging parameters according to $n$ and formulae
(22),(23),(24), averaging windows $\Delta N_{\tau},\Delta N_{H}$ are chosen
with thew following the conditions taken into consideration:
$k_{\max}<\Delta N_{\tau}<\Delta N_{H}<N.$ (26)
According to the definitions (23),(24) for coefficients
$\tau\left(\bar{n}\right)$ and $H_{i}\left(k,\bar{n}\right)$ following
conditions of the normalization take place:
$<\tau\left(\bar{n}\right)>_{\bar{n},N}=1;\begin{array}[]{c}{}\hfil\end{array}H_{i}\left(1,\bar{n}\right)=1,$
(27)
and the multiplier ${1\mathord{\left/{\vphantom{1\Delta
t_{\min}}}\right.\kern-1.2pt}\Delta t_{\min}}$ on the right side (21) can be
considered as an invariant component of an economic Planck’s constant $h$:
$\bar{h}={1\mathord{\left/{\vphantom{1\Delta
t_{\min}}}\right.\kern-1.2pt}\Delta t_{\min}},$ (28)
As you can see, $\bar{h}$ has a natural dimension ¡time¿ to the negative first
power.
It is also worth noting that average economic mass of the whole set of series
(or any separate group of the series) can be introduced with the help of the
following formula:
$\frac{1}{m}=\frac{1}{M}\sum_{i=1}^{M}\frac{1}{m_{i}}.$ (29)
Acquired with the help of series (6) ratios (7,(19)-(28) also allow different
interpretations. For example, it can be considered that normalized series (7)
depict the trajectory of a certain hypothetical economic quantum quasi-
particle in an abstract $M$-dimensional space of economic indices, and
$m_{i}^{-1}$ are the main components of inverse mass tensor of the quasi-
particle (the analogy with quasi-particles as free carriers of the electric
charge in semiconductors), which has already been used in the previous
chapter.
In the final part of this chapter we would like to pay attention to the chosen
variant of the theory, which is probably the simplest one, because of the
following reasons.
Carrying out various $n$ (discrete time) and $i$ (series number) averagings,
we didn’t take into account at least two fairly important factors: 1) amounts
of financial and material resources (their movement is reflected by each
series) and 2) possible correlation between the series.
However generalization of the theory and introduced notions is not so
difficult in this case. It is enough to form a row
$\left(\alpha_{1},\alpha_{2},...\alpha_{M}\right)$ of positive weight
coefficients with the following condition of normalization:
$\sum_{i=1}^{M}\alpha_{i}=M,$ (30)
with each of them taking into account the importance of separate series in
terms of a certain criterion, while for random values
$\phi_{i}\left(n,k\right)=\sqrt{\frac{\alpha_{i}}{k}}\ln\frac{X_{i}(t_{n+k})}{X_{i}(t_{n})},\begin{array}[]{c}{}\hfil\end{array}i=1,2,...M$
(31)
instead of a one-dimensional massive (25) we should introduce a covariance
matrix:
$\Psi=\left[\psi_{ij}\right],$ (32)
where
$\psi_{ij}=\psi_{ij}\left(k,\tilde{N}\right)=\begin{array}[]{c}{}\hfil\end{array}<\left(\phi_{i}\left(n,k\right)-\bar{\phi}_{i}\left(n,k\right)\right)\cdot\left(\phi_{j}\left(n,k\right)-\bar{\phi}_{j}\left(n,k\right)\right)>_{n,\tilde{N}},$
(33)
$\bar{\phi}_{i}\left(j,k,\tilde{N}\right)=\begin{array}[]{c}{}\hfil\end{array}<\phi_{i}\left(n,k\right)>_{n,\tilde{N}}$
(34)
(with $\alpha_{i}=1$ and absence of correlations
$\psi_{ij}=\varphi_{i}\delta_{ij}$). Using a standard algorithm of
characteristic constants
$\lambda_{i},\begin{array}[]{c}{}\hfil\end{array}i=1,2,...M$ and corresponding
orthonormal vectors $C_{i}=\left(c_{i1},c_{i2},...c_{iM}\right)$ search in
$\Psi$ matrix, we proceed to the new basis, where “renormalized” series
$y_{i}(t_{n})=\sum_{j=1}^{M}c_{ij}x_{j}\left(t_{n}\right)$ (new basis vectors)
aren’t correlated any more. However the presence of zero characteristic
constants or $\lambda_{i}$, which are distinguished with relatively low values
in absolute magnitude, will mean that the real dimension of the set of series
(7) is in fact less than $M$ (initial series (7) or their parts are strongly
correlated). In this case renormalized series $y_{i}(t_{n})$ with zero or low
characteristic constants have to be discarded. The remaining renormalized
series will undergo all above-listed procedures.
## 4 Experimental results and their discussion
To test the suggested ratios and definitions we have chosen 9 economic series
with $\Delta t_{\min}$ in one day for the period from April 27, 1993 to March
31, 2010. The chosen series correspond to the following groups that differ in
their origin:
1) stock market indices: USA (S&P500), Great Britain (FTSE 100) and Brazil
(BVSP);
2) currency dollar cross-rates (chf, jpy, gbp);
3) commodity market (gold, silver, and oil prices).
On Fig. 1-3 normalized plots of the corresponding series, divided by groups,
are introduced, while $\Delta t_{\min}$is taken equal to the unit.
As you can see from the Fig. 1-3, all time series include visually noticeable
chaotic component and obviously differ from each other, which allows us to
hope for the successful application the afore-mentioned theory to the
interpretation and analysis of real series. Let us confine to its elementary
variant.
As an example on fig. 4 we suggest absolute values of immediate speeds (or
absolute returns according to the general terminology used in literature),
calculated with the help of the formula evaluation (9), and their variance
(volatility), calculated with the help of the formula evaluation (13) for the
series of Japanese yen (jpy) US dollar cross-rates.
As we can see from the plots, the dependence of immediate speed or returns on
time is of chaotic nature, while the dependence of volatility is smooth but
not monotonous. For the rest of initial series, the dependencies of volatility
and returns are similar to the depicted on the fig. 4 ones.
Figure 1: USA (S&P500), Great Britain (FTSE 100), and Brazil (BVSP) daily
stock indices from April 27, 1993 to March 31, 2010.
Figure 2: Daily currency dollar cross-rates (chf, jpy, gbp)
Figure 3: Commodity market. Daily gold, silver, and oil prices.
Figure 4: Absolute values of immediate speeds (abs returns) and their
dispersions (volatility).
Fig. 5 shows averaged coefficients of time $\tau(t)$ compression-expansion
(formula (23)) for three groups of incoming series: currency (forex), stock,
and commodity markets.
Figure 5: Coefficients of time compression-expansion, market “temperature”.
The explanation is in the text.
The formulae (11),(23),(25) show that $\tau$(t) exists in proportion to the
averaged square speed (according to the chosen time span and series), i.e.
average “energy” of the economical “particle” (as it is in our analogy), and
can be thus interpreted as the series “temperature”. Crises are distinguished
with the intensification of economic processes (the “temperature” is rising),
while during the crisis-free period their deceleration can be observed (the
“temperature” is falling), what can be interpreted as the heterogenous flow of
economic time. $\tau(t)$ dependences shown on the fig. 5 illustrate all afore-
mentioned. Note that local time acceleration-deceleration can be rather
significant.
Transition to heterogenous economic time allows to make the observed economic
series more homogenous, which can simplify both analysis and prediction r042 .
In table we give the values of a non-dimentional economic mass of the $m_{i}$
series, calculated using (22) for all 9 incoming series, as well as average
masses of each group (formula (29)).
Table. Economic series masses
Incoming series | Economic mass | Average economic mass of the group
---|---|---
Commodity market | gold | $2,816\cdot 10^{4}$ | $4,983\cdot 10^{3}$
silver | $4,843\cdot 10^{3}$
oil | $2,777\cdot 10^{3}$
Currency market | jpy | $2,148\cdot 10^{4}$ | $2,499\cdot 10^{4}$
gbp | $3,523\cdot 10^{4}$
chf | $2,180\cdot 10^{4}$
Stock market | S&P 500 | $6,251\cdot 10^{3}$ | $4,748\cdot 10^{3}$
FTSE 100 | $6,487\cdot 10^{3}$
BVSP | $1,507\cdot 10^{3}$
As you can see from the table, the stock market is distinguished with the
lowest mass value, while the currency one shows the maximum number. Oil price
series has the lowest mass on the commodity market, gold – the highest one. As
for the currency market, British pound (gbp) have the highest value and
Japanese yen rates (jpy) demonstrates the minimum mass of the group, although
the dispersion is lower than that of the commodity market. The smallest spread
is peculiar to the currency market. Dynamic and developing Brazilian market
(BVSP) has the lowest mass, while the maximum value, just like in the previous
case, corresponds to Great Britain (FTSE 100). It is explained by the well-
known fact: Britain has been always known for its relatively “closed” economy
as comraped with the rest of the European and non-European countries.
The last group of experimental data corresponds to the dependence of Planck’s
economic constant (calculated for different series) on time $\Delta t=k\Delta
t_{\min}$ (time between the neighbouring registered observations), which is
characterised by $H_{i}\left(k,\bar{n}\right)$ coefficient (see formula (24)).
On fig. 6-8 integral dependencies $H_{i}\left(k\right)$ are depicted. The
following are averaged on the whole period of time 1993-2010 and calculated
for commodity, stock, and currency markets. As you can see there are no
obvious regularities, which can be explained by various crises and recessions
of the world and national economies that took place during the investigated
period.
To decide whether it is possible for local regularities of Planck’s economic
constant dependence on $\Delta t$ to appear, we have chosen relatively small
averaging fragments, $\Delta N=500$, which approximately equals two years.
Corresponding results for some of these fragments on commodity, currency, and
stock markets are given on fig. 9-11. Evidently, all three figures show clear
tendencies of $H_{i}\left(k,\bar{n}\right)$ recession and rise for each type
of the market (unlike integral dependences $H_{i}\left(k\right)$).
Figure 6: Integral coefficient $H_{i}\left(k\right)$ dependences for commodity
market.
Figure 7: Integral coefficient $H_{i}\left(k\right)$ dependences for stock
market.
Figure 8: Integral coefficient $H_{i}\left(k\right)$ dependences for currency
market.
Figure 9: Local coefficient $H_{i}\left(k,\bar{n}\right)$ dependeces for
commodity market (averaging time span from 27.04.1993 to 12.06.1995, 500 daily
values).
Figure 10: Local coefficient $H_{i}\left(k,\bar{n}\right)$ for currency market
(averaging time span from 12.06.1995 to 15.07.1997, 500 daily values).
Figure 11: Local coefficient $H_{i}\left(k,\bar{n}\right)$ dependences for
stock market (averaging time span from 12.06.1995 to 15.07.1997, 500 daily
values).
## 5 Conclusions
We have conducted methodological and philosophical analysis of physical
notions and their formal and informal connections with real economic
measurements. Basic ideas of the general relativity theory and relativistic
qantum mechanics concerning spacetime properties and physical dimensions
peculiarities were used as well. We have suggested procedures of detecting
normalized economic coordinates, economic mass and heterogenous economic time.
The afore-mentioned procedures are based on socio-economic time series
analysis and economical interpretation of Heisenberg’s uncertainty principle.
The notion of economic Planck’s constant has also been introduced. The theory
has been tested on real economic time series, including stock indices,
currency rates, and commodity prices. Acquired results indicate the
availability of further investigations.
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|
arxiv-papers
| 2011-11-10T17:47:21 |
2024-09-04T02:49:24.588977
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Vladimir Soloviev and Vladimir Saptsin",
"submitter": "Vladimir Saptsin",
"url": "https://arxiv.org/abs/1111.5289"
}
|
1111.5340
|
# On the Expected Complexity of Random Convex Hulls††thanks: This work has
been supported by a grant from the U.S.–Israeli Binational Science Foundation.
This work is part of the author’s Ph.D. thesis, prepared at Tel-Aviv
University under the supervision of Prof. Micha Sharir.
Sariel Har-Peled Department of Computer Science; University of Illinois; 201
N. Goodwin Avenue; Urbana, IL, 61801, USA; sariel@uiuc.edu;
http://www.uiuc.edu/~sariel/.
###### Abstract
In this paper we present several results on the expected complexity of a
convex hull of $n$ points chosen uniformly and independently from a convex
shape.
(i) We show that the expected number of vertices of the convex hull of $n$
points, chosen uniformly and independently from a disk is $O(n^{1/3})$, and
$O(k\log{n})$ for the case a convex polygon with $k$ sides. Those results are
well known (see [RS63, Ray70, PS85]), but we believe that the elementary proof
given here are simpler and more intuitive.
(ii) Let ${\cal D}$ be a set of directions in the plane, we define a
generalized notion of convexity induced by ${\cal D}$, which extends both
rectilinear convexity and standard convexity.
We prove that the expected complexity of the ${\cal D}$-convex hull of a set
of $n$ points, chosen uniformly and independently from a disk, is
$O\left({n^{1/3}+\sqrt{n\alpha({\cal D})}}\right)$, where $\alpha({\cal D})$
is the largest angle between two consecutive vectors in ${\cal D}$. This
result extends the known bounds for the cases of rectilinear and standard
convexity.
(iii) Let ${\cal B}$ be an axis parallel hypercube in ${\rm I\\!R}^{d}$. We
prove that the expected number of points on the boundary of the quadrant hull
of a set $S$ of $n$ points, chosen uniformly and independently from ${\cal B}$
is $O(\log^{d-1}n)$. Quadrant hull of a set of points is an extension of
rectilinear convexity to higher dimensions. In particular, this number is
larger than the number of maxima in $S$, and is also larger than the number of
points of $S$ that are vertices of the convex hull of $S$.
Those bounds are known [BKST78], but we believe the new proof is simpler.
## 1 Introduction
Let $C$ be a fixed compact convex shape, and let $X_{n}$ be a random sample of
$n$ points chosen uniformly and independently from $C$. Let $Z_{n}$ denote the
number of vertices of the convex hull of $X_{n}$. Rényi and Sulanke [RS63]
showed that $E[Z_{n}]=O(k\log{n})$, when $C$ is a convex polygon with $k$
vertices in the plane. Raynaud [Ray70] showed that expected number of facets
of the convex hull is $O(n^{(d-1)/(d+1)})$, where $C$ is a ball in ${\rm
I\\!R}^{d}$, so $E[Z_{n}]=O(n^{1/3})$ when $C$ is a disk in the plane. Raynaud
[Ray70] showed that the expected number of facets of
${\mathop{\mathrm{CH}}}(X_{n})=ConvexHull(X_{n})$ is
$O\left({(\log(n))^{(d-1)/2}}\right)$, where the points are chosen from ${\rm
I\\!R}^{d}$ by a $d$-dimensional normal distribution. See [WW93] for a survey
of related results.
All these bounds are essentially derived by computing or estimating integrals
that quantify the probability of two specific points of $X_{n}$ to form an
edge of the convex hull (multiplying this probability by $\binom{n}{2}$ gives
$E[Z_{n}]$). Those integrals are fairly complicated to analyze, and the
resulting proofs are rather long, counter-intuitive and not elementary.
Efron [Efr65] showed that instead of arguing about the expected number of
vertices directly, one can argue about the expected area/volume of the convex
hull, and this in turn implies a bound on the expected number of vertices of
the convex hull. In this paper, we present a new argument on the expected
area/volume of the convex hull (this method can be interpreted as a discrete
approximation to the integral methods). The argument goes as follows:
Decompose $C$ the into smaller shapes (called tiles). Using the topology of
the tiling and the underlining type of convexity, we argue about the expected
number of tiles that are exposed by the random convex hull, where a tile is
exposed if it does not lie completely in the interior of the random convex
hull. Resulting in a lower bound on the area/volume of the random convex hull.
We apply this technique to the standard case, and also for more exotic types
of convexity.
In Section 2, we give a rather simple and elementary proofs of the
aforementioned bounds $E[Z_{n}]=O(n^{1/3})$ for $C$ a disk, and
$E[Z_{n}]=O\left({k\log{n}}\right)$ for $C$ a convex $k$-gon. We believe that
these new elementary proofs are indeed simpler and more intuitive111Preparata
and Shamos [PS85, pp. 152] comment on the older proof for the case of a disk:
“Because the circle has no corners, the expected number of hull vertices is
comparatively high, although we know of no elementary explanation of the
$n^{1/3}$ phenomenon in the planar case.” It is the author’s belief that the
proof given here remedies this situation. than the previous integral-based
proofs.
The question on the expected complexity of the convex hull remains valid, even
if we change our type of convexity. In Section 3, we define a generalized
notion of convexity induced by ${\cal D}$, a given set of directions. This
extends both rectilinear convexity, and standard convexity. We prove that the
expected complexity of the ${\cal D}$-convex hull of a set of $n$ points,
chosen uniformly and independently from a disk, is
$O\left({n^{1/3}+\sqrt{n\alpha({\cal D})}}\right)$, where $\alpha({\cal D})$
is the largest angle between two consecutive vectors in ${\cal D}$. This
result extends the known bounds for the cases of rectilinear and standard
convexity.
Finally, in Section 4, we deal with another type convexity, which is an
extension of the generalized convexity mentioned above for the higher
dimensions, where the set of the directions is the standard orthonormal basis
of ${\rm I\\!R}^{d}$. We prove that the expected number of points that lie on
the boundary of the quadrant hull of $n$ points, chosen uniformly and
independently from the axis-parallel unit hypercube in ${\rm I\\!R}^{d}$, is
$O(\log^{d-1}n)$. This readily imply $O(\log^{d-1}n)$ bound on the expected
number of maxima and the expected number of vertices of the convex hull of
such a point set. Those bounds are known [BKST78], but we believe the new
proof is simpler and more intuitive.
## 2 On the Complexity of the Convex Hull of a Random Point Set
In this section, we show that the expected number of vertices of the convex
hull of $n$ points, chosen uniformly and independently from a disk, is
$O(n^{1/3})$. Applying the same technique to a convex polygon with $k$ sides,
we prove that the expected number of vertices of the convex hull is
$O(k\log{n})$.222As already noted, these results are well known ([RS63, Ray70,
PS85]), but we believe that the elementary proofs given here are simpler and
more intuitive. The following lemma, shows that the larger the expected area
outside the random convex hull, the larger is the expected number of vertices
of the convex hull.
###### Lemma 2.1
Let $C$ be a bounded convex set in the plane, such that the expected area of
the convex hull of $n$ points, chosen uniformly and independently from $C$, is
at least $\left({1-f(n)}\right)Area(C)$, where $1\geq f(n)\geq 0$, for $n\geq
0$. Then the expected number of vertices of the convex hull is $\leq nf(n/2)$.
###### Proof.
Let $N$ be a random sample of $n$ points, chosen uniformly and independently
from $C$. Let $N_{1}$ (resp. $N_{2}$) denote the set of the first (resp. last)
$n/2$ points of $N$. Let $V_{1}$ (resp. $V_{2}$) denote the number of vertices
of $H={\mathop{CH}}(N_{1}\cup N_{2})$ that belong to $N_{1}$ (resp. $N_{2}$),
where ${\mathop{CH}}(N_{1}\cup N_{2})={\mathrm{ConvexHull}}(N_{1}\cup N_{2})$.
Clearly, the expected number of vertices of $C$ is $E[V_{1}]+E[V_{2}]$. On the
other hand,
$E\left[{V_{1}\,\left|\,{N_{2}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq\frac{n}{2}\left({\frac{Area(C)-Area({\mathop{CH}}(N_{2}))}{Area(C)}}\right),$
since $V_{1}$ is bounded by the expected number of points of $N_{1}$ falling
outside ${\mathop{CH}}(N_{2})$.
We have
$\displaystyle E[V_{1}]$ $\displaystyle=$ $\displaystyle
E_{N_{2}}\left[{E[V_{1}|N_{2}]{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq
E\left[{\frac{n}{2}\left({\frac{Area(C)-Area({\mathop{CH}}(N_{2}))}{Area(C)}}\right){\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]$
$\displaystyle\leq$ $\displaystyle\frac{n}{2}f(n/2),$
since $E[X]=E_{Y}[E[X|Y]]$ for any two random variables $X,Y$. Thus, the
expected number of vertices of $H$ is $E[V_{1}]+E[V_{2}]\leq nf(n/2)$. ∎
###### Remark 2.2
Lemma 2.1 is known as Efron’s Theorem. See [Efr65].
###### Theorem 2.3
The expected number of vertices of the convex hull of $n$ points, chosen
uniformly and independently from the unit disk, is $O(n^{1/3})$.
###### Proof.
We claim that the expected area of the convex hull of $n$ points, chosen
uniformly and independently from the unit disk, is at least
$\pi-O\left({n^{-2/3}}\right)$.
Indeed, let $D$ denote the unit disk, and assume without loss of generality,
that $n=m^{3}$, where $m$ is a positive integer. Partition $D$ into $m$
sectors, ${\cal S}_{1},\ldots,{\cal S}_{m}$, by placing $m$ equally spaced
points on the boundary of $D$ and connecting them to the origin. Let
$D_{1},\ldots,D_{m^{2}}$ denote the $m^{2}$ disks centered at the origin, such
that (i) $D_{1}=D$, and (ii) $Area(D_{i-1})-Area(D_{i})=\pi/m^{2}$, for
$i=2,\ldots,m^{2}$. Let $r_{i}$ denote the radius of $D_{i}$, for
$i=1,\ldots,m^{2}$.
Let $S_{i,j}=(D_{i}\setminus D_{i+1})\cap{\cal S}_{j}$, and
$S_{m^{2},j}=D_{m^{2}}\cap{\cal S}_{j}$, for $i=1,\ldots,m^{2}-1$,
$j=1,\ldots,m$. The set $S_{i,j}$ is called the $i$-th tile of the sector
${\cal S}_{j}$, and its area is $\pi/n$, for $i=1,\ldots,m^{2}$,
$j=1,\ldots,m$.
Let $N$ be a random sample of $n$ points chosen uniformly and independently
from $D$. Let $X_{j}$ denote the first index $i$ such that $N\cap
S_{i,j}\neq\emptyset$, for $j=1,\ldots,m$. For a fixed
$j\in\left\\{{1,\ldots,m}\right\\}$, the probability that $X_{j}=k$ is upper-
bounded by the probability that the tiles $S_{1,j},\ldots,S_{(k-1),j}$ do not
contain any point of $N$; namely, by $\left({1-\frac{k-1}{n}}\right)^{n}$.
Thus, $P[X_{j}=k]\leq\left({1-\frac{k-1}{n}}\right)^{n}\leq e^{-(k-1)}$, since
$1-x\leq e^{-x}$, for $x\geq 0$. Thus,
$E\left[{X_{j}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=\sum_{k=1}^{m^{2}}kP[X_{j}=k]\leq\sum_{k=1}^{m^{2}}ke^{-(k-1)}=O(1),$
for $j=1,\ldots,m$.
Let $K_{o}$ denote the convex hull of $N\cup\left\\{{o}\right\\}$, where $o$
is the origin. The tile $S_{i,j}$ is exposed by a set $K$, if
$S_{i,j}\setminus K\neq\emptyset$. We claim that at most
$X_{j-1}+X_{j+1}+O(1)$ tiles are exposed by $K_{o}$ in the sector ${\cal
S}_{j}$, for $j=1,\ldots,m$ (where we put $X_{0}=X_{m}$, $X_{m+1}=X_{1}$).
$s$${\cal S}_{j+1}$${\cal S}_{j-1}$${\cal S}_{j}$$p$$q$$o$$T$${\cal S}_{i,j}$
Figure 1: Illustrating the proof that bounds the number of tiles exposed by
$T$ inside ${\cal S}_{j}$
Indeed, let $w=w(N,j)=max(X_{j-1},X_{j+1})$, and let $p,q$ be the two points
in $S_{j-1,w},S_{j+1,w}$, respectively, such that the number of sets exposed
by the triangle $T=\triangle{opq}$, in the sector ${\cal S}_{i}$, is maximal.
Both $p$ and $q$ lie on ${\partial}{D_{w+1}}$ and on the external radii
bounding ${\cal S}_{j-1}$ and ${\cal S}_{j+1}$, as shown in Figure 1. Clearly,
any tile which is exposed in ${\cal S}_{j}$ by $K_{o}$ is also exposed by $T$.
Let $s$ denote the segment connecting the middle of the base of $T$ to its
closest point on ${\partial}{D_{w}}$. The number of tiles in ${\cal S}_{j}$
exposed by $T$ is bounded by $\max\left({X_{j-1},X_{j+1}}\right)$, plus the
number of tiles intersecting the segment $s$. The length of $s$ is
$|oq|-|oq|\cos\left({\frac{3}{2}\cdot\frac{2\pi}{m}}\right)\leq
1-\cos\left({\frac{3}{2}\cdot\frac{2\pi}{m}}\right)\leq\frac{1}{2}\left({\frac{3\pi}{m}}\right)^{2}=\frac{4.5\pi^{2}}{m^{2}},$
since $\cos(x)\geq 1-x^{2}/2$, for $x\geq 0$.
On the other hand, $r_{i+1}-r_{i}\geq r_{i}-r_{i-1}\geq 1/(2m^{2})$, for
$i=2,\ldots,m^{2}$. Thus, the segment $s$ intersects at most
$\left\lceil{||s||/(1/(2m^{2}))}\right\rceil=\left\lceil{9\pi^{2}}\right\rceil=89$
tiles, and we have that the number of tiles exposed in the sector ${\cal
S}_{i}$ by $K_{o}$ is at most $\max\left({X_{j-1},X_{j+1}}\right)+89\leq
X_{j-1}+X_{j+1}+89$, for $j=1,\ldots,m$.
Thus, the expected number of tiles exposed by $K_{o}$ is at most
$E\left[{\sum_{i=1}^{m}\left({X_{j-1}+X_{j+1}+89}\right){\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=O(m).$
The area of $K={\mathop{CH}}(N)$ is bounded from below by the area of tiles
which are not exposed by $K$. The probability that $K\subsetneq K_{o}$
(namely, the origin is not inside $K$, or, equivalently, all points of $N$ lie
in some semidisk) is at most $2\pi/2^{n}$, as easily verified. Hence,
$E[Area(K)]\geq E[Area(C)]-P\left[{C\neq
K{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\pi=\pi-O(m)\frac{\pi}{n}-\frac{2\pi}{2^{n}}\pi=\pi-O\left({n^{-2/3}}\right).$
The assertion of the theorem now follows from Lemma 2.1. ∎
###### Lemma 2.4
The expected number of vertices of the convex hull of $n$ points, chosen
uniformly and independently from the unit square, is $O(\log{n})$.
$R_{5}(1)$$R_{5}(2)$$R_{5}(3)$$R_{5}(4)$$R_{5}^{\prime}(1)$$R_{5}^{\prime}(2)$$R_{5}^{\prime}(3)$$R_{5}^{\prime}(4)$$R_{5}^{\prime}(5)$$R_{5}^{\prime}(6)$
Figure 2: Illustrating the proof that bounds the number of tiles exposed by
${\mathop{CH}}(N)$ inside the $j$-th column, by using a non-uniform tiling of
the strips to the left and to the right of the $j$-th column. The area of such
a larger tile is at least $1/n$.
###### Proof.
We claim that the expected area of the convex hull of $n$ points, chosen
uniformly and independently from the unit square, is at least
$1-O\left({\log(n)/n}\right)$.
Let $S$ denote the unit square. Partition $S$ into $n$ rows and $n$ columns,
such that $S$ is partitioned into $n^{2}$ identical squares. Let
$S_{i,j}=[(i-1)/n,i/n]\times[(j-1)/n,j/n]$ denote the $j$-th square in the
$i$-th column, for $1\leq i,j\leq n$. Let ${\cal S}_{i}=\cup_{j=1}^{n}S_{i,j}$
denote the $i$-th column of $S$, for $i=1,\ldots,n$, and let ${\cal
S}(l,k)=\cup_{i=l}^{k}{\cal S}_{i}$, for $1\leq l\leq k\leq n$.
Let $N$ be a random sample of $n$ points chosen uniformly and independently
from $S$. Let $X_{j}$ denote the first index $i$ such that
$N\cap(\cup_{l=1}^{j-1}S_{l,i})\neq\emptyset$, for $j=2,\ldots,n-1$; namely,
$X_{j}$ is the index of the first row in ${\cal S}(1,j-1)$ that contains a
point from $N$. Symmetrically, let $X_{j}^{\prime}$ be the index of the first
row in ${\cal S}(j+1,n)$ that contains a point of $N$. Clearly,
$E[X_{j}]=E[X_{n-j+1}^{\prime}]$, for $j=2,\ldots,n-1$.
Let $Z_{j}$ denote the number of squares $S_{i,j}$ in the bottom of the $j$-th
column that are exposed by ${\mathop{CH}}(N)$, for $j=2,\ldots,n-1$. Arguing
as in the proof of Theorem 2.3, we have that
$Z_{j}\leq\max(X_{j},X_{j}^{\prime})\leq X_{j}+X_{j}^{\prime}$. Thus, in order
to bound $E[Z_{j}]$, we first bound $E[X_{j}]$ by covering the strips ${\cal
S}(1,j-1),{\cal S}(j+1,n)$ by tiles of area $\geq 1/n$. In particular, let
$h(l)=\left\lceil{n/(l-1)}\right\rceil$, and let
$R_{j}(m)=[0,(j-1)/n]\times[h(n-j+1)(m-1)/n,h(j)m/n]$, and let
$R_{j}^{\prime}(m)=[(j+1)/n,1]\times[h(j)(m-1)/n,h(j)m/n]$, for
$j=2,\ldots,n-1$. See Figure 2.
Let $Y_{j}$ denote the minimal index $i$ such that $R_{j}(i)\cap
N\neq\emptyset$. The area of $R_{j}(i)$ is at least $1/n$, for any $i$ and
$j$. Arguing as in the proof of Theorem 2.3, it follows that $E[Y_{j}]=O(1)$.
On the other hand, $E[X_{j}]\leq h(j)E[Y_{j}]=O(n/(j-1))$. Symmetrically,
$E[X_{j}^{\prime}]=O(n/(n-j))$.
Thus, by applying the above argument to the four directions (top, bottom,
left, right), we have that the expected number of squares $S_{i,j}$ exposed by
${\mathop{CH}}(N)$ is bounded by
$4n-4+4{\sum_{j=2}^{n-1}E[Z_{j}]}<4n+4{\sum_{j=2}^{n-1}(E[X_{j}]+E[X_{j}^{\prime}])}=4n+8{\sum_{j=2}^{n-1}O\left({\frac{n}{j-1}}\right)}=O(n\log{n}),$
where $4n-4$ is the number of squares adjacent to the boundary of $S$.
Since the area of each square is $1/n^{2}$, it follows that the expected area
of ${\mathop{CH}}(N)$ is at least $1-O(\log(n)/n)$.
By Lemma 2.1, the expected number of vertices of the convex hull is $O(\log
n)$. ∎
###### Lemma 2.5
The expected number of vertices of the convex hull of $n$ points, chosen
uniformly and independently from a triangle, is $O(\log{n})$.
###### Proof.
We claim that the expected area of the convex hull of $n$ points, chosen
uniformly and independently from a triangle $T$, is at least
$(1-O\left({\log(n)/n}\right))Area(T)$. We adapt the tiling used in Lemma 2.4
to a triangle. Namely, we partition $T$ into $n$ equal-area triangles, by
segments emanating from a fixed vertex, each of which is then partitioned into
$n$ equal-area trapezoids by segments parallel to the opposite side, such that
each resulting trapezoid has area $1/n^{2}$. See Figure 3.
Notice that this tiling has identical topology to the tiling used in Lemma
2.4. Thus, the proof of Lemma 2.4 can be applied directly to this case,
repeating the tiling process three times, once for each vertex of $T$. This
readily implies the asserted bound. ∎
$R_{5}(1)$$R_{5}(2)$$R_{5}(3)$$R_{5}^{\prime}(1)$$R_{5}^{\prime}(2)$$R_{5}^{\prime}(3)$$R_{5}^{\prime}(4)$$R_{5}^{\prime}(5)$
Figure 3: Illustrating the proof of Lemma 2.4 for the case of a triangle.
###### Theorem 2.6
The expected number of vertices of the convex hull of $n$ points, chosen
uniformly and independently from a polygon $P$ having $k$ sides, is
$O(k\log{n})$.
###### Proof.
We triangulate $P$ in an arbitrary manner into $k$ triangles
$T_{1},\ldots,T_{k}$. Let $N$ be a random sample of $n$ points, chosen
uniformly and independently from $P$. Let $Y_{i}=|T_{i}\cap N|$,
$N_{i}=T_{i}\cap N$, and $Z_{i}=|{\mathop{CH}}(N_{i})|$, for $i=1,\ldots,k$.
Notice that the distribution of the points of $N_{i}$ inside $T_{i}$ is
identical to the distribution of $Y_{i}$ points chosen uniformly and
independently from $T_{i}$. In particular, $E[Z_{i}|Y_{i}]=O(\log{Y_{i}})$, by
Lemma 2.5, and $E[Z_{i}]=E_{Y_{i}}[E[Z_{i}|Y_{i}]]=O(\log{n})$, for
$i=1,\ldots,k$.
Thus, $E[|{\mathop{CH}}(N)|]\leq
E\left[{\sum_{i=1}^{k}|{\mathop{CH}}(N_{i})|{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq\sum_{i=1}^{k}E[Z_{i}]=O(k\log{n})$.
∎
## 3 On the Expected Complexity of a Generalized Convex Hull Inside a Disk
In this section, we derive a bound on the expected complexity on a generalized
convex hull of a set of points, chosen uniformly and independently for the
unit disk. The new bound matches the known bounds, for the case of standard
convexity and maxima. The bound follows by extending the proof of Theorem 2.3.
We begin with some terminology and some initial observations, most of them
taken or adapted from [MP97]. A set ${\cal D}$ of vectors in the plane is a
set of directions, if the length of all the vectors in ${\cal D}$ is $1$, and
if $v\in{\cal D}$ then $-v\in{\cal D}$. Let ${\cal D}_{{\rm I\\!R}}$ denote
the set of all possible directions. A set $C$ is ${\cal D}$-convex if the
intersection of $C$ with any line with a direction in ${\cal D}$ is connected.
By definition, a set $C$ is convex (in the standard sense), if and only if it
is ${\cal D}_{{\rm I\\!R}}$-convex.
For a set $C$ in the plane, we denote by ${\mathop{\cal{CH}}}_{D}(C)$ the
${\cal D}$-convex hull of $C$; that is, the smallest ${\cal D}$-convex set
that contains $C$. While this seems like a reasonable extension of the regular
notion of convexity, its behavior is counterintuitive. For example, let ${\cal
D}_{Q}$ denote the set of all rational directions (the slopes of the
directions are rational numbers). Since ${\cal D}_{Q}$ is dense in ${\cal
D}_{{\rm I\\!R}}$, one would expect that ${\mathop{\cal{CH}}}_{{\cal
D}_{Q}}(C)={\mathop{\cal{CH}}}_{{\cal D}_{\rm
I\\!R}}(C)={\mathop{\mathrm{CH}}}(C)$. However, if $C$ is a set of points such
that the slope of any line connecting a pair of points of $C$ is irrational,
then ${\mathop{\cal{CH}}}_{D_{Q}}(C)=C$. See [OSSW85, RW88, RW87] for further
discussion of this type of convexity.
###### Definition 3.1
Let $f$ be a real function defined on a ${\cal D}$-convex set $C$. We say that
$f$ is ${\cal D}$-convex if, for any $x\in C$ and any $v\in{\cal D}$, the
function $g(t)=f(x+tv)$ is a convex function of the real variable $t$. (The
domain of $g$ is an interval in ${\rm I\\!R}$, as $C$ is assumed to be ${\cal
D}$-convex.)
Clearly, any convex function, in the standard sense, defined over the whole
plane satisfies this condition.
###### Definition 3.2
Let $C\subseteq{\rm I\\!R}^{2}$. The set ${\mathop{\cal{CH}}}^{\cal D}(C)$,
called the functional ${\cal D}$-convex hull of $C$, is defined as
${\mathop{\cal{CH}}}^{\cal D}(C)=\left\\{{x\in{\rm
I\\!R}^{2}\,\left|\,{f(x)\leq\sup_{y\in C}f(y)\text{ for all ${\cal D}$-convex
}f:{\rm I\\!R}^{2}\rightarrow{\rm
I\\!R}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}$
A set $C$ is functionally ${\cal D}$-convex if $C={\mathop{\cal{CH}}}^{{\cal
D}}(C)$.
###### Definition 3.3
Let ${\cal D}$ be a set of directions. A pair of vectors $v_{1},v_{2}\in{\cal
D}$, is a ${{\cal D}\text{-pair}}$, if $v_{2}$ is counterclockwise from
$v_{1}$, and there is no vector in ${\cal D}$ between $v_{1}$ and $v_{2}$. Let
${{\cal D}}_{\text{pairs}}$ denote the set of all ${{\cal D}\text{-pair}}$s.
Let $\mathop{pspan}(u_{1},u_{2})$ denote the portion of the plane that can be
represented as a positive linear combination of $u_{1},u_{2}\in{\cal D}$. Thus
$\mathop{pspan}(u_{1},u_{2})$ is the open wedge bounded by the rays emanating
from the origin in directions $u_{1},u_{2}$. We define by
$(v_{1},v_{2})_{L}=\mathop{pspan}(-v_{1},v_{2})$ and
$(v_{1},v_{2})_{R}=\mathop{pspan}(v_{1},-v_{2})$: these are two of the four
quadrants of the plane induced by the lines containing $v_{1}$ and $v_{2}$.
Similarly, for $v\in{\cal D}$ we denote by ${{v}_{L}}$ and ${{v}_{R}}$ the two
open half-planes defined by the line passing through $v$. Let
${\cal{Q}}({\cal D})=\left\\{{{{v}_{L}},{{v}_{R}}\,\left|\,{v\in{\cal
D}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}\cup\left\\{{(v_{1},v_{2})_{R},(v_{1},v_{2})_{L}\,\left|\,{(v_{1},v_{2})\in{D}_{\text{pairs}}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}.$
###### Definition 3.4
For a set $S\subseteq{\rm I\\!R}^{2}$ we denote by $T(S)$ the set of
translations of $S$ in the plane, that is
$T(S)=\left\\{{S+p\,\left|\,{p\in{\rm
I\\!R}^{2}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}$. Given a set
of directions ${\cal D}$, let ${\cal T}({\cal D})=\bigcup_{Q\in{\cal{Q}}({\cal
D})}T(Q)$.
For ${\cal D}_{\rm I\\!R}$, the set ${\cal T}({\cal D}_{\rm I\\!R})$ is the
set of all open half-planes. The standard convex hull of a planar point set
$S$ can be defined as follows: start from the whole plane, and remove from it
all the open half-planes $H^{+}$ such that $H^{+}\cap S=\emptyset$. We extend
this definition to handle ${\cal D}$-convexity for an arbitrary set of
directions ${\cal D}$, as follows:
$\mathop{{{\cal D}}\text{-}{\cal CH}}(S)={\rm
I\\!R}^{2}\setminus\left({\bigcup_{I\in{\cal T}({\cal D}),I\cap
S=\emptyset}I}\right);$
that is, we remove from the plane all the translations of quadrants and
halfplanes in ${\cal{Q}}({\cal D})$ that do not contain a point of $S$. See
Figures 4, 5.
(a)(b)(c)
Figure 4: (a) A set of directions ${\cal D}$, (b) the set of quadrants
${\cal{Q}}({\cal D})$ induced by ${\cal D}$, and (c) the $\mathop{{{\cal
D}}\text{-}{\cal CH}}$ of three points.
(a)(b)(c)
Figure 5: (a) A set of directions ${\cal D}$, such that $\alpha({\cal
D})>\pi/2$, (b) the set of quadrants ${\cal{Q}}({\cal D})$ induced by ${\cal
D}$, and (c) the $\mathop{{{\cal D}}\text{-}{\cal CH}}$ of a set of points
which is not connected.
For the case ${\cal D}_{xy}=\left\\{{(0,1),(1,0),(0,-1),(-1,0)}\right\\}$,
Matoušek and Plecháč [MP97] showed that if $\mathop{{{\cal
D}_{xy}}\text{-}{\cal CH}}(S)$ is connected, then ${\mathop{\cal{CH}}}^{{\cal
D}_{xy}}(S)=\mathop{{{\cal D}_{xy}}\text{-}{\cal CH}}(S)$.
###### Definition 3.5
For a set of directions ${\cal D}$, we define the density of ${\cal D}$ to be
$\alpha({\cal D})=\max_{(v_{1},v_{2})\in{{\cal
D}}_{\text{pairs}}}\alpha(v_{1},v_{2}),$
where $\alpha(v_{1},v_{2})$ denotes the counterclockwise angle from $v_{1}$ to
$v_{2}$.
See Figure 5, for an example of a set of directions with density larger than
$\pi/2$.
###### Corollary 3.6
Let ${\cal D}$ be a set of directions in the plane. Then:
* •
The set $\mathop{{{\cal D}}\text{-}{\cal CH}}(A)$ is ${\cal D}$-convex, for
any $A\subseteq{\rm I\\!R}^{2}$.
* •
For any $A\subseteq B\subseteq{\rm I\\!R}^{2}$, one has $\mathop{{{\cal
D}}\text{-}{\cal CH}}(A)\subseteq\mathop{{{\cal D}}\text{-}{\cal CH}}(B)$.
* •
For two sets of directions ${\cal D}_{1}\subseteq{\cal D}_{2}$ we have
$\mathop{{{\cal D}_{1}}\text{-}{\cal CH}}(S)\subseteq\mathop{{{\cal
D}_{2}}\text{-}{\cal CH}}(S)$, for any $S\subseteq{\rm I\\!R}^{2}$.
* •
Let $S$ be a bounded set in the plane, and let ${\cal D}_{1}\subseteq{\cal
D}_{2}\subseteq{\cal D}_{3}\cdots$ be a sequence of sets of directions, such
that $\lim_{i\rightarrow\infty}\alpha(D_{i})=0$. Then,
$\mathop{int}{{\mathop{\mathrm{CH}}}(S)}\subseteq\lim_{i\rightarrow\infty}\mathop{{{\cal
D}_{i}}\text{-}{\cal CH}}(S)\subseteq{\mathop{\mathrm{CH}}}(S)$.
###### Lemma 3.7
Let ${\cal D}$ a set of directions, and let $S$ be a finite set of points in
the plane. Then $C=\mathop{{{\cal D}}\text{-}{\cal CH}}(S)$ is a polygonal set
whose complexity is $O(|S\cap{\partial}{C}|)$.
###### Proof.
It is easy to show that $C$ is polygonal. We charge each vertex of $C$ to some
point of $S^{\prime}=S\cap{\partial}{C}$. Let $C^{\prime}$ be a connected
component of $C$. If $C^{\prime}$ is a single point, then this is a point of
$S^{\prime}$. Otherwise, let $e$ be an edge of $C^{\prime}$, and let $I$ be a
set in ${\cal T}({\cal D})$ such that $e\subseteq{\partial}{I}$, and $I\cap
S=\emptyset$.
Since $e$ is an edge of $C^{\prime}$, there is no $q\in{\rm I\\!R}^{2}$ such
that $e\subseteq q+I$, and $(q+I)\cap S=\emptyset$. This implies that there
must be a point $p$ of $S$ on ${\partial}{I}\cap l_{e}$, where $l_{e}$ is the
line passing through $e$. However, $C$ is a ${\cal D}$-convex set, and the
direction of $e$ belongs to ${\cal D}$. It follows that $l_{e}$ intersects $C$
along a connected set (i.e., the segment $e$), and $p\in l_{e}\cap C=e$. We
charge the edge $e$ to $p$. We claim that a point $p$ of $S^{\prime}$ can be
charged at most 4 times. Indeed, for each edge $e^{\prime}$ of $C$ incident to
$p$, there is a supporting set in ${\cal T}({\cal D})$, such that $p$ and
$e^{\prime}$ lie on its boundary. Only two of those sets can have angle less
than $\pi/2$ at $p$ (because such a set corresponds to a ${{\cal
D}\text{-pair}}(v_{1},v_{2})$ with $\alpha(v_{1},v_{2})>\pi/2$). Thus, a point
of $S^{\prime}$ is charged at most $\max(2\pi/(\pi/2),\pi/(\pi/2)+2)=4$ times.
∎
###### Lemma 3.8
Let ${\cal D}$ be a set of directions, and let $K$ be a bounded convex body in
the plane, such that the expected area of $\mathop{{{\cal D}}\text{-}{\cal
CH}}(N)$ of a set $N$ of $n$ points, chosen uniformly and independently from
$K$, is at least $\left({1-f(n)}\right)Area(K)$, where $1\geq f(n)\geq 0$, for
$n\geq 1$. Then, the expected number of vertices of $C=\mathop{{{\cal
D}}\text{-}{\cal CH}}(N)$ is $O(nf(n/2))$.
###### Proof.
By Lemma 3.7, the complexity of $C$ is proportional to the number of points of
$N$ on the boundary of $C$. Using this observation, it is easy to verify that
the proof of Lemma 2.1 can be extended to this case. ∎
We would like to apply the proof of Theorem 2.3 to bound the expected
complexity of a random ${\cal D}$-convex hull inside a disk. Unfortunately, if
we try to concentrate only on three consecutive sectors (as in Figure 1) it
might be that there is a quadrant $I$ of ${\cal T}({\cal D})$ that intersects
the middle the middle sector from the side (i.e. through the two adjacent
sectors). This, of course, can not happen when working with the regular
convexity. Thus, we first would like to decompose the unit disk into “safe”
regions, where we can apply a similar analysis as the regular case, and the
“unsafe” areas. To do so, we will first show that, with high probability, the
$\mathop{{{\cal D}}\text{-}{\cal CH}}$ of a random point set inside a disk,
contains a “large” disk in its interior. Next, we argue that this implies that
the random $\mathop{{{\cal D}}\text{-}{\cal CH}}$ covers almost the whole
disk, and the desired bound will readily follows from the above Lemma.
###### Definition 3.9
For $r\geq 0$, let $B_{r}$ denote the disk of radius of $r$ centered at the
origin.
###### Lemma 3.10
Let ${\cal D}$ be a set of directions, such that $0\leq\alpha({\cal
D})\leq\pi/2$. Let $N$ be a set of $n$ points chosen uniformly and
independently from the unit disk. Then, with probability $1-n^{-10}$ the set
$\mathop{{{\cal D}}\text{-}{\cal CH}}(N)$ contains $B_{r}$ in its interior,
where $r=1-c\sqrt{\log{n}/n}$, for an appropriate constant $c$.
###### Proof.
Let $r^{\prime}=1-c\sqrt{(\log{n})/n}$, where $c$ is a constant to be
specified shortly. Let $q$ be any point of $B_{r^{\prime}}$. We bound the
probability that $q$ lies outside $C=\mathop{{{\cal D}}\text{-}{\cal CH}}(N)$
as follows: Draw $8$ rays around $q$, such that the angle between any two
consecutive rays is $\pi/4$. This partitions $q+B_{r^{\prime\prime}}$, where
$r^{\prime\prime}=c\sqrt{(\log{n})/n}$, into eight portions
$R_{1},\ldots,R_{8}$, each having area $\pi c^{2}\log{n}/(8n)$. Moreover,
$R_{i}\subseteq q+B_{r^{\prime\prime}}\subseteq B_{1}$, for $i=1,\ldots,8$.
The probability of a point of $N$ to lie outside $R_{i}$ is
$1-c^{2}\log{n}/(8n)$. Thus, the probability that all the points of $N$ lie
outside $R_{i}$ is
$P\left[{N\cap
R_{i}=\emptyset{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq\left({1-\frac{c^{2}\log{n}}{8n}}\right)^{n}\leq
e^{-(c^{2}\log{n})/8}=n^{-c^{2}/8},$
since $1-x\leq e^{-x}$, for $x\geq 0$. Thus, the probability that one of the
$R_{i}$’s does not contain a point of $N$ is bounded by $8n^{-c^{2}/8}$. We
claim that if $R_{i}\cap N\neq\emptyset$, for every $i=1,\ldots,8$, then $q\in
C$. Indeed, if $q\notin C$ then there exists a set $Q\in{\cal{Q}}({\cal D})$,
such that $(q+Q)\cap N=\emptyset$. Since $\alpha({\cal D})\leq\pi/2$ there
exists an $i$, $1\leq i\leq 8$, such that $R_{i}\subseteq q+Q$; see Figure 6.
This is a contradiction, since $R_{i}\cap N\neq\emptyset$. Thus, the
probability that $q$ lies outside $C$ is $\leq 8n^{-c^{2}/8}$.
$R_{1}$$R_{2}$$F_{3}$$R_{4}$$R_{5}$$R_{6}$$R_{7}$$R_{8}$$q$$q+Q$
Figure 6: Since $\alpha({\cal D})\leq\pi/2$, any quadrant $Q\in{\cal{Q}}({\cal
D})$, when translated by $q$, must contain one of the $R_{i}$’s.
Let $N^{\prime}$ denote a set of $n^{10}$ points spread uniformly on the
boundary of $B_{r^{\prime}}$. By the above analysis, all the points of
$N^{\prime}$ lie inside $C$ with probability at least $1-8n^{10-c^{2}/8}$.
Furthermore, arguing as above, we conclude that $B_{r}\subseteq\mathop{{{\cal
D}}\text{-}{\cal CH}}(N^{\prime})$, where $r=1-2c\sqrt{(\log{n})/n}$. Hence,
with probability at least $1-8n^{10-c^{2}/8}$, $\mathop{{{\cal
D}}\text{-}{\cal CH}}(C)$ contains $B_{r}$. The lemma now follows by setting
$c=20$, say. ∎
Since the set of directions may contain large gaps, there are points in
$B_{1}\setminus B_{r}$ that are “unsafe”, in the following sense:
###### Definition 3.11
Let ${\cal D}$ be a set of directions, and let $0\leq r\leq 1$ be a prescribed
constant, such that $0\leq\alpha({\cal D})\leq\pi/2$. A point $p$ in $B_{1}$
is safe, relative to $B_{r}$, if $op\subseteq\mathop{{{\cal D}}\text{-}{\cal
CH}}(B_{r}\cup\left\\{{p}\right\\})$.
See Figure 7 for an example how the unsafe area looks like. The behavior of
the $\mathop{{{\cal D}}\text{-}{\cal CH}}$ inside the unsafe areas is somewhat
unpredictable. Fortunately, those areas are relatively small.
$B_{r}$$B_{1}$$T$$\overrightarrow{v_{1}}$$\overrightarrow{v_{2}}$$o$
Figure 7: The dark areas are the unsafe areas for a consecutive pairs of
directions $v_{1},v_{2}\in{\cal D}$.
###### Lemma 3.12
Let ${\cal D}$ be a set of directions, such that $0\leq\alpha({\cal
D})\leq\pi/2$, and let $r=1-O\left({\sqrt{(\log{n})/n}}\right)$. The unsafe
area in $B_{1}$, relative to $B_{r}$, can be covered by a union of $O(1)$
caps. Furthermore, the length of the base of such a cap is
$O(((\log{n})/n)^{1/4})$, and its height is
$O\left({\sqrt{(\log{n})/n}}\right)$.
###### Proof.
Let $p$ be an unsafe point of $B_{1}$. Let
$\overrightarrow{v_{1}},\overrightarrow{v_{2}}$ be the consecutive pair of
vectors in ${\cal D}$, such that the vector $\overrightarrow{po}$ lies between
them. If $\mathop{ray}(p,\overrightarrow{v_{1}})\cap B_{r}\neq\emptyset$, and
$\mathop{ray}(p,\overrightarrow{v_{2}})\cap B_{r}\neq\emptyset$ then
$po\subseteq{\mathop{\mathrm{CH}}}\left({\left\\{{p,o,p_{1},p_{2}}\right\\}}\right)\subseteq\mathop{{{\cal
D}}\text{-}{\cal CH}}(B_{r}\cup\left\\{{p}\right\\})$, for any pair of points
$p_{1}\in B_{r}\cap\mathop{ray}(p,\overrightarrow{v_{1}}),p_{2}\in
B_{r}\cap\mathop{ray}(p,\overrightarrow{v_{2}})$. Thus, $p$ is unsafe only if
one of those two rays miss $B_{r}$. Since $p$ is close to $B_{r}$, the angle
between the two tangents to $B_{r}$ emanating from $p$ is close to $\pi$. This
implies that the angle between $\overrightarrow{v_{1}}$ and
$\overrightarrow{v_{2}}$ is at least $\pi/4$ (provided $n$ is a at least some
sufficiently large constant), and the number of such pairs is at most $8$.
The area in the plane that sees $o$ in a direction between
$\overrightarrow{v_{1}}$ and $\overrightarrow{v_{2}}$, is a quadrant $Q$ of
the plane. The area in $Q$ which is is safe, is a parallelogram $T$. Thus, the
unsafe area in $B_{1}$ that induced by the pair $\overrightarrow{v_{1}}$ and
$\overrightarrow{v_{2}}$ is $(B_{1}\cap Q)\setminus T$. Since $\alpha({\cal
D})\leq\pi/2$, this set can covered with two caps of $B_{1}$ with their base
lying on the boundary of $B_{r}$. See Figure 7.
The height of such a cap is
$1-r=O\left({\sqrt{\frac{\log{n}}{n(\pi-\alpha)}}}\right)$, and the length of
the base of such a cap is
$2\sqrt{1-r^{2}}=O\left({\left({\frac{\log{n}}{n(\pi-\alpha)}}\right)^{1/4}}\right)$.
∎
The proof of Lemma 3.12 is where our assumption that $\alpha({\cal
D})\leq\pi/2$ plays a critical role. Indeed, if $\alpha({\cal D})>\pi/2$, then
the unsafe areas in $B_{1}\setminus B_{r}$ becomes much larger, as indicated
by the proof.
###### Theorem 3.13
Let ${\cal D}$ be a set of directions, such that $0\leq\alpha({\cal
D})\leq\pi/2$. The expected number of vertices of $\mathop{{{\cal
D}}\text{-}{\cal CH}}(N)$, where $N$ is a set of $n$ points, chosen uniformly
and independently from the unit disk, is $O\left({n^{1/3}+\sqrt{n\alpha({\cal
D})}}\right)$.
###### Proof.
We claim that the expected area $\mathop{{{\cal D}}\text{-}{\cal CH}}(N)$ is
at least $\pi-O\left({n^{-2/3}+\sqrt{\alpha/n}}\right)$, where
$\alpha=\alpha({\cal D})$. The theorem will then follow from Lemma 3.8.
Indeed, let $m$ be an integer to be specified later, and assume, without loss
of generality, that $m$ divides $n$. Partition $B$ into $m$ congruent sectors,
${\cal S}_{1},\ldots,{\cal S}_{m}$. Let $B^{1},\ldots,B^{\mu}$ denote the
$\mu=n/m$ disks centered at the origin, such that (i) $B^{1}=B_{1}$, and (ii)
$Area(B^{i-1})-Area(B^{i})=\pi/\mu$, for $i=2,\ldots,\mu$. Let $r_{i}$ denote
the radius of $B^{i}$, for $i=1,\ldots,\mu$. Note333
$Area(B^{1})-Area(B^{2})=\pi(1-r_{2}^{2})=\pi/\mu$, thus $r_{2}^{2}=1-1/\mu$.
We have $r_{2}\leq 1-1/(2\mu)$, and $r_{1}-r_{2}\geq
1-(1-1/(2\mu))=1/(2\mu)$., that $r_{i}-r_{i+1}\geq r_{i-1}-r_{i}\geq
1/(2\mu)$, for $i=2,\ldots,\mu-1$.
Let $r=1-O\left({\sqrt{(\log{n})/n}}\right)$, and let $U$ be the set of
sectors that either intersect an unsafe area of $B$ relative to $B_{r}$, or
their neighboring sectors intersect the unsafe area of $B$. By Lemma 3.12, the
number of sectors in $U$ is $O(1)\cdot
O\left({\frac{((\log{n})/n)^{1/4}}{(2\pi/m)}}\right)=O(m((\log{n})/n)^{1/4})$.
Let $S_{i,j}=(B^{i}\setminus B^{i+1})\cap{\cal S}_{j}$, and
$S_{\mu,j}=B^{\mu}\cap{\cal S}_{j}$, for $i=1,\ldots,\mu-1$, and
$j=1,\ldots,m$. The set $S_{i,j}$ is called the $i$-th tile of the sector
${\cal S}_{j}$, and its area is $\pi/n$, for $i=1,\ldots,\mu$, and
$j=1,\ldots,m$.
Let $X_{j}$ denote the first index $i$ such that $N\cap S_{i,j}\neq\emptyset$,
for $j=1,\ldots,m$. The probability that $X_{j}=k$ is upper-bounded by the
probability that the tiles $S_{1,j},\ldots,S_{(k-1),j}$ do not contain any
point of $N$; namely, by $\left({1-\frac{k-1}{n}}\right)^{n}$. Thus,
$P[X_{j}=k]\leq\left({1-\frac{k-1}{n}}\right)^{n}\leq e^{-(k-1)}$. Thus,
$E\left[{X_{j}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=\sum_{k=1}^{\mu}kP[X_{j}=k]\leq\sum_{k=1}^{\mu}ke^{-(k-1)}=O(1),$
for $j=1,\ldots,m$.
Let $C$ denote the set $\mathop{{{\cal D}}\text{-}{\cal CH}}(N\cup B_{r})$.
The tile $S_{i,j}$ is exposed by a set $K$, if $S_{i,j}\setminus
K\neq\emptyset$.
We claim that the expected number of tiles exposed by $C$ in a section
$S_{j}\notin U$ is at most $X_{j-1}+X_{j+1}+O(\mu/m^{2}+\alpha\mu/m)$, for
$j=1,\ldots,m$ (where we put $X_{0}=X_{m}$, $X_{m+1}=X_{1}$).
Indeed, let $w=max(X_{j-1},X_{j+1})$, and let $p,q$ be the two points in
$S_{j-1,w},S_{j+1,w}$, respectively, such that the number of sets exposed by
the triangle $T=\triangle{opq}$, in the sector ${\cal S}_{j}$, is maximal.
Both $p$ and $q$ lie on ${\partial}{B^{w+1}}$ and on the external radii
bounding ${\cal S}_{j-1}$ and ${\cal S}_{j+1}$, as shown in Figure 1. Let $s$
denote the segment connecting the midpoint $\rho$ of the base of $T$ to its
closest point on ${\partial}{B^{w}}$. The number of tiles in ${\cal S}_{j}$
exposed by $T$ is bounded by $w$, plus the number of tiles intersecting the
segment $s$. The length of $s$ is
$|oq|-|oq|\cos\left({\frac{3}{2}\cdot\frac{2\pi}{m}}\right)\leq
1-\cos\left({\frac{3\pi}{m}}\right)\leq\frac{1}{2}\left({\frac{3\pi}{m}}\right)^{2}=\frac{4.5\pi^{2}}{m^{2}},$
since $\cos{x}\geq 1-x^{2}/2$, for $x\geq 0$.
On the other hand, the segment $s$ intersects at most
$\left\lceil{||s||/(1/(2\mu))}\right\rceil=O(\mu/m^{2})$ tiles, and we have
that the number of tiles exposed in the sector ${\cal S}_{i}$ by $T$ is at
most $w+O(\mu/m^{2})$, for $j=1,\ldots,m$.
Since ${\cal S}_{j}\notin U$, the points $p,q$ are safe, and $op,oq\subseteq
C$. This implies that the only additional tiles that might be exposed in
${\cal S}_{j}$ by $C$, are exposed by the portion of the boundary of $C$
between $p$ and $q$ that lie inside $T$. Let $V$ be the circular cap
consisting of the points in $T$ lying between $pq$ and a circular arc
$\gamma\subseteq T$, connecting $p$ to $q$, such that for any point
$p^{\prime}\in\gamma$ one has $\angle{pp^{\prime}q}=\pi-\alpha$. See Figure 8.
$T$$V$$p$$q$$\geq\pi-\alpha$
Figure 8: The portion of $T$ that can be removed by a quadrant $Q$ of ${\cal
T}({\cal D})$, is covered by the darkly-shaded circular cap, such that any
point on its bounding arc creates an angle $\pi-\alpha$ with $p$ and $q$.
Let $Q\in{\cal T}({\cal D})$ be any quadrant of the plane induce by ${\cal
D}$, such that $Q\cap N=\emptyset$ (i.e. $C\cap Q=\emptyset$), and $Q\cap
T\neq\emptyset$. Then, $Q\cap op=\emptyset,Q\cap oq=\emptyset$ since $p$ and
$q$ are safe. Moreover, the angle of $Q$ is at least $\pi-\alpha$, which
implies that $Q\cap T\subseteq V$. See Figure 8.
Let $s^{\prime}$ be the segment $o\rho\cap V$, where $\rho$ is as above, the
midpoint of $pq$. The length of $s^{\prime\prime}$ is
$|s^{\prime}|\leq\sin\left({\frac{3}{2}\cdot\frac{2\pi}{m}}\right)\tan{\frac{\alpha}{2}}\leq\frac{3\pi}{m}\frac{\sqrt{2}\alpha}{2}\leq\frac{3\pi\alpha}{m},$
since $\sin{x}\leq x$, for $x\geq 0$, and $1/\sqrt{2}\leq\cos{(\alpha/2)}$
(because $0\leq\alpha\leq\pi/2$).
Thus, the expected number of tiles exposed by $C$, in a sector ${\cal
S}_{j}\notin U$, is bounded by
$X_{j-1}+X_{j+1}+O\left({\frac{\mu}{m^{2}}}\right)+O\left({\frac{3\pi\alpha/m}{1/(2\mu)}}\right)=X_{j-1}+X_{j+1}+O\left({\frac{\mu}{m^{2}}}\right)+O\left({\frac{\alpha\mu}{m}}\right).$
Thus, the expected number of tiles exposed by $C$, in sectors that do not
belong to $U$, is at most
$E\left[{\sum_{j=1}^{m}\left({X_{j-1}+X_{j+1}+O\left({\frac{\mu}{m^{2}}}\right)+O\left({\frac{\alpha\mu}{m}}\right)}\right){\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=O\left({m+\frac{\mu}{m}+\alpha\mu}\right).$
Adding all the tiles that lie outside $B_{r}$ in the sectors that belong to
$U$, it follows that the expected number of tiles exposed by $C$ is at most
$\displaystyle O$
$\displaystyle\left({m+\frac{\mu}{m}+\alpha\mu+|U|\cdot\frac{1-r}{1/2\mu}}\right)=O\left({m+\frac{\mu}{m}+\alpha\mu+m\left({\frac{\log{n}}{n}}\right)^{1/4}\cdot\mu\sqrt{\left({\frac{\log{n}}{n}}\right)}}\right)$
$\displaystyle=$ $\displaystyle O\left({m+\frac{n}{m^{2}}+\frac{\alpha
n}{m}+n\left({\frac{\log{n}}{n}}\right)^{3/4}}\right)=O\left({m+\frac{n}{m^{2}}+\frac{\alpha
n}{m}+n^{1/4}\log^{3/4}{n}}\right).$
Setting $m=\max{\left({n^{1/3},\sqrt{n\alpha}}\right)}$, we conclude that the
expected number of tiles exposed by $C$ is
$O\left({n^{1/3}+\sqrt{n\alpha}}\right)$.
The area of $C^{\prime}=\mathop{{{\cal D}}\text{-}{\cal CH}}(N)$ is bounded
from below by the area of the tiles which are not exposed by $C^{\prime}$. The
probability that $C^{\prime}\neq C$ (namely, that the disk $B_{r}$ is not
inside $C^{\prime}$) is at most $n^{-10}$, by Lemma 3.10. Hence the expected
area of $C^{\prime}$ is at least
$E[Area(C)]-Prob\left[{C\neq
C^{\prime}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\pi=\pi-O\left({n^{1/3}+\sqrt{n\alpha}}\right)\frac{\pi}{n}-n^{-10}\pi=\pi-O\left({n^{-2/3}+\sqrt{\frac{\alpha}{n}}\;}\right).$
The assertion of the theorem now follows from Lemma 3.8. ∎
The expected complexity of the $\mathop{{{\cal D}_{xy}}\text{-}{\cal CH}}$ of
$n$ points, chosen uniformly and independently from the unit square, is
$O(\log{n})$ (Lemma 2.4). Unfortunately, this is a degenerate case for a set
of directions with $\alpha({\cal D})=\pi/2$, as the following corollary
testifies:
###### Corollary 3.14
Let ${\cal D}_{xy}^{\prime}$ be the set of directions resulting from rotating
${\cal D}_{xy}$ by 45 degrees. Let $N$ be a set of $n$ points, chosen
independently and uniformly from the unit square ${S^{\prime}}$. The expected
complexity of $\mathop{{{\cal D}_{xy}^{\prime}}\text{-}{\cal CH}}(N)$ is
$\Omega\left({\sqrt{n}}\right)$.
###### Proof.
Without loss of generality, assume that $n=m^{2}$ for some integer $m$. Tile
${S^{\prime}}$ with $n$ translated copies of a square of area $1/n$. Let
${\cal S}_{1},\ldots,{\cal S}_{m}$ denote the squares in the top raw of this
tiling, from left to right. Let $A_{j}$ denote the event that ${\cal S}_{j}$
contains a point of $N$, and neither of the two adjacent squares
$S_{j-1},S_{j+1}$ contains a point of $N$, for $j=2,\ldots,m-1$.
We have
$Prob\left[{A_{j}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=Prob\left[{{\cal
S}_{j+1}\cap N=\emptyset\text{ and }{\cal S}_{j-1}\cap
N=\emptyset{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]-Prob\left[{({\cal
S}_{j-1}\cup{\cal S}_{j}\cup{\cal S}_{j+1})\cap
N=\emptyset{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right],$
for $j=2,\ldots,m-1$. Hence,
$\lim_{n\rightarrow\infty}Prob\left[{A_{j}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=\lim_{n\rightarrow\infty}\left({\left({1-\frac{2}{n}}\right)^{n}-\left({1-\frac{3}{n}}\right)^{n}}\right)=e^{-2}-e^{-3}\approx
0.0855$
$q$$q+Q_{top}$${\cal S}_{j-1}$${\cal S}_{j}$${\cal S}_{j+1}$
Figure 9: If $A_{j}$ happens, then the squares ${\cal S}_{j-1},{\cal S}_{j+1}$
do not contain a point of $N$. Thus, if $q$ is the highest point in ${\cal
S}_{j}$, then $q+Q_{top}$ can not contain a point of $N$, and $q$ is a vertex
of $\mathop{{{\cal D}_{xy}^{\prime}}\text{-}{\cal CH}}(N)$.
This implies, that for $n$ large enough,
$Prob\left[{A_{j}{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\geq 0.01$. Thus,
the expected value of $Y$ is $\Omega(m)=\Omega\left({\sqrt{n}}\right)$, where
$Y$ is the number of $A_{j}$’s that have occurred, for $j=2,\ldots,m-1$.
However, if $A_{j}$ occurs, then $C=\mathop{{{\cal
D}_{xy}^{\prime}}\text{-}{\cal CH}}(N)$ must have a vertex at ${\cal S}_{j}$.
Indeed, let $Q_{top}$ denote the quadrant of ${\cal{Q}}({\cal
D}_{xy}^{\prime})$ that contains the positive $y$-axis. If we translate
$Q_{top}$ to the highest point in $S_{j}\cap N$, then it does not contain a
point of $N$ in its interior, implying that this point is a vertex of $C$, see
Figure 9.
This implies that the expected complexity of $\mathop{{{\cal
D}_{xy}^{\prime}}\text{-}{\cal CH}}(N)$ is $\Omega\left({\sqrt{n}}\right)$ ∎
## 4 On the Expected Number of Points on the Boundary of the Quadrant Hull
Inside a Hypercube
In this section, we show that the expected number of points on the boundary of
the quadrant hull of a set $S$ of $n$ points, chosen uniformly and
independently from the unit cube is $O(\log^{d-1}n)$. Those bounds are known
[BKST78], but we believe the new proof is simpler.
###### Definition 4.1 ([MP97])
Let ${\cal Q}$ be a family of subsets of ${\rm I\\!R}^{d}$. For a set
$A\subseteq{\rm I\\!R}^{d}$, we define the ${\cal Q}$-hull of $A$ as
$\mathrm{\mathop{{{\cal Q}}\text{-}{co}}}(A)=\bigcap\left\\{{Q\in{\cal
Q}\,\left|\,{A\subseteq
Q}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}.$
###### Definition 4.2 ([MP97])
For a sign vector $s\in\left\\{{-1,+1}\right\\}^{d}$, define
$q_{s}=\left\\{{x\in{\rm
I\\!R}^{d}\,\left|\,{\mathop{\mathrm{sign}}(x_{i})=s_{i},\text{ for
}i=1,\ldots,d}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\},$
and for $a\in{\rm I\\!R}^{d}$, let $q_{s}(a)=q_{s}+a$. We set ${\cal
Q}_{sc}=\left\\{{{\rm I\\!R}^{d}\setminus q_{s}(a)\,\left|\,{a\in{\rm
I\\!R}^{d},s\in\left\\{{-1,+1}\right\\}^{d}}\rule[-5.69046pt]{0.0pt}{11.38092pt}\right.}\right\\}$.
We shall refer to $\mathrm{\mathop{{{\cal Q}_{sc}}\text{-}{co}}}(A)$ as the
quadrant hull of $A$. These are all points which cannot be separated from $A$
by any open orthant in space (i.e., quadrant in the plane).
###### Definition 4.3
Given a set of points $S\subseteq{\rm I\\!R}^{d}$, a point $p\in{\rm
I\\!R}^{d}$ is ${\cal Q}_{sc}$-exposed, if there is
$s\in\left\\{{-1,+1}\right\\}^{d}$, such that $q_{s}(p)\cap S=\emptyset$. A
set $C$ is ${\cal Q}_{sc}$-exposed, if there exists a point $p\in C$ which is
${\cal Q}_{sc}$-exposed.
###### Definition 4.4
For a set $S\subseteq{\rm I\\!R}^{d}$, let $n_{sc}(S)$ denote the number of
points of $S$ on the boundary of $\mathrm{\mathop{{{\cal
Q}_{sc}}\text{-}{co}}}(S)$.
###### Theorem 4.5
Let ${\cal C}$ be a unit axis parallel hypercube in ${\rm I\\!R}^{d}$, and let
$S$ be a set of $n$ points chosen uniformly and independently from ${\cal C}$.
Then, the expected number of points of $S$ on the boundary of
$H=\mathrm{\mathop{{{\cal Q}_{sc}}\text{-}{co}}}(S)$ is $O(\log^{d-1}(n))$.
###### Proof.
We partition ${\cal C}$ into equal size tiles, of volume $1/n^{d}$; that is
$C(i_{1},i_{2},\ldots,i_{d})=[(i_{1}-1)/n,i_{1}/n]\times\cdots\times[(i_{d}-1)/n,i_{d}/n]$,
for $1\leq i_{1},i_{2},\ldots,i_{d}\leq n$.
We claim that the expect number of tiles in our partition of ${\cal C}$ which
are exposed by $S$ is $O(n^{d-1}\log^{d-1}n)$.
Indeed, let $q=q_{(-1,-1,\ldots,-1)}$ be the “negative” quadrant of ${\rm
I\\!R}^{d}$. Let $X(i_{2},\ldots,i_{d})$ be the maximal integer $k$, for which
$C(k,i_{2},\ldots,i_{d})$ is exposed by $q$. The probability that
$X(i_{2},\ldots,i_{d})\geq k$ is bounded by the probability that the cubes
$C(l_{1},l_{2},\ldots,l_{d})$ does not contain a point of $S$, where
$l_{1}<k,l_{2}<i_{2},\ldots,l_{d}<i_{d}$. Thus,
$\displaystyle\Pr\left[{X(i_{2},\ldots,i_{d})\geq
k{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]$ $\displaystyle\leq$
$\displaystyle\left({1-\frac{(k-1)(i_{2}-1)\cdots(i_{d}-1)}{n^{d}}}\right)^{n}$
$\displaystyle\leq$
$\displaystyle\exp\left({-\frac{{(k-1)(i_{2}-1)\cdots(i_{d}-1)}}{n^{d-1}}}\right),$
since $1-x\leq e^{-x}$, for $x\geq 0$.
Hence, the probability that $\Pr\left[{X(i_{2},\ldots,i_{d})\geq i\cdot
m+1{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq e^{-i}$, where
$m=\left\lceil{\frac{n^{d-1}}{(i_{2}-1)\cdots(i_{d}-1)}}\right\rceil$. Thus,
$\displaystyle
E\left[{X(i_{2},\ldots,i_{d}){\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]$
$\displaystyle=$
$\displaystyle\sum_{i=1}^{\infty}i\Pr\left[{X(i_{2},\ldots,i_{d})=i{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]=\sum_{i=0}^{\infty}\sum_{j=im+1}^{(i+1)m}j\Pr\left[{X(i_{2},\ldots,i_{d})=j{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]$
$\displaystyle\leq$
$\displaystyle\sum_{i=0}^{\infty}(i+1)m\Pr\left[{X(i_{2},\ldots,i_{d})\geq
im+1{\rule[-5.69046pt]{0.0pt}{0.0pt}}}\right]\leq\sum_{i=0}^{\infty}(i+1)me^{-i}=O(m).$
Let $r$ denote the expected number of tiles exposed by $q$ in ${\cal C}$. If
$C(i_{1},\ldots,i_{d})$ is exposed by $q$, then $X(i_{2},\ldots,i_{d})\geq
i_{1}$. Thus, one can bound $r$ by the number of tiles on the boundary of
${\cal C}$, plus the sum of the expectations of the variables
$X(i_{2},\ldots,i_{d})$. We have
$\displaystyle r$ $\displaystyle=$ $\displaystyle
O(n^{d-1})+\sum_{i_{2}=2}^{n-1}\sum_{i_{3}=2}^{n-1}\cdots\sum_{i_{d}=2}^{n-1}O\left({\frac{n^{d-1}}{(i_{2}-1)(i_{3}-1)\cdots(i_{d}-1)}}\right)$
$\displaystyle=$ $\displaystyle
O\left({n^{d-1}}\right)\sum_{i_{2}=2}^{n-1}\frac{1}{i_{2}-1}\sum_{i_{3}=2}^{n-1}\frac{1}{i_{3}-1}\cdots\sum_{i_{d}=2}^{n-1}\frac{1}{i_{d}-1}=O\left({n^{d-1}\log^{d-1}{n}}\right).$
The set ${\cal Q}_{sc}$ contains translation of $2^{d}$ different quadrants.
This implies, by symmetry, that the expected number of tiles exposed in ${\cal
C}$ by $S$ is
$O\left({2^{d}n^{d-1}\log^{d-1}{n}}\right)=O\left({n^{d-1}\log^{d-1}{n}}\right)$.
However, if a tile is not exposed by any $q_{s}$, for
$s\in\left\\{{-1,+1}\right\\}^{d}$, then it lies in the interior of $H$.
Implying that the expected volume of $H$ is at least
$\frac{n^{d}-O\left({n^{d-1}\log^{d-1}{n}}\right)}{n^{d}}=1-O\left({\frac{\log^{d-1}n}{n}}\right).$
We now apply an argument similar to the one used in Lemma 2.1 (Efron’s
Theorem), and the theorem follows. ∎
###### Remark 4.6
A point $p$ of $S$ is a maxima, if there is no point $p^{\prime}$ in $S$, such
that $p_{i}\leq p^{\prime}_{i}$, for $i=1,\ldots,d$. Clearly, a point which is
a maxima, is also on the boundary of $\mathrm{\mathop{{{\cal
Q}_{sc}}\text{-}{co}}}(S)$. By Theorem 4.5, the expected number of maxima in a
set of $n$ points chosen independently and uniformly from the unit hypercube
in ${\rm I\\!R}^{d}$ is $O(\log^{d-1}n)$. This was also proved in [BKST78],
but we believe that our new proof is simpler.
Also, as noted in [BKST78], a vertex of the convex hull of $S$ is a point of
$S$ lying on the boundary of the $\mathrm{\mathop{{{\cal
Q}_{sc}}\text{-}{co}}}(S)$. Hence, the expected number of vertices of the
convex hull of a set of $n$ points chosen uniformly and independently from a
hypercube in ${\rm I\\!R}^{d}$ is $O(\log^{d-1}n)$.
### Acknowledgments
I wish to thank my thesis advisor, Micha Sharir, for his help in preparing
this manuscript. I also wish to thank Pankaj Agarwal, and Imre Bárány for
helpful discussions concerning this and related problems.
## References
* [BKST78] J. L. Bentley, H. T. Kung, M. Schkolnick, and C. D. Thompson. On the average number of maxima in a set of vectors and applications. Journal of the ACM, 25:536–543, 1978.
* [Efr65] B. Efron. The convex hull of a random set of points. Biometrika, 52(3):331–343, 1965.
* [MP97] J. Matoušek and P. Plecháč. On functional separately convex hulls. To appear in Discrete Comput. Geom., 1997.
* [OSSW85] T. Ottmann, E. Soisalon-Soininen, and D. Wood. On the definition and computation of rectilinear convex hulls. Inform. Sci., 33:157–171, 1985.
* [PS85] F. P. Preparata and M. I. Shamos. Computational Geometry: An Introduction. Springer-Verlag, 1985.
* [Ray70] H. Raynaud. Sur l’enveloppe convex des nuages de points aleatoires dans $R^{n}$. J. Appl. Probab., 7:35–48, 1970.
* [RS63] A. Rényi and R. Sulanke. Über die konvexe Hülle von $n$ zufällig gerwähten Punkten I. Z. Wahrsch. Verw. Gebiete, 2:75–84, 1963.
* [RW87] G. J. E. Rawlins and D. Wood. Optimal computation of finitely oriented convex hulls. Inform. Comput., 72:150–166, 1987.
* [RW88] G. J. E. Rawlins and D. Wood. Computational geometry with restricted orientations. In Proc. 13th IFIP Conf. System Modelling and Optimization, volume 113 of Lecture Notes in Control and Information Science, pages 375–384. Springer-Verlag, 1988.
* [WW93] W. Weil and J. A. Wieacker. Stochastic geometry. In P. M. Gruber and J. M. Wills, editors, Handbook of Convex Geometry, volume B, chapter 5.2, pages 1393–1438. North-Holland, 1993.
|
arxiv-papers
| 2011-11-22T21:17:34 |
2024-09-04T02:49:24.601234
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Sariel Har-Peled",
"submitter": "Sariel Har-Peled",
"url": "https://arxiv.org/abs/1111.5340"
}
|
1111.5385
|
# The elliptic flow in Au+Au collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ =
7.7, 11.5 and 39 GeV at STAR††thanks: Presented at the conference ‘Strangeness
in Quark Matter 2011’, Cracow, Poland, September 18-24, 2011
Shusu Shi (for the STAR collaboration) Institute of Particle Physics, Central
China Normal University, Wuhan, Hubei, 430079, China The Key Laboratory of
Quark and Lepton Physics (Central China Normal University), Ministry of
Education, Wuhan, Hubei, 430079, China
###### Abstract
We present elliptic flow, $v_{2}$, measurements for charged and identified
particles at midrapidity in Au+Au collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$
= 7.7, 11.5 and 39 GeV at STAR. We compare the inclusive charged hadron
$v_{2}$ to those from high energies at RHIC ($\sqrt{\mathrm{s}_{{}_{NN}}}$ =
62.4 and 200 GeV), at LHC ($\sqrt{\mathrm{s}_{{}_{NN}}}$ = 2.76 TeV). A
significant difference in $v_{2}$ between baryons and anti-baryons is observed
and the difference increases with decreasing beam energy. We observed the
$v_{2}$ of $\phi$ meson is systematically lower than other particles in Au+Au
collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 11.5 GeV.
25.75.Ld, 25.75.Dw
## 1 Introduction
The main physics motivation for the Beam Energy Scan at RHIC-STAR experiment
is searching for the phase boundary and critical point predicted by QCD
theory. The elliptic flow, $v_{2}$, which is generated by the initial
anisotropy in the coordinate space, is defined by
$v_{2}=\langle\cos 2(\phi-\Psi_{R})\rangle$ (1)
, where $\phi$ is azimuthal angle of an outgoing particle, $\Psi_{R}$ is the
azimuthal angle of the impact parameter, and angular brackets denote an
average over many particles and events. Due to the self-quenching effect, it
is sensitive to the early stage of the heavy ion collisions [1]. The number of
constituent quark (NCQ) scaling observed in the top energy of RHIC Au+Au and
Cu+Cu collisions [2, 3, 5] reflects that the collectivity has been built up at
the partonic stage. Especially, the NCQ scaling for multi-strange hadrons,
$\phi$ ($s\overline{s}$) and $\Omega$ ($sss$) supports the deconfinement and
partonic collectivity picture [4, 6]. A study based on a multi-phase transport
model (AMPT) indicates the NCQ scaling is related to the degrees of freedom in
the system [7]. The scaling and no scaling in $v_{2}$ reflects the partonic
and hadronic degrees of freedom respectively. The importance of $\phi$ meson
has been emphasized, where the $\phi$ meson $v_{2}$ could be small or zero
without partonic phase [8]. Thus, the measurements of elliptic flow with the
Beam Energy Scan data offer us the opportunity to investigate the QCD phase
boundary.
Figure 1: (Color online) The top panels show the $v_{2}\\{4\\}$ vs. $p_{T}$ at
midrapidity for various beam energies ($\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7 GeV
to 2.76 TeV). The results for $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7 to 200 GeV
are for $\mathrm{Au+Au}$ collisions and those for 2.76 TeV are for Pb + Pb
collisions. The red line shows the fit to the results from $\mathrm{Au+Au}$
collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 200 GeV. The bottom panels show
the ratio of $v_{2}\\{4\\}$ vs. $p_{T}$ for all $\sqrt{\mathrm{s}_{{}_{NN}}}$
with respect to this fitted line. The results are shown for three collision
centrality classes: $10-20\%$ (a1), $20-30\%$ (b1) and $30-40\%$ (c1).
In this proceedings, we present the $v_{2}$ results of charged and identified
hadrons by the STAR experiment from Au+Au collisions in
$\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7, 11.5 and 39 GeV. STAR’s Time Projection
Chamber (TPC) [9] is used as the main detector for event plane determination.
The centrality was determined by the number of tracks from the pseudorapidity
region $|\eta|\leq 0.5$. The particle identification for $\pi^{\pm}$,
$K^{\pm}$ and $p~{}(\overline{p})$ is achieved via the energy loss in the TPC
and the time of flight information from the multi-gap resistive plate chamber
detector [10]. Strange hadrons are reconstructed with the decay channels:
${K}^{0}_{S}$ $\rightarrow\pi^{+}+\pi^{-}$, $\phi\rightarrow K^{+}+K^{-}$,
$\Lambda$ $\rightarrow p+\pi^{-}$
($\overline{\Lambda}\rightarrow\overline{p}+\pi^{+}$), and
$\Xi^{-}\rightarrow$ $\Lambda$ $+\ \pi^{-}$ ($\overline{\Xi}^{+}\rightarrow$
$\overline{\Lambda}$\+ $\pi^{+}$)). The detailed description of the procedure
can be found in Refs. [2, 3, 11]. The event plane method [12] and cumulant
method [13, 14] are used for the $v_{2}$ measurement.
## 2 Results and Discussions
Figure 2: (Color online) The difference of $v_{2}$ for particles and anti-
particles ($v_{2}(X)-v_{2}(\overline{X})$) divided by particle $v_{2}$
($v_{2}(X)$) as a function of beam energy in Au+Au collisions (0-80%). Figure
3: (Color online) The number of constituent quark ($n_{\rm cq}$) scaled
$v_{2}$ as a function of transverse kinetic energy over $n_{\rm cq}$
($(m_{T}-m_{0})/n_{\rm cq}$) for various identified particles in Au+Au (0-80%)
collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 11.5 and 39 GeV.
The Beam Energy Scan data from RHIC-STAR experiment, offer a opportunity to
study the beam energy dependence of $v_{2}$ in a wide range of beam energy.
Figure 1 shows the results of transverse momentum ($p_{T}$) dependence of
$v_{2}\\{4\\}$ for charged hadrons from $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7
GeV to 2.76 TeV in $10-20\%$ (a1), $20-30\%$ (b1) and $30-40\%$ (c1)
centrality bins, where the ALICE results in Pb + Pb collisions at
$\sqrt{s_{NN}}$ = 2.76 TeV are taken from [15]. At low $p_{T}$
($p_{T}<2~{}\mbox{$\mathrm{GeV}/c$}$), the $v_{2}$ values increases with
increase in beam energy. Beyond $p_{T}=2~{}\mbox{$\mathrm{GeV}/c$}$ the
$v_{2}$ results show comparable values with in statistic errors. There is no
saturation signal of $v_{2}$ up to collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$
= 2.76 TeV. Figure 2 shows excitation function for the relative difference of
$v_{2}$ between particles and anti-particles. Here, to reduce the non-flow
effect, the $\eta$-sub event plane method is used to calculate $v_{2}$. The
$\eta$-sub event plane method is similar to the event plane method, except one
defines the flow vector for each particle based on particles measured in the
opposite hemisphere in pseudorapidity. An $\eta$ gap of $|\eta|<0.05$ is used
between negative/positive $\eta$ sub-event to guarantee that non-flow effects
are reduced by enlarging the separation between the correlated particles. The
difference for baryon and anti-baryon (protons and $\Lambda$s) could be
observed from 7.7 to 62.4 GeV. The difference of $v_{2}$ for baryons is within
10% at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 and 62.4 GeV, while a significant
difference is observed below $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 GeV. For
example, the difference of protons versus anti-protons is around 60%. There is
no obvious difference for $\pi^{+}$ versus $\pi^{-}$ (within 3%) and $K^{+}$
versus $K^{-}$ (within 2%) when $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 GeV. As the
decrease of beam energy, $\pi^{+}$ versus $\pi^{-}$ and $K^{+}$ versus $K^{-}$
start to show the difference. The $v_{2}$ of $\pi^{-}$ is larger than that of
$\pi^{+}$ and the $v_{2}$ of $K^{+}$ is larger than that of $K^{-}$. This
difference between particles and anti-particles might be due to the baryon
transport effect to midrapidity [16] or absorption effect in the hadronic
stage. The results could indicate the hadronic interaction become more
dominant in lower beam energy. The immediate consequence of the significant
difference between baryon and anti-baryon $v_{2}$ is that the NCQ scaling is
broken between particles and anti-particles at $\sqrt{\mathrm{s}_{{}_{NN}}}$ =
7.7 and 11.5 GeV. The transverse momentum dependence of $v_{2}$ for the
selected identified particles is shown in Fig. 3. The $v_{2}$ and
$m_{T}-m_{0}$ has been divided by number of constituent quark in each hadron.
In Au+Au collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 GeV, the similar
scaling behavior at $\sqrt{\mathrm{s}_{{}_{NN}}}$ =200 GeV is observed.
Especially, the $\phi$ mesons which is insensitive to the later hadronic
rescatterings follows the same trend of other particles. It suggests that the
partonic collectivity has been built up in collisions at
$\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 GeV. However, the $v_{2}$ for $\phi$ mesons
falls off from other particles at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 11.5 GeV.
The mean deviation to the $v_{2}$ pions is 2.6 $\sigma$. It indicates that the
hadronic interaction are dominant in collisions at
$\sqrt{\mathrm{s}_{{}_{NN}}}$ = 11.5 GeV.
## 3 Summary
In summary, we present the $v_{2}$ measurement for charged hadrons and
identified hadrons in Au+Au collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7,
11.5 and 39 GeV. The magitude of $v_{2}$ increases as increasing of the beam
energy from 7.7 GeV to 2.76 TeV. The difference between the $v_{2}$ of
particles and anti-particles is observed. The baryon and anti-baryon show
significant difference at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 7.7 and 11.5 GeV.
The ongoing analysis with 19.6 and 27 GeV data collected in 2011 will fill the
gap between 11.5 and 39 GeV. The pions and kaons are almost consistent at
$\sqrt{\mathrm{s}_{{}_{NN}}}$ = 39 GeV. This difference increase with
decreasing of the beam energy. The $v_{2}$ of $\phi$ meson falls off from
other particles in collisions at $\sqrt{\mathrm{s}_{{}_{NN}}}$ = 11.5 GeV.
Experimental data suggests the hadronic interactions are dominant when
$\sqrt{\mathrm{s}_{{}_{NN}}}$ $\leq$ 11.5 GeV.
## 4 Acknowledgments
This work was supported in part by the National Natural Science Foundation of
China under grant No. 11105060, 10775060 and 11135011.
## References
* [1] S. A. Voloshin, A. M. Poskanzer and R. Snellings, arXiv:0809.2949.
* [2] J. Adams et al. (STAR Collaboration), Phys. Rev. Lett. 92, 052302 (2004).
* [3] J. Adams et al. (STAR Collaboration), Phys. Rev. Lett. 95, 122301 (2005).
* [4] B. I. Abelev et al., (STAR Collaboration), Phys. Rev. Lett. 99 112301 (2007).
* [5] B. I. Abelev et al., (STAR Collaboration), Phys. Rev. C 81, 044902 (2010).
* [6] S. Shi (for the STAR collaboration), Nucl. Phys. A 830, 187c (2009).
* [7] F. Liu, K.J. Wu, and N. Xu, J. Phys. G 37 094029(2010).
* [8] B. Mohanty and N. Xu, J. Phys. G 36, 064022(2009).
* [9] K. H. Ackermann et al. (STAR Collaboration), Nucl. Instrum. Methods A 499, 624 (2003).
* [10] W. J. Llope (STAR TOF Group), Nucl. Instr. and Meth. B 241, 306 (2005).
* [11] C. Adler et al. (STAR Collaboration), Phys. Rev. Lett. 89, 132301 (2002).
* [12] A. M. Poskanzer and S. A. Voloshin, Phys. Rev. C 58 1671 (1998).
* [13] N. Borghini, P. M. Dinh, and J.-Y. Ollitrault, Phys. Rev. C 63, 054906 (2001).
* [14] N. Borghini, P. M. Dinh, and J.-Y. Ollitrault, Phys. Rev. C 64, 054901 (2001).
* [15] K. Aamodt et al. (ALICE Collaboration), Phys. Rev. Lett. 105, 252302 (2010).
* [16] J. Dunlop, M.A. Lisa and P. Sorensen, arXiv:1107.3078.
|
arxiv-papers
| 2011-11-23T02:07:41 |
2024-09-04T02:49:24.612319
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Shusu Shi (for the STAR collaboration)",
"submitter": "Shusu Shi",
"url": "https://arxiv.org/abs/1111.5385"
}
|
1111.5410
|
# Forward Jets and Forward-Central Jets at CMS
Niladri Sen, on behalf of the CMS Collaboration
###### Abstract
We report on cross section measurements for inclusive forward jet production
and for the simultaneous production of a forward and a central jet in
$\sqrt{s}=$7 TeV $pp$-collisions. Data collected in 2010, corresponding to an
integrated luminosity of 3.14 pb-1, is used for the measurements. Jets in the
transverse momentum range pT = 35 - 140 GeV/c are reconstructed with the anti-
kT (R = 0.5) algorithm. The extended coverage of large pseudo-rapidities is
provided by the Hadronic Forward calorimeter (3.2 $<|\eta|<$ 4.7), while
central jets are limited to $|\eta|<$ 2.8, covered by the main detector
components. The two differential cross sections are presented as a function of
the jet transverse momentum. Comparisons to next-to-leading order perturbative
calculations, and predictions from event generators based on different parton
showering mechanisms (pythia and herwig) and parton dynamics (cascade) are
shown.
###### Keywords:
CMS, forward physics, jets, QCD
###### :
13.60.Hb, 13.85.Fb,13.87.-a
## 1 Introduction
Jet production in hadron-hadron collisions is sensitive to underlying partonic
processes, initial- and final-state radiation (isr and fsr), and to the parton
density functions (pdfs) of the colliding hadrons. Measurements of jet cross
sections at previous colliders are well described over several orders of
magnitude by perturbative calculations Collaboration (2008a) Collaboration
(2008b). However, the jets were limited to central pseudo-rapidities
($|\eta|<$2.4), where the momentum fraction of the incoming partons ($x_{1}$,
$x_{2}$) were of the same order. Jets at large pseudo-rapidities (i.e.,
forward or backward jets, $|\eta|>$3 ) result from interactions between
colliding partons with differing momentum fractions (e.g. $x_{1}<<x_{2}$),
allowing us to investigate QCD-effects at small-$x$. At such small-$x$ values,
pdfs are well constrained by dis data, but there might be additional effects
that play a role. We expect signals of parton dynamics beyond the standard
dglap evolution (e.g., BFKL or CCFM), and saturation effects are foreseen.
Moreover, forward jets are of interest in vector-boson-fusion processes, which
is one of the mechanisms for Higgs boson production.
The Large Hadron Collidor (lhc) is a proton-proton collider with a beam energy
of 7 TeV and a design luminosity of $L=$34 cm-2s-1, designed to explore a new
energy scale. At the collider, momentum fraction of the proton ($x$) carried
by the partons can become very small and the parton densities become large.
Additionally, the probability of more than one partonic interaction per event
increases for collision energies produced at the lhc, and probing the forward
region opens up opens up a large phase space for qcd emissions. Thus, the lhc
provides the perfect opportunity to study small-$x$ physics and qcd effects
through measurements of forward and forward-central jet production. Here, a
measurement of the cross section of inclusive forward jets Collaboration
(2011a) and of the production of a forward jet in conjunction with a central
jet Collaboration (2011b), using data collected by the Compact Muon Solenoid
(cms), in $\sqrt{s}=$ 7 TeV $pp$-collisions at the lhc is presented.
## 2 The CMS Detector
The CMS experiment Collaboration (2008c), located on the French-side of CERN,
operates a multi-puprpse detector to study $pp$ collisions at the LHC. The
detector covers a solid angle approaching 4$\pi$ and incorporates vertex,
calorimeter and muon chambers. The tracking system covers the pseudo-rapidity
range -2.5 $<\eta<$ 2.5 and the calorimeter system covers the range
\- 5 $<\eta<$ 5\. The Hadronic Forward Calorimeter (HF), made of iron
absorbers embedded with radiation hard quartz fibres, is used for measuring
forward jets and energy, providing an almost hermetic coverage upto
$|\eta|\sim$ 5.
## 3 Cross Section Measurement For Inclusive Forward Jet Production
Jets are built from calorimeter information using the infrared and collinear
safe anti-kT (R=0.5) jet clustering algorithm M. Cacciari (2008). A single jet
trigger (p${}_{T}>$ 15 GeV), fully efficient in region of the measurement
(p${}_{T}>$ 35 GeV) selects jets from the data sample. Jet quality criteria
are applied to remove unphysical energy deposits, and the selection
requirements ensure that the jets are well contained within the fiducial
acceptance, i.e. their reconstructed axis is within 3.2 $<|\eta|<$ 4.7
Collaboration (2007).
The jet cross section is fully corrected for detector reconstruction effects
via a bin-by-bin correction factor calculated from simulated samples. The
number of jets, Njets, are binned in transverse momentum (pT) and pseudo-
rapidity ($\eta$). The final differential inclusive jet cross-section is:
$\frac{d^{2}\sigma}{dp_{T}d\eta}=\frac{C_{unfold}}{L}\cdot\frac{N_{jets}}{\Delta
p_{T}\cdot\Delta{\eta}}$ (1)
where Cunfold is the correction factor accounting for detector effects (e.g.
migrations, resolution), and $\Delta p_{T}$ and $\Delta{\eta}$ are the bin
widths in $p_{T}$ and $\eta$ respectively.
Both experimental and theoretical sources of uncertainties have been
considered. The dominant experimental systematic uncertainty is the accuracy
of the jet energy scale (jes), which is 20%–30% in all pT bins of the measured
cross section. Additional sources include the pT resolution for bin-by-bin
corrections (3%–6%) and the luminosity measurement (4%). The theoretical
uncertainties are estimated by checking non-perturbative effects through
pythia and herwig comparisons , ascertaining the impact of different pdfs and
the variation of $\mu_{r}$ and $\mu{f}$ by a factor of 2.
Figure 1 shows the fully corrected jet cross section as a function of pT, in
comparison to various theoretical models. The yellow-band indicates the
experimental uncertainty. Figure 2 shows the fractional difference between the
experimental jet cross-section and the theoretical predictions. Within the
uncertainties, the hadron-level predictions are in good agreement with the
data. The exception is cascade, where shape variations at high pT are observed
between the measured forward jet spectra and the prediction.
Figure 1: Inclusive jet cross section measured at forward pseudo-rapidities
(3.2 $<|\eta|<$ 4.7), fully corrected and unfolded, compared to various
hadron-level predictions. The error band represents the experimental
systematic uncertainty.
Figure 2: Fractional differences between forward jet spectra and theoretical
predictions. The error-bars on the data points show statistical uncertainties.
The error bands represent the systematical and theoretical uncertainties.
## 4 Cross Section Measurement for Simultaneous Production of a Forward and a
Central Jet
The jet clustering algorithm, widths in the pT spectrum, and selection
criteria are the same in this measurement as for the inclusive jet cross
section analysis. The only difference is the additional requirement of a well-
reconstructed central ($|\eta|<$ 2.8) jet . That is, an event is accepted if
there are at least two reconstructed jets with p${}_{T}>$ 35 GeV, one with its
axis within the central region ($|\eta|<$ 2.8), and the other within the hf
(3.2 $<|\eta|<$ 4.7). If there is more than one jet present in either of the
regions, the leading jet is considered.
The jet cross-sections are obtained using a bin-by-bin correction as for the
inclusive jet production analysis. In this case, however, the simulated
samples were re-weighted at hadron level to describe the measured data
distributions. The jes remains as the dominant systematic uncertainty ($\sim$
25%). The model dependence is the second largest uncertainty, estimated by the
variation of bin-by-bin correction factors from different mcs (5% – 15%).
The fully corrected cross section for simultaneous production of at least one
central and one forward jet is measured as a function of jet pT. Figures 3– 4
present the measurement with the corresponding hadron-level predictions in
contrast, for central jets and forward jets respectively. The yellow bands
indicate the experimental systematic uncertainties summed in quadrature and
the error bars represent the statistical uncertainty. The plots on the left
show some discrepancies between pythia and data, especially in the central
region. herwig, which uses angular-ordered parton showers, describes the data
better, for both regions of pseudo-rapidity.
Figure 3: Fully corrected pT-differential jet cross-section for the central
region ($|\eta|<$ 2.8) compared to various event generators, pythia and
cascade (left), herwig and hej (right). The error-bars on the data points show
statistical uncertainties. The error bands represent the systematical.
Figure 4: Fully corrected pT-differential jet cross-section for the forward
region (3.2 $<|\eta|<$ 4.7) compared to various event generators, pythia and
cascade (left), herwig and hej (right). The error-bars on the data points show
statistical uncertainties. The error bands represent the systematical
uncertainties.
## 5 Conclusion
We present a measurement of jet production in the forward pseudo-rapidity
range 3.2 $<|\eta|<$ 4.7, and the cross section for the simultaneous
production of one central and one forward jet, using 3.14 pb-1 of cms data
collected during the early $\sqrt{s}=$ 7 TeV pp-collisions. Within the current
experimental and theoretical uncertainties, perturbative calculations
reproduce the measured inclusive forward jet cross section well. The data-
model comparison of the forward-central jet measurement show that only some
calculations are in reasonable agreement with data. Both measurements are a
first test of perturbative qcd calculations in the forward region, providing
the basis for further investigation of this interesting region of phase space.
## 6 Acknowledgements
Copyright CERN for the benefit of the CMS Collaboration.
## References
* Collaboration (2008a) CDF Collaboration, _Phys. Rev._ D78 (2008a).
* Collaboration (2008b) D0 Collaboration, _Phys. Rev. Lett._ 101 (2008b).
* Collaboration (2011a) CMS Collaboration, Measurement of forward jets in proton-proton collisions at 7 TeV (2011a), CMS Physics Analysis Summary FWD-10-003.
* Collaboration (2011b) CMS Collaboration, Cross section measurement for simultaneous production of a central and a forward jet in proton-proton collisions at 7 TeV (2011b), CMS Physics Analysis Summary FWD-10-006.
* Collaboration (2008c) CMS Collaboration, _The CMS experiment at the CERN LHC_ , JINST 3:S08004, 2008c.
* M. Cacciari (2008) G.Soyez, M. Cacciari, G.P. Salam, _JHEP_ 04 (2008).
* Collaboration (2007) CMS Collaboration, Performance of jet algorithms in CMS (2007), CMS Physics Analysis Summary JME-07-003.
|
arxiv-papers
| 2011-11-23T06:31:44 |
2024-09-04T02:49:24.617544
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Niladri Sen (for the CMS Collaboration)",
"submitter": "Niladri Sen",
"url": "https://arxiv.org/abs/1111.5410"
}
|
1111.5550
|
11institutetext: National Laboratory of Solid State Microstructure,
Department of Physics, Nanjing University - Nanjing 210093, China
Center for Statistical and Theoretical Condensed Matter Physics, Zhejiang
Normal University - Jinhua 321004, China
Quantum Hall effects Spin polarized transport in semiconductors
Magnetoelectronics; spintronics: devices exploiting spin polarized transport
or integrated magnetic fields
# Transversal Propagation of Helical Edge States in Quantum Spin Hall Systems
Feng Lu 11 Yuan Zhou 11 Jin An and Chang-De Gong 1122111122
###### Abstract
The transversal propagation of the edge states in a two-dimensional quantum
spin Hall (QSH) system is classified by the characteristic parameter
$\lambda$. There are two different types of the helical edge states, the
normal and special edge states, exhibiting distinct behaviors. The penetration
depth of the normal edge state is momentum dependent, and the finite gap for
edge band decays monotonously with sample width, leading to the normal finite
size effect. In contrast, the penetration depth maintains a uniform minimal
value in the special edge states, and consequently the finite gap decays non-
monotonously with sample width, leading to the anomalous finite size effect.
To demonstrate their difference explicitly, we compared the real materials in
phase diagram. An intuitive way to search for the special edge states in the
two-dimensional QSH system is also proposed.
###### pacs:
73.43.-f
###### pacs:
72.25.Dc
###### pacs:
85.75.-d
## 1 Introduction
Owing to the linear dispersion and $Z_{2}$ topological invariant, the
anomalous transport properties of the helical edge states in a quantum spin
Hall (QSH)system are predicted[1, 2, 3, 4, 5, 6]. Well localized edge states
were usually treated as ideal one-dimensional channels to investigate the
exotic properties[2, 5]. However, their behaviors are significantly changed
due to the transversal broadening of these edge states in real samples. The
penetration depth of helical edge states had been discussed in both HgTe
quantum well system and Bi thin film[7]. In those two systems, its value are
determined by the inverse of momentum space distance between the edge state
and the absorption point into the bulk. The finite penetration depth also
leads to the so-called finite size effect in two-dimensional (2D) QSH
system[8]. A gap opens at $\Gamma$ point when the opposite helical edge
channels overlap each other, which had been used to confirm the intrinsic spin
Hall effect in HgTe quantum well system[9]. Recently, several electric devices
had been designed to manipulate the charge and spin transport with such finite
size effect[10, 11]. An anomalous finite size effect was further reported in
three-dimensional ($3$D) topological insulator $Bi_{2}Se_{3}$, much shorter
penetration depth and oscillatory finite size gap had been revealed[12]. This
gap oscillation had also been used to search new candidate of topological non-
trivial systems[13, 14]. These previous discussions were all based on the
specific materials. The general comprehension of the transversal propagation
behaviors is expected in QSH system, especially the intrinsic difference
between the normal and anomalous finite size effects.
In this paper, the transversal propagation behaviors of the helical edge
states in 2D QSH system are investigated. The helical edge states can be
classified into two modes, the normal and special edge states, according to
the decay characteristic quantity $\lambda$. In normal edge states, the
penetration depth shows clear momentum-dependence, and the finite gap for edge
states decays monotonously with the sample width. While in special edge
states, the penetration depth keeps unchanged in the momentum space, and its
finite gap decays oscilatorily with sample width. The normal and anomalous
finite size effect can be found in the respective edge states. These facts
give explicit explanations on the difference between the real materials. Based
on the theoretical calculations, the search of the special edge states in the
$2$D case is proposed.
The paper is organized as follows: In Sec.II, we specify two different
transversal propagation modes of the helical edge states without specific
boundary condition. A semi-finite boundary condition is adopted to show the
distinct evolution of the penetration depth in the normal and special edge
states in Sec.III. As a consequent effect, the normal and anomalous finite
size effects are discussed in Sec.IV. In Sec.V, the role of particle-hole
asymmetry and the comparison of real materials are further discussed. The
conclusion is drawn in Sec.VI.
## 2 Transversal Modes of 2D QSH Model
The QSH effect was theoretically predicted in $HgTe$ quantum well[15, 3], and
soon confirmed experimentally by König el at.[16]. We start from the effective
4$\times$4 model for a 2D QSH system proposed by Bernevig, Hughes, and
Zhang[5, 3]. Very recently, it was also adopted as an effective 2D model for
the 3D topological insulator in ultrathin limit[17]. The model Hamiltonian is
expressed as
$H(k)=\left[\begin{array}[]{cc}h(k)&0\\\ 0&h^{*}(-k)\end{array}\right],$ (1)
where $h(k)=\varepsilon_{k}\mathbf{I}_{2\times
2}+\mathbf{d}_{k}\cdot{\bm{\sigma}}$, with
$\varepsilon_{k}=C-D(k^{2}_{x}+k^{2}_{y})$. The k-dependent effective field
$\mathbf{d}_{k}=(Ak_{x},-Ak_{y},M_{k})$, where
$M_{k}=M-B(k^{2}_{x}+k^{2}_{y})$. $\bm{\sigma}$ is the Pauli matrices. $A$,
$B$, $C$, $D$, and $M$ are determined by the quantum well geometry in real
materials. Here, we treat them as independent parameters and study their
respective role first. Keep in mind that, the interested topological non-
trivial QSH phase[18] emerges only when $MB>0$. In HgTe quantum well system,
such condition is controlled by the thickness of the quantum well[3, 16]. The
properties of counter-part $h^{*}(-k)$ can be conveniently obtained by
applying the time reversal operation to $h_{k}$.
To focus on the edge properties, $k_{y}$ needs to be replaced by
$-i\partial_{y}$, while $k_{x}$ remains a good quantum number due to the
translational symmetry. A trial solution of
$\Psi_{k_{x}}(y)=C_{k_{x}}e^{-\lambda y}$ can be introduced, and the decay
characteristic quantities are subsequently obtained as[8]
$\lambda_{1,2}^{2}(k_{x},E)=k_{x}^{2}+F(E)\pm\sqrt{F^{2}(E)+\frac{E^{2}-M^{2}}{B^{2}-D^{2}}},$
(2)
where $F(E)=\frac{A^{2}-2(MB+DE)}{2(B^{2}-D^{2})}$ is a function of energy
$E$.
The transversal propagation behaviors of the states are determined by both of
$\lambda_{1,2}$. However, it is clear in Eqs. (2) that, $\lambda(k_{x},E)$ has
a definite distribution in momentum space, independent of the boundary
condition. Hence, we can directly discuss the transversal propagation
behaviors of the states from Eqs. (2). There are four modes of the states
specified by different combinations of $\lambda_{1,2}$, as illustrated in
TABLE 2.
[hbtp] Combination for $\lambda_{1,2}$ Mode* Condition $\lambda_{1}$
$\lambda_{2}$ Edge1 $\lambda_{1}^{2}>\lambda_{2}^{2}\geq 0$ Real Real Bulk1**
$\lambda_{1}^{2}\geq 0>\lambda_{2}^{2}$ Real Imaginary Bulk2
$0>\lambda_{1}^{2}>\lambda_{2}^{2}$ Imaginary Imaginary Edge2
$\lambda_{1}^{2}=(\lambda_{2}^{2})^{*}$ Complex Complex
* *
For bulk states, at least one of $\lambda_{1,2}$ is purely imaginary. While
for edge states, both of $\lambda_{1,2}$ should have a non-zero real part.
* **
Trivial edge states are also included due to the real $\lambda_{1}$.
Since only the non-trivial edge states are interested in this paper, it is
natural to ask whether there is something different between the two edge modes
in TABLE 2. For convenience, we specify the Edge1 state with both
$\lambda_{1,2}$ real as the normal edge state (NES); and the Edge2 state with
$\lambda_{1,2}$ complex conjugates as the special edge state (SES). To fulfill
the condition of conjugate, the term under square root in Eqs. (2) must be
negative, which restricts SES existing in a specific regime in momentum space
confined by
$E^{SES}_{\pm}=D\left(\frac{A^{2}-2MB}{2B^{2}}\right)\pm|\frac{A}{2B}|\sqrt{\gamma\left(4MB-A^{2}\right)}$
(3)
Here a screen factor denoting the particle-hole asymmetry
$\gamma=1-\frac{D^{2}}{B^{2}}$ is introduced. The system undergoes a phase
transition from an insulator to a semimetal when $D\geq B$[19], which is not
interested for us. Eqs. (3) naturally requires $4MB\geq A^{2}\geq 0$, implying
a non-trivial QSH state. There is no special restriction for the NES.
Present classification is a natrual consequence of the breaking of periodic
boundary condition, which leads to a definite distribution of transversal
propagation modes in momentum space. An explicit boundary condition just
creates a specific spectrum onto such distribution. The emergence of SES is
determined by the $k_{x}$-independent $E_{\pm}^{SES}$. For given parameters,
$E^{SES}_{\pm}$ will squish the bulk band, leading to a flat valence band top
(or conduct band bottom), as in Fig. 1(c)(e)(f). Conversely, such feature can
be viewed as a sufficient condition for SES in QSH system, even without the
explicit knowledge of material parameters. Moreover, the squished bulk band
also gives rise to the Bulk2 states in TABLE 2, which exhibit a larger density
of states than Bulk1 states in our numerical results. Such classification is
also significant in the problems of interference tunneling and restricted edge
transport[20, 11]. More differences of the edge modes will be discussed in the
following sections.
[width=13.5cm]spectra.eps
Figure 1: Helical edge spectra for different effective parameters. $A=0.4$ and
$B=-1.0$. The upper/lower panels are spectra with/without particle-hole
symmetry. $|M|$ increases from left to right. The upper, and lower part in
each panel is the edge spectrum, and corresponding $\lambda_{1,2}$
respectively. The bold solid line in black/blue represents the edge spectrum
of NES/SES. The lighter/darker gray regimes describe the Bulk1/Bulk2 states in
Table 2. The light grey dash line are the confines of SES given by Eqs. (3).
The red solid lines in lower part of each panel is $\Re\lambda_{2}$, which
give the inverse of $\ell$. The blue dash lines are the corresponding
$\Im\lambda_{1,2}$. In panel (c), three states are specified, and
corresponding transversal propagation behaviors in real space are shown in
Fig. 2.
## 3 Penetration Depth of Helical Edge States
The penetration depth distribution of the helical edge states is distinct from
that of the chiral edge states in integer quantum Hall effect[21, 22]. The
latter is determined by the universal magnetic length, which is related to the
external magnetic field. In contrast, the penetration depth in QSH system is
$k$-dependent, originated from the band structure[5]. To address this, the
semi-infinite boundary condition[7] is adopted here. We restrict
$\Re\lambda_{1,2}\geq 0$ to obtain an evanescent edge state localized near the
boundary. The boundary condition $\Psi_{k_{x}}(0)=0$ gives the linear
dispersion relation[7, 8]
$E=-A\sqrt{\gamma}k_{x}-\frac{MD}{B}$ (4)
The penetration depth $\ell=max\\{\Re\lambda_{1,2}^{-1}\\}$[6] behaves
differently in NES and SES. The NES situation had been discussed in previous
work as in Fig. 1(a) and (d). The $k_{x}$-dependent $\ell_{NES}$ reaches its
minimum at $k_{x}=\frac{A^{2}}{4}(1-\gamma^{-1})$. The edge state is absorbed
by the bulk when $\Re\lambda_{2}=0$. In contrast, the penetration depth in SES
maintains a uniform minimal value across the whole regime of SES, given by
$\ell_{SES}=2/(\lambda_{1}+\lambda_{2})=|\frac{2\sqrt{\gamma}B}{A}|,$ (5)
which is also independent of $E$, $k_{x}$ and $M$, as shown in Fig. 1. Here
$\Re\lambda_{1,2}$ governs the transversal decay behavior. In fact, although
the relation $\lambda_{1}+\lambda_{2}=\frac{A}{\sqrt{\gamma}B}$ keeps
unchanged even in NES, the penetration depth in NES is merely determined by
the minimum of real $\lambda_{1,2}$.
In Fig. 2, the transversal propagation behaviors of three selected helical
edge states are plotted in real space. The wave function of SES(1) and SES(2)
exhibit an evanescent oscillation with different periods. However, they share
the same penetration depth. In contrast, the NES(3) shows no oscillation but
much longer penetration depth. The non-monotonous decay behavior of edge state
was also reported in the lattice model[23], which can be naturally attributed
to the SES.
The HgTe quantum well[3] and the ultra-thin $Bi_{2}Se_{3}$ film[17] correspond
to the situation in Fig. 1(d), where the SES is absent. The penetration depth
is estimated to be about 50 nm[8]. In previous studies,[7] the $Bi\\{111\\}$
thin film was compared with the HgTe quantum well system, and remarkable
difference was found in the behaviors of the penetration depth. We notice
that, a flat valence band top emerges in $Bi\\{111\\}$ spectrum[7], implying
the existence of SES. Therefore, such difference can be well understood within
present discussion. Due to the similarity of $\lambda$ at $\Gamma$ point, the
topological surface states (TSS) of the 3D topological insulator can be
equivalently discussed within our framework, corresponding to the situation in
Fig. 1(f), where the edge states are SES dominated. The penetration depth of
TSS in 3D $Bi_{2}Se_{3}$ was also reported in previous work, with a shorter
$\ell$ of about 10 nm[12]. However, they concluded $\ell$ is proportional to
the inverse of $\left|M\right|$, distinguished from present discussion. In
fact, this situation does not belong to NES, but SES, since both the decay
characteristic quantities $\lambda_{1,2}$ have image part as they stressed,
too. Therefore, the penetration depth should be independent of
$\left|M\right|$. The difference between NES and SES will be further discussed
in the next section.
[width=8.5cm]propagation.eps
Figure 2: Transversal propagation behaviors of edge state wave function in
real space. Corresponding states are marked in Fig. 1(c). The dash line
roughly gives the penetration depth behaviors with function $exp(-y/\ell)$.
The black dash line takes $\ell_{SES}$ given by Eqs. (5), while the blue one
takes $\ell=\Re\lambda_{2}^{-1}$.
## 4 Normal And Anomalous Finite Size Effects
The finite size effect in QSH system arises from the overlap of the opposite
channel due to the decreasing sample width, leading to the finite energy gap
opening for the energy dispersion of edge state near Dirac point[8]. Since the
penetration depth of NES and SES is quite different, the consequent finite
size effect is also expected to be distinct. We now turn to the ribbon
geometry with the boundary condition of
$\Psi_{k_{x}}(-L/2)=\Psi_{k_{x}}(L/2)=0$, where $L$ is the width of the
ribbon. Our numerical results reveal that, the relative gap
$\delta\Delta(k_{x})=\Delta(k_{x})-2|Ak_{x}|$ reaches its maximum at $\Gamma$
point and decays exponentially with $|k_{x}|$. Hence, we just focus on the
situation at $\Gamma$ point where $\delta\Delta(0)=\Delta(0)$. We follow the
previous discussions[8, 14] to evaluate the finite size gap in different
situations.
[width=8cm]finite-size.eps
Figure 3: Normal and anomalous finite size effects. (a), and (b) show the edge
bands and gap behaviors varying with $L$ in NES, and SES, respectively. (c)
gives the $M$-dependent $\Delta(0)$ with the same parameter adopted in Fig. 1
(in logarithmic scale). The normal (left), and anomalous (right) finite size
effect is divided by the critical $M_{c}$ (black dash line). (d), and (e)
present the conductance at finite temperature corresponding to the situation
of (a), and (b), respectively, with $1/k_{B}T=400$ and
$G_{0}=\frac{e^{2}}{h}$.
When the NES dominates the Dirac point, as discussed previously in the HgTe
quantum well[8], the gap was estimated to be
$\Delta(0)\simeq\frac{4|AM\gamma|}{\sqrt{A^{2}-4MB\gamma}}e^{-\lambda_{2}L}.$
(6)
Here we assume $\lambda_{1}L>>1$ and $\lambda_{1}>>\lambda_{2}$. This is the
normal finite size effect as shown in Fig. 3(a). When Dirac point locates
inside the SES regime, the gap turns to be
$\Delta(0)\simeq\frac{8|AM\gamma\sin(\Im\lambda_{2}L)|}{\sqrt{A^{2}-4MB\gamma}}e^{-\ell_{SES}^{-1}L},$
(7)
here $\Im\lambda_{2}=\sqrt{\frac{M}{B}-\frac{A^{2}}{4\gamma B^{2}}}$ is the
imaginary part of $\lambda_{2}$. The gap exhibits an oscillatory behavior with
$L$, as described in Fig. 3(b). This oscillation was also predicted in 3D
topological insulator[12, 14], referred as the anomalous finite size effect.
We numerically investigate the $M$-dependent evolution of $\Delta(0)$ to
distinguish the difference between NES and SES as shown in Fig. 3(c). Here
$D=0$ is applied to avoid the mismatch between the Dirac point and the regime
of SES. For small $|M|$, the Dirac point is NES, and the corresponding
$\Delta(0)$ evolves monotonously with $|M|$. For large $L$, the Dirac point
turns to be SES, $\Delta(0)$ is oscillatory. A critical $|M_{c}|=0.04$ is
obtained with the same parameters taken in Fig. 1. The number of oscillatory
periods increase with $|M|$, owing to a decreasing $\Im\lambda$. Considerable
gap always opens at $L\sim 35$ (arb. units) for all $|M|>|M_{c}|$, which
coincides with the uniform minimal $\ell_{SES}$ discussed above. Here we
emphasize that, the uniform minimal $\ell_{SES}$ found in Fig. 1 is protected
by the linear dispersion of Eqs. (4). This linear relation is not preserved
when the finite size gap opens, then $\ell$ becomes momentum-dependent again
even in SES.
The essence of such differences can be understood based on present results.
For SES, a $y$-dependent phase factor emerges due to the finite $\Im\lambda$,
which is absent in NES. The edge band is renormalized, together with gap
opening, due to the overlap of opposite edge states. Meanwhile, the
transversal phase coherence of the opposite edge states contributes to the
oscillation of $\Delta$ for SES. Hence, the interference-fringe-like picture
can be obtained as shown in Fig. 3(c).
Such effect can be detected in transport measurements at low temperature[9].
The conductance at finite temperature is simply given by[8]
$G(\mu,T)=(2e^{2}/h)\left[f(\Delta/2-\mu)-f(-\Delta/2-\mu)+1\right]$ (8)
when $\mu$ locates inside the bulk gap. Here $f(E)$ is the Fermi distribution
function and $\Delta$ is the finite size gap. Fig. 3(d) and Fig. 3(e) present
the conductance for NES and SES respectively. Recently, a non-monotonous gap
evolution had been observed in ultra-thin $Bi_{2}Se_{3}$, which is noted as a
possible anomalous finite size effect for TSS[24].
## 5 Discussion
[width=8.5cm]comparison.eps
Figure 4: Phase diagrams in parameter space. (a), and (b) are with, and
without the particle-hole symmetry respectively. The black solid line
($A^{2}=4MB$) divides the region of edge state into SES and NES. The dash line
($A^{2}=4\gamma MB$) in (b) indicates the mismatch between SES and the Dirac
point as described in the text, and moves along the direction indicated by
arrow when $\gamma$ decreases. The intensity stands for the inverse of
penetration depth $\ell^{-1}$ at Dirac point. Different materials are compared
in (c), where $A$, $B$ and $M$ are unified into the units of $meV$ and $nm$,
so that $\ell$ has a common unit of $nm$. Here the logarithmic scale is
adopted. Materials with/without SES Dirac point are marked with
cycles/squares. The intensity in (c) is for $\gamma=1$. For $Bi_{2}Se_{3}$
TSS, and $Sb_{2}Te_{3}$, $|D/B|$ is $0.13$, and $0.63$, respectively.
Up to now, we have discussed the penetration depth and the finite size effect
in NES and SES. The particle-hole asymmetry factor $\gamma$ also plays a
subtle role on these properties. The Dirac point moves upward, and the SES
regime also shifts, leading to the possible mismatch as shown in Fig. 1(e).
The existence of SES Dirac point requires
$\frac{A^{2}}{B^{2}}\leq\gamma\frac{4M}{B}$. In Fig. 4(b), the regime between
the solid line and the dash line describes the mismatch: the SES exists, but
the Dirac point moves outside. It should be pointed out that such mismatch is
not sensitive with selected $|D/B|$ unless it approaches to $1$.
The penetration depth of the edge states at Dirac point and finite size effect
are discussed together in the phase space of relative parameters $\frac{A}{B}$
and $\frac{M}{B}$, as shown in Fig. 4. In these phase diagrams, since the
solid line divides the parameter space into two regimes, $\ell^{-l}$ at Dirac
point reveals two distinct evolutions. $\ell^{-1}$ at Dirac point increases
with $M/B$ and decreases with $|A/B|$ in NES. In contrast, it remains
unchanged with $M/B$ but increases with $|A/B|$ in SES. The recently
discovered topological non-trivial systems: 2D HgTe quantum well[8], Ultrathin
$Bi_{2}Se_{3}$ film[17], 3D $Bi_{2}Se_{3}$, $Sb_{2}Te_{3}$ and $Bi_{2}Te_{3}$
[25, 6] are compared in the same phase space. The former two effective 2D
systems are described by the same model of Eqs. (1). Although the 3D
topological insulators have a different effective model[25], the situation at
Dirac point is equivalent to the two 2D systems under proper parameter
substitution[14, 12]. As in Fig. 4(c), the helical edge states in the two 2D
systems contain only NES, therefore, large size is required to avoid the
normal finite size effect. In contrast, the 2D TSS of 3D $Bi_{2}Se_{3}$ and
$Sb_{2}Te_{3}$ implies a shorter penetration depth and a possible anomalous
finite size effect[12, 24]. Interestingly, the SES exits in $Bi_{2}Te_{3}$,
however, its Dirac point moves into bulk state due to strong particle-hole
asymmetry mentioned above. This may be true as compared with the angle
resolved photoemission spectroscopy measurements[26, 6]. Similar behavior may
also can be found in the bulk $HgTe$ under uniaxial strain[27]. We expect
that, these special effects can be electrically detected in other QSH systems
with smaller $A/B$ or larger $M/B$ as shown in Fig. 4. Recently, several
designs had been performed based on the finite size effect[10, 11]. The future
applications could be quite sensitive to these properties. In this sense,
present work provides a theoretical prediction on the possible finite size
effects in new materials.
## 6 Conclusion
In conclusion, two different transversal propagation modes of the helical edge
states, i.e., NES and SES, in the QSH system are specified by the decay
characteristic quantities $\lambda$. The emergence of the flat bulk band
implies the special edge state, which gives a sufficient criterion to
distinguish the two modes. The penetration depth of SES keeps a uniform
minimal value, independent of the selected $E$, $k_{x}$ and $M$. In contrast,
it is much larger and shows clear momentum dependence in NES. Different finite
size effects are studied in respective edge states. Especially, the
oscillatory gap for edge band is found in SES. Some real materials are
compared in the phase diagram to demonstrate the difference between NES and
SES. We also give clues to search possible QSH materials with SES for future
applications.
###### Acknowledgements.
We would like to thank Y. F. Wang, and L. Xu for helpful discussions. This
work is supported by NSFC Project No. 10804047, and A Project Funded by the
Priority Academic Program Development of Jiangsu Higher Education
Institutions. J. An acknowledges NSFC Project No. 10804073. C. D. Gong also
acknowledges 973 Projects No. 2009CB929504.
## References
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* [5] König M., Buhmann H., Molenkamp L. W., Hughes T., Liu C.-X., Qi X.-L. Zhang S.-C. Journal of Physical Society of Japan772008031007.
* [6] Qi X.-L. Zhang S.-C. Rev. Mod. Phys.8320111057.
* [7] Wada M., Murakami S., Freimuth F. Bihlmayer G. Phys. Rev. B832011121310.
* [8] Zhou B., Lu H.-Z., Chu R.-L., Shen S.-Q. Niu Q. Phys. Rev. Lett.1012008246807.
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* [16] König M., Wiedmann S., Brüne C., Roth A., Buhmann H., Molenkamp L. W., Qi X.-L. Zhang S.-C. Science3182007766.
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* [19] Mao S. Kuramoto Y. Phys. Rev. B832011085114.
* [20] Guigou M., Recher P., Cayssol J. Trauzettel B. Phys. Rev. B842011094534.
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* [22] Thouless D. J., Kohmoto M., Nightingale M. P. den Nijs M. Phys. Rev. Lett.491982405.
* [23] Ohyama Y., Tsuchiura H. Sakuma A. Journal of Physics: Conference Series2662011012103.
* [24] Sakamoto Y., Hirahara T., Miyazaki H., Kimura S. Hasegawa S. Phys. Rev. B812010165432.
* [25] Zhang H., Liu C.-X., Qi X.-L., Dai X., Fang Z. Zhang S.-C. Nat. Phys.52009438.
* [26] Li Y.-Y., Wang G., Zhu X.-G., Liu M.-H., Ye C., Chen X., Wang Y.-Y., He K., Wang L.-L., Ma X.-C., Zhang H.-J., Dai X., Fang Z., Xie X.-C., Liu Y., Qi X.-L., Jia J.-F., Zhang S.-C., Xue Q.-K. Adv. Mat.2220104002.
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|
arxiv-papers
| 2011-11-23T16:54:56 |
2024-09-04T02:49:24.625528
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Feng Lu, Yuan Zhou, Jin An, and Chang-De Gong",
"submitter": "Yuan Zhou",
"url": "https://arxiv.org/abs/1111.5550"
}
|
1111.5836
|
LPT Orsay 11-60
Enhanced Higgs Mediated Lepton Flavour Violating Processes in the
Supersymmetric Inverse Seesaw Model
Asmaa Abada, Debottam Das and Cédric Weiland
Laboratoire de Physique Théorique, CNRS – UMR 8627,
Université Paris-Sud 11, F-91405 Orsay Cedex, France
We study the impact of the inverse seesaw mechanism on several low-energy
flavour violating observables such as $\tau\rightarrow\mu\mu\mu$ in the
context of the Minimal Supersymmetric Standard Model. As a consequence of the
inverse seesaw, the contributions of the right-handed sneutrinos significantly
enhance the Higgs-mediated penguin diagrams. We find that different flavour
violating branching ratios can be enhanced by as much as two orders of
magnitude. We also comment on the impact of the Higgs-mediated processes on
the leptonic $B$-meson decays and on the Higgs flavour violating decays.
KEYWORDS: Supersymmetry, Lepton Flavour Violation, Inverse Seesaw
## 1 Introduction
Neutrino oscillations have provided indisputable evidence for flavour
violation in the neutral lepton sector. In the absence of any fundamental
principle that prevents charged lepton flavour violation, one expects that
extensions of the Standard Model (SM) accommodating neutrino masses and
mixings should also allow for lepton flavour violation (LFV) in the charged
lepton sector. Indeed, the additional new flavour dynamics and new field
content present in many extensions of the SM may provide contributions to
charged LFV (cLFV) processes such as radiative (e.g. $\mu\to e\gamma$) and
three-body lepton decays (for instance $\tau\to\mu\mu\mu$). These decays
generally arise from higher order processes, and so their branching ratios
(Brs) are expected to be small, making them difficult to observe. Thus, any
cLFV signal would provide clear evidence for new physics: mixings in the
lepton sector and probably the presence of new particles, possibly shedding
light on the origin of neutrino mass generation.
The search for manifestations of charged LFV constitutes the goal of several
experiments [1, 2, 3, 4, 5, 6, 7, 8, 9, 10, 11, 12, 13], dedicated to look for
signals of processes such as rare radiative decays, three-body decays and
muon-electron conversion in nuclei. Despite the fact that a cLFV signal could
provide clear evidence for new physics, the underlying mechanism of lepton
mixing might be difficult to unravel. In parallel to the low-energy searches
for new physics, i.e. via indirect effects of possible new particles, the LHC
has started to search directly for these new particles in its quest to unveil
the mechanism of electroweak symmetry breaking, thus possibly providing a
solution to the SM hierarchy problem.
Among the many possible extensions of the SM, supersymmetry (SUSY) is a well
motivated solution for the hierarchy problem, providing many other appealing
aspects such as gauge coupling unification and dark matter candidates. If the
LHC experiments indeed discover SUSY, it is then extremely interesting to
consider supersymmetric models that can also explain neutrino masses and
mixings. Furthermore, it is only natural to expect that such models might also
give rise to cLFV. If SUSY is indeed realised in nature, cLFV (mediated by new
sparticles) would provide a new probe to explore the origin of lepton mixings,
playing a complementary rôle to other searches of new physics, i.e. LHC direct
searches and neutrino dedicated experiments.
One of the most economical possibilities is to embed a seesaw mechanism [14,
15, 16] in the framework of SUSY models (i.e. the SUSY seesaw) [17]. For any
seesaw realisation, the neutrino Yukawa couplings could leave their imprints
in the SUSY soft-breaking slepton mass matrices, and consequently induce
flavour violation at low energies due to the renormalisation group (RG)
evolution of the SUSY soft-breaking parameters. The caveat of these scenarios
is that, in order to have sufficiently large Yukawa couplings (as required to
account for large cLFV Brs), the typical scale of the extra particles (such as
right handed neutrinos, scalar or fermionic isospin triplets) is in general
very high, potentially very close to the gauge coupling unification scale.
However, such a high (seesaw) scale would be impossible to probe
experimentally.
On the other hand, the so-called inverse seesaw [18] constitutes a very
appealing alternative to the ”standard” seesaw realisations. Embedding an
inverse seesaw mechanism in the Minimal Supersymmetric extension of the SM
(MSSM) requires the inclusion of two additional gauge singlet superfields,
with opposite lepton numbers ($+1$ and $-1$). When compared to other SUSY
seesaw realisations, cLFV observables are enhanced in this framework , and
such an enhancement can be attributed to large neutrino Yukawa couplings
($Y_{\nu}\sim O(1)$), compatible with a seesaw scale $M$, close to the
electroweak one, thus within LHC reach.
The differences between the inverse seesaw and the standard one can be
conceptually understood from an effective point of view and linked to the
distinct properties of the lepton number violating dimension-5 (Weinberg)
operator (responsible for neutrino masses and mixings) and the total lepton
number conserving dimension-6 operator, which is at the origin of cLFV.
Contrary to what occurs in the standard seesaw, these two operators are de-
correlated in the inverse seesaw, implying that the suppression of the
coefficient of the dimension-5 operator will not affect the size of the
coefficient of the dimension-6 operator. In both seesaws, the latter operator
is proportional to
$\left(Y_{\nu}^{\dagger}\frac{1}{\left|M\right|^{2}}Y_{\nu}\right)$; however,
in the case of a type I seesaw, the dimension-5 operator is proportional to
$\left(Y_{\nu}^{\dagger}\frac{1}{M}Y_{\nu}\right)$, while in the case of an
inverse seesaw, it has a further suppression of $\frac{\mu}{M}$ ($\mu$ being a
dimensionful parameter, linked to the mass of the sterile singlets). The
dimension-6 operator will thus be extremely suppressed in the case of a type I
seesaw, since in this case $M$ is very large to accommodate natural $Y_{\nu}$.
In contrast, in the inverse seesaw, small neutrino masses can easily be
accommodated via tiny values of $\mu$, which will not affect the dimension 6
operator. Furthermore, such small values of $\mu$ are natural in the sense of
’t Hooft since in the limit where $\mu\to 0$, the total lepton number symmetry
is restored [19].
In view of the strong potential of the inverse seesaw mechanism regarding
cLFV, several phenomenological studies have recently been carried out [20, 21,
22, 23, 24, 25, 26, 27, 28, 29]. While a non-supersymmetric inverse seesaw
requires two pairs of singlets to explain neutrino oscillation data [26], the
supersymmetric generalization can accommodate neutrino data [28] with just one
pair of singlets. The latter scenario is also known as the minimal
supersymmetric inverse seesaw model (MSISM). This model can also comply with
the dark matter relic abundance of the Universe [23].
The extra TeV scale singlet neutrinos may significantly contribute to cLFV
observables, irrespective of the supersymmetric states [30]. Supersymmetric
realisations of the inverse seesaw may enhance these cLFV rates even further
(e.g. the contributions to $l_{i}\rightarrow l_{j}\gamma$, which has been
analysed in [20]). Furthermore, this seesaw model can have LHC signatures: the
extra singlets can participate in the decay chains, leading to effects which
can be important, particularly in the case in which one of the singlets is the
lightest supersymmetric particle (LSP) [28].
In this paper, we focus on contributions to cLFV observables, such as
$\tau\to\mu\mu\mu$, arising from a Higgs-mediated effective vertex. We explore
the contributions which are due to the presence of comparatively light right-
handed neutrinos and sneutrinos (which are usually negligible in the framework
of a type I SUSY-seesaw), while still having large neutrino Yukawa couplings.
We find that all these contributions can lead to a significant enhancement of
several cLFV observables.
The paper is organised as follows. In Section 2, we define the model,
providing a brief overview on the implementation of the inverse seesaw in the
MSSM. In Section 3, we discuss the implications of this model regarding low-
energy cLFV observables, particular emphasis being given to the Higgs-mediated
processes. In Section 4, we study the Higgs-mediated contributions to several
lepton flavour violating observables and compare our results to present bounds
and to future experimental sensitivities in Section 5. Then we draw some
generic conclusions on the viability of an inverse seesaw as the underlying
mechanism of LFV. We finally conclude in Section 6.
## 2 Inverse Seesaw Mechanism in the MSSM
The model consists of the MSSM extended by three pairs of singlet superfields,
$\widehat{\nu}^{c}_{i}$ and $\widehat{X}_{i}$ ($i=1,2,3$)111
$\widetilde{\nu}^{c}=\widetilde{\nu}_{R}^{*}$ with lepton numbers assigned to
be $-1$ and $+1$, respectively. The supersymmetric inverse seesaw model is
defined by the superpotential
$\displaystyle{\mathcal{W}}$ $\displaystyle=$
$\displaystyle\varepsilon_{ab}\left[Y^{ij}_{d}\widehat{D}_{i}\widehat{Q}_{j}^{b}\widehat{H}_{d}^{a}+Y^{ij}_{u}\widehat{U}_{i}\widehat{Q}_{j}^{a}\widehat{H}_{u}^{b}+Y^{ij}_{e}\widehat{E}_{i}\widehat{L}_{j}^{b}\widehat{H}_{d}^{a}\right.$
(2.1) $\displaystyle+$
$\displaystyle\left.Y^{ij}_{\nu}\widehat{\nu}^{c}_{i}\widehat{L}^{a}_{j}\widehat{H}_{u}^{b}-\mu\widehat{H}_{d}^{a}\widehat{H}_{u}^{b}\right]+M_{R_{i}}\widehat{\nu}^{c}_{i}\widehat{X}_{i}+\frac{1}{2}\mu_{X_{i}}\widehat{X}_{i}\widehat{X}_{i}~{},$
where $i,j=1,2,3$ denote generation indices. In the above, $\widehat{H}_{d}$
and $\widehat{H}_{u}$ are the down- and up-type Higgs superfields,
$\widehat{L}_{i}$ denotes the SU(2) doublet lepton superfields. $M_{R_{i}}$
represents the right-handed neutrino mass term which conserves lepton number.
Due to the presence of non-vanishing $\mu_{X_{i}}$, the total lepton number
$L$ is no longer a good quantum number; nevertheless, notice that in our
formulation $(-1)^{L}$ is still a good symmetry. Without loss of generality,
the terms $\widehat{\nu}^{c}_{i}\widehat{X}_{i}$ and
$\widehat{X}_{i}\widehat{X}_{i}$ are taken to be diagonal in generation space.
Clearly, as $\mu_{X_{i}}\rightarrow 0$, lepton number conservation is
restored, since $M_{R}$ does not violate lepton number. Although in the
present study we consider three generations of $\widehat{\nu}^{c}$ and
$\widehat{X}$, we recall that in the minimal version of the SUSY inverse
seesaw (where only one generation of $\widehat{\nu}^{c}$ and $\widehat{X}$ is
included), neutrino data can be accommodated [28].
The soft SUSY breaking Lagrangian can be written as
$-{\mathcal{L}}_{\rm soft}=-{\mathcal{L}}^{\rm MSSM}_{\rm
soft}+m^{2}_{\widetilde{\nu}^{c}}\widetilde{\nu}^{c\dagger}_{i}\widetilde{\nu}^{c}_{i}+m^{2}_{X}\widetilde{X}^{\dagger}_{i}\widetilde{X}_{i}+\left(A_{\nu}Y^{ij}_{\nu}\varepsilon_{ab}\widetilde{\nu}^{c}_{i}\widetilde{L}^{a}_{j}H_{u}^{b}+B_{M_{R_{i}}}\widetilde{\nu}^{c}_{i}\widetilde{X}_{i}+\frac{1}{2}B_{\mu_{X_{i}}}\widetilde{X}_{i}\widetilde{X}_{i}+{\rm
h.c.}\right),$ (2.2)
where ${\mathcal{L}}^{\rm MSSM}_{\rm soft}$ denotes the soft SUSY breaking
terms of the MSSM. In the above, the singlet scalar states $\widetilde{X}_{i}$
and $\widetilde{\nu}^{c}_{i}$ are assumed to have flavour universal masses,
i.e. $m^{2}_{X_{i}}=m^{2}_{X}$ and
$m^{2}_{\widetilde{\nu}^{c}_{i}}=m^{2}_{\widetilde{\nu}^{c}}$. The parameters
$B_{M_{R_{i}}}$ and $B_{\mu_{X_{i}}}$ are the new terms involving the scalar
partners of the sterile neutrino states (notice that while the former
conserves lepton number, the latter gives rise to a lepton number violating
$\Delta L=2$ term). Working under the assumption of a flavour-blind mechanism
for SUSY breaking, we will assume universal boundary conditions222In our
subsequent numerical analysis, we will relax some of these universality
conditions, considering non-universal soft breaking terms for the Higgs
sector. In what concerns the right-handed sneutrino sector, we will assume
that the corresponding soft-breaking masses hardly run between the GUT and the
low-energy scale. for the soft SUSY breaking parameters at some very high
energy scale (e.g. the gauge coupling unification scale $\sim 10^{16}$ GeV),
$m_{\phi}=m_{0}\,,M_{\text{gaugino}}=M_{1/2}\,,A_{i}=A_{0}\,.$ (2.3)
Before addressing neutrino mass generation, a few comments on the nature of
the superpotential are in order. As can be seen from Eq. (2.1), the two
singlets $\widehat{\nu}^{c}_{i}$ and $\widehat{X}_{i}$ are differently treated
in the sense that a $\Delta L=2$ Majorana mass term is present for
$\widehat{X}_{i}$ ($\mu_{X_{i}}\widehat{X}_{i}\widehat{X}_{i}$), while no
$\mu_{\nu^{c}_{i}}\widehat{\nu}^{c}_{i}\widehat{\nu}^{c}_{i}$ is present in
${\mathcal{W}}$. Although a generic superpotential with $(-1)^{L}$ should
contain the latter term, let us notice that similar to what occurs for
$\mu_{X_{i}}$, the absence of $\mu_{\nu^{c}_{i}}$ also enhances the symmetry
of the model; moreover, we emphasise that it is the magnitude of $\mu_{X_{i}}$
(and not that of $\mu_{\nu^{c}_{i}}$) which controls the size of the light
neutrino mass [24, 29]. In view of this, and for the sake of simplicity, we
have assumed $\mu_{\nu^{c}_{i}}=0$ (considering non-vanishing, yet small
values of $\mu_{\nu^{c}_{i}}$ would not change the qualitative features of the
model). Although in our formulation we treat $\mu_{X_{i}}$ as an effective
parameter, its origin can be explained either dynamically or in a framework of
SUSY Grand Unified Theories (GUT) [24, 29, 25]. Furthermore
$\mu_{\nu^{c}_{i}}\ll\mu_{X_{i}}$ can also be realised in extended frameworks
[24].
In order to illustrate the pattern of light neutrino masses in the inverse
seesaw model and how it is related to the lepton number violating parameter
$\mu_{X_{i}}$, we consider the one-generation case. In the
$\\{\nu,{\nu^{c}},X\\}$ basis the $(3\times 3)$ neutrino mass matrix can be
written as
$\displaystyle{\cal M}$ $\displaystyle=$
$\displaystyle\left(\begin{array}[]{ccc}0&m_{D}&0\\\ m_{D}&0&M_{R}\\\
0&M_{R}&\mu_{X}\\\ \end{array}\right)\ ,$ (2.7)
with $m_{D}=Y_{\nu}v_{u}$, yielding the mass eigenvalues ($m_{1}\ll m_{2,3}$):
$\displaystyle
m_{1}=\frac{m_{D}^{2}\mu_{X}}{m_{D}^{2}+M_{R}^{2}}\,,~{}~{}~{}~{}m_{2,3}=\mp\sqrt{M_{R}^{2}+ùm_{D}^{2}}+\frac{M_{R}^{2}\mu_{X}}{2(m_{D}^{2}+M_{R}^{2})}\,.$
(2.8)
The above equation clearly reveals that the lightness of the smallest
eigenvalue $m_{1}$ is indeed due to the smallness of $\mu_{X}$ ($\mu_{X}\simeq
m_{1}$). Thus the lepton number conserving mass parameters ($m_{D}$ and
$M_{R}$) are completely unconstrained in this model. Finally, it is worth
noticing that the effective right-handed sneutrino mass term (Dirac-like) is
given by
$M^{2}_{\widetilde{\nu}^{c}_{i}}=m^{2}_{\widetilde{\nu}^{c}}+M_{R_{i}}^{2}+\sum_{j}{|Y^{ij}_{\nu}|^{2}v_{u}^{2}}$.
Assuming $M_{R}\sim{\mathcal{O}}$(TeV), the effective mass term will not be
very large, in clear contrast to what occurs in the standard (type I) SUSY
seesaw. In our analysis, we will be particularly interested in the rôle of
such a light sneutrino (i.e. $M^{2}_{\widetilde{\nu}^{c}}\sim
M^{2}_{\text{SUSY}}$) in the enhancement of Higgs mediated contributions to
lepton flavour violating observables.
## 3 Lepton flavour violation: Higgs-mediated contributions
In the SUSY seesaw framework, the only source of flavour violation is encoded
in the neutrino Yukawa couplings (which are necessarily non-diagonal to
account for neutrino oscillation data); even under the assumption of universal
soft breaking terms at the GUT scale, radiative effects proportional to
$Y_{\nu}$ induce flavour violation in the slepton mass matrices, which in turn
give rise to slepton mediated cLFV observables [31, 32]. As an example, in the
leading logarithmic approximation, the RGE corrections to the left-handed
slepton soft-breaking masses are given by
$\displaystyle(\Delta m_{\widetilde{L}}^{2})_{ij}$ $\displaystyle\simeq$
$\displaystyle-\frac{1}{8\pi^{2}}(3m_{0}^{2}+A_{0}^{2})(Y_{\nu}^{\dagger}LY_{\nu})_{ij}\,,~{}~{}L=\ln\frac{M_{GUT}}{M_{R}}\,$
(3.1) $\displaystyle=$ $\displaystyle\xi(Y^{\dagger}_{\nu}Y_{\nu})_{ij}.$
(For simplicity, in the above we are implicitly assuming a degenerate right-
handed neutrino spectrum, $M_{R_{i}}=M_{R}$.) The RGE-induced flavour
violating entries, $(\Delta m_{\widetilde{L}}^{2})_{ij}$, give rise to the
dominant contributions to low-energy flavour violating observables in the
charged lepton sector, such as $\ell_{i}\to\ell_{j}\gamma$ (mediated by
chargino-sneutrino and neutralino-slepton loops) and
$\ell_{i}\to\ell_{j}\ell_{k}\ell_{m}$ (from photon, $Z$ and Higgs mediated
penguin diagrams).
Compared to the standard (type I) SUSY seesaw, where $M_{R}\sim 10^{14}$ GeV,
the inverse seesaw is characterised by a right handed neutrino mass scale
$M_{R}\sim\mathcal{O}(\text{TeV})$ and this in turn leads to an enhancement of
the factor $\xi$, (see Eq. (3.1)), and hence to all low-energy cLFV
observables, in the latter framework. Furthermore, having right-handed
sneutrinos whose mass is of the same order of the other sfermions, i.e.
$M^{2}_{\widetilde{\nu}^{c}}\sim M^{2}_{\text{SUSY}}$, the
$\widetilde{\nu}^{c}$-mediated processes are no longer suppressed, and might
even significantly contribute to the low-energy flavour violating observables.
Here, we focus on the impact of such a light $\widetilde{\nu}^{c}$ in the
Higgs mediated processes which are expected to be important in the large
$\tan\beta$ regime.
Although at tree level Higgs-mediated neutral currents are flavour conserving,
non-holomorphic Yukawa interactions of the type $\bar{D}_{R}Q_{L}H_{u}^{*}$
can be induced at the one-loop level, as first noticed in [33]. In the large
$\tan\beta$ regime, in addition to providing significant corrections to the
masses of the $b$-quark, these non-holomorphic couplings have an impact on
$B^{0}-\bar{B}^{0}$ mixing and flavour violating decays, in particular
$B_{s}\rightarrow\mu^{+}\mu^{-}$ [34, 35, 36, 37, 38]. Similarly, in the
lepton sector, the origin of the Higgs-mediated flavour violating couplings
can be traced to a non-holomorphic Yukawa term of the form
$\bar{E}_{R}LH_{u}^{*}$ [39]. Other than the corrections to the $\tau$ lepton
mass, these new couplings give rise to additional contributions to several
cLFV processes mediated by Higgs exchange. In particular
$B_{s}\rightarrow\mu\tau$, $B_{s}\rightarrow e\tau$ (the so-called double
penguin processes) were considered in [38], while $\tau\rightarrow\mu\eta$ was
studied in [40]. A detailed analysis of the several $\mu-\tau$ lepton flavour
violating processes, namely $\tau\rightarrow\mu X$
($X=\gamma,e^{+}e^{-},\mu^{+}\mu^{-},\rho,\pi,\eta,\eta^{\prime}$) can be
found in [41].
Even though the flavour violating processes in the quark and lepton sectors
have a similar diagrammatic origin, the source of flavour violation is
different in each case. In the quark sector, trilinear soft SUSY breaking
couplings involving up-type squarks provide the dominant source of flavour
violation [35], while in the lepton case, LFV stems from the radiatively
induced non-diagonal terms in the slepton masses (see Eq. (3.1)) [39].
In the standard SUSY seesaw (type I), the term
${\widetilde{\nu}^{c}_{i}}H_{u}{\widetilde{L}_{Lj}}$ is usually neglected, as
it is suppressed by the very heavy right handed sneutrino masses
(${M_{\widetilde{\nu}^{c}_{i}}}\sim 10^{14}$GeV). However, in scenarios such
as the inverse SUSY seesaw, where
${M_{\widetilde{\nu}^{c}_{i}}}\sim\mathcal{O}$(TeV), this term may provide the
dominant contributions to Higgs mediated lepton flavour violation.
The effective Lagrangian describing the couplings of the neutral Higgs fields
to the charged leptons is given by
$\displaystyle-{\cal
L}^{\text{eff}}=\bar{E}^{i}_{R}Y_{e}^{ii}\left[\delta_{ij}H_{d}^{0}+\left(\epsilon_{1}\delta_{ij}+\epsilon_{2ij}(Y_{\nu}^{\dagger}Y_{\nu})_{ij}\right)H_{u}^{0\ast}\right]E^{j}_{L}+\text{h.c.}\,.$
(3.2)
In the above, the first term corresponds to the usual Yukawa interaction,
while the coefficient $\epsilon_{1}$ encodes the corrections to the charged
lepton Yukawa couplings. In the basis where the charged lepton Yukawa
couplings are diagonal, the last term in Eq. (3.2), i.e.
$\epsilon_{2ij}(Y_{\nu}^{\dagger}Y_{\nu})_{ij}$, is in general non-diagonal,
thus providing a new source of charged lepton flavour violation through Higgs
mediation. Its origin can be diagrammatically understood from Fig.1, where
flavour violation is parametrized via a mass insertion $(\Delta
m_{\widetilde{L}}^{2})_{ij}$ (see Eq. (3.1)).
|
---|---
|
Figure 1: Diagrams contributing to $\epsilon_{2}$. Crosses on scalar lines
represent LFV mass insertions $(\Delta m_{\widetilde{L}}^{2})_{ij}$, while
those on fermion lines denote chirality flips.
The coefficient $\epsilon_{2}$ can be estimated as
$\displaystyle\epsilon_{2ij}$ $\displaystyle=$
$\displaystyle\frac{\alpha^{\prime}}{8\pi}\xi\mu
M_{1}\left[2F_{2}\left(M_{1}^{2},m_{\widetilde{E}_{Lj}}^{2},m_{\widetilde{E}_{Li}}^{2},m_{\widetilde{E}_{Ri}}^{2}\right)-F_{2}\left(\mu^{2},m^{2}_{\widetilde{E}_{Lj}},m^{2}_{\widetilde{E}_{Li}},M_{1}^{2}\right)\right]+$
(3.4) $\displaystyle\frac{\alpha_{2}}{8\pi}\xi\mu
M_{2}\left[F_{2}\left(\mu^{2},m^{2}_{\widetilde{E}_{Lj}},m^{2}_{\widetilde{E}_{Li}},M_{2}^{2}\right)+2F_{2}\left(\mu^{2},m_{\widetilde{\nu}_{Lj}}^{2},m_{\widetilde{\nu}_{Li}}^{2},M_{2}^{2}\right)\right]\,,$
where
$\displaystyle F_{2}\left(x,y,z,w\right)=-\frac{x\ln
x}{(x-y)(x-z)(x-w)}-\frac{y\ln y}{(y-x)(y-z)(y-w)}+(x\leftrightarrow
z,y\leftrightarrow w)\,.$ (3.5)
Here, $M_{1}$ and $M_{2}$ are the masses of the electroweak gauginos at low
energies. On the other hand, the flavour conserving loop-induced form factor
$\epsilon_{1}$ (notice that the diagrams of Fig.1 contribute to this form
factor, but without the slepton flavour mixings in the internal lines) can be
expressed as [39, 38]
$\displaystyle\epsilon_{1}$ $\displaystyle=$
$\displaystyle\frac{\alpha^{\prime}}{8\pi}\mu
M_{1}\left[2F_{1}\left(M_{1}^{2},m_{\widetilde{E}_{L}}^{2},m_{\widetilde{E}_{R}}^{2}\right)-F_{1}\left(M_{1}^{2},\mu^{2},m^{2}_{\widetilde{E}_{L}}\right)+2F_{1}\left(M_{1}^{2},\mu^{2},m^{2}_{\widetilde{E}_{R}}\right)\right]$
(3.6) $\displaystyle+\frac{\alpha_{2}}{8\pi}\mu
M_{2}\left[F_{1}\left(\mu^{2},m^{2}_{\widetilde{E}_{L}},M_{2}^{2}\right)+2F_{1}\left(\mu^{2},m_{\widetilde{\nu}_{L}}^{2},M_{2}^{2}\right)\right],$
with
$\displaystyle F_{1}\left(x,y,z\right)$ $\displaystyle=$
$\displaystyle-\frac{xy\ln(x/y)+yz\ln(y/z)+zx\ln(z/x)}{(x-y)(y-z)(z-x)}\,.$
(3.7)
In the standard seesaw mechanism, the diagrams of Fig. 1 provide the only
source for Higgs-mediated lepton flavour violation. However, in the framework
of the inverse SUSY seesaw, there is an additional diagram that may even
account for the dominant Higgs-mediated lepton flavour violation contribution:
the sneutrino-chargino mediated loop, depicted in Fig. 2. (Due to the large
masses of $\widetilde{\nu}^{c}$ in the standard (type I) seesaw, this process
provides negligible contributions, and is hence not taken into account.)
Figure 2: Right-handed sneutrino contribution to $\epsilon^{\prime}_{2}$. This
contribution is particularly relevant when $\widetilde{\nu}^{c}$ is light.
The effective Lagrangian terms encoding lepton flavour violation is
accordingly modified as
$\displaystyle-{\cal
L}^{\text{LFV}}=\bar{E}^{i}_{R}Y_{e}^{ii}\epsilon^{\text{tot}}_{2ij}(Y_{\nu}^{\dagger}Y_{\nu})_{ij}H_{u}^{0\ast}E^{j}_{L}+\text{h.c.}\,,$
(3.8)
where $\epsilon^{\text{tot}}_{2}=\epsilon_{2}+\epsilon_{2}^{\prime}$,
$\epsilon_{2}^{\prime}$ being the contribution from the new diagram. This
contribution can be expressed as
$\displaystyle\epsilon^{\prime}_{2ij}=\frac{1}{16\pi^{2}}\mu
A_{\nu}F_{1}(\mu^{2},m^{2}_{\widetilde{\nu}_{i}},M^{2}_{\widetilde{\nu}^{c}_{j}}).$
(3.9)
In the above, we have parametrized the soft trilinear term for the neutral
leptons as $A_{\nu}Y_{\nu}$, where $A_{\nu}$ is a flavour independent real
mass term.
Below, we provide an approximate estimate of the relative contributions of the
terms $\epsilon_{2}$ and $\epsilon^{\prime}_{2}$: for simplicity we take
$M_{\widetilde{\nu}^{c}}\sim\mathcal{O}$(TeV) and assume common values for the
masses of all SUSY particles and dimensionful terms $A_{\nu}$ at low energies,
symbolically denoted by $A_{\nu}\sim\langle\widetilde{m}\rangle\sim
M_{\text{SUSY}}$. In this limit, the loop functions are given by
$F_{2}\left(x,x,x,x\right)=\frac{1}{6x^{2}}$ and
$F_{1}\left(x,x,x\right)=\frac{1}{2x}$. This leads to
$\displaystyle\epsilon_{2}=\frac{1}{8\pi}\xi{\widetilde{m}}^{2}\left(\frac{\alpha^{\prime}}{6{\widetilde{m}}^{4}}+3\frac{\alpha_{2}}{{6\widetilde{m}}^{4}}\right)\simeq-0.0007\,,$
(3.10)
while
$\displaystyle\epsilon^{\prime}_{2}=\frac{1}{16\pi^{2}}{\widetilde{m}}^{2}\frac{1}{2{\widetilde{m}}^{2}}\simeq
0.003\,.$ (3.11)
In this illustrative (leading order) calculation, we have assumed that at
$M_{\text{GUT}}$, one has $A_{0}=0$, taking for the gauge couplings
$\alpha_{2}=0.03$ and $\alpha^{\prime}=0.008$. Following Eq. (3.1), and
assuming $M_{R}=10^{3}~{}$GeV, one gets $\xi\sim-1.1\,m_{0}^{2}$. Thus, at the
leading order in the inverse seesaw, the lepton flavour violation coefficient
becomes
$|\epsilon^{\text{tot}}_{2}|=|\epsilon_{2}+\epsilon^{\prime}_{2}|\simeq
2\times 10^{-3}$.
For completeness, let us notice that in the standard seesaw model (where
sizable Yukawa couplings are typically associated to a right-handed neutrino
mass scale $\sim 10^{14}$ GeV), assuming the same amount of flavour violation
as parametrized by $\xi$, one finds
$|\epsilon^{\text{tot}}_{2}|=|\epsilon_{2}|\simeq 2\times 10^{-4}$. This
clearly reveals that in the inverse SUSY seesaw, $\epsilon^{\text{tot}}_{2}$
is enhanced by a factor of order $\sim 10$ compared to the standard seesaw.
The large enhancement of $\epsilon^{\text{tot}}_{2}$ will have an impact
regarding all Higgs-mediated lepton flavour violating observables. The
computation of the cLFV observables requires specifying the couplings of the
physical Higgs bosons to the leptons, in particular
$\bar{E}^{i}_{R}E^{j}_{L}H_{k}$ (where $H_{k}=h,H,A$). The effective
Lagrangian describing this interaction can be derived from Eq. (3.2), and
reads [39, 38] as
$\displaystyle-{\cal L}^{\text{eff}}_{i\neq
j}=(2G_{F}^{2})^{1/4}\,\frac{m_{E_{i}}\kappa^{E}_{ij}}{\cos^{2}\beta}\left(\bar{E}^{i}_{R}\,E^{j}_{L}\right)\left[\cos(\alpha-\beta)h+\sin(\alpha-\beta)H-iA\right]+\text{h.c.}\,,\,\,\,$
(3.12)
where $\alpha$ is the CP-even Higgs mixing angle and $\tan\beta=v_{u}/v_{d}$,
and
$\displaystyle\kappa^{E}_{ij}$ $\displaystyle=$
$\displaystyle\frac{\epsilon^{\text{tot}}_{2ij}(Y^{\dagger}_{\nu}Y_{\nu})_{ij}}{\left[1+\left(\epsilon_{1}+\epsilon^{\text{tot}}_{2ii}(Y^{\dagger}_{\nu}Y_{\nu})_{ii}\right)\tan\beta\right]^{2}}\
.$ (3.13)
As clear from the above equation, large values of $\epsilon^{\text{tot}}_{2}$
lead to an augmentation of $\kappa^{E}_{ij}$. Given that the cLFV branching
ratios are proportional to $({\kappa^{E}_{ij}})^{2}$, a sizeable enhancement,
as large as two orders of magnitude, is expected for all Higgs-mediated LFV
observables.
## 4 Higgs-mediated lepton flavour violating observables
Here we focus our attention on the cLFV observables where the dominant
contribution to flavour violation arises from the Higgs penguin diagrams, in
particular those involving $\tau$-leptons (due to the comparatively large
value of $Y_{\tau}$).
In what follows, we discuss some of these LFV decays in detail.
* •
$\tau\rightarrow 3\mu$
In the large $\tan\beta$ regime, Higgs-mediated flavour violating diagrams
would be particularly important in this decay mode. The branching ratio can be
expressed as [39, 38]
$\displaystyle\text{Br}(\tau\to 3\mu)$ $\displaystyle=$
$\displaystyle\frac{G_{F}^{2}\,m_{\mu}^{2}\,m_{\tau}^{7}\,\tau_{\tau}}{1536\,\pi^{3}\cos^{6}\beta}\,|\kappa_{\tau\mu}^{E}|^{2}\left[\left(\frac{\sin(\alpha-\beta)\cos\alpha}{M_{H}^{2}}-\frac{\cos(\alpha-\beta)\sin\alpha}{M_{h}^{2}}\right)^{2}+\frac{\sin^{2}\beta}{M_{A}^{4}}\right]$
(4.1) $\displaystyle\approx$
$\displaystyle\frac{G_{F}^{2}\,m_{\mu}^{2}\,m_{\tau}^{7}\,\tau_{\tau}}{768\,\pi^{3}\,M_{A}^{4}}|\kappa_{\tau\mu}^{E}|^{2}\tan^{6}\beta\,.$
(4.2)
In the above, $\tau_{\tau}$ is the $\tau$ life time and the approximate result
has been obtained in the large $\tan\beta$ regime. For other Higgs-mediated
lepton flavour violating 3-body decays, $\tau\rightarrow e\mu\mu$,
$\tau\rightarrow 3e$ or $\mu\rightarrow 3e$, their corresponding branching
ratios can easily be obtained with the appropriate kinematic factors and the
flavour changing factor $\kappa$. While $\text{Br}(\tau\rightarrow e\mu\mu)$
can be as large as $\text{Br}(\tau\rightarrow 3\mu)$ when
$(Y^{\dagger}_{\nu}Y_{\nu})_{13}\sim O(1)$ (which is possible in the case of
an inverted hierarchical light neutrino spectrum), other flavour violating
decays with final state electrons such as $\mu\rightarrow 3e$ are considerably
suppressed due to the smallness of the Yukawa couplings.
* •
$B_{s}\to\ell_{i}\ell_{j}$
$B$ mesons can also have Higgs-mediated LFV decays, which are significantly
enhanced in the large $\tan\beta$ regime. The branching fraction is given by
$\displaystyle\text{Br}(B_{s}\to\ell_{i}\ell_{j})$ $\displaystyle=$
$\displaystyle\frac{G_{F}^{4}\,M^{4}_{W}}{8\,\pi^{5}}\,|V_{tb}^{*}V_{ts}|^{2}\,M_{B_{s}}^{5}\,f_{B_{s}}^{2}\,\tau_{B_{s}}\biggl{(}\frac{m_{b}}{m_{b}+m_{s}}\biggr{)}^{2}$
(4.3) $\displaystyle\times$
$\displaystyle\sqrt{\biggl{[}1-\frac{(m_{\ell_{i}}+m_{\ell_{j}})^{2}}{M_{B_{s}}^{2}}\biggr{]}\biggl{[}1-\frac{(m_{\ell_{i}}-m_{\ell_{j}})^{2}}{M_{B_{s}}^{2}}\biggr{]}}$
$\displaystyle\times$
$\displaystyle\Biggl{\\{}\biggl{(}1-\frac{(m_{\ell_{i}}+m_{\ell_{j}})^{2}}{M_{B_{s}}^{2}}\biggr{)}|c^{ij}_{S}|^{2}+\biggl{(}1-\frac{(m_{\ell_{i}}-m_{\ell_{j}})^{2}}{M_{B_{s}}^{2}}\biggr{)}|c_{P}^{ij}|^{2}\Biggr{\\}}\,,$
where $V_{ij}$ represents the Cabibbo-Kobayashi-Maskawa (CKM) matrix,
$M_{B_{s}}$ and $\tau_{B_{s}}$ respectively denote the mass and lifetime of
the $B_{s}$ meson, while $f_{B_{s}}=230\pm 30$ MeV [42] is the ${B_{s}}$ meson
decay constant and $c_{P}^{ij}$, $c_{S}^{ij}$ are the form factors. As an
example, the lepton flavour violating (double-penguin) $B_{s}\to\mu\tau$ decay
can be computed with the following form factors [38]:
$\displaystyle c_{S}^{\mu\tau}=c_{P}^{\mu\tau}$ $\displaystyle=$
$\displaystyle\frac{\sqrt{2}\,\pi^{2}}{G_{F}\,M^{2}_{W}}\frac{m_{\tau}\,\kappa_{bs}^{d}\,\kappa_{\tau\mu}^{E\ast}}{\cos^{4}\beta\,\bar{\lambda}^{t}_{bs}}\left[\frac{\sin^{2}(\alpha-\beta)}{M^{2}_{H}}+\frac{\cos^{2}(\alpha-\beta)}{M^{2}_{h}}+\frac{1}{M^{2}_{A}}\right]$
(4.4) $\displaystyle\approx$
$\displaystyle\frac{8\,\pi^{2}\,m_{\tau}\,m_{t}^{2}}{M^{2}_{W}}\frac{\epsilon_{Y}~{}\kappa_{\tau\mu}^{E}~{}\tan^{4}\beta}{\left[1+(\epsilon_{0}+\epsilon_{Y}Y^{2}_{t})\tan\beta\right]\left[1+\epsilon_{0}\tan\beta\right]}\frac{1}{M^{2}_{A}}\,.$
Here, $\kappa_{bs}^{d}$ represents the flavour mixing in the quark sector
while $\bar{\lambda}^{t}_{bs}=V^{*}_{tb}V_{ts}$. Similarly, $\epsilon_{0}$ and
$\epsilon_{Y}$ are the down type quark form factors mediated by gluino and
squark exchange diagrams. The final result was, once again, derived in the
large $\tan\beta$ regime. The branching fractions of other flavour violating
decays such as $\text{Br}(B_{d,s}\rightarrow\tau e)$, would receive identical
contribution from the Higgs penguins. Likewise, the
$\text{Br}(B_{d,s}\rightarrow\mu e)$ can be calculated using the appropriate
form factors and lepton masses; as expected, these will be suppressed when
compared to $\text{Br}(B_{d,s}\rightarrow\tau\mu)$.
* •
$\tau\to\mu P$
Similar to what occurred in the previous processes, virtual Higgs exchange
could also induce decays such as $\tau\to\mu P$, where $P$ denotes a neutral
pseudoscalar meson ($P=\pi,\eta,\eta^{\prime})$. In the large $\tan\beta$
limit, where the pseudoscalar Higgs couplings to down-type quarks are
enhanced, CP-odd Higgs boson exchanges provide the dominant contribution to
the $\tau\to\mu P$ decay. The coupling can be written as
$\displaystyle-i(\sqrt{2}\,G_{F})^{1/2}\tan\beta~{}A(\xi_{d}\,m_{d}\,\bar{d}\,d+\xi_{s}\,m_{s}\,\bar{s}\,s+\xi_{b}\,m_{b}\,\bar{b}\,b)+{\rm
h.c.}.$ (4.5)
Here, the parameters $\xi_{d},\,\xi_{s},\,\xi_{b}$ are of order
$\mathcal{O}(1)$. Since we are mostly interested in the Higgs-mediated
contributions, we estimate the amplitude of these processes in the limit when
both $\tau\rightarrow 3\mu$ and $\tau\to\mu P$ are indeed dominated by the
exchange of the scalar fields. Accordingly, and following [41], one can write
$\displaystyle\frac{\text{Br}(\tau\to\mu\eta)}{\text{Br}(\tau\to 3\mu)}$
$\displaystyle\simeq$ $\displaystyle
36\,\pi^{2}\left(\frac{f^{8}_{\eta}\,m^{2}_{\eta}}{m_{\mu}\,m^{2}_{\tau}}\right)^{2}(1-x_{\eta})^{2}\left[\xi_{s}+\frac{\xi_{b}}{3}\left(1+\sqrt{2}\,\frac{f^{0}_{\eta}}{f^{8}_{\eta}}\right)\right]^{2},$
(4.6)
$\displaystyle\frac{\text{Br}(\tau\to\mu\eta^{\prime})}{\text{Br}(\tau\to\mu\eta)}$
$\displaystyle\simeq$
$\displaystyle\frac{2}{9}\left(\frac{f^{0}_{\eta^{\prime}}}{f^{8}_{\eta}}\right)^{2}\frac{m^{4}_{\eta^{\prime}}}{m^{4}_{\eta}}\left(\frac{1-x_{\eta^{\prime}}}{1-x_{\eta}}\right)^{2}\left[\frac{1+\frac{3}{\sqrt{2}}\,\frac{f^{8}_{\eta^{\prime}}}{f^{0}_{\eta^{\prime}}}\left(\frac{\xi_{s}}{\xi_{b}}+\frac{1}{3}\right)}{\frac{\xi_{s}}{\xi_{b}}+\frac{1}{3}+\frac{\sqrt{2}}{3}\,\frac{f^{0}_{\eta}}{f^{8}_{\eta}}}\right]^{2},$
(4.7)
$\displaystyle\frac{\text{Br}(\tau\to\mu\pi)}{\text{Br}(\tau\to\mu\eta)}$
$\displaystyle\simeq$
$\displaystyle\frac{4}{3}\left(\frac{f_{\pi}}{f^{8}_{\eta}}\right)^{2}\,\frac{m^{4}_{\pi}}{m^{4}_{\eta}}~{}(1-x_{\eta})^{-2}\left[\frac{\frac{\xi_{d}}{\xi_{b}}\,\frac{1}{1+z}+\frac{1}{2}\,(1+\frac{\xi_{s}}{\xi_{b}})\frac{1-z}{1+z}}{\frac{\xi_{s}}{\xi_{b}}+\frac{1}{3}+\frac{\sqrt{2}}{3}\,\frac{f^{0}_{\eta}}{f^{8}_{\eta}}}\right]^{2}\,,$
(4.8)
where $z=m_{u}/m_{d}$, $m_{\pi},\ f_{\pi}$ are the pion mass and decay
constant, $m_{\eta,\eta^{\prime}}$ are the masses of $\eta,\ \eta^{\prime}$,
$x_{\eta,\eta^{\prime}}=m_{\eta,\eta^{\prime}}^{2}/m_{\pi}^{2}$, and
$f^{8}_{\eta,\eta^{\prime}}$ and $f^{0}_{\eta,\eta^{\prime}}$ are evaluated
from the corresponding matrix elements. As first discussed in [40], and taking
$\xi_{s},\xi_{b}\sim 1$ and fixing the other parameters as in [41], one finds
$\frac{\text{Br}(\tau\to\mu\eta)}{\text{Br}(\tau\to 3\mu)}\simeq 5$. The other
branching fractions such as $\text{Br}(\tau\to\mu\eta^{\prime},\mu\pi)$ are
considerably suppressed compared to $\text{Br}(\tau\to\mu\eta)$. While the
ratio $\frac{\text{Br}(\tau\to\mu\eta^{\prime})}{\text{Br}(\tau\to\mu\eta)}$
can be as large as $6\times 10^{-3}$,
$\frac{\text{Br}(\tau\to\mu\pi)}{\text{Br}(\tau\to\mu\eta)}$ would
approximately lie in the range $10^{-3}-4\times 10^{-3}$ [41]. Since all these
ratios are independent of $\kappa_{\tau\mu}^{E}$, the above quoted numbers can
also be applied to the present framework. However, an enhancement in the
$\text{Br}(\tau\to 3\mu)$, due to the large values of $\kappa_{\tau\mu}^{E}$,
would also imply sizeable values of $\text{Br}(\tau\to\mu\eta)$.
* •
$H_{k}\to\mu\tau\;(H_{k}=h,H,A)$
The branching ratios of flavour violating Higgs decays provide another
interesting probe of lepton flavour violation. Following [43], the branching
fraction $H_{k}\to\mu\tau$ (normalised to the flavour conserving one
$H_{k}\to\tau\tau$) can be cast as:
$\displaystyle{\text{Br}(H_{k}\to\mu\tau)}=\tan^{2}\beta~{}(|\kappa_{\tau\mu}^{E}|^{2})~{}C_{\Phi}~{}{\text{Br}(H_{k}\to\tau\tau)}\,,$
(4.9)
where we approximated $1/\cos^{2}\beta\simeq\tan^{2}\beta$. The coefficients
$C_{\Phi}$ are given by:
$\displaystyle
C_{h}=\left[\frac{\cos(\beta-\alpha)}{\sin\alpha}\right]^{2},~{}~{}~{}~{}C_{H}=\left[\frac{\sin(\beta-\alpha)}{\cos\alpha}\right]^{2},~{}~{}~{}~{}C_{A}=1.$
(4.10)
## 5 Results and Discussion
As discussed in Section 3, in the inverse supersymmetric seesaw, Higgs-
mediated contributions can lead to an enhancement of several LFV observables
by as much as two orders of magnitude, compared to what is expected in the
standard SUSY seesaw.
As expected from the analytical study of Section 4, $m_{A}$ and $\tan\beta$
are the most relevant parameters in the Higgs-mediated flavour violating
processes. To better illustrate this, in Fig. 3 we study the dependence of
Br($\tau\rightarrow 3\mu$) on the aforementioned parameters. We have assumed a
common value for the squark masses, $m_{\widetilde{q}}\sim\text{TeV}$, while
for left- and right-handed sleptons we take $m_{\widetilde{\ell}}\sim
400~{}\text{GeV}$ and $M_{\widetilde{\nu}^{c}}\sim 3~{}\text{TeV}$ for the
right handed sneutrinos. The contours correspond to different values of the
branching ratios (the purple region has already been experimentally excluded).
From this figure one can easily identify the regimes for $m_{A}$ and
$\tan\beta$ which are associated to values of the LFV observables within reach
of the present and future experiments.
Figure 3: Branching ratio of the process $\tau\rightarrow 3\mu$ as a function
of $m_{A}$ (GeV) and $\tan\beta$. From left to right, the contours correspond
to $\text{Br}(\tau\rightarrow 3\mu)=2.1\times 10^{-8}$, $10^{-9}$, $10^{-10}$,
$10^{-11}$. The purple region has already been experimentally excluded[44].
In what follows, we numerically evaluate some LFV observables. Concerning the
mSUGRA parameters (and instead of scanning over the parameter space), we have
selected a few benchmark points [45] that already take into account the most
recent LHC constraints [46]. We have also considered the case in which the GUT
scale universality conditions are relaxed for the Higgs sector, i.e. scenarios
of Non-Universal Higgs Masses (NUHM), as this allows to explore the impact of
a light CP-odd Higgs boson. In Table 1, we list the chosen points: CMSSM-A and
CMSSM-B respectively correspond to the 10.2.2 and 40.1.1 benchmark points in
[45], while NUHM-C is an example of a non-universal scenario.
Point | $\tan\beta$ | $m_{1/2}$ | $m_{0}$ | $m^{2}_{H_{U}}$ | $m^{2}_{H_{D}}$ | $A_{0}$ | $\mu$ | $m_{A}$
---|---|---|---|---|---|---|---|---
CMSSM-A | 10 | 550 | 225 | $(225)^{2}$ | $(225)^{2}$ | 0 | 690 | 782
CMSSM-B | 40 | 500 | 330 | $(330)^{2}$ | $(330)^{2}$ | -500 | 698 | 604
NUHM-C | 15 | 550 | 225 | $(652)^{2}$ | $-(570)^{2}$ | 0 | 478 | 150
Table 1: Benchmark points used in the numerical analysis (dimensionful
parameters in GeV). CMSSM-A and CMSSM-B correspond to 10.2.2 and 40.1.1
benchmark points of [45].
For each point considered, the low-energy SUSY parameters were obtained using
SuSpect [47]. In what concerns the evolution of the soft-breaking right-handed
sneutrino masses $m_{\tilde{\nu}^{c}}^{2}$, we have assumed that the latter
hardly run between the GUT scale and the low-energy one. The flavour-violating
charged slepton parameters (e.g. $(\Delta m_{\widetilde{L}}^{2})_{ij}$ or
$\xi$), were estimated at the leading order using Eq. (3.1). Concerning NUHM,
we use the same value of $\xi$ as for CMSSM-A. Here, we are particularly
interested to study the effect of light CP-odd Higgs boson and this naive
approximation will serve our purpose. Furthermore, we use the mass insertion
approximation, assuming that mixing between left and right chiral slepton
states are relatively small. In computing the branching fractions and the
flavour violating factor $\kappa^{E}_{ij}$ we have assumed (physical) right-
handed sneutrino masses $M_{\widetilde{\nu}^{c}}\approx 3$ TeV and
$\left(Y_{\nu}^{\dagger}Y_{\nu}\right)=0.7$, in agreement with low-energy
neutrino data as well as other low-energy constraints, which are particularly
relevant in the inverse seesaw case such as Non-Standard Neutrino Interactions
bounds [48]. Moreover, in our numerical analysis, we have fixed the trilinear
soft breaking parameter $A_{\nu}=-500$ GeV (at the SUSY scale).
We now proceed to present our results for the flavour violating observables
discussed in Section 4. In Table 2, we collect the values of the different
branching ratios, as obtained for the considered benchmark points of Table 1.
We have also presented the corresponding current experimental bounds and
future sensitivity.
LFV Process | Present Bound | Future Sensitivity | CMSSM-A | CMSSM-B | NUHM-C
---|---|---|---|---|---
$\tau\rightarrow\mu\mu\mu$ | $2.1\times 10^{-8}$[44] | $8.2\times 10^{-10}$ [52] | $1.4\times 10^{-15}$ | $3.9\times 10^{-11}$ | $8.0\times 10^{-12}$
$\tau^{-}\rightarrow e^{-}\mu^{+}\mu^{-}$ | $2.7\times 10^{-8}$[44] | $\sim 10^{-10}$ [52] | $1.4\times 10^{-15}$ | $3.4\times 10^{-11}$ | $8.0\times 10^{-12}$
$\tau\rightarrow eee$ | $2.7\times 10^{-8}$[44] | $2.3\times 10^{-10}$ [52] | $3.2\times 10^{-20}$ | $9.2\times 10^{-16}$ | $1.9\times 10^{-16}$
$\mu\rightarrow eee$ | $1.0\times 10^{-12}$[1] | | $6.3\times 10^{-22}$ | $1.5\times 10^{-17}$ | $3.7\times 10^{-18}$
$\tau\rightarrow\mu\eta$ | $2.3\times 10^{-8}$[49] | $\sim 10^{-10}$ [52] | $8.0\times 10^{-15}$ | $3.3\times 10^{-10}$ | $4.6\times 10^{-11}$
$\tau\rightarrow\mu\eta^{\prime}$ | $3.8\times 10^{-8}$[49] | $\sim 10^{-10}$ [52] | $4.3\times 10^{-16}$ | $1.1\times 10^{-10}$ | $3.1\times 10^{-12}$
$\tau\rightarrow\mu\pi^{0}$ | $2.2\times 10^{-8}$[49] | $\sim 10^{-10}$ [52] | $1.8\times 10^{-17}$ | $8.5\times 10^{-13}$ | $1.0\times 10^{-13}$
$B^{0}_{d}\rightarrow\mu\tau$ | $2.2\times 10^{-5}$[50] | | $2.7\times 10^{-15}$ | $8.5\times 10^{-10}$ | $2.7\times 10^{-11}$
$B^{0}_{d}\rightarrow e\mu$ | $6.4\times 10^{-8}$[51] | $1.6\times 10^{-8}$[53] | $1.2\times 10^{-17}$ | $3.1\times 10^{-12}$ | $1.2\times 10^{-13}$
$B^{0}_{s}\rightarrow\mu\tau$ | | | $7.7\times 10^{-14}$ | $2.5\times 10^{-8}$ | $7.8\times 10^{-10}$
$B^{0}_{s}\rightarrow e\mu$ | $2.0\times 10^{-7}$[51] | $6.5\times 10^{-8}$[53] | $3.4\times 10^{-16}$ | $8.9\times 10^{-11}$ | $3.4\times 10^{-12}$
$h\rightarrow\mu\tau$ | | | $1.3\times 10^{-8}$ | $2.6\times 10^{-7}$ | $2.3\times 10^{-6}$
$A,H\rightarrow\mu\tau$ | | | $3.4\times 10^{-6}$ | $1.3\times 10^{-4}$ | $5.0\times 10^{-6}$
Table 2: Higgs-mediated contributions to the branching ratios of several
lepton flavour violating processes, for the different benchmark points of
Table 1. We also present the current experimental bounds and future
sensitivities for the LFV observables.
From Table 2, one can verify that from an experimental point of view, the most
promising channel in the supersymmetric inverse seesaw is
$\tau\rightarrow\mu\eta$ which could be tested at the next generation of $B$
factories. The $B^{0}_{d,s}\rightarrow\mu\tau$ decay might also be
interesting, but little conclusions can be drawn due to lack of information
concerning the future sensitivities.
It is important to stress that the numerical results summarised in Table 2
correspond to considering only Higgs-mediated contributions. In the low
$\tan\beta$ regime, photon- and $Z$-penguin diagrams may induce comparable or
even larger contributions to the observables, and potentially enhance the
branching fractions. Thus, the results for small $\tan\beta$ should be
interpreted as conservative estimates, representing only partial
contributions. For large $\tan\beta$ values, Higgs penguins do indeed provide
the leading contributions. Comparing our results with those obtained for a
type I SUSY seesaw at high scales (or even with a TeV scale SUSY seesaw), we
find a large enhancement of the branching fractions in the inverse seesaw
framework.
Another interesting property of the Higgs-mediated processes is that the
corresponding amplitude strongly depends on the chirality of the heaviest
lepton (be it the decaying lepton, or the heaviest lepton produced in $B$
decays). Considering the decays of a left-handed lepton
$\ell^{i}_{L}\rightarrow\ell^{j}_{R}X$, one finds that the corresponding
branching ratios would be suppressed by a factor
$\frac{m_{\ell^{j}}^{2}}{m_{\ell^{i}}^{2}}$ compared to those of the right-
handed lepton $\ell^{i}_{R}\rightarrow\ell^{j}_{L}X$. This can induce an
asymmetry that potentially allows to identify if Higgs mediation is the
dominant contribution to the LFV observables. Furthermore this asymmetry would
be more pronounced in the inverse-seesaw framework.
Due to its strong enhancement of the Higgs-penguin contributions, if realised
in Nature, the inverse seesaw offers a unique framework to test Higgs effects
in LFV processes. In fact, and as discussed in [32], if photon penguins
provide the dominant contribution to both $\text{Br}(\tau\rightarrow 3\mu)$
and $\text{Br}(\tau\rightarrow\mu\gamma)$, then the latter observables are
strongly correlated, $\frac{\text{Br}(\tau\rightarrow
3\mu)}{\text{Br}(\tau\rightarrow\mu\gamma)}\sim 0.003$ (see [32]). On the
other hand, if the dominant contribution to the three-body decays arises from
Higgs penguins, the correlation no longer holds, and the latter ratio can be
significantly enhanced. This would be the case of the present framework.
## 6 Conclusions
If observed, charged lepton flavour violation clearly signals the presence of
new physics. In this work, we have studied Higgs-mediated LFV processes in the
framework of the supersymmetric inverse seesaw. TeV scale right-handed
neutrinos (and hence light right-handed sneutrinos) offer the possibility to
enhance the Higgs-mediated contributions to several LFV processes. As shown in
this work, in the inverse SUSY seesaw, LFV branching ratios can be enhanced by
as much as two orders of magnitude when compared to the standard (type I) SUSY
seesaw.
### Acknowledgements
The authors are thankful to A. Vicente for many enlightening discussions. D.D.
acknowledges financial support from the CNRS. This work has been partly done
under the ANR project CPV-LFV-LHC NT09-508531.
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|
arxiv-papers
| 2011-11-24T19:06:55 |
2024-09-04T02:49:24.638009
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Asmaa Abada, Debottam Das, C\\'edric Weiland",
"submitter": "C\\'edric Weiland",
"url": "https://arxiv.org/abs/1111.5836"
}
|
1111.5994
|
# Possibility and Impossibility of the Entropy Balance in Lattice Boltzmann
Collisions
Alexander N. Gorban and Dave Packwood Department of Mathematics, University
of Leicester, United Kingdom
###### Abstract
We demonstrate that in the space of distributions operated on by lattice
Boltzmann methods that there exists a vicinity of the equilibrium where
collisions with entropy balance are possible and, at the same time, there
exist an area of nonequilibrium distributions where such collisions are
impossible. We calculate and graphically represent these areas for some simple
entropic equilibria using single relaxation time models. Therefore it is shown
that the definition of an entropic LBM is incomplete without a strategy to
deal with certain highly nonequilibrium states. Such strategies should be
explicitly stated as they may result in the production of additional entropy.
## I Introduction
Lattice Boltzmann schemes are a type of discrete algorithm which can be used
to simulate fluid dynamics and more Succi ; Benzi . Although such a method can
be derived as a discretization of the fully continuous Boltzmann equation,
some thermodynamics properties may be lost in this process. The Entropic
lattice Boltzmann method (ELBM) was invented first in 1998 as a tool for the
construction of single relaxation time lattice Boltzman models which respects
a $H$-theorem Htheorem ; SucciRevModPhys . For this purpose, instead of the
mirror image with a local equilibrium as the reflection center, the entropic
involution was proposed, which preserves the entropy value. Later, it was
called the Karlin-Succi involution gorban06 .
Nevertheless, controlling the proper entropy balance remained until recently a
challenging problem for many lattice Boltzmann models Nonexist . Some
discussions of modern ELBM implementations and results were published recently
KarlinSucciComment .
The distribution functions at the centre of lattice Boltzmann methods are
often referred to and understood as particle densities. Of course for such an
interpretation to be meaningful the distribution function should be strictly
positive. Despite this some lattice Boltzmann implementations may, as a
numerical scheme, tolerate negative population values. An ELBM usually
involves an evaluation of a Boltzmann type entropy function, which does not
exist for negative populations, hence such an ELBM cannot ever tolerate a
negative population value. Due to this there are population values for which
an entropic involution cannot be performed. A complete definition of an ELBM
must include a strategy for what to do in such a situation. The choice of such
a strategy should be explicitly given in any definition of an ELBM as it may
have side-effects with modification of dissipation which should be understood
separately from the influence of the proper entropy balance.
In this paper we study the regions in the spaces of distributions
(populations) where collisions with entropy preservation are possible (near
the equilibrium) and where they are impossible (sufficiently far from the
equilibrium) and demonstrate that both such areas always exist apart some
trivial degenerated cases.
## II Single Relaxation Time LB Schemes
For fluids, LB systems can be derived as a discretization of the Boltzmann
Equation
${\partial}_{t}f+\mathbf{v}\cdot{\partial}_{\mathbf{x}}f=Q(f)$ (1)
where $f\equiv f(\mathbf{x},\mathbf{v},t)$ is a one particle distribution
function over space, velocity space and time and $Q(f)$ represents the
interaction between particles, sometimes called a collision operation. A
particular example of the interaction $Q(f)$ is the Bhatnagar-Gross-Krook
equation
$Q(f)=-\frac{1}{\tau}(f-f^{\mathrm{eq}}).$ (2)
The BGK operation represents a relaxation towards the local equilibrium
$f^{\mathrm{eq}}$ with rate $1/\tau$. The distribution $f^{\mathrm{eq}}$ is
given by the Maxwell Boltzmann distribution,
$f^{\mathrm{eq}}=\frac{\rho}{(2\pi
T)^{D/2}}\exp\left(\frac{-(\mathbf{v}-\mathbf{u})^{2}}{2T}\right).$ (3)
The macroscopic quantities are available as integrals over velocity space of
the distribution function,
$\rho=\int
f\;\mathrm{d}\mathbf{v},\;\rho\mathbf{u}=\int\mathbf{v}f\;\mathrm{d}\mathbf{v},\;\rho\mathbf{u}^{2}+\rho
T=\int\mathbf{v}^{2}f\;\mathrm{d}\mathbf{v}.$
A discrete approximation to these integrals is the first ingredient to
discretize this system. The scalar field of the population function (over
space, vector space and time) becomes a sequence of vector fields (over space)
in time $f_{i}(\mathbf{x},n_{t}\;\epsilon),n_{t}\in\mathbb{Z}$, where the
elements of the vector each correspond with an element of the quadrature.
Explicitly the macroscopic moments are given by,
$\rho=\sum_{i=1}^{n}f_{i},\;\rho\mathbf{u}=\sum_{i=1}^{n}\mathbf{v}_{i}f_{i},\;\rho\mathbf{u}^{2}+\rho
T=\sum_{i=1}^{n}\mathbf{v}_{i}^{2}f_{i}.$
The complete discrete scheme is given by
$f_{i}(\mathbf{x}+\epsilon\mathbf{v}_{i},t+\epsilon)=f_{i}(\mathbf{x},t)+\omega(f^{\mathrm{eq}}_{i}(\mathbf{x},t)-f_{i}(\mathbf{x},t))$
(4)
where $\epsilon$ is the time step. For this system a discrete equilibrium must
be used. The choice of the velocity set
$\\{\mathbf{v}_{1},\ldots,\mathbf{v}_{n}\\}$ and the discrete equilibrium
distribution $f^{\mathrm{eq}}_{i}$ should provide the best approximation of
the transport equations for the moments by the discrete scheme (4).
## III ELBM
In the continuous case the Maxwellian distribution maximizes entropy, as
measured by the Boltzmann $H$ function, and therefore also has zero entropy
production. In the context of lattice Boltzmann methods a discrete form of the
$H$-theorem has been suggested as a way to introduce thermodynamic control to
the system Htheorem ; Boghosian .
A variation on the LBGK is the ELBGK ELBM . In this family of methods, the
equilibria are defined as the conditional entropy maximizers under given
values of macroscopic variables (entropic equilibria). The entropies have been
constructed in a lattice dependent fashion in LatticeEntropies . A slightly
different notation is used for the lattice Boltzmann algorithm,
$f_{i}(\mathbf{x}+\epsilon\mathbf{v}_{i},t+\epsilon)=f_{i}(\mathbf{x},t)+\alpha\beta(f^{\mathrm{eq}}_{i}(\mathbf{x},t)-f_{i}(\mathbf{x},t)).$
(5)
The single parameter $\omega$ is replaced by a composite parameter
$\alpha\beta$. In this case $\beta$ controls the viscosity and $\alpha$ is
varied to ensure a constant entropy condition according to the discrete
$H$-theorem. With knowledge of the entropy function $S$, $\alpha$ is found as
the non-trivial root of the equation
$S(\mathbf{f})=S(\mathbf{f}+\alpha(\mathbf{f}^{\ast}-\mathbf{f})).$ (6)
The trivial root $\alpha=0$ returns the entropy value of the original
populations. ELBGK then finds the non-trivial $\alpha$ such that (6) holds.
This version of the BGK collision one calls entropic BGK (or EBGK) collision.
A solution of (6) must be found at every time step and lattice site. The EBGK
collision obviously respects the Second Law (if $\beta\leq 1$), and simple
analysis of entropy dissipation gives the proper evaluation of viscosity.
In general the entropy function is based upon the lattice. For example, in the
case of the simple one dimensional lattice with velocities
$\mathbf{v}=(-c,0,c)$ and corresponding populations
$\mathbf{f}=(f_{-},f_{0},f_{+})$ an explicit Boltzmann style entropy function
is known LatticeEntropies :
$S(\mathbf{f})=-f_{-}\log(f_{-})-f_{0}\log(f_{0}/4)-f_{+}\log(f_{+}).$ (7)
## IV Regions of Existence and Non-existence of Entropic Involution
Let us study the entropic involution in the distribution simplex $\Sigma$
given by $\sum f_{i}=const>0$, $f_{i}\geq 0$.
Figure 1: The simplex $\Sigma$ is given by the white background. (1)
Populations relax through the equilibrium given by the single point to an
equal entropy point, if possible. The boundary of this possibility is given.
(2) The regions $A$ (the entropic involution is possible) and $B$ (the
involution is impossible) as subsets of the simplex divided by this boundary
are presented. Figure 2: The simplex $\Sigma$ is given by the white
background. (1) Populations relax through the their corresponding equilibrium
point along the line given by constant $u$ to an equal entropy point, if
possible. The boundary of this possibility is given. (2) The regions $A$ and
$B$ separated by this boundary are presented.
Let us prove that under very natural assumptions about some properties of the
entropy that the simplex of distributions can be split into two subsets $A$
and $B$: in the set $A$ the entropic involution exists, and for distributions
from the set $B$ equation (6) has no non-trivial solutions. Both sets $A$ and
$B$ have non-empty interior (apart of a trivial symmetric degenerated case).
Let the entropy $S$ be a strictly concave continuous function in the
distribution simplex $\Sigma$. We assume also that $S$ is twice
differentiable, the Hessian of $S$, $\partial^{2}S/\partial f_{i}\partial
f_{j}$, is negative definite in the interior of the simplex, $\Sigma_{+}$,
where $\sum f_{i}=const$, $f_{i}>0$ and the global maximizer of $S$, the
equilibrium, belongs to the interior of the simplex.
For example, the relative Boltzmann entropy, $S=-\sum
f_{i}(\ln(f_{i}/W_{i})-1)$, $W_{i}>0$, satisfies these conditions, because
$f\ln f\to 0$ when $f\to 0$ and $\partial^{2}S/\partial f_{i}\partial
f_{j}=-\delta_{ij}/f_{i}$, whereas the relative Burg entropy $S=\sum
W_{i}(\ln(f_{i}/W_{i}))$ does not satisfy these conditions because it does not
exist on the border of the simplex.
Macroscopic variables are linear functions of $\mathbf{f}$. The sets with
given values of the macroscopic variables in the simplex $\Sigma$ are
polyhedra, intersections of $\Sigma$ with linear manifolds with the given
values of moments. We assume that in any such a polyhedron the entropy
achieves its (conditionally) global maximum at an internal point. This
assumption holds for the Boltzmann relative entropy because of the logarithmic
singularity of the “chemical potentials” $\mu_{i}=\ln(f_{i}/W_{i})$ on the
border of positivity. These maximizers are equilibria. If $\mathbf{f}$ is
sufficiently close to a positive equilibrium then, due to the implicit
function theorem, the nontrivial solution to equation (6) exists and it gives
$\alpha=2+o(\mathbf{f}-\mathbf{f}^{*})$. The value $\alpha=2$ corresponds to
the mirror image, the small term $o(\mathbf{f}-\mathbf{f}^{*})$ gives the
corrections to the value $\alpha=2$. Therefore, in some vicinity of the
equilibrium the entropic involution exists.
To prove the existence of the area where entropic involution is impossible,
let us consider one polyhedron with given values of the macroscopic variables
and a positive equilibrium. The local minima of the entropy in this polyhedron
are situated at the vertices. At least one of them is a global minimum. Let
this vertex be $\mathbf{f}^{\mathbf{v}}$. Let us draw a straight line $l$
through points $\mathbf{f}^{\mathbf{v}}$ and $\mathbf{f}^{*}$. The
intersection $l\cap\Sigma$ is an interval and $S$ achieves its global minimum
on this interval at the point $\mathbf{f}^{\mathbf{v}}$. If the dimension of
the polyhedron is more than one then the opposite end of this interval is not
even a local minimum of $S$ in the polyhedron and the entropic involution does
not exists for $\mathbf{f}^{\mathbf{v}}$ and some vicinity around it.
A special degeneration is possible when the polyhedra are one-dimensional,
i.e. intervals, and the values of the entropy at both ends of each interval
coincide. For example, for two-dimensional distributions, $f_{+},f_{-}$, the
entropy $S==f_{+}\ln f_{+}-f_{-}\ln f_{-}$ and the macroscopic variable
$\rho=f_{+}+f_{-}$. Apart from such symmetric one-dimensional cases there
exists an area near the maximally non-equilibrium vertex
$\mathbf{f}^{\mathbf{v}}$ where the entropic involution cannot be defined.
Such an area may also exist near some other vertices, where local entropy
minima are reached.
For the Burg entropy, the entropic involution is always possible Boghosian
because it tends to $-\infty$ at the border of positivity. The same is true,
for the relative entropy of the form $S=-\beta^{-1}\sum
W_{i}((W_{i}/f_{i})^{\beta}-1)$ that tends to the Burg entropy when $\beta\to
0$ GoGoJudge2011 . This negative brunch of the relative Tsallis entropy is
less known. The standard Tsallis entropy Tsallis is finite at the border of
positivity, hence, collisions with entropy preservation are not always
possible for it.
We now demonstrate the population function values where the involution cannot
be performed for some simple examples. We use the standard 1-D lattice
described in Section III with the discrete equilibrium given in Eq 7. We begin
with an LBM with only one conserved moment in collision, namely density. The
equilibrium is
$f_{-}^{*}=\frac{\rho}{6},\;\;f_{0}^{*}=\frac{2\rho}{3},\;\;f_{+}^{*}=\frac{\rho}{6}.$
In Fig. 1, the simplex $\Sigma$ of positive populations with a fixed density
$\rho=1$ is the triangle given by the intersection of three half-planes,
$f^{+}>0,\,f^{-}>0$, and $1-f^{+}-f^{-}>0$. Within that region we plot several
entropy level contours $S(\mathbf{f})=c$ and the unique equilibrium point. The
region is divided into the parts where the entropic involution is possible
(around the equilibrium) and where it is impossible.
A more common use of lattice Boltzmann involves a second fixed moment,
momentum. The entropic equilibria used by the ELBGK are available explicitly
as the maximum of the entropy function (7),
$f_{\mp}^{*}=\frac{\rho}{6}(\mp
3u-1+2\sqrt{1\\!+\\!3u^{2}}),\;f_{0}^{*}=\frac{2\rho}{3}(2-\sqrt{1\\!+\\!3u^{2}}).$
In this case the dimension of the equilibrium is one greater. In Fig. 2 all
relaxation occurs parallel to the lines of constant $u$. The region where
entropic involution is possible is again given.
In each experiment the region is discretized into many individual points. For
each point a value for $\alpha$ is attempted to be found. The method used is
simply to begin with a guess of $\alpha=1$ and then add increments of
$10^{-3}$ until a solution of Eq. 6 occurs, or the edge of the positivity
domain is reached. This method would be inappropriate to use in a usual ELBM,
due to the very large computational cost, but it is very robust and hence
useful for this experiment with many higly non-equilibrium distributions.
Another approach (with the same result) implies calculation of the entropic
involution for all the boundary points where it exists. In this method we draw
a straight line $l$ through a boundary point $\mathbf{f}$ and the equilibrium
and find the intersection $l\cap\Sigma$ which consists of all points on $l$
with non-negative coordinates. One end of this interval is $\mathbf{f}$,
another end is also a boundary point, $\mathbf{f}^{\prime}$. The entropic
involution for $\mathbf{f}$ exits if and only if $S(\mathbf{f})^{\prime}\leq
S(\mathbf{f})$. After we check this inequality, we can solve Eq. (6). The
images of these involutions form the border that separates sets $A$ and $B$
(see Figs).
## V Conclusion
The entropic involution is not always possible to perform. We have
demonstrated that apart some special one-dimensional spaces of distributions
with additional symmetry there exist domains where collisions with the
preservation of entropy are not possible. We illustrated this statement by
some simple and well known examples of ELBGK systems for which we directly
calculated the areas where entropic collisions exist and where they do not
exist.
Such phenomena should be observable in all ELBM schemes with the classical
entropies: there exists a vicinity of the equilibrium where the entropic
involution is possible but for some areas of non-equilibrium distributions
there exists no non-trivial root of equation (6). A collision which preserves
entropy does not exist for this area. Therefore, for the regimes close to
equilibrium (the vicinities $A$ of equilibria, Figs 1,2), ELBM schemes
guaranty the precise balance of the entropy and for more nonequilibrium
regimes, when at some sites the distribution belonges to sets $B$, ELBM
schemes work as limiters Limiters . with additional dissipation. It is
necessary for any complete definition of an ELBM algorithm to prescribe what
to do when the involution is not possible. A reasonable choice would be to
over-relax the maximum amount possible while maintaining positive population
values. Such a technique is independently in use as a stabilizer for lattice
Boltzmann schemes, sometimes called the ‘positivity limiter’ BGJ ; Limiters ;
Li ; Servan ; Tosi . An effect of this operation is a local increase in
viscosity/entropy production. Hence, if an ELBM were to apply such a scheme it
would necessarily break the proper entropy balance. In this sense, ELBM
belongs to a large family of add-ons that regularise LBM by the management of
the addtional dissipation Add-ons .
## References
* (1) S. Ansumali and I. V. Karlin, Phys. Rev. E, 62 (6), 7999–8003 (2000).
* (2) R. Benzi, S. Succi, and M. Vergassola, Phys. Reports, 222, 145–197 (1992).
* (3) B. M. Boghosian, J. Yepez, P. V. Coveney, A. Wagner, Proc. R. Soc. Lond. A, 457, 717–766 (2001).
* (4) R. A. Brownlee, A. N. Gorban, and J. Levesley, Phys. Rev. E, 75, 036711 (2007).
* (5) R. A. Brownlee, A. N. Gorban, and J. Levesley, Physica A, 387 (2-3), 385–406 (2008).
* (6) R. A. Brownlee, J. Levesley, D. Packwood, A.N. Gorban, arXiv:1110.0270 [physics.comp-ph].
* (7) A. N. Gorban, in _Model Reduction and Coarse-Graining Approaches for Multiscale Phenomena_ , (Springer, Berlin-Heidelberg, New York 2006), 117–176; arXiv:cond-mat/0602024 [cond-mat.stat-mech].
* (8) A. N. Gorban, P. A. Gorban, G. Judge, arXiv:1003.1377 [physics.data-an].
* (9) I. V. Karlin, A. Ferrante, and H. C. Öttinger, Europhys. Lett. 47, 182–188 (1999).
* (10) I. V. Karlin, A. N. Gorban, S. Succi, and V. Boffi, Phys. Rev. Lett., 81, 6–9 (1998).
* (11) I. V. Karlin, S. Succi, arXiv:1107.3025 [cond-mat.stat-mech]
* (12) Y. Li, R. Shock, R. Zhang, H. Chen. J. Fluid Mech., 519, 273–300 (2004)
* (13) B. Servan-Camas, FT-C. Tsai. J Comput Physics 228 (1), 236–256 (2009).
* (14) S. Succi, _The lattice Boltzmann equation for fluid dynamics and beyond_ (Oxford University Press, New York 2001).
* (15) S. Succi, I. V. Karlin, and H. Chen, Rev. Mod. Phys. 74, 1203 (2002).
* (16) F. Tosi, S. Ubertini, S. Succi, H. Chen, I.V. Karlin, Math Comput Simulation, 72, 227–231 (2006).
* (17) C. Tsallis, J. Stat. Phys., 52, 479–487 (1988).
* (18) W.-A. Yong and L.-S. Luo, Phys. Rev. E, 67, 051105 (2003).
|
arxiv-papers
| 2011-11-25T14:18:54 |
2024-09-04T02:49:24.652497
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "A. N. Gorban, D. Packwood",
"submitter": "Alexander Gorban",
"url": "https://arxiv.org/abs/1111.5994"
}
|
1111.6213
|
# version1.4
Study of triangular flow $v_{3}$ in Au+Au and Cu+Cu collisions with a
multiphase transport model
Kai Xiao Institute of Particle Physics, Central China Normal University,
Wuhan, Hubei, 430079, China The Key Laboratory of Quark and Lepton Physics
(Central China Normal University), Ministry of Education, Wuhan, Hubei,
430079, China Na Li nli@mail.hust.edu.cn Department of Physics, Huazhong
University of Science and Technology, Wuhan 430074, China Shusu Shi
sss@iopp.ccnu.edu.cn Institute of Particle Physics, Central China Normal
University, Wuhan, Hubei, 430079, China The Key Laboratory of Quark and
Lepton Physics (Central China Normal University), Ministry of Education,
Wuhan, Hubei, 430079, China Feng Liu Institute of Particle Physics, Central
China Normal University, Wuhan, Hubei, 430079, China The Key Laboratory of
Quark and Lepton Physics (Central China Normal University), Ministry of
Education, Wuhan, Hubei, 430079, China
###### Abstract
We studied the relation between the initial geometry anisotropy and the
anisotropic flow in a multiphase transport model (AMPT) for both Au+Au and
Cu+Cu collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV. It is found that unlike the
elliptic flow $v_{2}$, little centrality dependence of the triangular flow
$v_{3}$ is observed. After removing the initial geometry effect,
$v_{3}/\varepsilon_{3}$ increases with the transverse particle density, which
is similar to $v_{2}/\varepsilon_{2}$. The transverse momentum ($p_{T}$)
dependence of $v_{3}$ from identified particles is qualitatively similar to
the $p_{T}$ dependence of $v_{2}$.
###### pacs:
25.75.Ld, 25.75.Dw
## I Introduction
A novel state of matter called quark-gluon plasma (QGP) composed by deconfined
partons is believed to be created experimentally in heavy ion collisions at
RHIC QGP . The discovery of large elliptic flow indicates that the partonic
collectivity is built up during the collisions, and the number-of-quark
scaling suggests that the partonic degrees of freedom are active nqscaling .
The anisotropic flow is usually described by a Fourier decomposition of the
azimuthal distribution with respect to the reaction plane ArtPRC . The second
harmonic coefficient, $v_{2}$, so called elliptic flow, has the biggest
magnitude at high energy collisions largeflow . It is believed that the
observed anisotropy in the momentum space is caused by the anisotropy in the
coordinate space in the initial condition. Lots of attention has been put on
the relation between $v_{2}$ and spatial eccentricity to see the hydrodynamics
behavior of the created system v2ecc ; PRCrun4 ; cucu_STAR .
Recent studies show that the event-by-event fluctuation of the initial
geometry v3ecc may play an important role in the study of collective flow.
The triangular shape in the initial geometry will be transferred to the
momentum space as the system expands, and finally leads to the none zero value
of the third harmonic coefficient, $v_{3}$. It is found that the triangular
flow $v_{3}$ is responsible for the ridge and shoulder structures and the
broad away-side of two-particle azimuthal correlation v3ridge . Besides, it is
also considered to be a good probe to study the viscous hydrodynamics behavior
of the colliding system v3hydro .
Lots of properties of the triangular flow $v_{3}$ have been studied in
hydrodynamic and transport models v3hydro ; v3ampt . However, since the
reaction plane can not be directly measured in the experiment, those
anisotropic parameters can not be directly obtained. It is found that
different methods may cause up to $20\%$ discrepancy on $v_{2}$ Artv2Review ,
thus it should also be carefully evaluated for the $v_{3}$ study. Besides,
$v_{3}$ is directly related with the initial fluctuation, it is interesting to
see its system size dependence.
In this paper, we will study the triangular flow $v_{3}$ in both Au+Au and
Cu+Cu collisions in a multiphase transport model (AMPT) AMPT . The relation
between $v_{3}$ and $\varepsilon_{3}$ is studied as a function of number of
participants and transverse momentum. The paper is organized as follows: In
Sec. II, the observables and technical methods are introduced. A brief
description of AMPT model is given in Sec. III. The results and discussions
are presented in Sec. IV. Finally, a summary is given in Sec. V.
## II OBSERVABLES
In a non-central collisions, the overlap region of two nuclei is an almond
shape. Since the position of nucleons may fluctuate event by event, as
discussed in Ref epart2 ; definition ; partecc , those initial geometric
irregularities of the colliding system can be described by $\varepsilon_{n}$:
$\varepsilon_{n}=\frac{{\sqrt{\left\langle{r^{2}\cos(n\varphi)}\right\rangle^{2}+\left\langle{r^{2}\sin(n\varphi)}\right\rangle^{2}}}}{{\left\langle{r^{2}}\right\rangle}},$
(1)
where $r$ and $\varphi$ are the polar coordinate position of participating
nucleons and $\langle\cdots\rangle$ is the average over all the participants
in an event. $n$ refers to the $n$-th harmonic, i.e., $\varepsilon_{2}$
describes the elliptic shape and $\varepsilon_{3}$ describes the triangular
shape.
As the system evolves, the anisotropy in the coordinate space is transferred
to the anisotropy in the momentum space due to the pressure gradient. The
particle distribution then can be written as
$\frac{dN}{d\phi}\propto 1+2\sum_{n=1}v_{n}\cos[n(\phi-\Psi_{n})],$ (2)
where $\phi$ is the azimuthal angle, and $\Psi_{n}$ is the $n$-th event plane
angle reconstructed by the final state particles:__
$\Psi_{n}=\frac{1}{n}\left[\tan^{-1}\frac{\sum\limits_{i}\sin(n\phi_{i})}{\sum\limits_{i}\cos(n\phi_{i})}\right].$
(3)
The observed anisotropic flow is defined as the $n$-th Flourier coefficient
$v_{n}$:
$v^{\rm obs}_{n}=\left\langle\cos[n(\phi-\Psi_{n})]\right\rangle.$ (4)
Here $\langle\cdots\rangle$ is taking the average over all the particles in
the sample.
This is the so-called event plane method of calculating $v_{n}$. The
reconstructed event plane fluctuates around the reaction plane. The observed
signals need to be revised by the corresponding resolution ArtPRC :
$v_{n}=\frac{v^{\rm obs}_{n}}{\mathscr{R}_{n}}.$ (5)
Due to the finite multiplicity of final state particles, the resolution
$\mathscr{R}_{n}=\langle\cos[n(\Psi_{n}-\Psi_{\rm nR})]\rangle$ (6)
is usually smaller than 1. $\Psi_{\rm nR}$ represents the nth real event plane
angle.
## III AMPT Model
There are four main components in AMPT model: the initial conditions, parton
interactions, hadronization and hadron interactions. The initial conditions
are obtained from the HIJING model HIJING , which includes the spatial and
momentum distributions of minijet partons from hard processes and strings from
soft processes. The time evolution of partons is then treated according to the
ZPC ZPC parton cascade model. After partons stop interacting, a combined
coalescence and string fragmentation model are used for the hadronization of
partons. The scattering among the resulting hadrons is described by a
relativistic transport (ART) model ART which includes baryon-baryon, baryon-
meson and meson-meson elastic and inelastic scattering.
In our study, we analyzed the events from AMPT with the parton cross section
equals to 3 mb and 10 mb. As all the conclusions are independent on the parton
cross section, only the results from 3 mb are shown in this paper. There are
about 8 million events in $\mathrm{Au}+\mathrm{Au}$ collisions and 19 million
events in $\mathrm{Cu}+\mathrm{Cu}$ collisions at $\sqrt{s_{{}_{NN}}}$= 200
used. The string melting AMPT version is used since the previous study shows
that the string melting AMPT version agrees with the experimental results
better AMPT . The centrality is defined by the impact parameter.
## IV Results and Discussions
Figure 1: The second and third harmonic event plane resolution calculated by
the particles with pseudo-rapidty region of $|\eta|>2$ as a function of
centrality in both Au+Au and Cu+Cu collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV
in AMPT model. Figure 2: $v_{n}$ as a function of centrality in both Au+Au and
Cu+Cu collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV in AMPT model.
In order to be comparable with the experimental data, the event plane method
is used to calculate $v_{n}$. The procedure is slightly different between ours
and Ref. v3ampt , in which the event plane was reconstructed by initial
partons. Charged particles with $p_{T}\leq 2$ GeV/$c$, $|\eta|>2$ are chosen
to reconstruct the event plane according to the Eq. 3. The particles used for
the $v_{n}$ measurements are within the $|\eta|<1$. The $\eta$ gap used here
is to reduce the auto-correlation between the particles used to reconstruct
the event plane and the particles of interest. In the following, the observed
$v_{n}$ are all corrected by the corresponding resolution. Fig. 1 shows the
resolution of $v_{2}$ and $v_{3}$ in both Au+Au and Cu+Cu collisions. The
resolution of $v_{2}$ shows a peak in mid-central collisions which is
consistent with the experimental result PRCrun2 . This is because the
resolution of $v_{2}$ is affected by both of the $v_{2}$ signal and the
multiplicity used to reconstruct the event plane. While the resolution of
$v_{3}$ only depends on the multiplicity, and keeps decreasing as the
multiplicity drops.
In Fig. 2, $v_{2}$ and $v_{3}$ are shown as functions of centrality in both
Au+Au and Cu+Cu collisions. We can see that $v_{2}$ shows strong centrality
dependence since it is mainly coming from the elliptic anisotropy in the
initial geometry. Unlike $v_{2}$, the dependence of $v_{3}$ on centrality and
system size are much smaller. The trend of $v_{3}$ observed is the same as
that in Ref v3ampt . However, the event plane angle $\Psi_{n}$ in Ref v3ampt
is obtained from the initial parton distribution, which is not observed in the
experiment, while in our study it is from the final state particle
distribution. The results indicate that the triangular flow is less sensitive
to the centrality and system size compared with the elliptic flow. It could be
understood as a result of combined effects from initial geometrical
fluctuation and collective dynamics which requires the size of bulk to
interact among themselves.
Figure 3: (Color online) $v_{2}$ as a function of $\varepsilon_{2}$ in Au+Au
collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV in AMPT model. The black points are
the average $v_{2}$ in the selected $\varepsilon_{2}$ bin. The pad in the
right down corner is the average of $v_{2}$ with smaller scale. Figure 4:
(Color online) $v_{3}$ as a function of $\varepsilon_{3}$ in Au+Au collisions
at $\sqrt{s_{{}_{NN}}}$=200 GeV in AMPT model. The black points are the
average $v_{3}$ in the selected $\varepsilon_{3}$ bin. The pad in the right
down corner is the average of $v_{3}$ with smaller scale.
It is commonly assumed that the harmonic flow coefficients $v_{n}$ linearly
depends on the $\varepsilon_{n}$. This assumption is supported by hydrodynamic
simulations v3hydro as long as one probes deformed initial profiles with only
a single non-vanishing harmonic eccentricity coefficient. In Fig. 3 and Fig.
4, we investigate the feasibility of this assumption for $v_{2}$ and $v_{3}$
respectively. The relations between $v_{n}$ and $\varepsilon_{n}$ are drawn
event by event in the two-dimension plots. The black points are the average
values of $v_{n}$ in an selected $\varepsilon_{n}$ bin, and the curves are the
connection of points to guide our eyes. The pads in the right down corners are
the average values of $v_{n}$ with smaller scale. In Fig. 3, the $v_{2}$
increases with $\varepsilon_{2}$, which is consistent with the ideal
hydrodynamic calculation Ulrichv3hydro . While in Fig. 4, the triangular flow
$v_{3}$ firstly increases with $\varepsilon_{3}$ up to 0.17, and then
decreases. Based on our study, the higher $\varepsilon_{3}$ bin corresponds to
the more peripheral collisions. It is known that $v_{3}$ is caused by initial
geometrical fluctuation, and built up by the interactions of constituents. The
less interactions in the higher $\varepsilon_{3}$ bin may cause the less
converting efficiency from $\varepsilon_{3}$ to $v_{3}$. That could be the
reason of decreasing trend of $v_{3}$ when $\varepsilon_{3}$ is larger than
0.17. Both the trend and the value of $v_{3}$ show discrepancy to the ideal
hydrodynamic calculations Ulrichv3hydro . As discussed in Ref. v3hydro , the
viscosity causes the decrease of $v_{3}$, however, the effects to $v_{3}$
versus $\varepsilon_{3}$ is not shown in the viscous hydrodynamic calculation.
Figure 5: $v_{n}/\varepsilon_{n}$ as a function of transverse particle density
in both Au+Au and Cu+Cu collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV in AMPT
model. Figure 6: $v_{3}$ as a function of transverse momentum in (a) Au+Au and
(b) Cu+Cu collisions at $\sqrt{s_{{}_{NN}}}$=200 GeV in AMPT model.
The ratio of elliptic flow to eccentricity $v_{2}/\varepsilon_{2}$ gains lots
of interests by comparing with the hydrodynamic model definition ;
Ollitrault_etaOverS . Recently, the behavior of triangular flow $v_{3}$ in
ideal hydrodynamics is also discussed v3hydro . In Fig. 5, we study the
$v_{n}/\varepsilon_{n}$ as a function of transverse particle density. From the
plot we can see that $v_{n}/\varepsilon_{n}$ from Au+Au and Cu+Cu are
consistent with each other very well. As the transverse particle density
increases, $v_{3}/\varepsilon_{3}$ rises with smaller value than
$v_{2}/\varepsilon_{2}$. It implies that as the particle density increases,
the initial geometry asymmetry transfers to momentum asymmetry more
efficiently while the system expands. Besides, the second order harmonic is
more efficient than the third order.
At last, the transverse momentum dependence of $v_{3}$ for $\pi$, $K$, p and
$\Lambda$ is also studied in Au+Au and Cu+Cu collisions. In Fig. 6, we can see
that $v_{3}$ shows quite similar trend to $v_{2}$. At low $p_{T}$, the mass
ordering phenomena is observed. The lighter particles are found with larger
$v_{3}$. It indicates that although $v_{3}$ is driven by $\varepsilon_{3}$,
its transverse momentum dependence is dominated by the hydrodynamics behavior
of the system. While when $p_{T}\geq 1.5$ GeV/$c$, baryons and mesons are
separated into two groups. The $p_{T}$ dependence of $v_{3}$ from identified
particles is qualitatively similar to the $p_{T}$ dependence of $v_{2}$
PRCrun4 ; cucu_STAR . The $v_{3}$ results of identified particles from AMPT
model are similar to the STAR preliminary results Yadav .
## V Summary
In summary, we studied the relation between initial geometry parameter
$\varepsilon_{n}$ and anisotropic flow $v_{n}$ in Au+Au and Cu+Cu collisions
using the AMPT Monte-Carlo model. We find that the triangular flow $v_{3}$ is
less sensitive to the centrality and system size compared with the elliptic
flow $v_{2}$. The $v_{2}$ displays an increasing trend as a function of
$\varepsilon_{2}$, which is qualitatively consistent with hydrodynamic
calculation. We found that $v_{3}$ shows an increasing trend when
$\varepsilon_{3}$ is less than 0.17, and then decreases beyond
$\varepsilon_{3}$ = 0.17. It may be because of the lower converting efficiency
from $\varepsilon_{3}$ to $v_{3}$ in the higher $\varepsilon_{3}$ bin. This
decreasing trend is in contrast to the results of ideal hydrodynamic
calculation. Both $v_{2}/\varepsilon_{2}$ and $v_{3}/\varepsilon_{3}$ increase
with the transverse particle density, and the second harmonic asymmetry in the
initial geometry seems to transfer to the momentum asymmetry more efficiently
than the third harmonic. The triangular flow $v_{3}$ of identified particles
shows a mass ordering in low $p_{T}$ and meson-baryon splitting at
intermediate $p_{T}$ in both Au+Au and Cu+Cu collisions which is similar to
the $p_{T}$ dependence of $v_{2}$.
## VI Acknowledgments
We wish to thank Prof. Fuqiang Wang for useful suggestions, and Dr. Kejun Wu
for useful discussions on the AMPT model. This work was supported in part by
the National Natural Science Foundation of China under grant No. 10775060,
11105060, 11135011, 11147146 and ‘the Fundamental Research Funds for the
Central Universities’, Grant No. HUST: 2011QN195.
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* (16) S. A. Voloshin, A. M. Poskanzer, R. Snellings arXiv:0809.2949 (2008) and Voloshin, Sergei A., Poskanzer, Arthur M., Snellings, Raimond: Collective Phenomena in Non-Central Nuclear Collisions. Stock, R. (ed.). SpringerMaterials - The Landolt-Bornstein Database (http://www.springermaterials.com). Springer-Verlag Berlin Heidelberg, 2010. $DOI:10.1007/978-3-642-01539-7_{1}0$
* (17) X.-N. Wang, Phys. Rev. D 43 104 (1991).
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* (19) B.A. Li and C.M. Ko, Phys. Rev. C 52 2037 (1995).
* (20) J. Adams et al (STAR Collaboration), Phys. Rev. C 72, 14904 (2005).
* (21) Z. Qiu and U. Heinz, Phys. Rev. C 84, 024911 (2005).
* (22) H. J. Drescher, A. Dumitru, C. Gombeaud and J. Y. Ollitrault, Phys. Rev. C76 024905 (2007).
* (23) Y. Pandit for the STAR collaboration, poster, RHIC and AGS Annual Users’ Meeting, Upton, NY, Jun. 20-24, 2011.
|
arxiv-papers
| 2011-11-27T01:31:48 |
2024-09-04T02:49:24.662529
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kai Xiao, Na Li, Shusu Shi, Feng Liu",
"submitter": "Shusu Shi",
"url": "https://arxiv.org/abs/1111.6213"
}
|
1111.6515
|
# On the interplay of direct and indirect $C\\!P$ violation in the charm
sector
M Gersabeck1, M Alexander2, S Borghi2,3, VV Gligorov1 and C Parkes3,1 1
European Organization for Nuclear Research (CERN), Geneva, Switzerland 2
School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom 3 School of Physics and Astronomy, University of Manchester,
Manchester, United Kingdom marco.gersabeck@cern.ch
###### Abstract
Charm mixing and $C\\!P$ violation observables are examined in the light of
the recently reported evidence from LHCb for $C\\!P$ violation in the charm
sector. If the result is confirmed as being due to direct $C\\!P$ violation at
the $1\%$ level, its effect will need to be taken into account in the
interpretation of $C\\!P$ violation observables. The contributions of direct
and indirect $C\\!P$ violation to the decay rate asymmetry difference $\Delta
A_{C\\!P}$ and the ratios of effective lifetimes $A_{\Gamma}$ and $y_{CP}$ are
considered here. Terms relevant to the interpretation of future high precision
measurements which have been neglected in previous literature are identified.
###### pacs:
13.25.Ft
††: J. Phys. G: Nucl. Part. Phys.
Charm, after strange and beauty, is the last system of neutral flavoured
mesons where $C\\!P$ violation remains to be discovered. While neutral
$\mathrm{B}$ mesons are characterised by their mass splitting leading to fast
oscillations and neutral kaons by their width splitting which results in a
short-lived and a long-lived state, neutral charm mesons have a very small
splitting in both mass and width. For charm, compared to the beauty sector,
this leads to rather subtle mixing-related effects in time-dependent as well
as in time-integrated charm measurements, which are examined in detail here.
First evidence for $C\\!P$ violation in the charm sector has recently been
reported by the LHCb collaboration in the study of the difference of the time-
integrated asymmetries of $D^{0}\to K^{+}K^{-}$ and $D^{0}\to\pi^{+}\pi^{-}$
decay rates through the parameter $\Delta A_{C\\!P}$ [1]. This measurement is
primarily sensitive to the difference in direct $C\\!P$ violation between the
two final states as discussed further below. Direct $C\\!P$ violation depends
on the final state and is the asymmetry of the rates of particle and
antiparticle decays. It can be caused by a difference in the magnitude of the
decay rates or by a difference in their phase. Indirect $C\\!P$ violation is
considered universal, i.e. final-state independent, and is an asymmetry in the
mixing rate or in its weak phase.
Indirect $C\\!P$ violation can be measured in time-dependent analyses. To
date, two types of measurements were used to search for indirect $C\\!P$
violation in the charm sector. One uses the asymmetry of the lifetimes,
$A_{\Gamma}$, measured in $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays to the $C\\!P$ eigenstates
$K^{+}K^{-}$ or $\pi^{+}\pi^{-}$ [2, 3, 4]. The other is a time-dependent
Dalitz plot analysis of $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays to $K^{0}_{\rm\scriptstyle
S}\pi^{+}\pi^{-}$ or $K^{0}_{\rm\scriptstyle S}K^{+}K^{-}$ [5, 6]. Another
observable studied related to $A_{\Gamma}$ is $y_{CP}$, which is given by the
deviation from one of the ratio of the lifetimes measured in decays to a
Cabibbo-allowed, $C\\!P$ averaged, and a Cabibbo-suppressed, $C\\!P$
eigenstate, final state. Any deviation of a measurement of $y_{CP}$ from that
of the mixing parameter $y$ would signal $C\\!P$ violation.
In the interpretation of $A_{\Gamma}$ and $y_{CP}$, direct $C\\!P$ violation
is commonly neglected [7]. In the light of the new evidence this assumption is
no longer justified. The relevance of direct $C\\!P$ violation to a
measurement of $A_{\Gamma}$ has previously been pointed out in [8]. However, a
closer look at both $A_{\Gamma}$ and $y_{CP}$ is necessary to examine the
contribution of direct and indirect $C\\!P$ violation in these observables as
well as their connection to $\Delta A_{C\\!P}$.
The mass eigenstates of neutral $D$ mesons, $|D_{1,2}\rangle$, with masses
$m_{1,2}$ and widths $\Gamma_{1,2}$ can be written as linear combinations of
the flavour eigenstates $|D_{1,2}\rangle=p|D^{0}\rangle\pm{}q|\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rangle$, with complex coefficients
$p$ and $q$ which satisfy $|p|^{2}+|q|^{2}=1$. The average mass and width are
defined as $m\equiv(m_{1}+m_{2})/2$ and
$\Gamma\equiv(\Gamma_{1}+\Gamma_{2})/2$. The $D$ mixing parameters are defined
using the mass and width difference as $x\equiv(m_{2}-m_{1})/\Gamma$ and
$y\equiv(\Gamma_{2}-\Gamma_{1})/2\Gamma$. The phase convention of $p$ and $q$
is chosen such that $C\\!P|D^{0}\rangle=-|\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\rangle$.
According to [8], the time dependent decay rates of $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays to the final state $f$,
which is a $C\\!P$ eigenstate with eigenvalue $\eta_{C\\!P}$, can be expressed
as
$\displaystyle\Gamma(D^{0}(t)\to
f)=\frac{1}{2}\rme^{-\tau}\left|A_{f}\right|^{2}\Big{\\{}$
$\displaystyle\left(1+|\lambda_{f}|^{2}\right)\cosh(y\tau)+\left(1-|\lambda_{f}|^{2}\right)\cos(x\tau)$
$\displaystyle+2\Re(\lambda_{f})\sinh(y\tau)-2\Im(\lambda_{f})\sin(x\tau)\Big{\\}},$
$\displaystyle\Gamma(\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}(t)\to
f)=\frac{1}{2}\rme^{-\tau}\left|\bar{A}_{f}\right|^{2}\Big{\\{}$
$\displaystyle\left(1+|\lambda^{-1}_{f}|^{2}\right)\cosh(y\tau)+\left(1-|\lambda^{-1}_{f}|^{2}\right)\cos(x\tau)$
(1)
$\displaystyle+2\Re(\lambda^{-1}_{f})\sinh(y\tau)-2\Im(\lambda^{-1}_{f})\sin(x\tau)\Big{\\}},$
where $\tau\equiv\Gamma t$, $\kern-1.99997pt\stackrel{{\scriptstyle\kern
1.39998pt\textsf{(---)}}}{{A}}_{\kern-2.10002ptf}\kern-3.00003pt$ are the
decay amplitudes and $\lambda_{f}$ is given by
$\lambda_{f}=\frac{q\bar{A}_{f}}{pA_{f}}=-\eta_{C\\!P}\left|\frac{q}{p}\right|\left|\frac{\bar{A}_{f}}{A_{f}}\right|\rme^{i\phi},$
(2)
where $\eta_{C\\!P}$ is the $C\\!P$ eigenvalue of the final state $f$ and
$\phi$ is the $C\\!P$ violating relative phase between $q/p$ and
$\bar{A}_{f}/A_{f}$. Introducing $|q/p|^{\pm 2}\approx 1\pm A_{m}$ and
$|\bar{A}_{f}/A_{f}|^{\pm 2}\approx 1\pm A_{d}$, one can write
$|\lambda_{f}^{\pm 1}|^{2}\approx(1\pm A_{m})(1\pm A_{d}),$ (3)
where $A_{m}$ represents a $C\\!P$ violation contribution from mixing and
$A_{d}$ from direct $C\\!P$ violation and where both $A_{m}$ and $A_{d}$ are
assumed to be small.
Expanding (On the interplay of direct and indirect $C\\!P$ violation in the
charm sector) up to second order in $\tau$ one can write the effective
lifetimes, i.e. those measured as a single exponential, as
$\displaystyle\hat{\Gamma}(\kern-1.00006pt\stackrel{{\scriptstyle\kern
0.70004pt\textsf{(---)}}}{{D}}\kern-3.00003pt(t)\to f)\approx$
$\displaystyle\Gamma\Bigg{\\{}1+\left[1\pm\frac{1}{2}(A_{m}+A_{d})-\frac{1}{8}(A_{m}^{2}-2A_{m}A_{d})\right]\eta_{C\\!P}(y\cos\phi\mp
x\sin\phi)$ (4) $\displaystyle\quad\mp A_{m}(x^{2}+y^{2})\pm
2A_{m}y^{2}\cos^{2}\phi\mp 4xy\cos\phi\sin\phi\Bigg{\\}},$
where terms below order $10^{-5}$ have been ignored. The experimental
constraints [9] give $x$, $y$, and $A_{d}$ for the final states $K^{+}K^{-}$
and $\pi^{+}\pi^{-}$ of order $10^{-2}$, and $A_{m}$ and $\sin\phi$ of order
$10^{-1}$. The sum of measurements of $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays leads to the definition of
the observable $y_{CP}$ which is given by
$y_{CP}=\frac{\hat{\Gamma}+\hat{\bar{\Gamma}}}{2\Gamma}-1\approx\eta_{C\\!P}\Bigg{\\{}\left[1-\frac{1}{8}(A_{m}^{2}-2A_{m}A_{d})\right]y\cos\phi-\frac{1}{2}(A_{m}+A_{d})x\sin\phi\Bigg{\\}}.$
(5)
The difference of measurements of $D^{0}$ and $\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ decays leads to the parameter
$A_{\Gamma}$ which is defined as
$\displaystyle
A_{\Gamma}=(\hat{\Gamma}-\hat{\bar{\Gamma}})(\hat{\Gamma}+\hat{\bar{\Gamma}})^{-1}\approx$
$\displaystyle\bigg{[}\frac{1}{2}(A_{m}+A_{d})y\cos\phi-\left(1-\frac{1}{8}A_{m}^{2}\right)x\sin\phi-
A_{m}(x^{2}+y^{2})$ (6)
$\displaystyle+2A_{m}y^{2}\cos^{2}\phi-4xy\cos\phi\sin\phi\bigg{]}\frac{\eta_{C\\!P}}{1+y_{CP}}.$
The weak phase $\phi$ has not been assumed to be universal. When averaging
measurements from different channels, a potential decay-dependent weak phase
of the amplitude ratio has to be taken into account [8]. Expanding only up to
order $10^{-4}$ leads to
$y_{CP}\approx\eta_{C\\!P}\left[\left(1-\frac{1}{8}A_{m}^{2}\right)y\cos\phi-\frac{1}{2}(A_{m})x\sin\phi\right],$
(7)
and
$A_{\Gamma}\approx\bigg{[}\frac{1}{2}(A_{m}+A_{d})y\cos\phi-x\sin\phi\bigg{]}\frac{\eta_{C\\!P}}{1+y_{CP}}\approx\eta_{C\\!P}\left[\frac{1}{2}(A_{m}+A_{d})y\cos\phi-x\sin\phi\right].$
(8)
The difference of $y_{CP}$ evaluated in (7) from the expression used in
literature [7] so far is the term $\eta_{C\\!P}\frac{1}{8}A_{m}^{2}y\cos\phi$,
which can be of similar order as $\frac{1}{2}A_{m}x\sin\phi$ and should
therefore not be ignored. Equation (8) shows that there can be a significant
contribution to $A_{\Gamma}$ from direct $C\\!P$ violation. Assuming $y=1\%$
and $\cos\phi=1$, direct $C\\!P$ violation at the level of $A_{d}/2=1\%$ would
lead to a contribution to $A_{\Gamma}$ of $10^{-4}$. Current measurements
yield a sensitivity of a few $10^{-3}$ [2, 3, 4]. Future measurements at LHCb
and future $\mathrm{B}$ factory experiments are expected to reach
uncertainties of the level of $10^{-4}$, i.e. that of the direct $C\\!P$
violation contribution. More precise measurements may well change the
approximations made in (7) and (8), in particular a measurement of
$A_{m}\lesssim A_{d}$.
In time-integrated measurements the rate asymmetry is measured which is
defined as
$A_{C\\!P}\equiv\frac{\Gamma(D^{0}\to f)-\Gamma(\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}\to f)}{\Gamma(D^{0}\to
f)+\Gamma(\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}\to f)}.$ (9)
Introducing
$a_{C\\!P}^{dir}\equiv\frac{|A_{f}|^{2}-|\bar{A}_{f}|^{2}}{|A_{f}|^{2}+|\bar{A}_{f}|^{2}}=\frac{1-\left|\frac{\bar{A}_{f}}{A_{f}}\right|^{2}}{1+\left|\frac{\bar{A}_{f}}{A_{f}}\right|^{2}}=\frac{-A_{d}}{2+A_{d}}\approx-\frac{1}{2}A_{d},$
(10)
and using (On the interplay of direct and indirect $C\\!P$ violation in the
charm sector), (9) becomes
$A_{C\\!P}\approx
a_{C\\!P}^{dir}-A_{\Gamma}(1-(a_{C\\!P}^{dir})^{2})\frac{\langle
t\rangle}{\tau}\approx a_{C\\!P}^{dir}-A_{\Gamma}\frac{\langle
t\rangle}{\tau},$ (11)
where $\langle t\rangle$ denotes the average decay time of the observed
candidates. Terms in $\langle t\rangle^{2}$ are below order $10^{-4}$, given
current experimental constraints, and have been ignored.
A common way to reduce experimental systematic uncertainties is to measure the
difference in time-integrated asymmetries in related final states. For the
two-body final states $K^{+}K^{-}$ and $\pi^{+}\pi^{-}$, this difference is
given by
$\displaystyle\Delta A_{C\\!P}$ $\displaystyle\equiv
A_{C\\!P}(K^{+}K^{-})-A_{C\\!P}(\pi^{+}\pi^{-})$ (12)
$\displaystyle=a_{C\\!P}^{dir}(K^{+}K^{-})-a_{C\\!P}^{dir}(\pi^{+}\pi^{-})$
$\displaystyle\quad-A_{\Gamma}(K^{+}K^{-})\frac{\langle
t(K^{+}K^{-})\rangle}{\tau}+A_{\Gamma}(\pi^{+}\pi^{-})\frac{\langle
t(\pi^{+}\pi^{-})\rangle}{\tau}.$
Assuming the $C\\!P$ violating phase $\phi$ to be universal [10] this can be
rewritten as
$\displaystyle\Delta A_{C\\!P}\approx\Delta
a_{C\\!P}^{dir}\left(1+y\cos\phi\frac{\overline{\langle
t\rangle}}{\tau}\right)+\left(a_{C\\!P}^{ind}+\overline{a_{C\\!P}^{dir}}y\cos\phi\right)\frac{\Delta\langle
t\rangle}{\tau}$ (13)
where $\Delta X\equiv X(K^{+}K^{-})-X(\pi^{+}\pi^{-})$, $\Delta X\equiv
X(K^{+}K^{-})-X(\pi^{+}\pi^{-})$, and
$a_{C\\!P}^{ind}=-(A_{m}/2)y\cos\phi+x\sin\phi$. The ratio $\overline{\langle
t\rangle}/\tau$ is equal to one for the lifetime-unbiased $\mathrm{B}$ factory
measurements [11, 12] and is $2.083\pm 0.001$ for LHCb [1] and $2.53\pm 0.02$
for CDF [13], thus leading to a correction of $\Delta a_{C\\!P}^{dir}$ of the
order of $10^{-2}$. The factor $\Delta\langle t\rangle/\tau$ multiplying the
indirect $C\\!P$ violation is zero for the $\mathrm{B}$ factory measurements
and ranges from $0.098\pm 0.003$ to $0.26\pm 0.01$ for LHCb and CDF,
respectively. Therefore, $\Delta A_{C\\!P}$ is largely a measure of direct
$C\\!P$ violation while an obvious contribution from indirect $C\\!P$
violation exists. The contribution from direct $C\\!P$ violation to
$A_{\Gamma}$ pointed out in (8) leads to a term proportional to $y$. This term
may be of similar size as the term proportional to $\Delta\langle t\rangle$
and should therefore be taken into account.
In summary, the mixing and $C\\!P$ violation parameters $y_{CP}$, $A_{\Gamma}$
and $\Delta A_{C\\!P}$ have been discussed in the light of the recent evidence
for $C\\!P$ violation in the $D^{0}$ sector. The parameter $y_{CP}$ is least
affected by direct $C\\!P$ violation, however, it contains a term which has
been neglected in the literature so far and which can be of the same order as
the constribution proportional to $x$. A measurement of $A_{\Gamma}$ can
exhibit a contribution of direct $C\\!P$ violation at the level of $10^{-4}$,
comparable to the expected future experimental sensitivity. The direct $C\\!P$
violation term in the $\Delta A_{C\\!P}$ measurement contains a contribution
proportional to $y$. The interpretation of future high precision measurements
of these observables will need to take account of these contributions.
The authors would like to thank Tim Gershon, Bostjan Golob, Alex Kagan and
Vincenzo Vagnoni for helpful discussions and comments. Furthermore, the
authors thank the LHCb collaboration whose work inspired this paper. MG and
VVG are supported by a Marie Curie Action: “Cofunding of the CERN Fellowship
Programme (COFUND-CERN)” of the European Community’s Seventh Framework
Programme under contract number (PCOFUND-GA-2008-229600). MA, SB and CP
acknowledge the support of the STFC (United Kingdom).
## References
## References
* [1] Aaij R et al. [LHCb collaboration]. Evidence for $C\\!P$ violation in time-integrated $D^{0}\to h^{-}h^{+}$ decay rates. 2011\. Preprint hep-ex 1112.0938.
* [2] Staric M et al. [Belle collaboration]. Evidence for $D^{0}$\- $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ Mixing. Phys. Rev. Lett., 98:211803, 2007.
* [3] Aubert B et al. [BaBar collaboration]. Measurement of $D^{0}$\- $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ Mixing using the Ratio of Lifetimes for the Decays $D^{0}\to K^{-}\pi^{+}$ and $K^{+}K^{-}$. Phys. Rev., D80:071103, 2009.
* [4] Aaij R et al. [LHCb collaboration]. Measurement of mixing and $C\\!P$ violation parameters in two-body charm decays. 2011\. Preprint hep-ex 1112.4698.
* [5] Abe K et al. [Belle collaboration]. Measurement of $D^{0}$-$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ Mixing Parameters in $D^{0}\to K^{0}_{\rm\scriptstyle S}\pi^{+}\pi^{-}$ decays. Phys. Rev. Lett., 99:131803, 2007.
* [6] Del Amo Sanchez P et al. [BaBar collaboration]. Measurement of $D^{0}$-$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mixing parameters using $D^{0}\to K^{0}_{\rm\scriptstyle S}\pi^{+}\pi^{-}$ and $D^{0}\to K^{0}_{\rm\scriptstyle S}K^{+}K^{-}$ decays. Phys. Rev. Lett., 105:081803, 2010.
* [7] Bergmann S, Grossman Y, Ligeti Z, Nir Y, and Petrov A A. Lessons from CLEO and FOCUS measurements of $D^{0}$-$\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ mixing parameters. Phys. Lett., B486:418–425, 2000.
* [8] Kagan A L and Sokoloff M D. On Indirect $C\\!P$ Violation and Implications for $D^{0}$\- $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$ and $\mathrm{B}_{(s)}$ \- $\kern 1.99997pt\overline{\kern-1.99997pt\mathrm{B}}{}_{(s)}$ mixing. Phys.Rev., D80:076008, 2009.
* [9] Asner D et al. [Heavy Flavor Averaging Group]. Averages of b-hadron, c-hadron, and $\tau$-lepton Properties. 2010\. online update at http://www.slac.stanford.edu/xorg/hfag/index.html.
* [10] Grossman Y, Kagan A L, and Nir Y. New physics and $C\\!P$ violation in singly Cabibbo suppressed D decays. Phys.Rev., D75:036008, 2007.
* [11] Aubert B et al. [BaBar collaboration]. Search for $C\\!P$ violation in the decays $D^{0}\to K^{-}K^{+}$ and $D^{0}\to\pi^{-}\pi^{+}$. Phys. Rev. Lett., 100:061803, 2008.
* [12] Staric M et al. [Belle collaboration]. Measurement of $C\\!P$ asymmetry in Cabibbo suppressed $D^{0}$ decays. Phys. Lett., B670:190–195, 2008.
* [13] Aaltonen T et al. [CDF collaboration]. Measurement of $C\\!P$–violating asymmetries in $D^{0}\to\pi^{+}\pi^{-}$ and $D^{0}\to K^{+}K^{-}$ decays at CDF. Preprint hep-ex 1111.5023.
|
arxiv-papers
| 2011-11-28T17:03:29 |
2024-09-04T02:49:24.677351
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Marco Gersabeck (1), Michael Alexander (2), Silvia Borghi (2 and 3),\n Vladimir V Gligorov (1), Chris Parkes (3 and 1) ((1) CERN, Geneva,\n Switzerland, (2) School of Physics and Astronomy, University of Glasgow,\n Glasgow, UK, (3) School of Physics and Astronomy, University of Manchester,\n Manchester, UK)",
"submitter": "Marco Gersabeck",
"url": "https://arxiv.org/abs/1111.6515"
}
|
1111.6550
|
# Sensitivity to eV-scale Neutrinos of Experiments at a
Very Low Energy Neutrino Factory
J.H. Cobb C.D. Tunnell Corresponding author: tunnell@fnal.gov Subdepartment
of Particle Physics, University of Oxford, Denys Wilkinson Building, Keble
Road, Oxford, OX1 3RH, UK A.D. Bross Fermilab, P.O. Box 500, Batavia, IL
60510-0500, USA
###### Abstract
The results of LSND have yet to be confirmed at the $5\sigma$-level. An
experiment is proposed utilizing a 3 GeV muon storage ring that would allow
for both disappearance and appearance channels to be explored at short-
baselines. The appearance channel could provide well over $5\sigma$
confirmation or rejection of the LSND result. Other physics could also be
performed at such a facility such as the measurement of electron-neutrino
cross sections. The sensitivity of experiments at a Very Low Energy Neutrino
Factory (VLENF) to neutrinos at the eV-scale is presented.
Suggested keywords
###### pacs:
14.60.Pq, 14.60.St
††preprint: APS/123-QED
Despite the experimental success of the past decade in establishing that
neutrinos are massive and mix, whether they be produced in the Sun Aharmim
_et al._ (2008), reactors Ahn _et al._ (2006), accelerators Adamson _et al._
(2008), or the atmosphere Wendell _et al._ (2010), there exist nevertheless
some anomalies. All oscillation experiments to date can be explained with
three neutrinos except LSND — and various experiments could be made to agree
with either. The LEP collider experiments lep (2006) showed that additional
light neutrinos cannot couple to the $Z$. The MiniBooNe experiment has been
unable to deny or confirm the LSND result since their data that uses the same
anti-neutrino beam polarity as LSND agrees with the background-only hypothesis
at 0.5% and agrees with short-baseline oscillations at 8.7% Aguilar-Arevalo
_et al._ (2010). An experiment with a sensitivity greater than $5\sigma$ is
needed in order to refute or confirm evidence of neutrinos at the eV-scale.
The LSND data suggests a new mass splitting Athanassopoulos _et al._ (1998)
that is potentially observed elsewhere. Recently, a recalculation of reactor
fluxes resulted in the _reactor neutrino anomaly_ Mueller _et al._ (2011);
Huber (2011) which, along with the _Gallium anomaly_ Acero _et al._ (2008),
are possibly caused by sterile neutrinos. Global fits reconcile these
anomalies into a framework that introduces a sterile neutrino at the eV-scale,
_e.g._ , Giunti and Laveder (2011a); Kopp _et al._ (2011). Additional
indications come from cosmology where WMAP favors more than three neutrinos
Komatsu _et al._ (2011).
The Very Low Energy Neutrino Factory (VLENF) would utilize a muon storage ring
to study eV-scale oscillation physics and measure cross sections (including
$\nu_{e}$). Pions are collected from a target and injected into a storage ring
where they decay to muons. The storage ring is optimized for 2 GeV muons where
the energy is optimized for the needs of both oscillation and cross section
physics. The muons decay according to $\mu^{+}\to
e^{+}\bar{\nu}_{\mu}\nu_{e}$. Straight sections in the storage ring result in
neutrinos directed at a near- and far-detector.
The storage ring could be a fixed field alternating gradient (FFAG) lattice
for large momentum acceptance, which is important given the large momentum
spectrum after the target. The pricing and engineering of FFAGs are well-
understood because of experience building them c:f (2005). The design would
require only normal conducting magnets which simplifies construction,
commissioning, and operations. Design work for the injection into the storage
ring and the particle collection downstream of the target are underway.
The near detector will be placed at 20-50 meters from the end of the straight
and will measure neutrino-nucleon cross sections of interest to future long-
baseline experiments including the first precision measurement of $\nu_{e}$
cross sections. The far detector at 800 m would measure disappearance and
wrong-sign muon appearance channels. The detector would need to be magnetized
for the wrong-sign muon appearance channel. Numerous possibilities exist for
detector technologies that include liquid argon, MINOS inspired, and totally
active scintillating detectors. For the purposes of this study, a detector
inspired by MINOS but with thinner plates is assumed. The experiment will take
advantage of the “golden channel” of oscillation appearance
$\nu_{e}\to\nu_{\mu}$ where the resulting final state is a wrong-sign muon.
The probability $\nu_{e}\to\nu_{\mu}$ depends on the mixing matrix, $U$. Let
$R_{ij}$ be a rotation between the $i$th and $j$th mass eigenstates without CP
violation. For $N$ neutrinos, $R_{ij}$ has dimension $N\times N$. By
convention, the three neutrino mixing matrix is
$U_{\text{PMNS}}=R_{23}R_{13}R_{12}$. In the (3+1) model of neutrino
oscillations, extra rotations can be introduced such that the mixing matrix is
$U_{\text{(3+1)}}=R_{34}R_{24}R_{14}U_{\text{PMNS}}$. Given that $\Delta
m^{2}_{41}>>\Delta m^{2}_{31}$, $U_{\text{PMNS}}$ can be approximated by the
identity matrix (_ie._ the “short-baseline approximation”). It then follows
that $|U_{e4}|^{2}=\sin(\theta_{14})$ and $|U_{\mu
4}|^{2}=\sin(\theta_{24})\cos(\theta_{14})$.
The oscillation probabilities for appearance and disappearance, respectively,
are:
$\displaystyle\text{P}_{\nu_{e}\to\nu_{\mu}}=$ $\displaystyle
4|U_{e4}|^{2}|U_{\mu 4}|^{2}\sin^{2}\left(\frac{\Delta
m^{2}_{41}L}{4E}\right),$ (1)
$\displaystyle\text{P}_{\nu_{\alpha}\to\nu_{\alpha}}=$ $\displaystyle
1-\left[4|U_{\alpha 4}|^{2}(1-|U_{\alpha
4}|^{2})\right]\sin^{2}\left(\frac{\Delta m^{2}_{41}L}{4E}\right).$ (2)
Disappearance measurements will constrain $|U_{\mu 4}|^{2}$ and $|U_{e4}|^{2}$
while the appearance channel measures their product $|U_{e4}|^{2}|U_{\mu
4}|^{2}$. These parameters are over-constrained. Only the appearance channel
will be assumed since the physics sensitivity of disappearance measurements is
well-understood for short-baseline experiments. Without oscillations there
would be $\sim 10^{5}$ charge current (CC) events in the far detector. The
disappearance oscillation probability, if LSND is correct, would be on the
order of a few percent for this experiment; the statistical errors and
systematic errors are expected to be about half a percent from the presence of
a near detector and beam instrumentation.
The oscillation sensitivities have been computed using the GLoBES software
(version 3.1.10) Huber _et al._ (2005, 2007). Since GLoBES, by default, only
allows for a $3\times 3$ mixing matrix, the SNU (version 1.1) add-on Kopp
(2008); Kopp _et al._ (2008) is used to extend computations in GLoBES to
$4\times 4$ mixing matrices.
Unlike previous analyses, the far detector approximation of the source and
detector being treated as point-sources cannot be used since the size of the
detector and accelerator straight are comparable to the baseline. This study
computes the neutrino flux by integrating the phase space of the stored muons
using Monte Carlo (MC) integration. The beam occupies a 6D phase space ($x$,
$y$, $z$, $p_{x}$, $p_{y}$, $p_{z}$) and the detector has a $6\text{ m}\times
6\text{ m}$ cross section. A random point is chosen within the beam phase
space and within the detector volume. The transverse phase space is
represented by the Twiss parameters $\alpha=0$ and $\beta=25\text{ m}$ where
the $1\sigma$ Gaussian emittance is assumed to be $15\text{ mm}$. It follows
that the spread in, for example, $x$ is $\sigma_{x}=\sqrt{\beta\epsilon}$ and
the angular divergence in $x$ is $\sigma_{x^{\prime}}=\sqrt{\epsilon/\beta}$.
The longitudinal phase space ($z$ and $p_{z}$) is described by assuming a
uniform distribution in $z$ and $p_{z}=(2\pm 20\%)\text{ GeV}$. The code for
the analysis herein is available online Tunnell (2011) under the GPL license
gpl .
This analysis assumes $2\times 10^{17}$ decays of $\mu^{+}$ in each straight,
normalized to an exposure similar to MiniBooNe of $10^{21}$ protons on target
(POT). The far detector at 800 m consists of a kilotonne of fiducial target
mass. The background rejection for the wrong-sign muons is assumed to be
$10^{-4}$ for the $\bar{\nu}_{\mu}\to\bar{\nu}_{\mu}$ neutral current (NC)
events and $10^{-5}$ for the charge misidentification of
$\bar{\nu}_{\mu}\to\bar{\nu}_{\mu}$ CC events. The backgrounds from
$\nu_{e}\to\nu_{e}$ CC and NC are negligible. The $\nu_{e}\to\nu_{\mu}$ CC
events have a 90% detection efficiency. These efficiencies were extrapolated
from previous studies Laing (2010) and reconfirmation of these calculations
for energies below 2 GeV is still required.
Figure 1: The oscillation probability for the “golden channel”
$\nu_{e}\to\nu_{\mu}$ using the (3+1) oscillation parameters in TABLE 1. A
baseline of 800 meters is assumed. This numerical treatment of the binned
oscillation probability agrees with Eq. 1.
To initially gauge the sensitivity of this experimental setup, some
oscillation parameters must be assumed that include the LSND effect. The best-
fit data for (3+1) sterile neutrinos is used Giunti and Laveder (2011b) where
the MB $\bar{\nu}$ and LSND $\bar{\nu}$ data was fit (TABLE 1). Lower energy
neutrinos contain the most information about the mixing matrix (FIG. 1).
Table 1: Best-fit oscillation parameters for the (3+1) sterile neutrino scenario for MB $\bar{\nu}$ and LSND $\bar{\nu}$ data Giunti and Laveder (2011b). Parameter | Value
---|---
$\Delta m^{2}_{41}$ [$\text{eV}^{2}$] | 0.89
$|U_{e4}|^{2}$ | 0.025
$|U_{\mu 4}|^{2}$ | 0.023
Figure 2: The spectrum of signal and background events at the far-detector.
The number of wrong-sign muon appearance signal events is 27 which arises
through the “golden channel” oscillation probability $\nu_{e}\to\nu_{\mu}$.
Roughly 2 background events are in the signal box and correspond to both
$\bar{\nu}_{\mu}$ CC with charge misidentification and $\bar{\nu}_{\mu}$ NC
events where a short-lived pion has faked a wrong-sign muon. An exposure of
$10^{21}$ POT is assumed.
The best-fit values allow for an estimation of the event rates. The true
number of events without efficiencies is shown in TABLE 2. Applying the
assumed efficiencies of $10^{-5}$ and $10^{-4}$ for
$\bar{\nu}_{\mu}\to\bar{\nu}_{\mu}$ CC and $\bar{\nu}_{X}\to\bar{\nu}_{X}$ NC,
respectively, reveals the event spectrum (FIG. 2). After cuts, there are 27
signal events and 2 background events.
Table 2: A table of the raw event rates. The first row corresponds to the “golden channel” appearance signal. The other rows are potential backgrounds to the signal. The background that drives the analysis is $\bar{\nu}_{\mu}\to\bar{\nu}_{\mu}$. Channel | Interaction | Pre-cuts
---|---|---
$\nu_{e}\to\nu_{\mu}$ | CC | 30
$\bar{\nu}_{X}\to\bar{\nu}_{X}$ | NC | 16850
$\bar{\nu}_{\mu}\to\bar{\nu}_{\mu}$ | CC | 42545
$\nu_{e}\to\nu_{e}$ | CC | 78974
Figure 3: Optimization of the mean stored muon energy and far detector
baseline. Above muon energies of 1 GeV, appearance physics is not sensitive to
the muon energy. For a fixed neutrino energy, the boosted muon frame will
result in more events the more it is boosted. The optimal baseline is around
800 meters.
In practice the muon energy is fixed by the needs of the cross section
physics. However, given a fixed fiducial mass, the baseline and muon energy
can be optimized for short-baseline physics. The metric for comparison is the
$\chi^{2}$ between simulations for the (3+1) best-fit mentioned earlier and
the background-only hypothesis. The short-baseline oscillation physics reach
is comparable at energies above 2 GeV and the optimal baseline is around 800
meters (FIG. 3).
Having chosen a baseline and energy, it is possible to investigate the physics
reach of this facility if we appropriately define the desired $\chi^{2}$
value. The design requirement is to measure the LSND effect at “$5\sigma$” but
that can be ambiguously defined. Herein “$n\sigma$” is defined according to
the Gaussian probability of the respective deviation. Let $\text{CDF}_{X}$ be
the cumulative distribution function of some distribution $X$. If $G$ is a
Gaussian with $\mu=0$ and $\sigma=1$, then the $p$-value of a $5\sigma$ effect
is $1-\text{CDF}_{G}(5)\simeq 3\times 10^{-7}$. Similarly, if $\chi^{2}_{2}$
is the chi-squared distribution with two degrees of freedom, then the required
value of the $\chi^{2}$ to be an $n\sigma$ effect is
$\text{CDF}^{-1}_{\chi^{2}_{2}}(\text{CDF}_{G}(n))$.
Figure 4: Sensitivity to sterile parameters. The orange band corresponds to
the LSND $\bar{\nu}$ and MiniBooNe $\bar{\nu}$ sterile neutrino fit performed
by Giunti and Laveder Giunti and Laveder (2011a) where they assumed only one
extra massive neutrino. The contours are the sensitivities of an experiment at
a 3 GeV FFAG storage ring using parameters defined in the text. The entire 99%
confidence interval for the LSND $\bar{\nu}$ and MiniBooNe $\bar{\nu}$ results
would be confirmed or excluded at $7\sigma$.
The $\chi^{2}$ is computed using the pull-method, which allows for systematics
to be included. The signal and background normalization errors are assigned to
be 2% and 20%, respectively, and these systematic uncertainties are
marginalized over. Spectral information is also used. The sensitivity to new
physics is shown in FIG. 4. The 99% confidence interval of a fit in a (3+1)
scheme to the MiniBooNe $\bar{\nu}$ and LSND $\bar{\nu}$ data is shown and
excluded at $7\sigma$. This meets and exceeds the requested criterion for
future eV-scale experiments.
Some effects are ignored that could reduce the significance of the result. For
example, not included are backgrounds associated with cosmic muons. This could
be problematic since muons fill the storage ring which results in a large duty
factor. One solution, if it is an issue, would be to put RF cavities in the
storage ring to bunch the beam. This merits further study.
It has been demonstrated that a VLENF has sensitivity to eV-scale
oscillations. Further optimizations of the target horn, collection, injection,
and the dynamic aperture of the storage ring will increase the neutrino flux.
But previously unthought-of backgrounds may arise and detector charge
identification performance is difficult at 1 GeV and below. A design study
must be performed to ensure enough contingency in the sensitivity to guarantee
a firm confirmation or refutation of the LSND effect. Further study is
required.
_Thanks_ for guidance and useful discussions with Joachim Kopp, David Neuffer,
Kenneth Long, and Alain Blondel. Also thanks to Nick Ryder for a careful
reading.
## References
* Aharmim _et al._ (2008) B. Aharmim _et al._ (SNO), Phys. Rev. Lett. 101, 111301 (2008), arXiv:0806.0989 [nucl-ex] .
* Ahn _et al._ (2006) M. H. Ahn _et al._ (K2K), Phys. Rev. D74, 072003 (2006), arXiv:hep-ex/0606032 .
* Adamson _et al._ (2008) P. Adamson _et al._ (MINOS), Phys. Rev. Lett. 101, 131802 (2008), arXiv:0806.2237 [hep-ex] .
* Wendell _et al._ (2010) R. Wendell _et al._ (Kamiokande), Phys. Rev. D81, 092004 (2010), arXiv:1002.3471 [hep-ex] .
* lep (2006) Phys. Rept. 427, 257 (2006), arXiv:hep-ex/0509008 .
* Aguilar-Arevalo _et al._ (2010) A. A. Aguilar-Arevalo _et al._ (MiniBooNE Collaboration), Phys. Rev. Lett. 105, 181801 (2010).
* Athanassopoulos _et al._ (1998) C. Athanassopoulos _et al._ (LSND Collaboration), Phys. Rev. Lett. 81, 1774 (1998).
* Mueller _et al._ (2011) T. A. Mueller _et al._ , Phys. Rev. C83, 054615 (2011), arXiv:1101.2663 [hep-ex] .
* Huber (2011) P. Huber, Phys. Rev. C84, 024617 (2011), arXiv:1106.0687 [hep-ph] .
* Acero _et al._ (2008) M. A. Acero, C. Giunti, and M. Laveder, Phys. Rev. D78, 073009 (2008), arXiv:0711.4222 [hep-ph] .
* Giunti and Laveder (2011a) C. Giunti and M. Laveder, (2011a), arXiv:1107.1452 [hep-ph] .
* Kopp _et al._ (2011) J. Kopp, M. Maltoni, and T. Schwetz, Phys. Rev. Lett. 107, 091801 (2011), arXiv:1103.4570 [hep-ph] .
* Komatsu _et al._ (2011) E. Komatsu _et al._ (WMAP), Astrophys. J. Suppl. 192, 18 (2011), arXiv:1001.4538 [astro-ph.CO] .
* c:f (2005) http://hadron.kek.jp/FFAG/FFAG05_HP/index.htm (2005), see talks including the summary talk.
* Huber _et al._ (2005) P. Huber, M. Lindner, and W. Winter, Comput.Phys.Commun. 167, 195 (2005), arXiv:hep-ph/0407333 [hep-ph] .
* Huber _et al._ (2007) P. Huber, J. Kopp, M. Lindner, M. Rolinec, and W. Winter, Comput.Phys.Commun. 177, 432 (2007), arXiv:hep-ph/0701187 [hep-ph] .
* Kopp (2008) J. Kopp, Int. J. Mod. Phys. C19, 523 (2008), erratum ibid. C19 (2008) 845, arXiv:physics/0610206 .
* Kopp _et al._ (2008) J. Kopp, M. Lindner, T. Ota, and J. Sato, Phys. Rev. D77, 013007 (2008), arXiv:0708.0152 [hep-ph] .
* Tunnell (2011) C. D. Tunnell, https://code.launchpad.net/~c-tunnell1/+junk/vlenf_tools (2011), questions should be directed to the corresponding author.
* (20) http://www.gnu.org/licenses/gpl-3.0.html.
* Laing (2010) A. B. Laing, “Optimisation of detectors for the golden channel at a neutrino factory,” (2010).
* Giunti and Laveder (2011b) C. Giunti and M. Laveder, (2011b), arXiv:1109.4033 [hep-ph] .
|
arxiv-papers
| 2011-11-28T18:57:56 |
2024-09-04T02:49:24.684118
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Christopher D. Tunnell and John H. Cobb and Alan D. Bross",
"submitter": "Christopher Tunnell",
"url": "https://arxiv.org/abs/1111.6550"
}
|
1111.6671
|
# The dynamics of the 3D radial NLS with the combined terms
Changxing Miao Changxing Miao:
Institute of Applied Physics and Computational Mathematics,
P. O. Box 8009, Beijing, China, 100088, miao_changxing@iapcm.ac.cn ,
Guixiang Xu Guixiang Xu
Institute of Applied Physics and Computational Mathematics,
P. O. Box 8009, Beijing, China, 100088, xu_guixiang@iapcm.ac.cn and Lifeng
Zhao Lifeng Zhao
University of Science and Technology of China,
Hefei, China, zhaolifengustc@yahoo.cn
###### Abstract.
In this paper, we show the scattering and blow-up result of the radial
solution with the energy below the threshold for the nonlinear Schrödinger
equation (NLS) with the combined terms
$\displaystyle iu_{t}+\Delta u=-|u|^{4}u+|u|^{2}u$ (CNLS)
in the energy space $H^{1}({\mathbb{R}}^{3})$. The threshold is given by the
ground state $W$ for the energy-critical NLS: $iu_{t}+\Delta u=-|u|^{4}u$.
This problem was proposed by Tao, Visan and Zhang in [37]. The main difficulty
is the lack of the scaling invariance. Illuminated by [17], we need give the
new radial profile decomposition with the scaling parameter, then apply it
into the scattering theory. Our result shows that the defocusing,
$\dot{H}^{1}$-subcritical perturbation $|u|^{2}u$ does not affect the
determination of the threshold of the scattering solution of (CNLS) in the
energy space.
###### Key words and phrases:
Blow up; Dynamics; Nonlinear Schrödinger Equation; Scattering; Threshold
Energy.
###### 2000 Mathematics Subject Classification:
Primary: 35L70, Secondary: 35Q55
## 1\. Introduction
We consider the dynamics of the radial solutions for the nonlinear Schrödinger
equation (NLS) with the combined nonlinearities in $H^{1}({\mathbb{R}}^{3})$
$\left\\{\begin{aligned} iu_{t}+\Delta
u=&\;f_{1}(u)+f_{2}(u),,\quad(t,x)\in{\mathbb{R}}\times{\mathbb{R}}^{3},\\\
u(0)=&\;u_{0}(x)\in H^{1}({\mathbb{R}}^{3}).\end{aligned}\right.$ (1.1)
where $u:{\mathbb{R}}\times{\mathbb{R}}^{3}\mapsto{\mathbb{C}}$ and
$f_{1}(u)=-|u|^{4}u$, $f_{2}(u)=|u|^{2}u$. As we known, $f_{1}$ has the
$\dot{H}^{1}$-critical growth, $f_{2}$ has the $\dot{H}^{1}$-subcritical
growth.
The equation has the following mass and Hamiltonian quantities
$\displaystyle M(u)(t)=$
$\displaystyle\frac{1}{2}\int_{{\mathbb{R}}^{3}}|u(t,x)|^{2}\;dx;\quad
E(u)(t)=\int_{{\mathbb{R}}^{3}}\frac{1}{2}|\nabla
u(t,x)|^{2}\;dx+F_{1}(u(t))+F_{2}(u(t))$
where
$F_{1}(u(t))=\displaystyle-\frac{1}{6}\int_{{\mathbb{R}}^{3}}|u(t,x)|^{6}\;dx,\;\;F_{2}(u(t))=\frac{1}{4}\int_{{\mathbb{R}}^{3}}|u(t,x)|^{4}\;dx.$
They are conserved for the sufficient smooth solutions of (1.1).
In [37], Tao, Visan and Zhang made the comprehensive study of
$\displaystyle iu_{t}+\Delta u=|u|^{4}u+|u|^{2}u$
in the energy space. They made use of the interaction Morawetz estimate
established in [6] and the stability theory for the scattering solution. Their
result is based on the scattering result of the defocusing, energy-critical
NLS in the energy space, which is established by Bourgain [3, 4] for the
radial case, I-team [7], Ryckman-Visan [34] and Visan [38] for the general
data. Since the classical interaction Morawetz estimate in [6] fails for
(1.1), Tao, et al., leave the scattering and blow-up dichotomy of (1.1) below
the threshold as an open problem in [37]. For other results, please refer to
[15, 16, 30, 31, 32, 39, 40].
For the focusing, energy-critical NLS
$\displaystyle iu_{t}+\Delta u=-|u|^{4}u.$ (1.2)
Kenig and Merle first applied the concentration compactness in [2, 21, 22]
into the scattering theory of the radial solution of (1.2) in [19] with the
energy below that of the ground state of
$\displaystyle-\Delta W=|W|^{4}W.$ (1.3)
In this paper, we will also make use of the concentration compactness argument
and the stability theory to study the dichotomy of the radial solution of
(1.1) with the energy below the threshold, which will be shown to be the
energy of the ground state $W$ for (1.2). For the applications of the
concentration compactness in the scattering theory and rigidity theory of the
critical NLS, NLW, NLKG and Hartree equations, please see [8, 9, 10, 11, 12,
13, 17, 20, 23, 24, 25, 26, 27, 28, 29].
We now show the differences between (1.1) and (1.2). On one hand, there is an
explicit solution $W$ for (1.2), which is the ground state of (1.3) and does
not scatter. The threshold of the scattering solution of (1.2) is determined
by the energy of $W$. While for (1.1), there is no such explicit solution,
whose energy is the threshold of the scattering solution of (1.1). We need
look for a mechanism to determine the threshold of the scattering solution of
(1.1). It turns out that the constrained minimization of the energy as (1.5)
is appropriate111The similar constrained minimization of the energy as (1.5)
is not appropriate for the focusing perturbation: $iu_{t}+\Delta
u=-|u|^{4}u-|u|^{2}u$, since the threshold $m$ in this way equals to $0$ and
it is not the desired result.. On the other hand, for (1.2), it is
$\dot{H}^{1}$-scaling invariant, which gives us many conveniences, especially
in the nonlinear profile decomposition about (1.2). While for (1.1), it is the
lack of scaling invariance. We need give the new profile decomposition with
the scaling parameter of (1.1) in $H^{1}(R^{3})$, take care of the role of the
scaling parameter in the linear and nonlinear profile decompositions, then
apply them into the scattering theory.
Now for $\varphi\in H^{1}$, we denote the scaling quantity
$\varphi^{\lambda}_{3,-2}$ by
$\displaystyle\varphi^{\lambda}_{3,-2}(x)=e^{3\lambda}\varphi(e^{2\lambda}x).$
We denote the scaling derivative of $E$ by $K(\varphi)$
$\displaystyle
K(\varphi)=\mathcal{L}E(\varphi):=\dfrac{d}{d\lambda}\Big{|}_{\lambda=0}E(\varphi^{\lambda}_{3,-2})=\int_{{\mathbb{R}}^{3}}\left(\frac{4}{2}|\nabla\varphi|^{2}-\frac{12}{6}|\varphi|^{6}+\frac{6}{4}|\varphi|^{4}\right)\;dx,$
(1.4)
which is connected with the Virial identity, and then plays the important role
in the blow-up and scattering of the solution of (1.1).
Now the threshold $m$ is determined by the following constrained
minimization222In fact, the following minimization of the static energy
$\inf\\{M(\varphi)+E(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)=0\\}$ also equals to $m$.
of the energy $E(\varphi)$
$\displaystyle m=\inf\\{E(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)=0\\}.$ (1.5)
Since we consider the $\dot{H}^{1}$-critical growth with the
$\dot{H}^{1}$-subcritical perturbation, we will use the modified energy later
$\displaystyle E^{c}(u)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\frac{1}{2}|\nabla
u(t,x)|^{2}-\frac{1}{6}|u(t,x)|^{6}\right)\;dx.$
As the nonlinearity $|u|^{2}u$ is the defocusing, $\dot{H}^{1}$-subcritical
perturbation, one think that the focusing, $\dot{H}^{1}$-critical term plays
the decisive role of the threshold of the scattering solution of (1.1) in the
energy space. The first result is to characterize the threshold energy $m$ as
following
###### Proposition 1.1.
There is no minimizer for (1.5). But for the threshold energy $m$, we have
$\displaystyle m=E^{c}(W),$
where $W\in\dot{H}^{1}({\mathbb{R}}^{3})$ is the ground state of the massless
equation
$\displaystyle-\Delta W=|W|^{4}W.$
As the dynamics of the solution of (1.1) with the energy less than the
threshold $m$, the conjecture is
###### Conjecture 1.2.
Let $u_{0}\in H^{1}({\mathbb{R}}^{3})$ with
$E(u_{0})<m,$ (1.6)
and $u$ be the solution of (1.1) and $I$ be its maximal interval of existence.
Then
1. (a)
If $K(u_{0})\geq 0$, then $I={\mathbb{R}}$, and $u$ scatters in both time
directions as $t\rightarrow\pm\infty$ in $H^{1}$;
2. (b)
If $K(u_{0})<0$, then $u$ blows up both forward and backward at finite time in
$H^{1}$.
In this paper, we verify the conjecture in the radial case.
###### Theorem 1.3.
Conjecture 1.2 holds whenever $u$ is spherically symmetric.
###### Remark 1.4.
Our consideration of the radial case is based on the following facts:
1. (1)
It is an open problem that the scattering result of (1.2) in dimension three,
except for the radial case in [19]. Our result is based on the corresponding
scattering result of (1.2).
2. (2)
It seems to be hard to lower the regularity of the critical element to
$L^{\infty}\dot{H}^{s}$ for some $s<0$ by the double Duhamel argument in
dimension three to obtain the compactness of the critical element in $L^{2}$,
which is used to control the spatial center function $x(t)$ of the critical
element.
###### Remark 1.5.
We can remove the radial assumption under the stronger constraint that
$\displaystyle M(u_{0})+E(u_{0})<m,$
which can help us to obtain the compactness of the critical element in $L^{2}$
and control the spatial center function $x(t)$ of the critical element. Of
course, we need the precondition333By the relation between the sharp Sobolev
constant and the ground state $W$, we know that the constrained condition
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\big{|}\nabla
u_{0}\big{|}^{2}-\big{|}u_{0}\big{|}^{6}\right)\;dx\geq
0,\quad\int_{{\mathbb{R}}^{3}}\left(\frac{1}{2}\big{|}\nabla
u_{0}\big{|}^{2}-\frac{1}{6}\big{|}u_{0}\big{|}^{6}\right)\;dx<E^{c}(W)$ is
equivalent to the constrained condition $\displaystyle\big{\|}\nabla
u_{0}\big{\|}^{2}_{L^{2}}\leq\big{\|}\nabla
W\big{\|}^{2}_{L^{2}},\quad\int_{{\mathbb{R}}^{3}}\left(\frac{1}{2}\big{|}\nabla
u_{0}\big{|}^{2}-\frac{1}{6}\big{|}u_{0}\big{|}^{6}\right)\;dx<E^{c}(W).$ We
use the former in this paper while the latter is given by Kenig-Merle in [19].
that the global wellposedness and scattering result of (1.2) holds for
$u_{0}\in\dot{H}^{1}({\mathbb{R}}^{3})$ with
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\big{|}\nabla
u_{0}\big{|}^{2}-\big{|}u_{0}\big{|}^{6}\right)\;dx\geq
0,\quad\int_{{\mathbb{R}}^{3}}\left(\frac{1}{2}\big{|}\nabla
u_{0}\big{|}^{2}-\frac{1}{6}\big{|}u_{0}\big{|}^{6}\right)\;dx<$
$\displaystyle m.$
###### Remark 1.6.
From the assumption in Theorem 1.3, we know that the solution starts from the
following subsets of the energy space,
$\displaystyle\mathcal{K}^{+}=$ $\displaystyle\Big{\\{}\varphi\in
H^{1}({\mathbb{R}}^{3})\;\;\Big{|}\;\;\varphi\;\text{is
radial},\;E(\varphi)<m,\;K(\varphi)\geq 0\Big{\\}},$
$\displaystyle\mathcal{K}^{-}=$ $\displaystyle\Big{\\{}\varphi\in
H^{1}({\mathbb{R}}^{3})\;\;\Big{|}\;\;\varphi\;\text{is
radial},\;E(\varphi)<m,\;K(\varphi)<0\Big{\\}}.$
By the scaling argument, we know that $\mathcal{K}^{\pm}\not=\emptyset$ (we
can also know that $\mathcal{K}^{+}\not=\emptyset$ by the small data theory).
In fact, let $\chi(x)$ be a radial smooth cut-off function satisfying
$0\leq\chi\leq 1$, $\chi(x)=1$ for $|x|\leq 1$ and $\chi(x)=0$ for $|x|\geq
2$. If we take $\chi_{R}(x)=\chi(x/R)$ and
$\displaystyle\varphi(x)=\theta\lambda^{-1/2}\chi_{R}(x/\lambda)W(x/\lambda),$
where $\theta,\lambda,R$ is determined later and the cutoff function
$\chi_{R}$ is not needed for dimension $d\geq 5$ since $W\in H^{1}$. Then we
have
$\displaystyle\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}=$
$\displaystyle\theta^{2}\left(\big{\|}\nabla
W\big{\|}^{2}_{L^{2}}+\int\left((\chi_{R}^{2}-1)\big{|}\nabla
W\big{|}^{2}+|\nabla\chi_{R}|^{2}|W|^{2}+2\chi_{R}\nabla\chi_{R}\cdot W\nabla
W\right)\;dx\right),$ $\displaystyle\big{\|}\varphi\big{\|}^{6}_{L^{6}}=$
$\displaystyle\theta^{6}\left(\big{\|}W\big{\|}^{6}_{L^{6}}+\int(\chi_{R}^{6}-1)|W|^{6}\;dx\right),\quad\big{\|}\varphi\big{\|}^{4}_{L^{4}}=\lambda\cdot\theta^{4}\big{\|}\chi_{R}W\big{\|}^{4}_{L^{4}},$
$\displaystyle\big{\|}\varphi\big{\|}^{2}_{L^{2}}=$
$\displaystyle\lambda^{2}\cdot\theta^{2}\big{\|}\chi_{R}W\big{\|}^{2}_{L^{2}}.$
Therefore, taking $R$ sufficiently large, $\theta=1+\epsilon$ and
$\lambda=\epsilon^{3}$ , we have
$\displaystyle E(\varphi)=$ $\displaystyle\frac{\theta^{2}}{2}\big{\|}\nabla
W\big{\|}^{2}_{L^{2}}-\frac{\theta^{6}}{6}\big{\|}W\big{\|}^{6}_{L^{6}}$
$\displaystyle+\frac{\theta^{2}}{2}\int\left((\chi_{R}^{2}-1)\big{|}\nabla
W\big{|}^{2}+|\nabla\chi_{R}|^{2}|W|^{2}+2\chi_{R}\nabla\chi_{R}\cdot W\nabla
W\right)\;dx$
$\displaystyle-\frac{\theta^{6}}{6}\int(\chi_{R}^{6}-1)|W|^{6}\;dx+\lambda\cdot\frac{\theta^{4}}{4}\big{\|}\chi_{R}W\big{\|}^{4}_{L^{4}}$
$\displaystyle=$ $\displaystyle m-6\epsilon^{2}m+o(\epsilon^{2}),$
$\displaystyle K(\varphi)=$ $\displaystyle 2\theta^{2}\big{\|}\nabla
W\big{\|}^{2}_{L^{2}}-2\theta^{6}\big{\|}W\big{\|}^{6}_{L^{6}}$
$\displaystyle+2\theta^{2}\int\left((\chi_{R}^{2}-1)\big{|}\nabla
W\big{|}^{2}+|\nabla\chi_{R}|^{2}|W|^{2}+2\chi_{R}\nabla\chi_{R}\cdot W\nabla
W\right)\;dx$
$\displaystyle-2\theta^{6}\int(\chi_{R}^{6}-1)|W|^{6}\;dx+\lambda\cdot\frac{3\theta^{4}}{2}\big{\|}\chi_{R}W\big{\|}^{4}_{L^{4}}$
$\displaystyle=$ $\displaystyle-24\epsilon m+o(\epsilon^{2}).$
If taking $\epsilon<0$ and $|\epsilon|$ sufficient small, then we have
$\varphi\in\mathcal{K}^{+}$; If taking $\epsilon>0$ and sufficient small, then
we have $\varphi\in\mathcal{K}^{-}$.
### Acknowledgements.
The authors are partly supported by the NSF of China (No. 10801015, No.
10901148, No. 11171033). The authors would like to thank Professor K.
Nakanishi for his valuable communications. ∎
## 2\. Preliminaries
In this section, we give some notation and some wellknown results.
### 2.1. Littlewood-Paley decomposition and Besov space
Let $\Lambda_{0}(x)\in\mathcal{S}({\mathbb{R}}^{3})$ such that its Fourier
transform $\widetilde{\Lambda}_{0}(\xi)=1$ for $|\xi|\leq 1$ and
$\widetilde{\Lambda}_{0}(\xi)=0$ for $|\xi|\geq 2$. Then we define
$\Lambda_{k}(x)$ for any $k\in{\mathbb{Z}}\backslash\\{0\\}$ and
$\Lambda_{(0)}(x)$ by the Fourier transforms:
$\displaystyle\widetilde{\Lambda}_{k}(\xi)=\widetilde{\Lambda}_{0}(2^{-k}\xi)-\widetilde{\Lambda}_{0}(2^{-k+1}\xi),\quad\widetilde{\Lambda}_{(0)}(\xi)=\widetilde{\Lambda}_{0}(\xi)-\widetilde{\Lambda}_{0}(2\xi).$
Let $s\in{\mathbb{R}}$, $1\leq p,q\leq\infty$. The inhomogeneous Besov space
$B^{s}_{p,q}$ is defined by
$\displaystyle
B^{s}_{p,q}=\left\\{f\;\big{|}\;f\in\mathcal{S}^{\prime},\Big{\|}2^{ks}\big{\|}\Lambda_{k}*f\big{\|}_{L^{p}_{x}}\Big{\|}_{l^{q}_{k\geq
0}}<\infty\right\\},$
where $\mathcal{S}^{\prime}$ denotes the space of tempered distributions. The
homogeneous Besov space $\dot{B}^{s}_{p,q}$ can be defined by
$\displaystyle\dot{B}^{s}_{p,q}=\left\\{f\;\Big{|}\;f\in\mathcal{S}^{\prime},\left(\sum_{k\in{\mathbb{Z}}\backslash\\{0\\}}2^{qks}\big{\|}\Lambda_{k}*f\big{\|}^{q}_{L^{p}_{x}}+\big{\|}\Lambda_{(0)}*f\big{\|}_{L^{p}_{x}}\Big{\|}^{q}\right)^{1/q}<\infty\right\\}.$
### 2.2. Linear estimates
We say that a pair of exponents $(q,r)$ is Schröidnger
$\dot{H}^{s}$-admissible in dimension three if
$\displaystyle\dfrac{2}{q}+\dfrac{3}{r}=\dfrac{3}{2}-s$
and $2\leq q,r\leq\infty$. If $I\times{\mathbb{R}}^{3}$ is a space-time slab,
we define the $\dot{S}^{0}(I\times{\mathbb{R}}^{3})$ Strichartz norm by
$\displaystyle\big{\|}u\big{\|}_{\dot{S}^{0}(I\times{\mathbb{R}}^{3})}:=\sup\big{\|}u\big{\|}_{L^{q}_{t}L^{r}_{x}(I\times{\mathbb{R}}^{3})}$
where the sup is taken over all $L^{2}$-admissible pairs $(q,r)$. We define
the $\dot{S}^{s}(I\times{\mathbb{R}}^{3})$ Strichartz norm to be
$\displaystyle\big{\|}u\big{\|}_{\dot{S}^{s}(I\times{\mathbb{R}}^{3})}:=\big{\|}D^{s}u\big{\|}_{\dot{S}^{0}(I\times{\mathbb{R}}^{3})}.$
We also use $\dot{N}^{0}(I\times{\mathbb{R}}^{3})$ to denote the dual space of
$\dot{S}^{0}(I\times{\mathbb{R}}^{3})$ and
$\displaystyle\dot{N}^{k}(I\times{\mathbb{R}}^{3}):=\\{u;D^{k}u\in\dot{N}^{0}(I\times{\mathbb{R}}^{3})\\}.$
By definition and Sobolev’s inequality, we have
###### Lemma 2.1.
For any $\dot{S}^{1}$ function $u$ on $I\times{\mathbb{R}}^{3}$, we have
$\displaystyle\big{\|}\nabla
u\big{\|}_{L^{\infty}_{t}L^{2}_{x}}+\big{\|}u\big{\|}_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}(I\times{\mathbb{R}}^{3})}+\big{\|}u\big{\|}_{L^{\infty}_{t}L^{6}_{x}}+\big{\|}u\big{\|}_{L^{12}_{t}L^{9}_{x}}+\big{\|}u\big{\|}_{L^{10}_{t,x}}\lesssim\big{\|}u\big{\|}_{\dot{S}^{1}}.$
For any $\dot{S}^{1/2}$ function $u$ on $I\times{\mathbb{R}}^{3}$, we have
$\displaystyle\big{\|}u\big{\|}_{L^{\infty}_{t}\dot{H}^{1/2}_{x}}+\big{\|}u\big{\|}_{L^{6}_{t}\dot{B}^{1/2}_{18/7,2}(I\times{\mathbb{R}}^{3})}+\big{\|}u\big{\|}_{L^{\infty}_{t}L^{3}_{x}}+\big{\|}u\big{\|}_{L^{6}_{t}L^{9/2}_{x}}+\big{\|}u\big{\|}_{L^{5}_{t,x}}\lesssim\big{\|}u\big{\|}_{\dot{S}^{1/2}}.$
Now we state the standard Strichartz estimate.
###### Lemma 2.2 ([5, 18, 36]).
Let $I$ be a compact time interval, $k\in[0,1]$, and let
$u:I\times{\mathbb{R}}^{3}\rightarrow{\mathbb{C}}$ be an $\dot{S}^{k}$
solution to the forced Schrödinger equation
$\displaystyle iu_{t}+\Delta u=F$
for a function $F$. Then we have
$\displaystyle\big{\|}u\big{\|}_{\dot{S}^{k}(I\times{\mathbb{R}}^{3})}\lesssim\big{\|}u(t_{0})\big{\|}_{\dot{H}^{k}({\mathbb{R}}^{d})}+\big{\|}F\big{\|}_{\dot{N}^{k}(I\times{\mathbb{R}}^{3})},$
for any time $t_{0}\in I$.
We shall also need the following exotic Strichartz estimate, which is
important in the application of the stability theory.
###### Lemma 2.3 ([14]).
For any $F\in L^{2}_{t}\left(I;\dot{B}^{1/3}_{18/11,2}\right)$, we have
$\displaystyle\left\|\int^{t}_{0}e^{i(t-s)\Delta}F(s)\;ds\right\|_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}}\lesssim\big{\|}F\big{\|}_{L^{2}_{t}\dot{B}^{1/3}_{18/11,2}}.$
### 2.3. Local wellposedness and Virial identity
Let
$ST(I):=L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap L^{12}_{t}L^{9}_{x}\cap
L^{6}_{t}\dot{B}^{1/2}_{18/7,2}\cap L^{5}_{t,x}(I\times{\mathbb{R}}^{3}).$
By the definition of admissible pair, we know that
$L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap L^{12}_{t}L^{9}_{x}$ is the
$\dot{H}^{1}$-admissible space, $L^{6}_{t}\dot{B}^{1/2}_{18/7,2}\cap
L^{5}_{t,x}$ is the $\dot{H}^{1/2}$-admissible space. Now we have
###### Theorem 2.4 ([37]).
Let $u_{0}\in H^{1}$, then for every $\eta>0$, there exists $T=T(\eta)$ such
that if
$\displaystyle\big{\|}e^{it\Delta}u_{0}\big{\|}_{ST([-T,T])}\leq\eta,$
then (1.1) admits a unique strong $H^{1}_{x}$-solution $u$ defined on
$[-T,T]$. Let $(-T_{min},T_{max})$ be the maximal time interval on which $u$
is well-defined. Then, $u\in S^{1}(I\times{\mathbb{R}}^{d})$ for every compact
time interval $I\subset(-T_{min},T_{max})$ and the following properties hold:
1. (1)
If $T_{max}<\infty$, then
$\displaystyle\big{\|}u\big{\|}_{ST((0,T_{max})\times{\mathbb{R}}^{d})}=\infty.$
Similarly, if $T_{min}<\infty$, then
$\displaystyle\big{\|}u\big{\|}_{ST((-T_{min},0)\times{\mathbb{R}}^{d})}=\infty.$
2. (2)
The solution $u$ depends continuously on the initial data $u_{0}$ in the
following sense: The functions $T_{min}$ and $T_{max}$ are lower
semicontinuous from $\dot{H}^{1}_{x}\cap\dot{H}^{1/2}_{x}$ to $(0,+\infty]$.
Moreover, if $u^{(m)}_{0}\rightarrow u_{0}$ in
$\dot{H}^{1}_{x}\cap\dot{H}^{1/2}_{x}$ and $u^{(m)}$ is the maximal solution
to (1.1) with initial data $u^{(m)}_{0}$, then $u^{(m)}\rightarrow u$ in
$ST(I\times{\mathbb{R}}^{3})$ and every compact subinterval
$I\subset(-T_{min},T_{max})$.
###### Proof.
The proof is based on the Strichartz estimate and exotic Strichartz estimate
and the following nonlinear estimates.
$\displaystyle\big{\|}|u|^{4}u\big{\|}_{L^{2}\dot{B}^{1/3}_{18/11,2}}\lesssim\big{\|}u\big{\|}_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}}\big{\|}u\big{\|}^{4}_{L^{10}_{t,x}},\quad$
$\displaystyle\big{\|}|u|^{2}u\big{\|}_{L^{2}\dot{B}^{1/3}_{18/11,2}}\lesssim\big{\|}u\big{\|}_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}}\big{\|}u\big{\|}^{2}_{L^{5}_{t,x}},$
$\displaystyle\big{\|}|u|^{4}u\big{\|}_{L^{2}\dot{B}^{1/2}_{6/5,2}}\lesssim\big{\|}u\big{\|}_{L^{6}_{t}\dot{B}^{1/2}_{18/7,2}}\big{\|}u\big{\|}^{4}_{L^{12}_{t}L^{9}_{x}},\quad$
$\displaystyle\big{\|}|u|^{2}u\big{\|}_{L^{2}\dot{B}^{1/3}_{6/5,2}}\lesssim\big{\|}u\big{\|}_{L^{6}_{t}\dot{B}^{1/2}_{18/7,2}}\big{\|}u\big{\|}^{2}_{L^{6}_{t}L^{9/2}_{x}}.$
∎
###### Lemma 2.5.
Let $\phi\in C^{\infty}_{0}({\mathbb{R}}^{3})$, radially symmetric and $u$ be
the radial solution of (1.1). Then we have
$\displaystyle\partial_{t}\int_{{\mathbb{R}}^{3}}\phi(x)\big{|}u(t,x)\big{|}^{2}\;dx=$
$\displaystyle-2\Im\int_{{\mathbb{R}}^{3}}\nabla\phi\cdot\nabla\bar{u}\;u\;dx$
$\displaystyle\partial^{2}_{t}\int_{{\mathbb{R}}^{3}}\phi(x)\big{|}u(t,x)\big{|}^{2}\;dx=$
$\displaystyle 4\int_{{\mathbb{R}}^{3}}\phi^{\prime\prime}(r)\big{|}\nabla
u\big{|}^{2}\;dx-\int_{{\mathbb{R}}^{3}}\Delta^{2}\phi\big{|}u(t,x)\big{|}^{2}\;dx$
$\displaystyle-\frac{4}{3}\int_{{\mathbb{R}}^{3}}\Delta\phi\big{|}u(t,x)\big{|}^{6}\;dx+\int_{{\mathbb{R}}^{3}}\Delta\phi\big{|}u(t,x)\big{|}^{4}\;dx,$
where $r=|x|$.
###### Proof.
By the simple computation, we have
$\displaystyle\partial^{2}_{t}\int_{{\mathbb{R}}^{3}}\phi(x)\big{|}u(t,x)\big{|}^{2}\;dx=$
$\displaystyle
4\int_{{\mathbb{R}}^{3}}\phi_{jk}\cdot\Re(u_{k}\overline{u}_{j})\;dx-\int_{{\mathbb{R}}^{3}}\Delta^{2}\phi\cdot\big{|}u(t,x)\big{|}^{2}\;dx$
$\displaystyle-\frac{4}{3}\int_{{\mathbb{R}}^{3}}\Delta\phi\cdot\big{|}u(t,x)\big{|}^{6}\;dx+\int_{{\mathbb{R}}^{3}}\Delta\phi\cdot\big{|}u(t,x)\big{|}^{4}\;dx.$
Then the result comes from the following fact
$\displaystyle\partial^{2}_{jk}\phi(x)=\phi^{\prime\prime}(r)\frac{x_{j}x_{k}}{r^{2}}+\frac{\phi^{\prime}(r)}{r}\left(\delta_{jk}-\frac{x_{j}x_{k}}{r^{2}}\right)$
holds for any radial symmetric function $\phi(x)$. ∎
### 2.4. Variational characterization
In this subsection, we give the threshold energy $m$ (Proposition 1.1) by the
variational method, and various estimates for the solutions of (1.1) with the
energy below the threshold. There is no the radial assumption on the solution.
We first give some notation before we show the behavior of $K$ near the
origin. Let us denote the quadratic and nonlinear parts of $K$ by $K^{Q}$ and
$K^{N}$, that is,
$\displaystyle K(\varphi)=K^{Q}(\varphi)+K^{N}(\varphi),$
where $K^{Q}(\varphi)=\displaystyle
2\;\int_{{\mathbb{R}}^{3}}|\nabla\varphi|^{2}\;dx,$ and
$K^{N}(\varphi)=\displaystyle\int_{{\mathbb{R}}^{3}}\left(-2|\varphi|^{6}+\frac{3}{2}\varphi|^{4}\right)\;dx$.
###### Lemma 2.6.
For any $\varphi\in H^{1}({\mathbb{R}}^{3})$, we have
$\displaystyle\lim_{\lambda\rightarrow-\infty}K^{Q}(\varphi^{\lambda}_{3,-2})=0.$
(2.1)
###### Proof.
It is obvious by the definition of $K^{Q}$. ∎
Now we show the positivity of $K$ near 0 in the energy space.
###### Lemma 2.7.
For any bounded sequence $\varphi_{n}\in
H^{1}({\mathbb{R}}^{3})\backslash\\{0\\}$ with
$\displaystyle\lim_{n\rightarrow+\infty}K^{Q}(\varphi_{n})=0,$
then for large $n$, we have
$\displaystyle K(\varphi_{n})>0.$
###### Proof.
By the fact that $K^{Q}(\varphi_{n})\rightarrow 0$, we know that
$\displaystyle\lim_{n\rightarrow+\infty}\big{\|}\nabla\varphi_{n}\big{\|}^{2}_{L^{2}}=0.$
Then by the Sobolev and Gagliardo-Nirenberg inequalities, we have for large
$n$
$\displaystyle\big{\|}\varphi_{n}\big{\|}^{6}_{L^{6}_{x}}\lesssim\big{\|}\nabla\varphi_{n}\big{\|}^{6}_{L^{2}_{x}}=$
$\displaystyle o(\big{\|}\nabla\varphi_{n}\big{\|}^{2}_{L^{2}}),$
$\displaystyle\big{\|}\varphi_{n}\big{\|}^{4}_{L^{4}_{x}}\lesssim\big{\|}\varphi_{n}\big{\|}_{L^{2}}\big{\|}\nabla\varphi_{n}\big{\|}^{3}_{L^{2}}$
$\displaystyle=o(\big{\|}\nabla\varphi_{n}\big{\|}^{2}_{L^{2}}),$
where we use the boundedness of $\big{\|}\varphi_{n}\big{\|}_{L^{2}}$. Hence
for large $n$, we have
$\displaystyle K(\varphi_{n})=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(2|\nabla\varphi_{n}|^{2}-2|\varphi_{n}|^{6}+\frac{3}{2}|\varphi_{n}|^{4}\right)\;dx\thickapprox\int_{{\mathbb{R}}^{3}}|\nabla\varphi_{n}|^{2}\;dx>0.$
This concludes the proof. ∎
By the definition of $K$, we denote two real numbers by
$\displaystyle\bar{\mu}=\max\\{4,0,6\\}=6,\quad\underline{\mu}=\min\\{4,0,6\\}=0.$
Next, we show the behavior of the scaling derivative functional $K$.
###### Lemma 2.8.
For any $\varphi\in H^{1}$, we have
$\displaystyle\left(\bar{\mu}-\mathcal{L}\right)E(\varphi)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\big{|}\nabla\varphi\big{|}^{2}+\big{|}\varphi\big{|}^{6}\right)\;dx,$
$\displaystyle\mathcal{L}\left(\bar{\mu}-\mathcal{L}\right)E(\varphi)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(4\big{|}\nabla\varphi\big{|}^{2}+12\big{|}\varphi\big{|}^{6}\right)\;dx.$
###### Proof.
By the definition of $\mathcal{L}$, we have
$\displaystyle\mathcal{L}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}=4\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}},\quad\mathcal{L}\big{\|}\varphi\big{\|}^{6}_{L^{6}}=12\big{\|}\varphi\big{\|}^{6}_{L^{6}},\quad\mathcal{L}\big{\|}\varphi\big{\|}^{4}_{L^{4}}=6\big{\|}\varphi\big{\|}^{4}_{L^{4}},$
which implies that
$\displaystyle\left(\bar{\mu}-\mathcal{L}\right)E(\varphi)=6E(\varphi)-K(\varphi)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\big{|}\nabla\varphi\big{|}^{2}+\big{|}\varphi\big{|}^{6}\right)\;dx,$
$\displaystyle\mathcal{L}\left(\bar{\mu}-\mathcal{L}\right)E(\varphi)=\mathcal{L}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}+\mathcal{L}\big{\|}\varphi\big{\|}^{6}_{L^{6}}$
$\displaystyle=\int_{{\mathbb{R}}^{3}}\left(4\big{|}\nabla\varphi\big{|}^{2}+12\big{|}\varphi\big{|}^{6}\right)\;dx.$
This completes the proof. ∎
According to the above analysis, we will replace the functional $E$ in (1.5)
with a positive functional $H$, while extending the minimizing region from
“$K(\varphi)=0$” to “$K(\varphi)\leq 0$”. Let
$\displaystyle
H(\varphi):=\left(1-\frac{\mathcal{L}}{\bar{\mu}}\right)E(\varphi)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\frac{1}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{1}{6}\big{|}\varphi\big{|}^{6}\right)\;dx,$
then for any $\varphi\in H^{1}\backslash\\{0\\}$, we have
$\displaystyle H(\varphi)>0,\quad\mathcal{L}H(\varphi)\geq 0.$
Now we can characterization the minimization problem (1.5) by use of $H$.
###### Lemma 2.9.
For the minimization $m$ in (1.5), we have
$\displaystyle m=$ $\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)\leq 0\\}$
$\displaystyle=$ $\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}.$ (2.2)
###### Proof.
For any $\varphi\in H^{1}$, $\varphi\not=0$ with $K(\varphi)=0$, we have
$E(\varphi)=H(\varphi)$, this implies that
$\displaystyle m=$ $\displaystyle\inf\\{E(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)=0\\}$ $\displaystyle\geq$
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)\leq 0\\}.$ (2.3)
On the other hand, for any $\varphi\in H^{1}$, $\varphi\not=0$ with
$K(\varphi)<0$, by Lemma 2.6, Lemma 2.7 and the continuity of $K$ in
$\lambda$, we know that there exists a $\lambda_{0}<0$ such that
$\displaystyle K(\varphi^{\lambda_{0}}_{3,-2})=0,$
then by $\mathcal{L}H\geq 0$, we have
$\displaystyle
E(\varphi^{\lambda_{0}}_{3,-2})=H(\varphi^{\lambda_{0}}_{3,-2})\leq
H(\varphi^{0}_{3,-2})=H(\varphi).$
Therefore,
$\displaystyle\inf\\{E(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)=0\\}$
$\displaystyle\leq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}.$ (2.4)
By (2.3) and (2.4), we have
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)\leq 0\\}$
$\displaystyle\leq m\leq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}.$
In order to show (2.2), it suffices to show that
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)\leq 0\\}$
$\displaystyle\geq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}.$ (2.5)
For any $\varphi\in H^{1}$, $\varphi\not=0$ with $K(\varphi)\leq 0$. By Lemma
2.8, we know that
$\displaystyle\mathcal{L}K(\varphi)=\bar{\mu}K(\varphi)-\int_{{\mathbb{R}}^{3}}\left(4\big{|}\nabla\varphi\big{|}^{2}+12\big{|}\varphi\big{|}^{6}\right)\;dx<0,$
then for any $\lambda>0$ we have
$\displaystyle K(\varphi^{\lambda}_{3,-2})<0,$
and as $\lambda\rightarrow 0$
$\displaystyle
H(\varphi^{\lambda}_{3,-2})=\int_{{\mathbb{R}}^{3}}\left(\frac{e^{4\lambda}}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{e^{12\lambda}}{6}\big{|}\varphi\big{|}^{6}\right)\;dx\longrightarrow
H(\varphi).$
This shows (2.5), and completes the proof. ∎
Next we will use the ($\dot{H}^{1}$-invariant) scaling argument to remove the
$L^{4}$ term (the lower regularity quantity than $\dot{H}^{1}$) in $K$, that
is, to replace the constrained condition $K(\varphi)<0$ with
$K^{c}(\varphi)<0$, where
$\displaystyle
K^{c}(\varphi):=\int_{{\mathbb{R}}^{3}}\left(2|\nabla\varphi|^{2}-2|\varphi|^{6}\right)\;dx.$
In fact, we have
###### Lemma 2.10.
For the minimization $m$ in (1.5), we have
$\displaystyle m=$ $\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)<0\\}$
$\displaystyle=$ $\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)\leq 0\\}.$
###### Proof.
Since $K^{c}(\varphi)\leq K(\varphi)$, it is obvious that
$\displaystyle m=$ $\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}$ $\displaystyle\geq$
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)<0\\}.$
Hence in order to show the first equality, it suffices to show that
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)<0\\}$
$\displaystyle\leq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)<0\\}.$ (2.6)
To do so, for any $\varphi\in H^{1}$, $\varphi\not=0$ with $K^{c}(\varphi)<0$,
taking
$\displaystyle\varphi^{\lambda}_{1,-2}(x)=e^{\lambda}\varphi(e^{2\lambda}x),$
we have $\varphi^{\lambda}_{1,-2}\in H^{1}$ and
$\varphi^{\lambda}_{1,-2}\not=0$ for any $\lambda>0$. In addition, we have
$\displaystyle
K(\varphi^{\lambda}_{1,-2})=\int_{{\mathbb{R}}^{3}}\left(2\big{|}\nabla\varphi\big{|}^{2}-2\big{|}\varphi\big{|}^{6}+\frac{3}{2}e^{-2\lambda}\big{|}\varphi\big{|}^{4}\right)\;dx$
$\displaystyle\longrightarrow K^{c}(\varphi),$ $\displaystyle
H(\varphi^{\lambda}_{1,-2})=\int_{{\mathbb{R}}^{3}}\left(\frac{1}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{1}{6}\big{|}\varphi\big{|}^{6}\right)\;dx=$
$\displaystyle H(\varphi),$
as $\lambda\rightarrow+\infty$. This gives (2.6), and completes the proof of
the first equality.
For the second equality, it is obvious that
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)<0\\}$
$\displaystyle\geq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)\leq 0\\},$
hence we only need to show that
$\displaystyle\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)<0\\}$
$\displaystyle\leq\inf\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K^{c}(\varphi)\leq 0\\}.$ (2.7)
To do this, we use the ($L^{2}$-invariant) scaling argument. For any
$\varphi\in H^{1}$, $\varphi\not=0$ with $K^{c}(\varphi)\leq 0$, we have
$\varphi^{\lambda}_{3,-2}\in H^{1}$, $\varphi^{\lambda}_{3,-2}\not=0$. In
addition, by
$\displaystyle\mathcal{L}K^{c}(\varphi)=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(8\big{|}\nabla\varphi\big{|}^{2}-24\big{|}\varphi\big{|}^{6}\right)\;dx=4K^{c}(\varphi)-16\big{\|}\varphi\big{\|}^{6}_{L^{6}}<0,$
$\displaystyle
H(\varphi^{\lambda}_{3,-2})=\int_{{\mathbb{R}}^{3}}\left(\frac{e^{4\lambda}}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{e^{12\lambda}}{6}\big{|}\varphi\big{|}^{6}\right)\;dx,$
we have $K^{c}(\varphi^{\lambda}_{3,-2})<0$ for any $\lambda>0$, and
$\displaystyle H(\varphi^{\lambda}_{3,-2})\rightarrow
H(\varphi),\;\;\text{as}\;\;\lambda\rightarrow 0.$
This implies (2.7) and completes the proof. ∎
After these preparations, we can now make use of the sharp Sobolev constant in
[1, 35] to compute the minimization $m$ of (1.5), which also shows Proposition
1.1.
###### Lemma 2.11.
For the minimization $m$ in (1.5), we have
$\displaystyle m=E^{c}(W).$
###### Proof.
By Lemma 2.10, we have
$\displaystyle m=$
$\displaystyle\inf\left\\{\frac{1}{6}\int_{{\mathbb{R}}^{3}}\left(|\nabla\varphi|^{2}+|\varphi|^{6}\right)\;dx\;\Big{|}\;\varphi\in
H^{1},\;\varphi\not=0,\;\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}\leq\big{\|}\varphi\big{\|}^{6}_{L^{6}}\right\\}$
$\displaystyle\geq$
$\displaystyle\inf\left\\{\int_{{\mathbb{R}}^{3}}\frac{1}{6}\left(|\nabla\varphi|^{2}+|\varphi|^{6}\right)+\frac{1}{6}\left(|\nabla\varphi|^{2}-|\varphi|^{6}\right)\;dx\;\Big{|}\;\varphi\in
H^{1},\;\varphi\not=0,\;\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}\leq\big{\|}\varphi\big{\|}^{6}_{L^{6}}\right\\}$
where the equality holds if and only if the minimization is taken by some
$\varphi$ with
$\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}=\big{\|}\varphi\big{\|}^{6}_{L^{6}}$.
While
$\displaystyle\inf\left\\{\int_{{\mathbb{R}}^{3}}\frac{1}{3}|\nabla\varphi|^{2}\;dx\;\big{|}\;\varphi\in
H^{1},\;\varphi\not=0,\;\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}\leq\big{\|}\varphi\big{\|}^{6}_{L^{6}}\right\\}$
$\displaystyle=\inf\left\\{\frac{1}{3}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}\left(\frac{\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}}{\big{\|}\varphi\big{\|}^{6}_{L^{6}}}\right)^{1/2}\;\Big{|}\;\varphi\in
H^{1},\;\varphi\not=0\right\\}$
$\displaystyle=\inf\left\\{\frac{1}{3}\left(\frac{\big{\|}\nabla\varphi\big{\|}_{L^{2}}}{\big{\|}\varphi\big{\|}_{L^{6}}}\right)^{3}\;\Big{|}\;\varphi\in
H^{1},\;\varphi\not=0\right\\}$
$\displaystyle=\inf\left\\{\frac{1}{3}\left(\frac{\big{\|}\nabla\varphi\big{\|}_{L^{2}}}{\big{\|}\varphi\big{\|}_{L^{6}}}\right)^{3}\;\Big{|}\;\varphi\in\dot{H}^{1},\;\varphi\not=0\right\\}=\frac{1}{3}\big{(}C^{*}_{3}\big{)}^{-3}.$
where we use the density property $H^{1}\hookrightarrow\dot{H}^{1}$ in the
last second equality and that $C^{*}_{3}$ is the sharp Sobolev constant in
${\mathbb{R}}^{3}$, that is,
$\displaystyle\big{\|}\varphi\big{\|}_{L^{6}_{x}}\leq
C^{*}_{3}\big{\|}\nabla\varphi\big{\|}_{L^{2}_{x}},\;\;\forall\;\varphi\in\dot{H}^{1}({\mathbb{R}}^{3}),$
and the equality can be attained by the ground state $W$ of the following
elliptic equation
$\displaystyle-\Delta W=|W|^{4}W.$
This implies that $\frac{1}{3}\big{(}C^{*}_{3}\big{)}^{-3}=E^{c}(W)$. The
proof is completed. ∎
After the computation of the minimization $m$ in (1.5), we next give some
variational estimates.
###### Lemma 2.12.
For any $\varphi\in H^{1}$ with $K(\varphi)\geq 0$, we have
$\displaystyle\int_{{\mathbb{R}}^{3}}\left(\frac{1}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{1}{6}\big{|}\varphi\big{|}^{6}\right)dx\leq
E(\varphi)\leq\int_{{\mathbb{R}}^{3}}\left(\frac{1}{2}\big{|}\nabla\varphi\big{|}^{2}+\frac{1}{4}\big{|}\varphi\big{|}^{4}\right)dx.$
(2.8)
###### Proof.
On one hand, the right hand side of (2.8) is trivial. On the other hand, by
the definition of $E$ and $K$, we have
$\displaystyle
E(\varphi)=\int_{{\mathbb{R}}^{3}}\left(\frac{1}{6}\big{|}\nabla\varphi\big{|}^{2}+\frac{1}{6}\big{|}\varphi\big{|}^{6}\right)\;dx+\frac{1}{6}K(\varphi),$
which implies the left hand side of (2.8). ∎
At the last of this section, we give the uniform bounds on the scaling
derivative functional $K(\varphi)$ with the energy $E(\varphi)$ below the
threshold $m$, which plays an important role for the blow-up and scattering
analysis in Section 3 and Section 6.
###### Lemma 2.13.
For any $\varphi\in H^{1}$ with $E(\varphi)<m$.
1. (1)
If $K(\varphi)<0$, then
$\displaystyle K(\varphi)\leq-6\big{(}m-E(\varphi)\big{)}.$ (2.9)
2. (2)
If $K(\varphi)\geq 0$, then
$\displaystyle
K(\varphi)\geq\min\left(6(m-E(\varphi)),\frac{2}{3}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}+\frac{1}{2}\big{\|}\varphi\big{\|}^{4}_{L^{4}}\right).$
(2.10)
###### Proof.
By Lemma 2.8, for any $\varphi\in H^{1}$, we have
$\displaystyle\mathcal{L}^{2}E(\varphi)=\bar{\mu}\mathcal{L}E(\varphi)-4\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}-12\big{\|}\varphi\big{\|}^{6}_{L^{6}}.$
Let $j(\lambda)=E(\varphi^{\lambda}_{3,-2})$, then we have
$\displaystyle
j^{\prime\prime}(\lambda)=\bar{\mu}j^{\prime}(\lambda)-4e^{4\lambda}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}-12e^{12\lambda}\big{\|}\varphi\big{\|}^{6}_{L^{6}}.$
(2.11)
Case I: If $K(\varphi)<0$, then by (2.1), Lemma 2.7 and the continuity of $K$
in $\lambda$, there exists a negative number $\lambda_{0}<0$ such that
$K(\varphi^{\lambda_{0}}_{3,-2})=0$, and
$\displaystyle
K(\varphi^{\lambda}_{3,-2})<0,\;\;\forall\;\;\lambda\in(\lambda_{0},0).$
By (1.5), we obtain $j(\lambda_{0})=E(\varphi^{\lambda_{0}}_{3,-2})\geq m$.
Now by integrating (2.11) over $[\lambda_{0},0]$, we have
$\displaystyle\int^{0}_{\lambda_{0}}j^{\prime\prime}(\lambda)\;d\lambda\leq\bar{\mu}\int^{0}_{\lambda_{0}}j^{\prime}(\lambda)\;d\lambda,$
which implies that
$\displaystyle
K(\varphi)=j^{\prime}(0)-j^{\prime}(\lambda_{0})\leq\bar{\mu}\left(j(0)-j(\lambda_{0})\right)\leq-\bar{\mu}\big{(}m-E(\varphi)\big{)},$
which implies (2.9).
Case II: $K(\varphi)\geq 0$. We divide it into two subcases:
When $2\bar{\mu}K(\varphi)\geq 12\big{\|}\varphi\big{\|}^{6}_{L^{6}}$. Since
$\displaystyle
12\int_{{\mathbb{R}}^{3}}\big{|}\varphi\big{|}^{6}\;dx=-6K(\varphi)+\int_{{\mathbb{R}}^{3}}\left(12\big{|}\nabla\varphi\big{|}^{2}+9\big{|}\varphi\big{|}^{4}\right)\;dx,$
then we have
$\displaystyle
2\bar{\mu}K(\varphi)\geq-6K(\varphi)+\int_{{\mathbb{R}}^{3}}\left(12\big{|}\nabla\varphi\big{|}^{2}+9\big{|}\varphi\big{|}^{4}\right)\;dx,$
which implies that
$\displaystyle
K(\varphi)\geq\frac{2}{3}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}+\frac{1}{2}\big{\|}\varphi\big{\|}^{4}_{L^{4}}.$
When $2\bar{\mu}K(\varphi)\leq 12\big{\|}\varphi\big{\|}^{6}_{L^{6}}$. By
(2.11), we have for $\lambda=0$
$\displaystyle 0<$ $\displaystyle
2\bar{\mu}j^{\prime}(\lambda)<12e^{12\lambda}\big{\|}\varphi\big{\|}^{6}_{L^{6}},$
$\displaystyle j^{\prime\prime}(\lambda)=\bar{\mu}j^{\prime}(\lambda)-$
$\displaystyle
4e^{4\lambda}\big{\|}\nabla\varphi\big{\|}^{2}_{L^{2}}-12e^{12\lambda}\big{\|}\varphi\big{\|}^{6}_{L^{6}}\leq-\bar{\mu}j^{\prime}(\lambda).$
(2.12)
By the continuity of $j^{\prime}$ and $j^{\prime\prime}$ in $\lambda$, we know
that $j^{\prime}$ is an accelerating decreasing function as $\lambda$
increases until $j^{\prime}(\lambda_{0})=0$ for some finite number
$\lambda_{0}>0$ and (2.12) holds on $[0,\lambda_{0}]$.
By $K(\varphi^{\lambda_{0}}_{3,-2})=j^{\prime}(\lambda_{0})=0,$ we know that
$\displaystyle E(\varphi^{\lambda_{0}}_{3,-2})\geq m.$
Now integrating (2.12) over $[0,\lambda_{0}]$, we obtain that
$\displaystyle-K(\varphi)=j^{\prime}(\lambda_{0})-j^{\prime}(0)\leq-\bar{\mu}\big{(}j(\lambda_{0})-j(0)\big{)}\leq-\bar{\mu}(m-E(\varphi)).$
This completes the proof. ∎
## 3\. Part I: Blow up for $\mathcal{K}^{-}$
In this section, we prove the blow-up result of Theorem 1.3. We can also refer
to [33]. Now let $\phi$ be a smooth, radial function satisfying
$\partial^{2}_{r}\phi(r)\leq 2$, $\phi(r)=r^{2}$ for $r\leq 1$, and $\phi(r)$
is constant for $r\geq 3$. For some $R$, we define
$\displaystyle
V_{R}(t):=\int_{{\mathbb{R}}^{3}}\phi_{R}(x)|u(t,x)|^{2}\;dx,\quad\phi_{R}(x)=R^{2}\phi\left(\frac{|x|}{R}\right).$
By Lemma 2.5, $\Delta\phi_{R}(r)=6$ for $r\leq R,$ and
$\Delta^{2}\phi_{R}(r)=0$ for $r\leq R,$ we have
$\displaystyle\partial^{2}_{t}V_{R}(t)=$
$\displaystyle\;4\int_{{\mathbb{R}}^{3}}\phi_{R}^{\prime\prime}(r)\big{|}\nabla
u(t,x)\big{|}^{2}\;dx-\int_{{\mathbb{R}}^{3}}(\Delta^{2}\phi_{R})(x)|u(t,x)|^{2}\;dx$
$\displaystyle-\frac{4}{3}\int_{{\mathbb{R}}^{3}}(\Delta\phi_{R})|u(t,x)|^{6}\;dx+\int_{{\mathbb{R}}^{3}}(\Delta\phi_{R})|u(t,x)|^{4}\;dx$
$\displaystyle\leq$ $\displaystyle\;4\int_{{\mathbb{R}}^{3}}\left(2|\nabla
u(t)|^{2}-2|u(t)|^{6}+\frac{3}{2}|u(t)|^{4}\right)\;dx$
$\displaystyle+\frac{c}{R^{2}}\int_{R\leq|x|\leq
3R}\big{|}u(t)\big{|}^{2}\;dx+c\int_{R\leq|x|\leq
3R}\left(\big{|}u(t)\big{|}^{4}+\big{|}u(t)\big{|}^{6}\right)\;dx.$
By the Gagliardo-Nirenberg and radial Sobolev inequalities, we have
$\displaystyle\big{\|}f\big{\|}^{4}_{L^{4}(|x|\geq R)}\leq$
$\displaystyle\frac{c}{R^{2}}\big{\|}f\big{\|}^{3}_{L^{2}(|x|\geq
R)}\big{\|}\nabla f\big{\|}_{L^{2}(|x|\geq R)},$
$\displaystyle\big{\|}f\big{\|}_{L^{\infty}(|x|\geq R)}\leq$
$\displaystyle\frac{c}{R}\big{\|}f\big{\|}^{1/2}_{L^{2}(|x|\geq
R)}\big{\|}\nabla f\big{\|}^{1/2}_{L^{2}(|x|\geq R)}.$
Therefore, by mass conservation and Young’s inequality, we know that for any
$\epsilon>0$ there exist sufficiently large $R$ such that
$\displaystyle\partial^{2}_{t}V_{R}(t)\leq$ $\displaystyle
4K(u(t))+\epsilon\big{\|}\nabla u(t,x)\big{\|}^{2}_{L^{2}}+\epsilon^{2}.$
$\displaystyle=$ $\displaystyle 48E(u)-\big{(}16-\epsilon\big{)}\big{\|}\nabla
u(t)\big{\|}^{2}_{L^{2}}-6\big{\|}u(t)\big{\|}^{4}_{L^{4}}+\epsilon^{2}$ (3.1)
By $K(u)<0$, mass and energy conservations, Lemma 2.13 and the continuity
argument, we know that for any $t\in I$, we have
$\displaystyle K(u(t))\leq-6\left(m-E(u(t))\right)<0.$
By Lemma 2.9, we have
$\displaystyle m\leq H(u(t))<\frac{1}{3}\big{\|}u(t)\big{\|}^{6}_{L^{6}}.$
where we have used the fact that $K(u(t))<0$ in the second inequality. By the
fact $m=\frac{1}{3}\left(C^{*}_{3}\right)^{-3}$ and the Sharp Sobolev
inequality, we have
$\displaystyle\big{\|}\nabla
u(t)\big{\|}^{6}_{L^{2}}\geq\left(C^{*}_{3}\right)^{-6}\big{\|}u(t)\big{\|}^{6}_{L^{6}}>\left(C^{*}_{3}\right)^{-9},$
which implies that $\big{\|}\nabla u(t)\big{\|}^{2}_{L^{2}}>3m$.
In addition, by $E(u_{0})<m$ and energy conservation, there exists
$\delta_{1}>0$ such that $E(u(t))\leq(1-\delta_{1})m$. Thus, if we choose
$\epsilon$ sufficiently small, we have
$\displaystyle\partial^{2}_{t}V_{R}(t)\leq
48(1-\delta_{1})m-3\big{(}16-\epsilon\big{)}m+\epsilon^{2}\leq-24\delta_{1}m,$
which implies that $u$ must blow up at finite time. ∎
## 4\. Perturbation theory
In this part, we give the perturbation theory of the solution of (1.1) with
the global space-time estimate. First we denote the space-time space $ST(I)$
on the time interval $I$ by
$\displaystyle ST(I):=\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\cap L^{6}_{t}\dot{B}^{1/2}_{18/7,2}\cap
L^{5}_{t,x}\right)$ $\displaystyle(I\times{\mathbb{R}}^{3}),$ $\displaystyle
ST^{*}(I):=\left(L^{2}_{t}\dot{B}^{1/3}_{18/11,2}\cap
L^{2}_{t}\dot{B}^{1/2}_{6/5,2}\right)(I\times{\mathbb{R}}^{3})$ .
The main result in this section is the following.
###### Proposition 4.1.
Let $I$ be a compact time interval and let $w$ be an approximate solution to
(1.1) on $I\times{\mathbb{R}}^{3}$ in the sense that
$\displaystyle i\partial_{t}w+\Delta w=-|w|^{4}w+|w|^{2}w+e$
for some suitable small function $e$. Assume that for some constants
$L,E_{0}>0$, we have
$\displaystyle\big{\|}w\big{\|}_{ST(I)}\leq
L,\quad\big{\|}w(t_{0})\big{\|}_{H^{1}_{x}({\mathbb{R}}^{3})}\leq E_{0}$
for some $t_{0}\in I$. Let $u(t_{0})$ close to $w(t_{0})$ in the sense that
for some $E^{\prime}>0$, we have
$\displaystyle\big{\|}u(t_{0})-w(t_{0})\big{\|}_{H^{1}_{x}}\leq E^{\prime}.$
Assume also that for some $\varepsilon$, we have
$\displaystyle\left\|e^{i(t-t_{0})\Delta}\big{(}u(t_{0})-w(t_{0})\big{)}\right\|_{ST(I)}$
$\displaystyle\leq\varepsilon,\quad\big{\|}e\big{\|}_{ST^{*}(I)}\leq\varepsilon,$
(4.1)
where $0<\varepsilon\leq\varepsilon_{0}=\varepsilon_{0}(E_{0},E^{\prime},L)$
is a small constant. Then there exists a solution $u$ to (1.1) on
$I\times{\mathbb{R}}^{3}$ with initial data $u(t_{0})$ at time $t=t_{0}$
satisfying
$\displaystyle\big{\|}u-w\big{\|}_{ST(I)}\leq$ $\displaystyle
C(E_{0},E^{\prime},L)\;\varepsilon,\quad\text{and}\quad\big{\|}u\big{\|}_{ST(I)}\leq
C(E_{0},E^{\prime},L).$
###### Proof.
Since $w\in ST(I)$, there exists a partition of the right half of $I$ at
$t_{0}$:
$\displaystyle t_{0}<t_{1}<\cdots<t_{N},\quad I_{j}=(t_{j},t_{j+1}),\quad
I\cap(t_{0},\infty)=(t_{0},t_{N}),$
such that $N\leq C(L,\delta)$ and for any $j=0,1,\ldots,N-1$, we have
$\displaystyle\big{\|}w\big{\|}_{ST(I_{j})}\leq\delta\ll 1.$ (4.2)
The estimate on the left half of $I$ at $t_{0}$ is analogue, we omit it.
Let
$\displaystyle\gamma(t,x)=$ $\displaystyle u(t,x)-w(t,x),$
$\displaystyle\gamma_{j}(t,x)=$ $\displaystyle
e^{i(t-t_{j})\Delta}\Big{(}u(t_{j},x)-w(t_{j},x)\Big{)},$
then $\gamma$ satisfies the following difference equation
$\displaystyle
i\gamma_{t}+\Delta\gamma=O(w^{4}\gamma+w^{3}\gamma^{2}+w^{2}\gamma^{3}+w\gamma^{4}+\gamma^{5}+w^{2}\gamma+w\gamma^{2}+\gamma^{3})-e,$
which implies that
$\displaystyle\gamma(t)=$
$\displaystyle\gamma_{j}(t)-i\int^{t}_{t_{j}}e^{i(t-s)\Delta}\Big{(}O(w^{4}\gamma+w^{3}\gamma^{2}+w^{2}\gamma^{3}+w\gamma^{4}+\gamma^{5}+w^{2}\gamma+w\gamma^{2}+\gamma^{3})-e\Big{)}\;ds,$
$\displaystyle\gamma_{j+1}(t)=$
$\displaystyle\gamma_{j}(t)-i\int^{t_{j+1}}_{t_{j}}e^{i(t-s)\Delta}\Big{(}O(w^{4}\gamma+w^{3}\gamma^{2}+w^{2}\gamma^{3}+w\gamma^{4}+\gamma^{5}+w^{2}\gamma+w\gamma^{2}+\gamma^{3})-e\Big{)}\;ds.$
By Lemma 2.2, we have
$\displaystyle\big{\|}\gamma-\gamma_{j}\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)\cap
L^{5}_{t,x}\left(I_{j}\right)}+\big{\|}\gamma_{j+1}-\gamma_{j}\big{\|}_{L^{6}_{t}\left({\mathbb{R}};\dot{B}^{1/2}_{18/7,2}\right)\cap
L^{5}_{t,x}\left({\mathbb{R}}\times{\mathbb{R}}^{3}\right)}$ (4.3)
$\displaystyle\lesssim\big{\|}O(w^{4}\gamma+w^{3}\gamma^{2}+w^{2}\gamma^{3}+w\gamma^{4}+\gamma^{5}\big{\|}_{L^{2}_{t}\left(I_{j};\dot{B}^{1/2}_{6/5,2}\right)}$
$\displaystyle\quad+\big{\|}w^{2}\gamma+w\gamma^{2}+\gamma^{3})\big{\|}_{L^{2}_{t}\left(I_{j};\dot{B}^{1/2}_{6/5,2}\right)}+\big{\|}e\big{\|}_{L^{2}_{t}\left(I_{j};\dot{B}^{1/2}_{6/5,2}\right)}$
$\displaystyle\lesssim\big{\|}w\big{\|}^{4}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}^{3}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}$
$\displaystyle\quad+\big{\|}w\big{\|}^{3}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}^{2}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}^{2}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}$
$\displaystyle\quad+\big{\|}w\big{\|}^{2}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}^{2}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}^{3}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}$
$\displaystyle\quad+\big{\|}w\big{\|}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}^{3}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}^{4}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}$
$\displaystyle\quad+\big{\|}\gamma\big{\|}^{4}_{L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}$
$\displaystyle\quad+\big{\|}w\big{\|}^{2}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}$
$\displaystyle\quad+\big{\|}w\big{\|}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}w\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}\big{\|}\gamma\big{\|}^{2}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}$
$\displaystyle\quad+\big{\|}\gamma\big{\|}^{2}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\big{\|}\gamma\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)}+\big{\|}e\big{\|}_{L^{2}_{t}\left(I_{j};\dot{B}^{1/2}_{6/5,2}\right)}.$
At the same time, by Lemma 2.3, we have
$\displaystyle\big{\|}\gamma-\gamma_{j}\big{\|}_{L^{10}_{t}\left(I_{j};\dot{B}^{1/3}_{90/19,2}\right)\cap
L^{12}_{t}\left(I_{j};L^{9}_{x}\right)}+\big{\|}\gamma_{j+1}-\gamma_{j}\big{\|}_{L^{10}_{t}\left({\mathbb{R}};\dot{B}^{1/3}_{90/19,2}\right)\cap
L^{12}_{t}\left({\mathbb{R}};L^{9}_{x}\right)}$ (4.4) $\displaystyle\lesssim$
$\displaystyle\left\|O(w^{4}\gamma+w^{3}\gamma^{2}+w^{2}\gamma^{3}+w\gamma^{4}+\gamma^{5}+w^{2}\gamma+w\gamma^{2}+\gamma^{3})\right\|_{L^{2}(I_{j};\dot{B}^{1/3}_{\frac{18}{11},2})}+\big{\|}e\big{\|}_{L^{2}(I_{j};\dot{B}^{1/3}_{\frac{18}{11},2})}$
$\displaystyle\lesssim$
$\displaystyle\big{\|}w\big{\|}^{4}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}^{3}_{L^{10}_{t,x}(I_{j})}\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}_{L^{10}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}w\big{\|}^{3}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}^{2}_{L^{10}_{t,x}(I_{j})}\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}^{2}_{L^{10}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}w\big{\|}^{2}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}^{2}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}_{L^{10}_{t,x}(I_{j})}\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}^{3}_{L^{10}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}w\big{\|}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}^{3}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}^{4}_{L^{10}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}\gamma\big{\|}^{4}_{L^{10}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}$
$\displaystyle+\big{\|}w\big{\|}^{2}_{L^{5}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}_{L^{5}_{t,x}(I_{j})}\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}_{L^{5}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}w\big{\|}_{L^{5}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{5}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}w\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}\big{\|}\gamma\big{\|}^{2}_{L^{5}_{t,x}(I_{j})}$
$\displaystyle+\big{\|}\gamma\big{\|}^{2}_{L^{5}_{t,x}(I_{j})}\big{\|}\gamma\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}+\big{\|}e\big{\|}_{L^{2}(I_{j};\dot{B}^{1/3}_{18//11,2})}.$
By the interpolation, we have
$\displaystyle\big{\|}f\big{\|}_{L^{6}\left(I_{j};L^{9/2}_{x}\right)}\lesssim\big{\|}f\big{\|}_{L^{6}_{t}\left(I_{j};\dot{B}^{1/2}_{18/7,2}\right)},\quad\big{\|}f\big{\|}_{L^{10}_{t,x}(I_{j})}\lesssim\big{\|}f\big{\|}_{L^{10}_{t}(I_{j};\dot{B}^{1/3}_{90/19,2})}.$
Therefore, assuming that
$\displaystyle\big{\|}\gamma\big{\|}_{ST(I_{j})}\leq\delta\ll
1,\quad\forall\;j=0,1,\ldots,N-1,$ (4.5)
then by (4.2), (4.3) and (4.4), we have
$\displaystyle\big{\|}\gamma\big{\|}_{ST(I_{j})}+\big{\|}\gamma_{j+1}\big{\|}_{ST(t_{j+1},t_{N})}\leq
C\big{\|}\gamma_{j}\big{\|}_{ST(t_{j},t_{N})}+\varepsilon,$
for some absolute constant $C>0$. By (4.1) and iteration on $j$, we get
$\displaystyle\big{\|}\gamma\big{\|}_{ST(I)}\leq(2C)^{N}\varepsilon\leq\frac{\delta}{2},$
if we choose $\varepsilon_{0}$ sufficiently small. Hence the assumption (4.5)
is justified by continuity in $t$ and induction on $j$. then repeating the
estimate (4.3) and (4.4) once again, we can obtain the $ST$-norm estimate on
$\gamma$, which implies the Strichartz estimate on $u$. ∎
## 5\. Profile decomposition
In this part, we will use the method in [2, 17, 21] to show the linear and
nonlinear profile decompositions of the sequences of radial, $H^{1}$-bounded
solutions of (1.1), which will be used to construct the critical element
(minimal energy non-scattering solution) and show its properties, especially
the compactness. In order to do it, we now introduce the complex-valued
function $\overrightarrow{v}(t,x)$ by
$\displaystyle\overrightarrow{v}(t,x)=\left<\nabla\right>v(t,x),\quad
v(t,x)=\left<\nabla\right>^{-1}\overrightarrow{v}(t,x).$
Given $(t^{j}_{n},h^{j}_{n})\in{\mathbb{R}}\times(0,1]$, let $\tau^{j}_{n}$,
$T^{j}_{n}$ denote the scaled time drift, the scaling transformation, defined
by
$\displaystyle\tau^{j}_{n}=-\frac{t^{j}_{n}}{\left(h^{j}_{n}\right)^{2}},\quad
T^{j}_{n}\varphi(x)=\frac{1}{(h^{j}_{n})^{3/2}}\varphi\left(\frac{x}{h^{j}_{n}}\right).$
We also introduce the set of Fourier multipliers on ${\mathbb{R}}^{3}$.
$\displaystyle\mathcal{MC}=\\{\mu=\mathcal{F}^{-1}\widetilde{\mu}\mathcal{F}\;|\;\widetilde{\mu}\in
C({\mathbb{R}}^{3}),\;\exists\lim_{|\xi|\rightarrow+\infty}\widetilde{\mu}(\xi)\in{\mathbb{R}}\\}.$
### 5.1. Linear profile decomposition
In this subsection, we show the profile decomposition with the scaling
parameter of a sequence of the radial, free Schrödinger solutions in the
energy space $H^{1}({\mathbb{R}}^{3})$, which implies the profile
decomposition of a sequence of radial initial data.
###### Proposition 5.1.
Let
$\overrightarrow{v}_{n}(t,x)=e^{it\Delta}\overrightarrow{v}_{n}(0)$
be a sequence of the radial solutions of the free Schrödinger equation with
bounded $L^{2}$ norm. Then up to a subsequence, there exist
$K\in\\{0,1,2,\ldots,\infty\\}$, radial functions
$\\{\varphi^{j}\\}_{j\in[0,K)}\subset L^{2}({\mathbb{R}}^{3})$ and
$\\{t^{j}_{n},h^{j}_{n}\\}_{n\in{\mathbb{N}}}\subset{\mathbb{R}}\times(0,1]$
satisfying
$\displaystyle\overrightarrow{v}_{n}(t,x)=\sum^{k-1}_{j=0}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x),$
(5.1)
where
$\overrightarrow{v}^{j}_{n}(t,x)=e^{i(t-t^{j}_{n})\Delta}T^{j}_{n}\varphi^{j}$,
and
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\big{\|}\overrightarrow{w}^{k}_{n}\big{\|}_{L^{\infty}_{t}({\mathbb{R}};B^{-3/2}_{\infty,\infty}({\mathbb{R}}^{3}))}=0,$
(5.2)
and for any Fourier multiplier $\mu\in\mathcal{MC}$, any $l<j<k\leq K$ and any
$t\in{\mathbb{R}}$,
$\displaystyle\lim_{n\rightarrow+\infty}\left(\log\left|\dfrac{h^{j}_{n}}{h^{l}_{n}}\right|+\left|\frac{t^{j}_{n}-t^{l}_{n}}{(h^{l}_{n})^{2}}\right|\right)=\infty,$
(5.3)
$\displaystyle\lim_{n\rightarrow+\infty}\left<\mu\overrightarrow{v}^{l}_{n}(t)\;,\;\mu\overrightarrow{v}^{j}_{n}(t)\right>_{L^{2}_{x}}=\lim_{n\rightarrow+\infty}\left<\mu\overrightarrow{v}^{j}_{n}(t)\;,\;\mu\overrightarrow{w}^{k}_{n}(t)\right>_{L^{2}_{x}}=0.$
(5.4)
Moreover, each sequence $\\{h^{j}_{n}\\}_{n\in{\mathbb{N}}}$ is either going
to $0$ or identically $1$ for all $n$.
###### Remark 5.2.
We call $\overrightarrow{v}^{j}_{n}$ and $\overrightarrow{w}^{k}_{n}$ the free
concentrating wave and the remainder, respectively. From (5.4), we have the
following asymptotic orthogonality
$\displaystyle\lim_{n\rightarrow+\infty}\left(\big{\|}\mu\overrightarrow{v}_{n}(t)\big{\|}^{2}_{L^{2}}-\sum^{k-1}_{j=0}\big{\|}\mu\overrightarrow{v}^{j}_{n}(t)\big{\|}^{2}_{L^{2}}-\big{\|}\mu\overrightarrow{w}^{k}_{n}(t)\big{\|}^{2}_{L^{2}}\right)$
$\displaystyle=0.$ (5.5)
###### Proof of Proposition 5.1.
Let
$\displaystyle\nu:=\varlimsup_{n\rightarrow\infty}\big{\|}\overrightarrow{v}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}=\varlimsup_{n\rightarrow\infty}\sup_{(t,x)\in{\mathbb{R}}\times{\mathbb{R}}^{3},\atop
k\geq 0}2^{-3k/2}\big{|}\Lambda_{k}*\overrightarrow{v}_{n}(t,x)\big{|}.$
If $\nu=0$, then we have done with $K=0$.
Otherwise,
$\displaystyle\nu=\varlimsup_{n\rightarrow\infty}\big{\|}\overrightarrow{v}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}>0$.
By the radial Gagliardo-Nirenberg inequality and the Bernstein inequality, we
have
$\displaystyle\sup_{t\in{\mathbb{R}},|2^{k}x|\geq R,\atop k\geq
0}2^{-3k/2}\big{|}\Lambda_{k}*\overrightarrow{v}_{n}(t,x)\big{|}\lesssim$
$\displaystyle\sup_{k\geq
0}\frac{2^{k}2^{-3k/2}}{R}\big{\|}\Lambda_{k}*\overrightarrow{v}_{n}(t,x)\big{\|}^{1/2}_{L^{\infty}_{t}L^{2}_{x}}\cdot\big{\|}\nabla\Lambda_{k}*\overrightarrow{v}_{n}(t,x)\big{\|}^{1/2}_{L^{\infty}_{t}L^{2}_{x}}$
$\displaystyle\lesssim$ $\displaystyle\sup_{k\geq
0}\frac{1}{R}\big{\|}\overrightarrow{v}_{n}(t,x)\big{\|}_{L^{\infty}_{t}L^{2}_{x}}\lesssim\frac{1}{R}.$
If taking $R$ sufficiently large, we have
$\displaystyle\sup_{t\in{\mathbb{R}},|2^{k}x|\geq R,k\geq
0}2^{-3k/2}\big{|}\Lambda_{k}*\overrightarrow{v}_{n}(t,x)\big{|}\leq\frac{1}{2}\nu.$
thus, there exists a sequence $(t_{n},x_{n},k_{n})$ with $k_{n}\geq 0$ and
$|2^{k_{n}}x_{n}|\leq R$ such that for large $n$,
$\displaystyle\frac{1}{2}\varlimsup_{n\rightarrow\infty}\big{\|}\overrightarrow{v}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}=\frac{1}{2}\nu\leq
2^{-3k_{n}/2}\big{|}\Lambda_{k_{n}}*\overrightarrow{v}_{n}(t_{n},x_{n})\big{|}.$
Now we define $h_{n}$ and $\psi_{n}$ by $h_{n}=2^{-k_{n}}\in(0,1]$ and
$\displaystyle\overrightarrow{v}_{n}(t_{n},x)=$
$\displaystyle\left(T_{n}\psi_{n}\right)(x-x_{n})=\frac{1}{(h_{n})^{3/2}}\psi_{n}\left(\frac{x-x_{n}}{h_{n}}\right)$
(5.6) $\displaystyle=$ $\displaystyle
T_{n}\left(\psi_{n}\Big{(}x-\frac{x_{n}}{h_{n}}\Big{)}\right).$
Since
$\big{\|}\psi_{n}\big{\|}_{L^{2}}=\big{\|}T_{n}\psi_{n}\big{\|}_{L^{2}}=\big{\|}\overrightarrow{v}_{n}(t_{n})\big{\|}_{L^{2}}\leq
C,$ then there exists some $\psi\in L^{2}$, such that, up to a subsequence, we
have as $n\rightarrow+\infty$
$\displaystyle\frac{x_{n}}{h_{n}}\rightarrow
x^{0},\;\;\text{and}\;\;\psi_{n}\rightharpoonup\psi\quad\text{weakly
in}\;\;L^{2}.$ (5.7)
On the other hand, if $k_{n}=0$, we have
$\displaystyle
2^{-3k_{n}/2}\big{|}\Lambda_{k_{n}}*\overrightarrow{v}_{n}(t_{n},x_{n})\big{|}=$
$\displaystyle\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{0}(y)\;\;2^{-3k_{n}/2}\overrightarrow{v}_{n}\left(t_{n},x_{n}-\frac{y}{2^{k_{n}}}\right)\;dy$
$\displaystyle=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{0}(y)\;\psi_{n}(-y)\;dy$
$\displaystyle\longrightarrow$
$\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{0}(y)\;\psi(-y)\;dy\lesssim\big{\|}\psi\big{\|}_{L^{2}}.$
By the same way, if $k_{n}\geq 1$, we have
$\displaystyle
2^{-3k_{n}/2}\big{|}\Lambda_{k_{n}}*\overrightarrow{v}_{n}(t_{n},x_{n})\big{|}=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{(0)}(y)\;2^{-3k_{n}/2}\overrightarrow{v}_{n}\left(t_{n},x_{n}-\frac{y}{2^{k_{n}}}\right)\;dy$
$\displaystyle=$
$\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{(0)}(y)\;\psi_{n}(-y)\;dy$
$\displaystyle\longrightarrow$
$\displaystyle\int_{{\mathbb{R}}^{3}}\Lambda_{(0)}(y)\;\psi(-y)\;dy\lesssim\big{\|}\psi\big{\|}_{L^{2}}.$
If $h_{n}\rightarrow 0$, then we take
$\displaystyle(t^{0}_{n},h^{0}_{n})=(t_{n},h_{n}),\quad\varphi^{0}(x)=\psi\left(x-x^{0}\right),$
otherwise, up to a subsequence, we may assume that $h_{n}\rightarrow
h_{\infty}$ for some $h_{\infty}\in(0,1]$, and take
$\displaystyle(t^{0}_{n},h^{0}_{n})=(t_{n},1),\quad\varphi^{0}(x)=\frac{1}{(h_{\infty})^{3/2}}\psi\left(\frac{x}{h_{\infty}}-x^{0}\right),$
then
$\displaystyle
T_{n}\left(\psi\Big{(}x-\frac{x_{n}}{h_{n}}\Big{)}\right)-T^{0}_{n}\varphi^{0}(x)\longrightarrow
0\quad\text{strongly in}\;\;L^{2}.$ (5.8)
In addition, since
$\overrightarrow{v}_{n}(t_{n},x)=\left(T_{n}\psi_{n}\right)(x-x_{n})$ is
radial, so is $\varphi^{0}(x)$.
Let
$\overrightarrow{v}^{0}_{n}(t,x)=e^{i(t-t^{0}_{n})\Delta}T^{0}_{n}\varphi^{0}$,
we define $\overrightarrow{w}^{1}_{n}$ by
$\displaystyle\overrightarrow{v}_{n}(t,x)=$
$\displaystyle\overrightarrow{v}^{0}_{n}(t,x)+\overrightarrow{w}^{1}_{n}(t,x),$
(5.9)
then by (5.7) and (5.8), we have
$\displaystyle(T^{0}_{n})^{-1}\overrightarrow{w}^{1}_{n}(t^{0}_{n})=(T^{0}_{n})^{-1}T_{n}\left(\psi_{n}\Big{(}x-\frac{x_{n}}{h_{n}}\Big{)}\right)-\varphi^{0}\rightharpoonup
0\quad\text{weakly in}\;\;L^{2},$
which implies that
$\displaystyle\left<\mu\overrightarrow{v}^{0}_{n}(t),\mu\overrightarrow{w}^{1}_{n}(t)\right>=$
$\displaystyle\left<\mu\overrightarrow{v}^{0}_{n}(t^{0}_{n}),\mu\overrightarrow{w}^{1}_{n}(t^{0}_{n})\right>=\left<\mu^{0}_{n}\varphi^{0},\mu^{0}_{n}(T^{0}_{n})^{-1}\overrightarrow{w}^{1}_{n}(t^{0}_{n})\right>\longrightarrow
0,$
where we used the conservation law in the first equality and the dominated
convergence theorem and $\mu^{0}_{n}(D)=\mu\left(\frac{D}{h^{0}_{n}}\right)$
in the last equality. It is the decomposition for $k=1$.
Next we apply the above procedure to the sequence $\overrightarrow{w}^{1}_{n}$
in place of $\overrightarrow{v}_{n}$, then either
$\displaystyle\varlimsup_{n\rightarrow\infty}\big{\|}\overrightarrow{w}^{1}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}=0$
or we can find the next concentrating wave $\overrightarrow{v}^{1}_{n}$ and
the remainder $\overrightarrow{w}^{2}_{n}$, such that for some
$(t^{1}_{n},h^{1}_{n})$ with $h^{1}_{n}\in(0,1]$ and radial function
$\varphi^{1}\in L^{2}({\mathbb{R}}^{3})$,
$\displaystyle\overrightarrow{w}^{1}_{n}(t,x)=\overrightarrow{v}^{1}_{n}(t,x)+$
$\displaystyle\overrightarrow{w}^{2}_{n}(t,x)=e^{i(t-t^{1}_{n})\Delta}T^{1}_{n}\varphi^{1}(x)+\overrightarrow{w}^{2}_{n}(t,x),$
(5.10)
and
$\displaystyle\varlimsup_{n\rightarrow+\infty}\big{\|}\overrightarrow{w}^{1}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}\lesssim$
$\displaystyle\;\big{\|}\varphi^{1}\big{\|}_{L^{2}}=\big{\|}\overrightarrow{v}^{1}_{n}\big{\|}_{L^{2}},$
(5.11)
$\displaystyle(T^{1}_{n})^{-1}\overrightarrow{w}^{2}_{n}(t^{1}_{n})\rightharpoonup
0\quad\text{weakly in}\;\;L^{2}$
$\displaystyle\Longrightarrow\left<\mu\overrightarrow{v}^{1}_{n}(t),\mu\overrightarrow{w}^{2}_{n}(t)\right>\longrightarrow
0.$
Iterating the above procedure, we can obtain the decomposition (5.1). It
remains to show the properties (5.2), (5.3) and (5.4).
We first assume that (5.4) holds, then by (5.5) and the Cauchy criterion, we
have
$\displaystyle\lim_{n\rightarrow+\infty}\big{\|}\overrightarrow{w}^{k}_{n}\big{\|}_{L^{\infty}_{t}B^{-3/2}_{\infty,\infty}}\lesssim\big{\|}\varphi^{k}\big{\|}_{L^{2}}=\big{\|}\overrightarrow{v}^{k}_{n}\big{\|}_{L^{2}}\longrightarrow
0\quad\text{as}\;\;k\rightarrow+\infty.$ (5.12)
which implies (5.2).
Now we show (5.3) by contradiction. Suppose that (5.3) fails, then there
exists a minimal $(l,j)$ which violates (5.3). By extracting a subsequence, We
may assume that $h^{l}_{n}\rightarrow h^{l}_{\infty}$ and
$h^{l}_{n}/h^{j}_{n}$ and $(t^{l}_{n}-t^{j}_{n})/(h^{l}_{n})^{2}$ all
converge.
Now consider
$\displaystyle\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{l+1}_{n}(t^{l}_{n})=$
$\displaystyle\sum^{j}_{m=l+1}\left(T^{l}_{n}\right)^{-1}\overrightarrow{v}^{m}_{n}(t^{l}_{n})+\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{l}_{n})$
$\displaystyle=$
$\displaystyle\sum^{j}_{m=l+1}\left(T^{l}_{n}\right)^{-1}e^{i(t^{l}_{n}-t^{m}_{n})\Delta}T^{m}_{n}\varphi^{m}+\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{l}_{n})$
$\displaystyle=$
$\displaystyle\sum^{j-1}_{m=l+1}S^{l,m}_{n}\varphi^{m}+S^{l,j}_{n}\varphi^{j}+\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{l}_{n}),$
where
$\displaystyle
S^{l,m}_{n}=\left(T^{l}_{n}\right)^{-1}e^{i(t^{l}_{n}-t^{m}_{n})\Delta}T^{m}_{n}=e^{i\frac{t^{l}_{n}-t^{m}_{n}}{(h^{l}_{n})^{2}}\Delta}\left(T^{l}_{n}\right)^{-1}T^{m}_{n}:=e^{it^{l,m}_{n}\Delta}T^{l,m}_{n}$
with the sequence
$\displaystyle t^{l,m}_{n}=\frac{t^{l}_{n}-t^{m}_{n}}{(h^{l}_{n})^{2}},\quad
h^{l,m}_{n}=\frac{h^{m}_{n}}{h^{l}_{n}}.$ (5.13)
By the procedure of constructing (5.1), as $n\rightarrow+\infty$, we have
$\displaystyle\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{l+1}_{n}(t^{l}_{n})\rightharpoonup
0$ $\displaystyle\quad\text{weakly in}\;L^{2},$
$\displaystyle\left(T^{j}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{j}_{n})\rightharpoonup
0$ $\displaystyle\quad\text{weakly in}\;L^{2},$
and by the asymptotic orthogonality (5.3) between $m$ and $l$ with
$m\in[l+1,j-1]$
$\displaystyle S^{l,m}_{n}\varphi^{m}\rightharpoonup
0,\;\;\forall\;m\in[l+1,j-1],$
and by the convergence of $h^{l}_{n}/h^{j}_{n}$ and
$(t^{l}_{n}-t^{j}_{n})/(h^{l}_{n})^{2}$, we have
$S^{l,j}_{n}\varphi^{j}\rightarrow S^{l,j}_{\infty}\varphi^{j}$ and
$\displaystyle\left(T^{l}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{l}_{n})=$
$\displaystyle
S^{l,j}_{n}\left(T^{j}_{n}\right)^{-1}\overrightarrow{w}^{j+1}_{n}(t^{j}_{n})\rightharpoonup
0\quad\text{weakly in}\;L^{2}.$
Then $\varphi^{j}=0$, it is a contradiction. Thus we obtain the orthogonality
(5.3).
Last we show (5.4). For $j\not=l$, we have
$\displaystyle\left<\mu\overrightarrow{v}^{l}_{n}(t)\;,\;\mu\overrightarrow{v}^{j}_{n}(t)\right>_{L^{2}_{x}}$
$\displaystyle=$
$\displaystyle\left<\mu\overrightarrow{v}^{l}_{n}(0)\;,\;\mu\overrightarrow{v}^{j}_{n}(0)\right>_{L^{2}_{x}}=\left<\mu
e^{-it^{l}_{n}\Delta}T^{l}_{n}\varphi^{l}\;,\;\mu
e^{-it^{j}_{n}\Delta}T^{j}_{n}\varphi^{j}\right>_{L^{2}_{x}}$ $\displaystyle=$
$\displaystyle\left<e^{-it^{l}_{n}\Delta}T^{l}_{n}\mu^{l}_{n}\varphi^{l}\;,\;e^{-it^{j}_{n}\Delta}T^{j}_{n}\mu^{j}_{n}\varphi^{j}\right>_{L^{2}_{x}}=\left<\left(T^{j}_{n}\right)^{-1}e^{i(t^{j}_{n}-t^{l}_{n})\Delta}T^{l}_{n}\mu^{l}_{n}\varphi^{l}\;,\;\mu^{j}_{n}\varphi^{j}\right>_{L^{2}_{x}}$
$\displaystyle=$
$\displaystyle\left<e^{i\frac{t^{j}_{n}-t^{l}_{n}}{(h^{j}_{n})^{2}}\Delta}\left(T^{j}_{n}\right)^{-1}T^{l}_{n}\mu^{l}_{n}\varphi^{l}\;,\;\mu^{j}_{n}\varphi^{j}\right>_{L^{2}_{x}}=\left<S^{j,l}_{n}\mu^{l}_{n}\varphi^{l}\;,\;\mu^{j}_{n}\varphi^{j}\right>_{L^{2}_{x}}\rightarrow
0\quad\text{as}\;\;n\rightarrow+\infty$
where $\widetilde{\mu}^{l}_{n}(\xi)=\widetilde{\mu}\left(\xi/h^{l}_{n}\right)$
and we used the fact that $S^{j,l}_{n}\rightharpoonup 0$ weakly in $L^{2}$ as
$n\rightarrow+\infty$ by (5.3). In addition, we have
$\displaystyle\left<\mu\overrightarrow{v}^{j}_{n}(t)\;,\;\mu\overrightarrow{w}^{k}_{n}(t)\right>_{L^{2}_{x}}=\left<\mu\overrightarrow{v}^{j}_{n}(t)\;,\;\mu\Big{(}\overrightarrow{w}^{j+1}_{n}(t)-\sum^{k-1}_{m=j+1}\overrightarrow{v}^{m}_{n}(t)\Big{)}\right>_{L^{2}_{x}}\longrightarrow
0$
as $n\rightarrow+\infty$. This completes the proof of (5.4). ∎
After the orthogonality’s proof of the linear energy, we begin with the
orthogonal analysis for the nonlinear energy.
###### Lemma 5.3.
Let $\overrightarrow{v}_{n}$ be a sequence of the radial solutions of the free
Schrödinger equation. Let
$\displaystyle\overrightarrow{v}_{n}(t,x)=\sum^{k-1}_{j=0}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x)$
be the linear profile decomposition given by Proposition 5.1. Then we have
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|M(v_{n}(0))-\sum^{k-1}_{j=0}M(v^{j}_{n}(0))-M(w^{k}_{n}(0))\right|=0,$
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|E(v_{n}(0))-\sum^{k-1}_{j=0}E(v^{j}_{n}(0))-E(w^{k}_{n}(0))\right|=0,$
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|K(v_{n}(0))-\sum^{k-1}_{j=0}K(v^{j}_{n}(0))-K(w^{k}_{n}(0))\right|=0.$
###### Proof.
We can show that the quadratic terms in $M$, $E$ and $K$ have the orthogonal
decomposition by taking $\mu=\frac{1}{\left<\nabla\right>}$ and
$\mu=\frac{|\nabla|}{\left<\nabla\right>}$ in Remark 5.2, thus it suffices to
show that
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|F_{i}\left(v_{n}(0)\right)-\sum_{j<k}F_{i}\left(v^{j}_{n}(0)\right)-F_{i}\left(w^{k}_{n}(0)\right)\right|=0,\quad
i=1,2,$
where $F_{1}$ and $F_{2}$ are denoted by
$\displaystyle
F_{1}(u(t))=\int_{{\mathbb{R}}^{3}}|u(t,x)|^{6}\;dx,\;\;F_{2}(u(t))=\int_{{\mathbb{R}}^{3}}|u(t,x)|^{4}\;dx.$
In order to do so, we need re-arrange the linear concentrating wave with
respect to its dispersive decay (whether $\tau^{j}_{n}$ goes to $\pm\infty$ or
not for all $j$). Let
$v^{<k}_{n}(0)=\displaystyle\sum_{j<k}v^{j}_{n}(0)=\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)+\sum_{j<k,\tau^{j}_{n}\rightarrow\pm\infty}v^{j}_{n}(0)$
for some finite numbers $\tau^{j}_{\infty}$’s, then we have
$\displaystyle\Big{|}F_{i}\left(v_{n}(0)\right)-$
$\displaystyle\sum_{j<k}F_{i}\left(v^{j}_{n}(0)\right)-F_{i}\left(w^{k}_{n}(0)\right)\Big{|}$
$\displaystyle\leq$
$\displaystyle\left|F_{i}\left(v_{n}(0)\right)-F_{i}\left(v^{<k}_{n}(0)\right)\right|+\left|F_{i}\left(v^{<k}_{n}(0)\right)-F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)\right|$
$\displaystyle+\left|F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(v^{j}_{n}(0)\right)\right|$
(5.14)
$\displaystyle+\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\pm\infty}F_{i}\left(v^{j}_{n}(0)\right)\right|+\left|F_{i}\left(w^{k}_{n}(0)\right)\right|.$
First, by (5.2) and interpolation, we have that
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\big{\|}w^{k}_{n}(0)\big{\|}_{L^{p}_{x}}=0,\quad\forall\;\;2<p\leq
6.$
which implies that
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|F_{i}\left(v_{n}(0)\right)-F_{i}\left(v^{<k}_{n}(0)\right)\right|$
$\displaystyle=0,$ $\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|F_{i}\left(w^{k}_{n}(0)\right)\right|=0.$
Second by the dispersive estimate for $v^{j}_{n}(0)$ with
$\tau^{j}_{n}\rightarrow\pm\infty$, we have
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|F_{i}(v^{<k}_{n}(0))-F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)\right|$
$\displaystyle=0,$ $\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\pm\infty}F_{i}\left(v^{j}_{n}(0)\right)\right|=0.$
Last we will use the approximation argument in [17] to show that every non-
dispersive concentrating wave will get away from the others, which contributes
to the orthogonality of (5.14). Let
$\psi^{j}:=e^{i\tau^{j}_{\infty}\Delta}\varphi^{j}\in L^{2}$, we have
$\displaystyle\left|F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(v^{j}_{n}(0)\right)\right|$
(5.15)
$\displaystyle\leq\left|F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)-F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)\right|$
$\displaystyle\quad+\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(v^{j}_{n}(0)\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|$
$\displaystyle\quad+\left|F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|.$
(5.16)
For those $v^{j}_{n}(0)$ with $\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}$, by
the continuity of the operator $e^{it\Delta}$ in $t$ in $H^{1}$, we have
$\displaystyle v^{j}_{n}(0)=$
$\displaystyle\left<\nabla\right>^{-1}e^{-it^{j}_{n}\Delta}T^{j}_{n}\varphi^{j}=\left<\nabla\right>^{-1}T^{j}_{n}e^{i\tau^{j}_{n}\Delta}\varphi^{j}\longrightarrow\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\quad\text{in}\;\;H^{1}({\mathbb{R}}^{3}),$
which implies that
$\displaystyle\left|F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}v^{j}_{n}(0)\right)-F_{i}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)\right|\rightarrow
0,$
$\displaystyle\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(v^{j}_{n}(0)\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{i}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|\rightarrow
0.$
Now we consider (5.16) for $i=1,2$, separately.
First for $i=2$, we compute as following,
$\displaystyle\left|F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|$
$\displaystyle\leq\left|F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)\right|$
$\displaystyle\quad+\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|$
$\displaystyle\quad+\left|F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|.$
For $h^{j}_{n}\rightarrow 0$, we have
$\displaystyle\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\rightarrow
0\quad\text{in}\;\;L^{p},\;\;\forall\;2\leq p<6,$
which implies that
$\displaystyle\left|F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j},h^{j}_{n}=1}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)\right|\rightarrow
0,$
$\displaystyle\left|\sum_{j<k,\tau^{j}n_{n}\rightarrow\tau^{j}}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j},h^{j}_{n}=1}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|\rightarrow
0.$
In addition, by the orthogonality (5.3), we know that there is at most one
term $\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}$ with
$\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1$, hence
$\displaystyle\left|F_{2}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty},h^{j}_{n}=1}F_{2}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|=0.$
Now we consider the case $i=1$, Let $\widehat{\psi}^{j}=|\nabla|^{-1}\psi^{j}$
if $h^{j}_{n}\rightarrow 0$, and
$\widehat{\psi}^{j}=\left<\nabla\right>^{-1}\psi^{j}$ if $h^{j}_{n}\equiv 1$,
then we have $\widehat{\psi}^{j}\in L^{6}_{x}$, and
$\displaystyle\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)\right|$
$\displaystyle\leq\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|$
$\displaystyle\quad+\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|$
$\displaystyle\quad+\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|.$
Since
$\displaystyle\big{\|}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}-h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\big{\|}_{L^{6}_{x}}=$
$\displaystyle\begin{cases}\big{\|}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}-h^{j}_{n}T^{j}_{n}|\nabla|^{-1}\psi^{j}\big{\|}_{L^{6}_{x}}\quad\text{if}\;\;h^{j}_{n}\rightarrow
0\\\
\big{\|}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}-h^{j}_{n}T^{j}_{n}\left<\nabla\right>^{-1}\psi^{j}\big{\|}_{L^{6}_{x}}\quad\text{if}\;\;h^{j}_{n}\equiv
1\end{cases}$ $\displaystyle=$
$\displaystyle\begin{cases}\big{\|}(\left<\nabla\right>^{j}_{n})^{-1}\psi^{j}-|\nabla|^{-1}\psi^{j}\big{\|}_{L^{6}_{x}}\quad\text{if}\;\;h^{j}_{n}\rightarrow
0\\\ 0\quad\text{if}\;\;h^{j}_{n}\equiv 1\end{cases}$
$\displaystyle\longrightarrow 0,\quad\text{as}\;\;n\rightarrow+\infty,$
which shows that
$\displaystyle\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j}\right)-F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|\longrightarrow
0,$
$\displaystyle\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(\left<\nabla\right>^{-1}T^{j}_{n}\psi^{j})\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|\longrightarrow
0.$
We further replace each $\widehat{\psi}^{j}$ by the non-overlap terms
$\widetilde{\psi}^{j}_{n}$ with each other
$\displaystyle\widetilde{\psi}^{j}_{n}=\widehat{\psi}^{j}\times\begin{cases}0;\quad\exists\;l<j,\;\text{such
that}\;h^{l}_{n}<h^{j}_{n}\;\;\text{and}\;\;\frac{x}{h^{j,l}_{n}}\in\text{supp}\;{\widehat{\psi}^{l}},\\\
1;\quad\text{otherwise},\end{cases}$
where $h^{j,l}_{n}$ is determined by (5.13). By (5.3), we know that
$h^{j,l}_{n}\rightarrow 0$, therefore as $n\rightarrow+\infty$
$\displaystyle\widetilde{\psi}^{j}_{n}\rightarrow\widehat{\psi}^{j},$
$\displaystyle\quad
a.e.\;x\in{\mathbb{R}}^{3},\quad\text{and}\quad\widetilde{\psi}^{j}_{n}\rightarrow\widehat{\psi}^{j},\quad\text{in}\;\;L^{6}_{x},$
which implies that
$\displaystyle\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right)\right|\longrightarrow
0,$
$\displaystyle\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right)\right|\longrightarrow
0.$
On the other hand, by the support property of $\widetilde{\psi}^{j}_{n}$, we
know that
$\displaystyle
F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right)=\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right).$
Therefore, we have
$\displaystyle\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)\right|$
$\displaystyle\leq\left|F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-F_{1}\left(\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right)\right|$
$\displaystyle\quad+\left|\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widehat{\psi}^{j}\right)-\sum_{j<k,\tau^{j}_{n}\rightarrow\tau^{j}_{\infty}}F_{1}\left(h^{j}_{n}T^{j}_{n}\widetilde{\psi}^{j}_{n}\right)\right|\longrightarrow
0.$
This completes the proof. ∎
###### Lemma 5.4.
Let $k\in{\mathbb{N}}$ and radial functions $\varphi_{0},\ldots,\varphi_{k}\in
H^{1}({\mathbb{R}}^{3})$, $m$ be determined by (1.5). Assume that there exist
some $\delta$, $\varepsilon>0$ with $4\varepsilon<3\delta$ such that
$\displaystyle\sum^{k}_{j=0}E(\varphi_{j})-\varepsilon\leq
E\left(\sum^{k}_{j=0}\varphi_{j}\right)<m-\delta,\;\;\text{and}\;\;-\varepsilon\leq
K\left(\sum^{k}_{j=0}\varphi_{j}\right)\leq\sum^{k}_{j=0}K(\varphi_{j})+\varepsilon.$
Then $\varphi_{j}\in\mathcal{K}^{+}$ for all $j=0,\ldots,k$.
###### Proof.
Suppose that $K(\varphi_{l})<0$ for some $l$. Then by Lemma 2.9, we have
$\displaystyle H(\varphi_{l})\geq\inf\left\\{H(\varphi)\;|\;\varphi\in
H^{1}({\mathbb{R}}^{3}),\;\varphi\not=0,\;K(\varphi)\leq 0\right\\}=m.$
By the nonnegativity of $H(\varphi_{j})$ for $j\geq 0$, we have
$\displaystyle m\leq$ $\displaystyle
H(\varphi_{l})\leq\sum^{k}_{j=0}H(\varphi_{j})=\sum^{k}_{j=0}\left(E(\varphi_{j})-\frac{1}{6}K(\varphi_{j})\right)$
$\displaystyle\leq$ $\displaystyle
E\left(\sum^{k}_{j=0}\varphi_{j}\right)+\varepsilon-\frac{1}{6}K\left(\sum^{k}_{j=0}\varphi_{j}\right)+\frac{1}{6}\varepsilon$
$\displaystyle\leq$ $\displaystyle
m-\delta+\varepsilon+\frac{1}{3}\varepsilon<m.$
It is a contradiction. Hence for any $j\in\\{0,\ldots,k\\}$, we have
$\displaystyle K(\varphi_{j})\geq 0,$
which implies that
$\displaystyle E(\varphi_{j})=H(\varphi_{j})+\frac{1}{6}K(\varphi_{j})\geq 0,$
and
$\displaystyle
E(\varphi_{j})\leq\sum^{k}_{i=0}E(\varphi_{i})<m-\delta+\varepsilon<m,$
which means that $\varphi_{j}\in\mathcal{K}^{+}$ for all $j$. ∎
According to the above results, we conclude as following.
###### Proposition 5.5.
Let $\overrightarrow{v}_{n}(t,x)$ be a sequence of the radial solutions of the
free Schrödinger equation satisfying
$\displaystyle v_{n}(0)\in\mathcal{K}^{+}\;\;\text{and}\;\;E(v_{n}(0))<m.$
Let
$\displaystyle\overrightarrow{v}_{n}(t,x)=\sum^{k-1}_{j=0}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x),$
be the linear profile decomposition given by Proposition 5.1. Then for large
$n$ and all $j<K$, we have
$\displaystyle
v^{j}_{n}(0)\in\mathcal{K}^{+},\;\;\;\;w^{K}_{n}(0)\in\mathcal{K}^{+},$
and
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|M(v_{n}(0))-\sum_{j<k}M(v^{j}_{n}(0))-M(w^{k}_{n}(0))\right|=0,$
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|E(v_{n}(0))-\sum_{j<k}E(v^{j}_{n}(0))-E(w^{k}_{n}(0))\right|=0,$
$\displaystyle\lim_{k\rightarrow
K}\varlimsup_{n\rightarrow+\infty}\left|K(v_{n}(0))-\sum_{j<k}K(v^{j}_{n}(0))-K(w^{k}_{n}(0))\right|=0.$
Moreover for all $j<K$, we have
$\displaystyle
0\leq\varliminf_{n\rightarrow+\infty}E(v^{j}_{n}(0))\leq\varlimsup_{n\rightarrow+\infty}E(v^{j}_{n}(0))\leq\varlimsup_{n\rightarrow+\infty}E(v_{n}(0)),$
where the last inequality becomes equality only if $K=1$ and
$w^{1}_{n}\rightarrow 0$ in $L^{\infty}_{t}\dot{H}^{1}_{x}$.
### 5.2. Nonlinear profile decomposition
After the linear profile decomposition of a sequence of initial data in the
last subsection, we now show the nonlinear profile decomposition of a sequence
of radial solutions of (1.1) with the same initial data in the energy space
$H^{1}({\mathbb{R}}^{3})$. First we introduce some notation
$\displaystyle\left<\nabla\right>^{j}_{n}=\sqrt{\left(h^{j}_{n}\right)^{2}-\Delta},\;\;\left<\nabla\right>^{j}_{\infty}=\sqrt{\left(h^{j}_{\infty}\right)^{2}-\Delta}\;.$
Now let $v_{n}(t,x)$ be a sequence of radial solutions for the free
Schrödinger equation with initial data in $\mathcal{K}^{+}$, that is,
$v_{n}\in H^{1}({\mathbb{R}}^{3})$ is radial and
$\displaystyle\left(i\partial_{t}+\Delta\right)v_{n}=0,\quad
v_{n}(0)\in\mathcal{K}^{+}.$
Let
$\displaystyle\overrightarrow{v}_{n}(t,x)=\left<\nabla\right>v_{n}(t,x),$
then by Proposition 5.1, we have a sequence of the radial, free concentrating
wave $\overrightarrow{v}^{j}_{n}(t,x)$ with
$\overrightarrow{v}^{j}_{n}(t^{j}_{n})=T^{j}_{n}\varphi^{j}$,
$v^{j}_{n}(0)\in\mathcal{K}^{+}$ for $j=0,\ldots,K$, such that
$\displaystyle\overrightarrow{v}_{n}(t,x)=$
$\displaystyle\sum^{k-1}_{j=0}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x)=\sum^{k-1}_{j=0}e^{i(t-t^{j}_{n})\Delta}T^{j}_{n}\varphi^{j}+\overrightarrow{w}^{k}_{n}$
$\displaystyle=$
$\displaystyle\sum^{k-1}_{j=0}T^{j}_{n}e^{i\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right)\Delta}\varphi^{j}+\overrightarrow{w}^{k}_{n}.$
Now for any concentrating wave $\overrightarrow{v}^{j}_{n}$, $j=0,\ldots,K$,
we undo the group action, i.e., the scaling transformation $T^{j}_{n}$, to
look for the linear profile $V^{j}$. Let
$\displaystyle\overrightarrow{v}^{j}_{n}(t,x)=$ $\displaystyle
T^{j}_{n}\overrightarrow{V}^{j}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right),$
then we have
$\displaystyle\left(i\partial_{t}+\Delta\right)\overrightarrow{V}^{j}=0,\quad\overrightarrow{V}^{j}(0)=\varphi^{j}.$
Now let $u^{j}_{n}(t,x)$ be the nonlinear solution of (1.1) with initial data
$v^{j}_{n}(0)$, that is
$\displaystyle\left(i\partial_{t}+\Delta\right)\overrightarrow{u}^{j}_{n}(t,x)=$
$\displaystyle\left<\nabla\right>f_{1}(\left<\nabla\right>^{-1}\overrightarrow{u}^{j}_{n})+\left<\nabla\right>f_{2}(\left<\nabla\right>^{-1}\overrightarrow{u}^{j}_{n}),$
$\displaystyle\quad\overrightarrow{u}^{j}_{n}(0)=$
$\displaystyle\overrightarrow{v}^{j}_{n}(0)=T^{j}_{n}\overrightarrow{V}^{j}(\tau^{j}_{n}),\quad
u^{j}_{n}(0)\in\mathcal{K}^{+},$
where $\tau^{j}_{n}=-t^{j}_{n}/(h^{j}_{n})^{2}$. In order to look for the
nonlinear profile $\overrightarrow{U}^{j}_{\infty}$ associated to the radial,
free concentrating wave
$\left(\overrightarrow{v}^{j}_{n};\;h^{j}_{n},t^{j}_{n}\right)$, we also need
undo the group action. We denote
$\displaystyle\overrightarrow{u}^{j}_{n}(t,x)=$ $\displaystyle
T^{j}_{n}\overrightarrow{U}^{j}_{n}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right),$
then we have
$\displaystyle\left(i\partial_{t}+\Delta\right)\overrightarrow{U}^{j}_{n}=$
$\displaystyle\left(\left<\nabla\right>^{j}_{n}\right)f_{1}\left(\left(\left<\nabla\right>^{j}_{n}\right)^{-1}\overrightarrow{U}^{j}_{n}\right)+h^{j}_{n}\cdot\left(\left<\nabla\right>^{j}_{n}\right)f_{2}\left(\left(\left<\nabla\right>^{j}_{n}\right)^{-1}\overrightarrow{U}^{j}_{n}\right),$
$\displaystyle\overrightarrow{U}^{j}_{n}(\tau^{j}_{n})=$
$\displaystyle\overrightarrow{V}^{j}(\tau^{j}_{n}).$
Up to a subsequence, we may assume that there exist
$h^{j}_{\infty}\in\\{0,1\\}$ and $\tau^{j}_{\infty}\in[-\infty,\infty]$ for
every $j=\\{0,\ldots,K\\}$, such that
$\displaystyle
h^{j}_{n}\rightarrow\;h^{j}_{\infty},\;\;\text{and}\;\;\tau^{j}_{n}\rightarrow\;\tau^{j}_{\infty}.$
As $n\rightarrow+\infty$, the limit equation of $\overrightarrow{U}^{j}_{n}$
is given by
$\displaystyle\left(i\partial_{t}+\Delta\right)\overrightarrow{U}^{j}_{\infty}=$
$\displaystyle\left(\left<\nabla\right>^{j}_{\infty}\right)f_{1}\left(\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{U}^{j}_{\infty}\right)+h^{j}_{\infty}\cdot\left(\left<\nabla\right>^{j}_{\infty}\right)f_{2}\left(\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{U}^{j}_{\infty}\right),$
$\displaystyle\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{\infty})=$
$\displaystyle\overrightarrow{V}^{j}(\tau^{j}_{\infty})\in
L^{2}({\mathbb{R}}^{3}).$
Let
$\displaystyle\widehat{U}^{j}_{\infty}:=\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{U}^{j}_{\infty},$
then
$\displaystyle\left(i\partial_{t}+\Delta\right)\widehat{U}^{j}_{\infty}=$
$\displaystyle f_{1}\left(\widehat{U}^{j}_{\infty}\right)+h^{j}_{\infty}\cdot
f_{2}\left(\widehat{U}^{j}_{\infty}\right),$ (5.17)
$\displaystyle\widehat{U}^{j}_{\infty}(\tau^{j}_{\infty})=$
$\displaystyle\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{V}^{j}(\tau^{j}_{\infty}).$
(5.18)
The unique existence of a local radial solution
$\overrightarrow{U}^{j}_{\infty}$ around $\tau^{j}_{\infty}$ is known in all
cases, including $h^{j}_{\infty}=0$ and $\tau^{j}_{\infty}=\pm\infty$.
$\overrightarrow{U}^{j}_{\infty}$ on the maximal existence interval is called
the nonlinear profile associated with the radial, free concentrating wave
$\left(\overrightarrow{v}^{j}_{n};\;h^{j}_{n},t^{j}_{n}\right)$.
The nonlinear concentrating wave $u^{j}_{(n)}$ associated with
$\left(\overrightarrow{v}^{j}_{n};\;h^{j}_{n},t^{j}_{n}\right)$ is defined by
$\displaystyle\overrightarrow{u}^{j}_{(n)}(t,x)=T^{j}_{n}\overrightarrow{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right),$
then we have
$\displaystyle\left(i\partial_{t}+\Delta\right)\overrightarrow{u}^{j}_{(n)}=$
$\displaystyle\left(\sqrt{|\nabla|^{2}+\left(\frac{h^{j}_{\infty}}{h^{j}_{n}}\right)^{2}}\right)f_{1}\left(\left(\sqrt{|\nabla|^{2}+\left(\frac{h^{j}_{\infty}}{h^{j}_{n}}\right)^{2}}\right)^{-1}\overrightarrow{u}^{j}_{(n)}\right)$
$\displaystyle+\frac{h^{j}_{\infty}}{h^{j}_{n}}\cdot\left(\sqrt{|\nabla|^{2}+\left(\frac{h^{j}_{\infty}}{h^{j}_{n}}\right)^{2}}\right)f_{2}\left(\left(\sqrt{|\nabla|^{2}+\left(\frac{h^{j}_{\infty}}{h^{j}_{n}}\right)^{2}}\right)^{-1}\overrightarrow{u}^{j}_{(n)}\right)$
$\displaystyle=$
$\displaystyle\left<\nabla\right>^{j}_{\infty}f_{1}\left(\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{u}^{j}_{(n)}\right)+h^{j}_{\infty}\cdot\left<\nabla\right>^{j}_{\infty}f_{2}\left(\left(\left<\nabla\right>^{j}_{\infty}\right)^{-1}\overrightarrow{u}^{j}_{(n)}\right),$
$\displaystyle\overrightarrow{u}^{j}_{(n)}(0)=$ $\displaystyle
T^{j}_{n}\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{n}),$
which implies that
$\displaystyle\big{\|}\overrightarrow{u}^{j}_{(n)}(0)-\overrightarrow{u}^{j}_{n}(0)\big{\|}_{L^{2}}=$
$\displaystyle\big{\|}T^{j}_{n}\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{n})-T^{j}_{n}\overrightarrow{V}^{j}(\tau^{j}_{n})\big{\|}_{L^{2}}=\big{\|}\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{n})-\overrightarrow{V}^{j}(\tau^{j}_{n})\big{\|}_{L^{2}}$
$\displaystyle\leq$
$\displaystyle\big{\|}\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{n})-\overrightarrow{U}^{j}_{\infty}(\tau^{j}_{\infty})\big{\|}_{L^{2}}+\big{\|}\overrightarrow{V}^{j}(\tau^{j}_{n})-\overrightarrow{V}^{j}(\tau^{j}_{\infty})\big{\|}_{L^{2}}\rightarrow
0.$
We denote
$\displaystyle\overrightarrow{u}^{j}_{(n)}=\left<\nabla\right>u^{j}_{(n)}.$
If $h^{j}_{\infty}=1$, we have $h^{j}_{n}\equiv 1$, then $u^{j}_{(n)}\in
H^{1}({\mathbb{R}}^{3})$ is radial and satisfies
$\displaystyle\left(i\partial_{t}+\Delta\right)u^{j}_{(n)}=f_{1}(u^{j}_{(n)})+f_{2}(u^{j}_{(n)}).$
If $h^{j}_{\infty}=0$, then $u^{j}_{(n)}\in H^{1}({\mathbb{R}}^{3})$ is radial
and satisfies
$\displaystyle\left(i\partial_{t}+\Delta\right)u^{j}_{(n)}=\frac{|\nabla|}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{|\nabla|}u^{j}_{(n)}\right).$
Let $u_{n}$ be a sequence of (local) radial solutions of (1.1) with initial
data in $\mathcal{K}^{+}$ at $t=0$, and let $v_{n}$ be the sequence of the
radial, free solutions with the same initial data. We consider the linear
profile decomposition given by Proposition 5.1
$\displaystyle\overrightarrow{v}_{n}(t,x)=\sum^{k-1}_{j=0}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x),\quad\overrightarrow{v}^{j}_{n}(t^{j}_{n})=T^{j}_{n}\varphi^{j},\quad
v^{j}_{n}(0)\in\mathcal{K}^{+}.$
With each free concentrating wave
$\\{\overrightarrow{v}^{j}_{n}\\}_{n\in{\mathbb{N}}}$, we associate the
nonlinear concentrating wave
$\\{\overrightarrow{u}^{j}_{(n)}\\}_{n\in{\mathbb{N}}}$. A nonlinear profile
decomposition of $u_{n}$ is given by
$\displaystyle\overrightarrow{u}^{<k}_{(n)}(t,x):=\sum^{k-1}_{j=0}\overrightarrow{u}^{j}_{(n)}(t,x)=\sum^{k-1}_{j=0}T^{j}_{n}\overrightarrow{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right).$
(5.19)
Since the smallness condition (5.2) and the orthogonality condition (5.3)
ensure that every nonlinear concentrating wave and the remainder interacts
weakly with the others, we will show that
$\overrightarrow{u}^{<k}_{(n)}+\overrightarrow{w}^{k}_{n}$ is a good
approximation for $\overrightarrow{u}_{n}$ provided that each nonlinear
profile has the finite global Strichartz norm.
Now we define the Strichartz norms for the nonlinear profile decomposition.
Let $ST(I)$ and $ST^{*}(I)$ be the function spaces on
$I\times{\mathbb{R}}^{3}$ defined as Section 4
$\displaystyle ST(I):=\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\cap L^{6}_{t}\dot{B}^{1/2}_{18/7,2}\cap
L^{5}_{t,x}\right)$ $\displaystyle(I\times{\mathbb{R}}^{3}),$ $\displaystyle
ST^{*}(I):=\left(L^{2}_{t}\dot{B}^{1/3}_{18/11,2}\cap
L^{2}_{t}\dot{B}^{1/2}_{6/5,2}\right)(I\times{\mathbb{R}}^{3})$ .
The Strichartz norm for the nonlinear profile $\widehat{U}^{j}_{\infty}$
depends on the scaling $h^{j}_{\infty}$.
$\displaystyle
ST^{j}_{\infty}(I):=\begin{cases}ST(I),\quad&\text{for}\;h^{j}_{\infty}=1,\\\
\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\right)(I\times{\mathbb{R}}^{3}),\quad&\text{for}\;h^{j}_{\infty}=0.\end{cases}$
###### Lemma 5.6.
In the nonlinear profile decomposition (5.19). Suppose that for each $j<K$, we
have
$\displaystyle\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{ST^{j}_{\infty}({\mathbb{R}})}+\big{\|}\overrightarrow{U}^{j}_{\infty}\big{\|}_{L^{\infty}_{t}L^{2}_{x}({\mathbb{R}}^{3})}<\infty.$
Then for any finite interval $I$, any $j<K$ and any $k\leq K$, we have
$\displaystyle\varlimsup_{n\rightarrow+\infty}\big{\|}u^{j}_{(n)}\big{\|}_{ST(I)}\lesssim$
$\displaystyle\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{ST^{j}_{\infty}({\mathbb{R}})},$
(5.20)
$\displaystyle\varlimsup_{n\rightarrow+\infty}\big{\|}u^{<k}_{(n)}\big{\|}^{2}_{ST(I)}\lesssim$
$\displaystyle\varlimsup_{n\rightarrow+\infty}\sum_{j<k}\big{\|}u^{j}_{(n)}\big{\|}^{2}_{ST(I)},$
(5.21)
where the implicit constants do not depend on $I,j$ or $k$. We also have
$\displaystyle\lim_{n\rightarrow+\infty}\left\|f_{1}\left(u^{<k}_{(n)}\right)-\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}(I)}=0,$
(5.22)
$\displaystyle\lim_{n\rightarrow+\infty}\left\|f_{2}\left(u^{<k}_{(n)}\right)-\sum_{j<k}h^{j}_{\infty}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{2}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}(I)}=0.$
(5.23)
###### Proof.
Proof of (5.20). By the definitions of $u^{j}_{(n)}$ and
$\widehat{U}^{j}_{\infty}$, we know that
$\displaystyle u^{j}_{(n)}(t,x)=$
$\displaystyle\left<\nabla\right>^{-1}\overrightarrow{u}^{j}_{(n)}(t,x)=\left<\nabla\right>^{-1}T^{j}_{n}\overrightarrow{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right)$
$\displaystyle=$
$\displaystyle\left<\nabla\right>^{-1}T^{j}_{n}\left<\nabla\right>^{j}_{\infty}\widehat{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right)=h^{j}_{n}T^{j}_{n}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>^{j}_{n}}\widehat{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right).$
For the case $h^{j}_{\infty}=1$, we have
$u^{j}_{(n)}(t,x)=\widehat{U}^{j}_{\infty}(t-t^{j}_{n},x)$, hence (5.20) is
trivial. For the case $h^{j}_{\infty}=0$, by the above relation between
$u^{j}_{(n)}$ and $\widehat{U}^{j}_{\infty}$, we have
$\displaystyle\big{\|}u^{j}_{(n)}\big{\|}_{\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\right)(I\times{\mathbb{R}}^{3})}\lesssim$
$\displaystyle\left\|\frac{|\nabla|}{\left<\nabla\right>^{j}_{n}}\widehat{U}^{j}_{\infty}\right\|_{\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\right)({\mathbb{R}}\times{\mathbb{R}}^{3})}$
$\displaystyle\lesssim$
$\displaystyle\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\right)({\mathbb{R}}\times{\mathbb{R}}^{3})},$
and
$\displaystyle\big{\|}u^{j}_{(n)}\big{\|}_{L^{6}_{t}\dot{B}^{1/2}_{18/7,2}(I\times{\mathbb{R}}^{3})}\lesssim$
$\displaystyle|I|^{\frac{1}{12}}\big{\|}u^{j}_{(n)}\big{\|}_{L^{12}_{t}\dot{B}^{1/2}_{18/7,2}}\lesssim|I|^{\frac{1}{12}}(h^{j}_{n})^{\frac{1}{3}}\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{L^{12}_{t}\dot{B}^{\frac{5}{6}}_{18/7,2}}\rightarrow
0,$
$\displaystyle\big{\|}u^{j}_{(n)}\big{\|}_{L^{5}_{t,x}(I\times{\mathbb{R}}^{3})}\lesssim$
$\displaystyle|I|^{\frac{7}{60}}\big{\|}u^{j}_{(n)}\big{\|}_{L^{12}_{t}L^{5}_{x}}\lesssim|I|^{\frac{7}{60}}(h^{j}_{n})^{\frac{4}{15}}\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{L^{12}_{t}\dot{B}^{\frac{4}{15}}_{5,2}}\rightarrow
0.$
where we use the fact that the boundedness of $\widehat{U}^{j}_{\infty}$ in
$L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap L^{12}_{t}L^{9}_{x}\cap
L^{\infty}\dot{H}^{1}$ implies its boundedness in
$L^{12}_{t}\dot{B}^{\frac{5}{6}}_{18/7,2}\cap
L^{12}_{t}\dot{B}^{\frac{4}{15}}_{5,2}$ by (5.17).
Proof of (5.21). We estimate the left hand side of (5.21) by
$\displaystyle\big{\|}u^{<k}_{(n)}\big{\|}^{2}_{ST(I)}=$
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n)}+\sum_{j<k,h^{j}_{\infty}=0}u^{j}_{(n)}\right\|^{2}_{ST(I)}$
$\displaystyle\lesssim$
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n)}\right\|^{2}_{ST(I)}+\left\|\sum_{j<k,h^{j}_{\infty}=0}u^{j}_{(n)}\right\|^{2}_{ST(I)}.$
For the case $h^{j}_{\infty}=1$. Define $\widehat{U}^{j}_{\infty,R}$ and
$u^{j}_{(n),R}$ by
$\displaystyle\widehat{U}^{j}_{\infty,R}(t,x)=\chi_{R}(t,x)\widehat{U}^{j}_{\infty}(t,x),\quad
u^{j}_{(n),R}(t,x)=T^{j}_{n}\widehat{U}^{j}_{\infty,R}(t-t^{j}_{n}),$
where $\chi_{R}$ is the cut-off function as in Remark 1.6. Then we have
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n)}\right\|^{2}_{ST(I)}\lesssim$
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n),R}\right\|^{2}_{ST(I)}+\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n)}-\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n),R}\right\|^{2}_{ST(I)}.$
On one hand, we know that
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n)}-\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n),R}\right\|_{ST(I)}\leq\sum_{j<k,h^{j}_{\infty}=1}\big{\|}(1-\chi_{R})u^{j}_{(n)}\big{\|}_{ST(I)}\rightarrow
0,$
as $R\rightarrow+\infty$. On the other hand, by (5.3) and the similar
orthogonality analysis as in [17], we know that
$\displaystyle\varlimsup_{n\rightarrow+\infty}\left\|\sum_{j<k,h^{j}_{\infty}=1}u^{j}_{(n),R}\right\|^{2}_{ST(I)}\lesssim\varlimsup_{n\rightarrow+\infty}\sum_{j<k,h^{j}_{\infty}=1}\left\|u^{j}_{(n),R}\right\|^{2}_{ST(I)}\lesssim\varlimsup_{n\rightarrow+\infty}\sum_{j<k,h^{j}_{\infty}=1}\left\|u^{j}_{(n)}\right\|^{2}_{ST(I)}.$
For the case $h^{j}_{\infty}=0$, On one hand, by $h^{j}_{n}\rightarrow 0$, we
have
$\displaystyle\varlimsup_{n\rightarrow+\infty}\left\|\sum_{j<k,h^{j}_{\infty}=0}u^{j}_{(n)}\right\|_{\left(L^{6}_{t}\dot{B}^{1/2}_{18/7,2}\cap
L^{5}_{t,x}\right)(I\times{\mathbb{R}}^{3})}=0.$
On the other hand, by (5.3) and the analogue approximation analysis as in
[17], we have
$\displaystyle\varlimsup_{n\rightarrow+\infty}\left\|\sum_{j<k,h^{j}_{\infty}=0}u^{j}_{(n)}\right\|^{2}_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}(I\times{\mathbb{R}}^{3})}\lesssim$
$\displaystyle\varlimsup_{n\rightarrow+\infty}\sum_{j<k,h^{j}_{\infty}=0}\left\|u^{j}_{(n)}\right\|^{2}_{L^{10}_{t}\dot{B}^{1/3}_{90/19,2}(I\times{\mathbb{R}}^{3})},$
$\displaystyle\varlimsup_{n\rightarrow+\infty}\left\|\sum_{j<k,h^{j}_{\infty}=0}u^{j}_{(n)}\right\|^{2}_{L^{12}_{t}L^{9}_{x}(I\times{\mathbb{R}}^{3})}\lesssim$
$\displaystyle\varlimsup_{n\rightarrow+\infty}\sum_{j<k,h^{j}_{\infty}=0}\left\|u^{j}_{(n)}\right\|^{2}_{L^{12}_{t}L^{9}_{x}(I\times{\mathbb{R}}^{3})}.$
Proof of (5.22). Let $\displaystyle
u^{<k}_{<n>}(t,x):=\sum_{j<k}u^{j}_{<n>}(t,x)$ where
$\displaystyle
u^{j}_{<n>}(t,x):=\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}=\frac{1}{\left<\nabla\right>^{j}_{\infty}}\overrightarrow{u}^{j}_{(n)}=\frac{1}{\left<\nabla\right>^{j}_{\infty}}T^{j}_{n}\overrightarrow{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right)=h^{j}_{n}T^{j}_{n}\widehat{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right),$
and
$\displaystyle
u^{j}_{(n)}(t,x)=h^{j}_{n}T^{j}_{n}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>^{j}_{n}}\widehat{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right).$
Then we have
$\displaystyle\left\|f_{1}\left(u^{<k}_{(n)}\right)-\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}}$
$\displaystyle\leq\left\|f_{1}\left(u^{<k}_{(n)}\right)-f_{1}\left(u^{<k}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|f_{1}\left(u^{<k}_{\left<n\right>}\right)-\sum_{j<k}f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}$
$\displaystyle\quad+\left\|\sum_{j<k}f_{1}\left(u^{j}_{\left<n\right>}\right)-\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}$
$\displaystyle\leq\left\|f_{1}\left(u^{<k}_{(n)}\right)-f_{1}\left(u^{<k}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|f_{1}\left(u^{<k}_{\left<n\right>}\right)-\sum_{j<k}f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}$
$\displaystyle\quad+\left\|\sum_{j<k,h^{j}_{\infty}=0}f_{1}\left(u^{j}_{\left<n\right>}\right)-\sum_{j<k,h^{j}_{\infty}=0}\frac{|\nabla|}{\left<\nabla\right>}f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}.$
By (5.3) and the approximation argument in [17], we have
$\displaystyle\left\|f_{1}\left(u^{<k}_{(n)}\right)-f_{1}\left(u^{<k}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|f_{1}\left(u^{<k}_{\left<n\right>}\right)-\sum_{j<k}f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}\longrightarrow
0$
as $n\rightarrow+\infty$. In addition, by $h^{j}_{n}\rightarrow 0$ as
$n\rightarrow+\infty$, we have
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=0}\left(1-\frac{|\nabla|}{\left<\nabla\right>}\right)f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{L^{2}_{t}\dot{B}^{1/3}_{18/11,2}}=$
$\displaystyle\left\|\left(1-\frac{|\nabla|}{\left<\nabla\right>^{j}_{n}}\right)\sum_{j<k,h^{j}_{\infty}=0}f_{1}\left(\widehat{U}^{j}_{\infty}\right)\right\|_{L^{2}_{t}\dot{B}^{1/3}_{18/11,2}}\longrightarrow
0,$
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=0}\left(1-\frac{|\nabla|}{\left<\nabla\right>}\right)f_{1}\left(u^{j}_{\left<n\right>}\right)\right\|_{L^{2}_{t}\dot{B}^{1/2}_{6/5,2}}=$
$\displaystyle\left(h^{j}_{n}\right)^{1/2}\left\|\left(1-\frac{|\nabla|}{\left<\nabla\right>^{j}_{n}}\right)\sum_{j<k,h^{j}_{\infty}=0}f_{1}\left(\widehat{U}^{j}_{\infty}\right)\right\|_{L^{2}_{t}\dot{B}^{1/2}_{6/5,2}}\longrightarrow
0,$
as $n\rightarrow+\infty$. Therefore, we have
$\displaystyle\lim_{n\rightarrow+\infty}\left\|f_{1}\left(u^{<k}_{(n)}\right)-\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}}=0.$
Proof of (5.23). Note that
$\displaystyle\left\|f_{2}\left(u^{<k}_{(n)}\right)-\sum_{j<k}h^{j}_{\infty}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{2}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}}$
$\displaystyle\leq\left\|f_{2}\left(u^{<k}_{(n)}\right)-f_{2}\left(u^{<k}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|f_{2}\left(u^{<k}_{\left<n\right>}\right)-\sum_{j<k}f_{2}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|\sum_{j<k,h^{j}_{\infty}=0}f_{2}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}.$
By the analogue analysis, we have
$\displaystyle\left\|f_{2}\left(u^{<k}_{(n)}\right)-f_{2}\left(u^{<k}_{\left<n\right>}\right)\right\|_{ST^{*}}+\left\|f_{2}\left(u^{<k}_{\left<n\right>}\right)-\sum_{j<k}f_{2}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}\longrightarrow
0,$
and
$\displaystyle\left\|\sum_{j<k,h^{j}_{\infty}=0}f_{2}\left(u^{j}_{\left<n\right>}\right)\right\|_{ST^{*}}\longrightarrow
0$
as $n\rightarrow+\infty$. Hence, we obtain
$\displaystyle\lim_{n\rightarrow+\infty}\left\|f_{2}\left(u^{<k}_{(n)}\right)-\sum_{j<k}h^{j}_{\infty}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{2}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)\right\|_{ST^{*}}=0.$
These complete the proof. ∎
After this preliminaries, we now show that
$\overrightarrow{u}^{<k}_{(n)}+\overrightarrow{w}^{k}_{n}$ is a good
approximation for $\overrightarrow{u}_{n}$ provided that each nonlinear
profile has finite global Strichartz norm.
###### Proposition 5.7.
Let $u_{n}$ be a sequence of local, radial solutions of (1.1) around $t=0$ in
$\mathcal{K}^{+}$ satisfying
$\displaystyle
M\left(u_{n}\right)<\infty,\quad\varlimsup_{n\rightarrow\infty}E(u_{n})<m.$
Suppose that in the nonlinear profile decomposition (5.19), every nonlinear
profile $\widehat{U}^{j}_{\infty}$ has finite global Strichartz and energy
norms we have
$\displaystyle\big{\|}\widehat{U}^{j}_{\infty}\big{\|}_{ST^{j}_{\infty}({\mathbb{R}})}+\big{\|}\overrightarrow{U}^{j}_{\infty}\big{\|}_{L^{\infty}_{t}L^{2}_{x}({\mathbb{R}}^{3})}<\infty.$
Then $u_{n}$ is bounded for large $n$ in the Strichartz and the energy norms
$\displaystyle\varlimsup_{n\rightarrow\infty}\big{\|}u_{n}\big{\|}_{ST({\mathbb{R}})}+\big{\|}\overrightarrow{u}_{n}\big{\|}_{L^{\infty}_{t}L^{2}_{x}({\mathbb{R}})}<\infty.$
###### Proof.
We only need to verify the condition of Proposition 4.1. Note that
$u^{<k}_{(n)}+w^{k}_{n}$ satisfies that
$\displaystyle\left(i\partial_{t}+\Delta\right)\left(u^{<k}_{(n)}+w^{k}_{n}\right)=f_{1}\left(u^{<k}_{(n)}+w^{k}_{n}\right)+f_{2}\left(u^{<k}_{(n)}+w^{k}_{n}\right)$
$\displaystyle\qquad+f_{1}\left(u^{<k}_{(n)}\right)-f_{1}\left(u^{<k}_{(n)}+w^{k}_{n}\right)+f_{2}\left(u^{<k}_{(n)}\right)-f_{2}\left(u^{<k}_{(n)}+w^{k}_{n}\right)$
$\displaystyle\qquad+\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)-f_{1}\left(u^{<k}_{(n)}\right)+\sum_{j<k}h^{j}_{\infty}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{2}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)-f_{2}\left(u^{<k}_{(n)}\right).$
First, by the construction of $\overrightarrow{u}^{<k}_{(n)}$, we know that
$\displaystyle\left\|\left(\overrightarrow{u}^{<k}_{(n)}(0)+\overrightarrow{w}^{k}_{n}(0)\right)-\overrightarrow{u}_{n}(0)\right\|_{L^{2}_{x}}\leq\sum_{j<k}\left\|\overrightarrow{u}^{j}_{(n)}(0)-\overrightarrow{u}^{j}_{n}(0)\right\|_{L^{2}_{x}}\rightarrow
0,$
as $n\rightarrow+\infty$, which also implies that for large $n$, we have
$\displaystyle\left\|\overrightarrow{u}^{<k}_{(n)}(0)+\overrightarrow{w}^{k}_{n}(0)\right\|_{L^{2}_{x}}\leq
E_{0}.$
Next, by the linear profile decomposition in Proposition 5.1, we know that
$\displaystyle\big{\|}u_{n}(0)\big{\|}^{2}_{L^{2}}=$
$\displaystyle\left\|v_{n}(0)\right\|^{2}_{L^{2}_{x}}=\sum_{j<k}\left\|v^{j}_{n}(0)\right\|^{2}_{L^{2}_{x}}+\left\|w^{k}_{n}(0)\right\|^{2}_{L^{2}_{x}}+o_{n}(1)$
$\displaystyle\geq$
$\displaystyle\sum_{j<k}\left\|v^{j}_{n}(0)\right\|^{2}_{L^{2}_{x}}+o_{n}(1)=\sum_{j<k}\left\|u^{j}_{(n)}(0)\right\|^{2}_{L^{2}_{x}}+o_{n}(1),$
$\displaystyle\left\|u_{n}(0)\right\|^{2}_{\dot{H}^{1}_{x}}=$
$\displaystyle\left\|v_{n}(0)\right\|^{2}_{\dot{H}^{1}_{x}}=\sum_{j<k}\left\|v^{j}_{n}(0)\right\|^{2}_{\dot{H}^{1}_{x}}+\left\|w^{k}_{n}(0)\right\|^{2}_{\dot{H}^{1}_{x}}+o_{n}(1)$
$\displaystyle\geq$
$\displaystyle\sum_{j<k}\left\|v^{j}_{n}(0)\right\|^{2}_{\dot{H}^{1}_{x}}+o_{n}(1)=\sum_{j<k}\left\|u^{j}_{(n)}(0)\right\|^{2}_{\dot{H}^{1}_{x}}+o_{n}(1),$
which means except for a finite set $J\subset{\mathbb{N}}$, the energy of
$u^{j}_{(n)}$ with $j\not\in J$ is smaller than the iteration threshold, hence
we have
$\displaystyle\big{\|}u^{j}_{(n)}\big{\|}_{ST({\mathbb{R}})}\lesssim\big{\|}\overrightarrow{u}^{j}_{(n)}(0)\big{\|}_{L^{2}_{x}},$
thus, for any finite interval $I$, by Lemma 5.6, we have
$\displaystyle\sup_{k}\varlimsup_{n\rightarrow+\infty}\big{\|}u^{<k}_{(n)}\big{\|}^{2}_{ST(I)}\lesssim$
$\displaystyle\sup_{k}\varlimsup_{n\rightarrow+\infty}\sum_{j<k}\big{\|}u^{j}_{(n)}\big{\|}^{2}_{ST(I)}$
$\displaystyle=$
$\displaystyle\sup_{k}\varlimsup_{n\rightarrow+\infty}\left[\sum_{j<k,j\in
J}\big{\|}u^{j}_{(n)}\big{\|}^{2}_{ST(I)}+\sum_{j<k,j\not\in
J}\big{\|}u^{j}_{(n)}\big{\|}^{2}_{ST(I)}\right]$ $\displaystyle\lesssim$
$\displaystyle\sum_{j<k,j\in
J}\big{\|}\widehat{U}^{j}_{\infty}\big{\|}^{2}_{ST^{j}_{\infty}(I)}+\sup_{k}\varlimsup_{n\rightarrow+\infty}\sum_{j<k,j\not\in
J}\big{\|}\overrightarrow{u}^{j}_{(n)}(0)\big{\|}^{2}_{L^{2}_{x}}$
$\displaystyle<$ $\displaystyle\infty.$
This together with the Strichartz estimate for $w^{k}_{n}$ implies that
$\displaystyle\sup_{k}\varlimsup_{n\rightarrow+\infty}\big{\|}u^{<k}_{(n)}+w^{k}_{n}\big{\|}^{2}_{ST(I)}<\infty.$
Last we need show the nonlinear perturbation is small in some sense. By
Proposition 5.1 and Lemma 5.6, we have
$\displaystyle\left\|f_{1}\left(u^{<k}_{(n)}\right)-f_{1}\left(u^{<k}_{(n)}+w^{k}_{n}\right)\right\|_{ST^{*}(I)}\rightarrow
0,$
$\displaystyle\left\|f_{2}\left(u^{<k}_{(n)}\right)-f_{2}\left(u^{<k}_{(n)}+w^{k}_{n}\right)\right\|_{ST^{*}(I)}\rightarrow
0,$
and
$\displaystyle\left\|\sum_{j<k}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{1}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)-f_{1}\left(u^{<k}_{(n)}\right)\right\|_{ST^{*}(I)}\rightarrow
0,$
$\displaystyle\left\|\sum_{j<k}h^{j}_{\infty}\frac{\left<\nabla\right>^{j}_{\infty}}{\left<\nabla\right>}f_{2}\left(\frac{\left<\nabla\right>}{\left<\nabla\right>^{j}_{\infty}}u^{j}_{(n)}\right)-f_{2}\left(u^{<k}_{(n)}\right)\right\|_{ST^{*}(I)}\rightarrow
0,$
as $n\rightarrow+\infty$. Therefore, by Proposition 4.1, we can obtain the
desired result, which concludes the proof. ∎
## 6\. Part II: GWP and Scattering for $\mathcal{K}^{+}$
After the stability analysis of the scattering solution of (1.1) and the
compactness analysis (linear and nonlinear profile decompositions) of a
sequence of the radial solutions of (1.1) in the energy space. We now use them
to show the scattering result of Theorem 1.3 by contradiction.
Let $E^{*}$ be the threshold for the uniform Strichartz norm bound, i.e.,
$\displaystyle E^{*}:=\sup\\{A>0,ST(A)<\infty\\}$
where $ST(A)$ denotes the supremum of $\big{\|}u\big{\|}_{ST(I)}$ for any
strong radial solution $u$ of (1.1) in $\mathcal{K}^{+}$ on any interval $I$
satisfying $E(u)\leq A$, $M(u)<\infty$.
The small solution scattering theory gives us $E^{*}>0$.
Now we are going to show that $E^{*}\geq m$ by contradiction. From now on,
suppose that $E^{*}\geq m$ fails, that is, we assume that
$\displaystyle E^{*}<m.$ (6.1)
### 6.1. Existence of a critical element
In this subsection, by the profile decomposition and the stability theory of
the scattering solution of (1.1), we show the existence of the critical
element, which is the radial, energy solution of (1.1) with the smallness
energy $E^{*}$ and infinite Strichartz norm.
By the definition of $E^{*}$ and the fact that $E^{*}<m$, there exist a
sequence of radial solutions $\\{u_{n}\\}_{n\in{\mathbb{N}}}$ of (1.1) in
$\mathcal{K}^{+}$, which have the maximal existence interval $I_{n}$ and
satisfy that
$\displaystyle M(u_{n})<\infty,\quad E(u_{n})\rightarrow
E^{*}<m,\quad\big{\|}u_{n}\big{\|}_{ST(I_{n})}\rightarrow+\infty,\quad\text{as}\;\;n\rightarrow+\infty,$
then we have $\big{\|}u_{n}\big{\|}_{H^{1}}<\infty$ by Lemma 2.12. By the
compact argument (profile decomposition) and the stability theory, we can show
that
###### Theorem 6.1.
Let $u_{n}$ be a sequence of radial solutions of (1.1) in $\mathcal{K}^{+}$ on
$I_{n}\subset{\mathbb{R}}$ satisfying
$\displaystyle M(u_{n})<\infty,\quad E(u_{n})\rightarrow
E^{*}<m,\quad\big{\|}u_{n}\big{\|}_{ST(I_{n})}\rightarrow+\infty,\quad\text{as}\;\;n\rightarrow+\infty.$
Then there exists a global, radial solution $u_{c}$ of (1.1) in
$\mathcal{K}^{+}$ satisfying
$\displaystyle E(u_{c})=E^{*}<m,\quad
K(u_{c})>0,\quad\big{\|}u_{c}\big{\|}_{ST({\mathbb{R}})}=\infty.$
In addition, there are a sequence $t_{n}\in{\mathbb{R}}$ and radial function
$\varphi\in L^{2}({\mathbb{R}}^{3})$ such that, up to a subsequence, we have
as $n\rightarrow+\infty$,
$\displaystyle\left\|\frac{|\nabla|}{\left<\nabla\right>}\Big{(}\overrightarrow{u}_{n}(0,x)-e^{-it_{n}\Delta}\varphi(x)\Big{)}\right\|_{L^{2}}\rightarrow
0.$ (6.2)
###### Proof.
By the time translation symmetry of (1.1), we can translate $u_{n}$ in $t$
such that $0\in I_{n}$ for all $n$. Then by the linear and nonlinear profile
decomposition of $u_{n}$, we have
$\displaystyle e^{it\Delta}\overrightarrow{u}_{n}(0,x)=$
$\displaystyle\sum_{j<k}\overrightarrow{v}^{j}_{n}(t,x)+\overrightarrow{w}^{k}_{n}(t,x),\quad\overrightarrow{v}^{j}_{n}(t,x)=e^{i(t-t^{j}_{n})\Delta}T^{j}_{n}\varphi^{j},$
$\displaystyle\overrightarrow{u}^{<k}_{(n)}(t,x)=$
$\displaystyle\sum_{j<k}\overrightarrow{u}^{j}_{(n)}(t,x),\quad\overrightarrow{u}^{j}_{(n)}(t,x)=T^{j}_{n}\overrightarrow{U}^{j}_{\infty}\left(\frac{t-t^{j}_{n}}{(h^{j}_{n})^{2}}\right),$
$\displaystyle\big{\|}\overrightarrow{u}^{j}_{(n)}(0)-\overrightarrow{v}^{j}_{n}(0)\big{\|}_{L^{2}}\rightarrow
0.$
By Proposition 5.5 and the following observations that
1. (1)
Every radial solution of (1.1) in $\mathcal{K}^{+}$ with the energy less than
$E^{*}$ has global finite Strichartz norm by the definition of $E^{*}$.
2. (2)
Lemma 5.7 precludes that all the nonlinear profiles
$\overrightarrow{U}^{j}_{\infty}$ have finite global Strichartz norm.
we deduce that there is only one radial profile and
$\displaystyle E(u^{0}_{(n)}(0))\rightarrow E^{*},\quad
u^{0}_{(n)}(0)\in\mathcal{K}^{+},\quad\big{\|}\widehat{U}^{0}_{\infty}\big{\|}_{ST^{0}_{\infty}(I)}=\infty,\quad\big{\|}w^{1}_{n}\big{\|}_{L^{\infty}_{t}\dot{H}^{1}_{x}}\rightarrow
0.$
If $h^{0}_{n}\rightarrow 0$, then
$\widehat{U}^{0}_{\infty}=|\nabla|^{-1}\overrightarrow{U}^{0}_{\infty}$ solves
the $\dot{H}^{1}$-critical NLS
$\displaystyle\left(i\partial_{t}+\Delta\right)\widehat{U}^{0}_{\infty}=f_{1}(\widehat{U}^{0}_{\infty})$
and satisfies
$\displaystyle
E^{c}\left(\widehat{U}^{0}_{\infty}(\tau^{0}_{\infty})\right)=E^{*}<m,\;\;K^{c}\left(\widehat{U}^{0}_{\infty}(\tau^{0}_{\infty})\right)\geq
0,\;\;\big{\|}\widehat{U}^{0}_{\infty}\big{\|}_{\left(L^{10}_{t}\dot{B}^{1/3}_{90/19,2}\cap
L^{12}_{t}L^{9}_{x}\right)(I\times{\mathbb{R}}^{3})}=\infty.$
However, it is in contradiction with Kenig-Merle’s result444By the global
$L^{10}_{t,x}$ estimate of solution $u$ of (1.2), we can obtain the global
$L^{q}_{t}\dot{W}^{1,r}_{x}$ estimate of $u$ for any Schrödinger
$L^{2}$-admissible pair $(q,r)$. in [19]. Hence $h^{0}_{n}\equiv 1$, which
implies (6.2).
Now we show that
$\widehat{U}^{0}_{\infty}=\left<\nabla\right>^{-1}\overrightarrow{U}^{j}_{\infty}$
is a global solution, which is the consequence of the compactness of (6.2).
Suppose not, then we can choose a sequence $t_{n}\in{\mathbb{R}}$ which
approaches the maximal existence time. Since
$\widehat{U}^{0}_{\infty}(t+t_{n})$ satisfies the assumption of this theorem,
then applying the above argument to it, we obtain that for some $\psi\in
L^{2}$ and another sequence $t^{\prime}_{n}\in{\mathbb{R}}$, as
$n\rightarrow+\infty$
$\displaystyle\left\|\frac{|\nabla|}{\left<\nabla\right>}\left(\overrightarrow{U}^{0}_{\infty}(t_{n})-e^{-it^{\prime}_{n}\Delta}\psi(x)\right)\right\|_{L^{2}}\rightarrow
0.$ (6.3)
Let $\overrightarrow{v}(t):=e^{it\Delta}\psi.$ For any $\varepsilon>0$, there
exist $\delta>0$ with $I=[-\delta,\delta]$ such that
$\displaystyle\big{\|}\left<\nabla\right>^{-1}\overrightarrow{v}(t-t^{\prime}_{n})\big{\|}_{ST(I)}\leq\varepsilon,$
which together with (6.3) implies that for sufficiently large $n$
$\displaystyle\big{\|}\left<\nabla\right>^{-1}e^{it\Delta}\overrightarrow{U}^{0}_{\infty}(t_{n})\big{\|}_{ST(I)}\leq\varepsilon.$
If $\varepsilon$ is small enough, this implies that the solution
$\widehat{U}^{0}_{\infty}$ exists on $[t_{n}-\delta,t_{n}+\delta]$ for large
$n$ by the small data theory. This contradicts the choice of $t_{n}$. Hence
$\widehat{U}^{0}_{\infty}$ is a global solution and it is just the desired
critical element $u_{c}$. By Proposition 1.1, we know that $K(u_{c})>0$. ∎
### 6.2. Compactness of the critical element
In order to preclude the critical element, we need obtain some useful
properties about the critical element. In the following subsections, we
establish some properties about the critical element by its minimal energy
with infinite Strichartz norm, especially its compactness and its consequence.
Since (1.1) is symmetric in $t$, we may assume that
$\displaystyle\big{\|}u_{c}\big{\|}_{ST(0,+\infty)}=\infty,$ (6.4)
we call it a forward critical element.
###### Proposition 6.2.
Let $u_{c}$ be a forward critical element. Then the set
$\displaystyle\\{u_{c}(t,x);0<t<\infty\\}$
is precompact in $\dot{H}^{s}$ for any $s\in(0,1]$.
###### Proof.
By the conservation of the mass, it suffices to prove the precompactness of
$u_{c}(t_{n})\\}$ in $\dot{H}^{1}$ for any positive time $t_{1},t_{2},\ldots$.
If $t_{n}$ converges, then it is trivial from the continuity in $t$.
If $t_{n}\rightarrow+\infty$. Applying Theorem 6.1 to the sequence of
solutions $\overrightarrow{u}_{c}(t+t_{n})$, we get another sequence
$t^{\prime}_{n}\in{\mathbb{R}}$ and radial function $\varphi\in L^{2}$ such
that
$\displaystyle\frac{|\nabla|}{\left<\nabla\right>}\left(\overrightarrow{u}_{c}(t_{n},x)-e^{-it^{\prime}_{n}\Delta}\varphi(x)\right)\rightarrow
0\quad\text{in}\;\;L^{2}.$
1. (1)
If $t^{\prime}_{n}\rightarrow-\infty$, then we have
$\displaystyle\big{\|}\left<\nabla\right>^{-1}e^{it\Delta}\overrightarrow{u}_{c}(t_{n})\big{\|}_{ST(0,+\infty)}=\big{\|}\left<\nabla\right>^{-1}e^{it\Delta}\varphi\big{\|}_{ST(-t^{\prime}_{n},+\infty)}+o_{n}(1)\rightarrow
0.$
Hence $u_{c}$ can solve (1.1) for $t>t_{n}$ with large $n$ globally by
iteration with small Strichartz norms, which contradicts (6.4).
2. (2)
If $t^{\prime}_{n}\rightarrow+\infty$, then we have
$\displaystyle\big{\|}\left<\nabla\right>^{-1}e^{it\Delta}\overrightarrow{u}_{c}(t_{n})\big{\|}_{ST(-\infty,0)}=\big{\|}\left<\nabla\right>^{-1}e^{it\Delta}\varphi\big{\|}_{ST(-\infty,-t^{\prime}_{n})}+o_{n}(1)\rightarrow
0$
Hence $u_{c}$ can solve (1.1) for $t<t_{n}$ with large $n$ with vanishing
Strichartz norms, which implies $u_{c}=0$ by taking the limit, which is a
contradiction.
Thus $t^{\prime}_{n}$ is bounded, which implies that $t^{\prime}_{n}$ is
precompact, so is $u_{c}(t_{n},x)$ in $\dot{H}^{1}$. ∎
As a consequence, the energy of $u_{c}$ stays within a fixed radius for all
positive time, modulo arbitrarily small rest. More precisely, we define the
exterior energy by
$\displaystyle E_{R}(u;t)=\int_{|x|\geq R}\Big{(}\big{|}\nabla
u(t,x)\big{|}^{2}+\big{|}u(t,x)\big{|}^{4}+\big{|}u(t,x)\big{|}^{6}\Big{)}\;dx$
for any $R>0$. Then we have
###### Corollary 6.3.
Let $u_{c}$ be a forward critical element. then for any $\varepsilon$, there
exist $R_{0}(\varepsilon)>0$ such that
$\displaystyle E_{R_{0}}(u_{c};t)\leq\varepsilon E(u_{c}),\;\text{for
any}\;t>0.$
### 6.3. Death of the critical element
We are in a position to preclude the soliton-like solution by a truncated
Virial identity.
###### Theorem 6.4.
The critical element $u_{c}$ of (1.1) cannot be a soliton in the sense of
Theorem 6.1.
###### Proof.
We still drop the subscript $c$. Now let $\phi$ be a smooth, radial function
satisfying $0\leq\phi\leq 1$, $\phi(x)=1$ for $|x|\leq 1$, and $\phi(x)=0$ for
$|x|\geq 2$. For some $R$, we define
$\displaystyle
V_{R}(t):=\int_{{\mathbb{R}}^{3}}\phi_{R}(x)|u(t,x)|^{2}\;dx,\quad\phi_{R}(x)=R^{2}\phi\left(\frac{|x|^{2}}{R^{2}}\right).$
On one hand, we have
$\displaystyle\partial_{t}V_{R}(t)=4\Im\int_{{\mathbb{R}}^{3}}\phi^{\prime}\left(\frac{|x|^{2}}{R^{2}}\right)x\cdot\nabla
u(t,x)\;\overline{u(t,x)}\;dx.$
Therefore, we have
$\displaystyle\big{|}\partial_{t}V_{R}(t)\big{|}\lesssim R$ (6.5)
for all $t\geq 0$ and $R>0$.
On the other hand, by Lemma 2.5 and Hölder’s inequality, we have
$\displaystyle\partial^{2}_{t}V_{R}(t)=4\int_{{\mathbb{R}}^{3}}\phi_{R}^{\prime\prime}(r)\big{|}\nabla
u(t,x)\big{|}^{2}\;dx-\int_{{\mathbb{R}}^{3}}(\Delta^{2}\phi_{R})(x)|u(t,x)|^{2}\;dx$
$\displaystyle\qquad\quad-\frac{4}{3}\int_{{\mathbb{R}}^{3}}(\Delta\phi_{R})(x)|u(t,x)|^{6}\;dx+\int_{{\mathbb{R}}^{3}}(\Delta\phi_{R})(x)|u(t,x)|^{4}\;dx$
$\displaystyle=$ $\displaystyle 4\int_{{\mathbb{R}}^{3}}\left(2|\nabla
u(t,x)|^{2}-2|u(t,x)|^{6}+\frac{3}{2}|u(t,x)|^{4}\right)\;dx$ $\displaystyle+$
$\displaystyle O\left(\int_{|x|\geq R}\left(|\nabla
u(t,x)|^{2}+|u(t,x)|^{6}+|u(t,x)|^{4}\right)\;dx+\left(\int_{R\leq|x|\leq
2R}|u(t,x)|^{6}\;dx\right)^{1/3}\right)$ $\displaystyle=$ $\displaystyle
4K\left(u(t)\right)+O\left(\int_{|x|\geq R}\left(|\nabla
u(t,x)|^{2}+|u(t,x)|^{4}\right)\;dx+\left(\int_{R\leq|x|\leq
2R}|u(t,x)|^{6}\;dx\right)^{1/3}\right).$
By Lemma 2.13, we have
$\displaystyle 4K\left(u(t)\right)=$
$\displaystyle\;4\int_{{\mathbb{R}}^{3}}\left(2|\nabla
u(t,x)|^{2}-2|u(t,x)|^{6}+\frac{3}{2}|u(t,x)|^{4}\right)\;dx$
$\displaystyle\gtrsim$
$\displaystyle\min\left(6(m-E(u(t))),\frac{2}{3}\big{\|}\nabla
u(t)\big{\|}^{2}_{L^{2}}+\frac{1}{2}\big{\|}u(t)\big{\|}^{4}_{L^{4}}\right)$
$\displaystyle\gtrsim$ $\displaystyle\big{\|}\nabla
u(t)\big{\|}^{2}_{L^{2}}+\big{\|}u(t)\big{\|}^{4}_{L^{4}}$
$\displaystyle\gtrsim$ $\displaystyle E(u(t)),$
Thus, choosing $\eta>0$ sufficiently small and $\displaystyle R:=C(\eta)$ and
by Corollary 6.3, we obtain
$\displaystyle\partial^{2}_{t}V_{R}(t)\gtrsim E(u(t))=E(u_{0}),$
which implies that for all $T_{1}>T_{0}$
$\displaystyle(T_{1}-T_{0})E(u_{0})\lesssim R=C(\eta).$
Taking $T_{1}$ sufficiently large, we obtain a contradiction unless $u\equiv
0$. But $u\equiv 0$ is not consistent with the fact that
$\big{\|}u\big{\|}_{ST({\mathbb{R}})}=\infty$. ∎
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|
arxiv-papers
| 2011-11-29T02:37:43 |
2024-09-04T02:49:24.694440
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Changxing Miao, Guixiang Xu and Lifeng Zhao",
"submitter": "Changxing Miao",
"url": "https://arxiv.org/abs/1111.6671"
}
|
1111.6746
|
F. Iwamuro et al.FMOS image-based reduction package
instrumentation: spectrographs — methods: data analysis
# FIBRE-pac: FMOS image-based reduction package
Fumihide Iwamuro11affiliation: Department of Astronomy, Kyoto University,
Kitashirakawa, Kyoto, Japan Yuuki Moritani11affiliationmark: Kiyoto
Yabe11affiliationmark: Masanao Sumiyoshi11affiliationmark: Kaori
Kawate11affiliationmark: Naoyuki Tamura22affiliation: Subaru Telescope,
National Astronomical Observatory of Japan, Hilo, HI, USA Masayuki
Akiyama33affiliation: Astronomical Institute, Tohoku University, Aramaki,
Sendai, Japan Masahiko Kimura22affiliationmark: Naruhisa
Takato22affiliationmark: Philip Tait22affiliationmark: Kouji
Ohta11affiliationmark: Tomonori Totani11affiliationmark: Yuji
Suzuki11affiliationmark: and Motonari Tonegawa11affiliationmark:
iwamuro@kusastro.kyoto-u.ac.jp
###### Abstract
The FIBRE-pac (FMOS image-based reduction package) is an IRAF-based reduction
tool for the fiber multiple-object spectrograph (FMOS) of the Subaru
telescope. To reduce FMOS images, a number of special techniques are necessary
because each image contains about 200 separate spectra with airglow emission
lines variable in spatial and time domains, and with complicated throughput
patterns for the airglow masks. In spite of these features, almost all of the
reduction processes except for a few steps are carried out automatically by
scripts in text format making it easy to check the commands step by step.
Wavelength- and flux-calibrated images together with their noise maps are
obtained using this reduction package.
## 1 Introduction
FMOS (Kimura et al., 2010) is the near-infrared fiber multi-object
spectrograph that has been in operation as one of the open-use instruments on
the Subaru telescope since 2009. It can configure 400 fibers of $1\farcs 2$
aperture in a 30′ diameter field of view at the primary focus. The 400
infrared spectra in two groups are taken by two spectrographs called IRS1 and
IRS2 (infrared spectrograph 1 and 2) in either of two modes: a low-resolution
mode with a spectral resolution of $\Delta\lambda\sim$ 20 Å in the 0.9-1.8
$\mu$m wavelength range, and a high-resolution mode of $\Delta\lambda\sim$ 5 Å
in one of the quarter wavelength ranges of 0.90-1.14, 1.01-1.25, 1.11-1.35,
1.40-1.62, 1.49-1.71, or 1.59-1.80 $\mu$m. The bright OH-airglow emission
lines are masked by the mask mirror (Iwamuro et al., 2001) installed in these
spectrographs.
The individual tools for image reduction in this package were developed during
engineering and guaranteed-time observations since 2008, in conjunction with
other operation and reduction software. Although these tools were not
developed as part of the public reduction software of FMOS, they are now
stable, thus they have earned the name of “FIBRE-pac” (FMOS Image-Based
REduction package). The basic concepts underlying the package are as follows:
1. 1.
Most of the reduction is processed by IRAF.
2. 2.
Several complicated steps are processed by original software written in C
using the cfitsio library (Pence, 1999).
3. 3.
Almost all of the reduction processes are automated using script files in text
format.
4. 4.
Modification of the text files is done by general UNIX commands.
5. 5.
The original 2-dimensional information is kept as far as possible throughout
the reduction processes.
These concepts enable easy implementation and open processing without the
inconvenience of licensed or black box parts and ensure traceable operation
with visual confirmation. The 2-dimensional information has advantages not
only in filtering out unexpected noise using their small sizes but also on
detection of faint emission-features. In this paper, we describe the reduction
process of the FMOS images based on its Aug.9-2011 version, taking care of
complex conditions in the infrared, using multiple fibers and the OH-
suppressed spectrograph.
## 2 The images
The FMOS images are acquired by a uniform interval non-destructive readout
technique called “Up-the-Ramp sampling” (hereafter “ramp sampling” for short).
A typical exposure of 900 seconds consists of 54 images. After an exposure is
finished, a final frame (treated as a “raw image” in the reduction process) is
created where the count in each pixel is calculated by performing a least
squares fit to the signal count of 54 images. In addition to suppressing
readout noise, the advantages of the ramp sampling are 1) saturated pixels can
be estimated from the counts prior to saturation, and 2) cosmic-ray events can
be detected as an unexpected jump in the counts and removed from the final
frame (Fixsen et al., 2000). The detection threshold of the cosmic-ray events
has been currently set to 10$\sigma$ in an empirical way based on the real
FMOS images. For IRS1, the fit and the cosmic-ray rejection are done during
the ramp sampling, so that the final frame is ready as the exposure finishes.
For IRS2, however, the nonlinear bias variation prevents fitting the slope
during the exposure. Instead, the ramp fitting is executed after the simple
background subtraction (cf. §3.2) has been performed for all of the images
taken during the sampling. For example, 54 background-subtracted ramp images
are prepared prior to fitting. Consequently the nonlinear bias component is
subtracted together with the background photons.
(80mm,80mm)Figure1.eps
Figure 1: Raw image taken in the IRS1 low-resolution mode with a ramp readout
of 54 times using the HAWAII-2 2k$\times$2k detector (corresponding to a 900-s
exposure). The left half of the image consists of spectra in the $J$ band,
while the right half is in the $H$ band. The typical FWHM of the each fiber
spectrum is 5 pixels with a pitch of 10 pixels between the spectra.
Figure 1 shows the resulting “raw image” after ramp sampling 54 times using
IRS1. Raw images for IRS2 do not exist for the reasons explained above. The
raw images have the following components:
1. 1.
Thermal background ($\sim$300 e/900 s)
2. 2.
Suppressed OH-airglow ($\sim$600 e/900 s on average)
3. 3.
Remaining cosmic ray events ($\sim$200 e/event)
4. 4.
Read noise ($\sim$10 e rms/54 ramp readout)
5. 5.
Object ($\sim$30 e/900 s for a 20 mag(AB) object)
6. 6.
Dark current ($<$10 e/900 s)
7. 7.
Bias offset ($<$10 e/exposure)
8. 8.
Cross-talk between quadrants ($\sim$0.15% of the count)
9. 9.
Bad pixels having no efficiency, unusual large dark current, or too large
noise ($\sim$0.2% of the pixels)
10. 10.
On-chip amplifier glow at the corners of the four quadrants ($\sim$2000 e/54
ramp readout)
11. 11.
Unknown external noise in a stripe pattern (quite rare and the amplitude is
$\sim$100 e/900 s at most)
12. 12.
Point spread function (PSF) of the fiber of 5 pixels with a pitch of 10 pixels
between spectra (having almost no cross talk of the spectra on either side)
13. 13.
Position of an individual fiber along the slit
14. 14.
Optical distortion
15. 15.
Quantum efficiency and its pixel-to-pixel variation (“flat” pattern)
16. 16.
System throughput variation with wavelength
17. 17.
Atmospheric transmittance
18. 18.
Intrinsic absorption and emission of the objects
The thermal background can be removed in the initial subtraction described in
subsection 3.2. The 2nd strongest component, residual airglow, can be
subtracted by interpolation of the background of other fibers after the
optical distortion is corrected using a dome-flat image taken in the same
observation mode. The remaining cosmic-ray and other strong noise features
having a size smaller than the PSF of a fiber are removed in three subsequent
stages with different threshold counts. Besides the science frames, the
following images are necessary for this reduction process.
Detector-flat image
: The homogeneous thermal image of the black board attached to the entrance
window of the camera dewar. This image gives the pixel-to-pixel quantum
efficiency ratio of the detector. More than 30 sets of the ramp sampling
frames are averaged to ensure the high signal-to-noise (S/N) ratio. This image
is included in the reduction package together with the other standard images
such as a bad pixel mask.
Bad pixel mask
: The distribution of bad pixels in the detector selected through the
reduction process of the detector-flat image.
Dome-flat image
: The dome-flat spectra taken before or after observation. This image is used
for measurements of the distortion parameters as well as the relative
throughput correction of the scientific images.
Th-Ar spectral image
: The emission-line spectra for wavelength calibration taken just before the
exposures of dome-flat frames.
The details of the reduction processes using these images are described in the
following section.
## 3 Reduction Process
### 3.1 Preparatory processing
Before performing a reduction of the science frames, a preparatory process is
applied to the dome-flat and Th-Ar spectral images to determine the optical
distortion and the wavelength calibration of the image.
First, Flat fielding: The dome-flat image is divided by the detector-flat
image to remove differences in the quantum efficiency between pixels.
Second, Bad pixel correction: The registered and temporally prominent bad
pixels, which are picked up by subtraction of 3$\times$3 median filtered
image, are replaced with an interpolated value from the surrounding pixels.
Third, Correction of the spatial distortion: The $y$-axis of the image is
converted using $y^{\prime}=y+(a_{1}y+a_{2})x^{2}+(b_{1}y+b_{2})x$. (Here, the
origin is at the center of the image.) The four parameters are chosen to make
the PSF amplitude in the image projected along the $x$-axis maximum (cf.
figure 2). In this modification process, the dome-flat spectrum from each
fiber with 9 pixels in width is converted into a parallel line to make a
“combined 1D image” (as in figure 3) that includes the pattern of the airglow
mask.
(80mm,79mm)Figure2a.eps (80mm,80mm)Figure2b.eps
Figure 2: Correction of the spatial distortion of the dome-flat image. The
$y$-axis of the left-hand image is converted to straighten the spectra, as
shown at right. The thick lines indicate the shape of the spectra at the top
and bottom edges of the detector. The left and right half regions are
separately corrected in this low-resolution mode, because the corresponding
mask mirrors are separated (Kimura et al., 2010). The projected images of each
half region along the $x$-axis are shown along the sides of the image.
Fourth, Correction of the spectral distortion: The $x$-axis of the combined 1D
image is converted using $x^{\prime}=x+(ax+b)$, in which the two parameters
are determined for each line so as to minimize the shift and the magnification
difference in the mask pattern. The 200 resulting $a$ and $b$ parameters are
fitted with fourth-order polynomials in $y$ to remove the local matching error
in the pattern. This conversion process makes the mask pattern straight in the
column direction (cf. figure 3), which is necessary for better subtraction of
the residual airglow lines in the science frames.
These parameters obtained from the dome-flat image are also applied to the Th-
Ar spectral image to make a combined 1D image, as shown in figure 4. Although
the slit image of each emission line and of the airglow masks should be
parallel, alignment errors cause small differences between them.
Finally, Wavelength calibration: The correlation between the observed
wavelength and the pixels of each spectrum are determined, as represented by
$\lambda=px^{3}+qx^{2}+rx+s$. The resulting $p$, $q$, $r$, and $s$
coefficients are used to calculate the corresponding wavelengths in each
spectrum without changing them, because the $x$ positions of the individual
fibers in the slit may exhibit some scatter due to imperfect fiber alignment
along the slit. The results of the wavelength calibration are confirmed by
comparing the reduced Th-Ar spectra with an artificial image based on the
known wavelengths of the Th-Ar emission lines (cf. figure 4). The typical
calibration error of the spectra is less than 1 pixel, corresponding to 5 Å or
1.2 Å in the low- or high-resolution mode, respectively.
(80mm,18mm)Figure3.eps
Figure 3: Correction of the spectral distortion of the combined 1D image of
the dome flat. The $x$-axis of the combined 1D image (at top) is converted
using $x^{\prime}=x+(ax+b)$, in which $a$ and $b$ are coefficients of a
fourth-order polynomial function in $y$. The images of the airglow masks are
consequently straightened in the corrected image (at bottom). The thick lines
indicate the shape of the airglow mask at the left and the right edges of the
detector.
(80mm,18mm)Figure4.eps
Figure 4: Combined 1D image of the Th-Ar spectra (at bottom) compared with an
artificial image based on the known wavelengths of the Th-Ar emission lines
(at top). The agreement in the position of the lines shows that the conversion
from the column numbers to wavelengths is correct.
### 3.2 Initial background subtraction
Since the usual FMOS observations are carried out using an $ABAB$ nodding
pattern of the telescope, $A-B$ simple sky subtractions can be performed using
two different sky images: $A_{n}-B_{n-1}$ and $A_{n}-B_{n}$ (where $A_{n}$
denotes the image taken at position $A$ of the $n$-th pair). For IRS2, these
sky-subtracted frames are calculated by the ramp fitting algorithm applied to
the sky-subtracted sub-frames. When the brightness of the OH airglow varies
monotonically, most of it can be canceled by merging these two images
according to $A_{n}-B=(A_{n}-B_{n})w+(A_{n}-B_{n-1})(1-w)$. The weight $w$ is
chosen to make the sum of the absolute count $|A_{n}-B|$ a minimum within the
range $-1<w<2$. Figure 5 shows a pair of sky-subtracted images and the merged
one. Typically, the weight $w$ is equal to about 0.5.
(50mm,50mm)Figure5a.eps (50mm,50mm)Figure5b.eps (50mm,50mm)Figure5c.eps
Figure 5: Canceling the residual OH-airglow emission. Left panel: Sky-
subtracted image using the previous image as the sky image ($A_{n}-B_{n-1}$).
Center panel: Sky-subtracted image using the next image as the sky image
($A_{n}-B_{n}$). Right panel: Merged image ($A_{n}-B$) with a weight of
$w=0.472$ in which the negative spectra cannot be used because they have
merged information with various weights.
### 3.3 Corrections of detector cross talk, bias difference, and bad pixels
After the initial background subtraction, cross talk is removed by subtracting
0.15% for IRS1 and 1% for IRS2 from each quadrant. Next, the bias difference
between the quadrants (as indicated in figure 5, the lower left quadrant tends
to show a higher bias level) is corrected to make the average over each
quadrant equal. After flat fielding using the detector flat image, the
registered and temporally prominent bad pixels are rejected, together with the
four adjacent pixels by interpolating the surrounding pixels.
### 3.4 Distortion correction and residual sky subtraction
The processed image is converted into a combined 1D image based on the
distortion parameters obtained in the preparatory reduction process. However,
only one line of each spectrum (9 pixels in width) is extracted instead of
summing the 9 lines. One can thereby make a set of 9 combined 1D images each
of which consists of a different part of each spectrum (cf. figure 6). In
other words, the PSF of a fiber is divided into 9 pieces, only one of which is
used in a combined 1D image. The flexure or temperature change in the
spectrograph causes a small difference between the dome-flat and the
scientific images along the vertical direction in position. This difference is
corrected using the vertical position of the spectra of bright stars. In this
way, the counts of the residual sky becomes smooth along the columns after the
relative throughput correction of the fibers. The residual airglow lines are
fitted and subtracted in these images, and then the relative throughput
difference is multiplied to restore the noise level back to the original
state. The residual subtracted images are recombined to form an image in which
the PSF of the fiber is determined. Medium-level bad pixels having smaller
size than the PSF of a fiber are replaced at the end of this process.
(50mm,50mm)Figure6a.eps (50mm,50mm)Figure6b.eps (50mm,50mm)Figure6c.eps
Figure 6: Residual background subtraction and recomposition process. Left
panel: A set of 9 combined 1D images from different part of the PSF. Center
panel: Residual subtracted images by performing fits along the columns in each
image. Right panel: Recomposed image from the 9 residual subtracted images.
### 3.5 Combine and residual background subtraction
There are two ways to observe scientific targets with the $ABAB$ nodding
pattern: “normal beam switching” (NBS), in which all targets are observed at
position $A$ (all the fibers are supposed to observe “blank” sky at position
$B$), and “cross beam switching” (CBS), in which less than half of the fibers
are allocated to the targets at position $A$ while the others are at position
$B$. In the CBS observation mode, the same targets appear in both images
collected in positions $A$ or $B$. However, the target spectra in position $B$
are merged to minimize the absolute flux of the residual airglow lines in the
initial background subtraction process. Thus the same reduction process has to
be applied to the images in position $B$, replacing the negative spectra of
position $A$ with the corresponding parts of the position $B$ image (as shown
in figure 7). The merged images (or all of the position $A$ images in the NBS
mode) are then combined into one averaged image.
Finally, the combined image is divided into a set of 9 combined 1D images
again, in order to perform a fine subtraction of the residual background. The
recomposed image also goes through a bad pixel rejection process.
(50mm,50mm)Figure7a.eps (50mm,50mm)Figure7b.eps (50mm,50mm)Figure7c.eps
Figure 7: Merging the CBS images from positions $A$ and $B$. Left panel:
Positive spectra in the reduced $A_{n}-B$ image in which the removed negative
(position $B$) spectra contain merged information (cf. figure 5). Center
panel: Negative spectra in the reduced $A-B_{n}$ image in which the removed
positive (position $A$) spectra are merged one. Right panel: Merged
$A_{n}-B_{n}$ image.
### 3.6 Mask edge correction and CBS combine
Although the correction for spectral distortions was performed to straighten
the images of the airglow masks (as in figure 3), there remain small
differences among the throughput patterns of the spectra because of the local
shape errors of the mask elements or because of the presence of dust on the
mask mirror. These differences are corrected by dividing the averaged image by
a “relative” dome-flat image in which the common spectral features have been
removed by normalizing the count along lines and columns (cf. figure 8). After
correcting the relative differences of the throughput patterns among the
spectra, the corresponding spectra from positions $A$ (positive) and $B$
(negative) taken in the CBS observation mode are ready to be combined. The
negative spectra in the image are then inverted and rearranged so that they
can be combined with the corresponding positive spectra (as shown in figure
9).
(80mm,18mm)Figure8.eps
Figure 8: Relative dome-flat image to correct the throughput differences
between spectra. Top panel: Flux-normalized dome-flat spectra in which the
throughput differences of the fiber have been removed. Bottom panel: Relative
dome-flat image created by normalizing the top image along the columns. This
image is displayed within a range of 0.8 (black) to 1.2 (white) to emphasize
the relative differences.
(50mm,50mm)Figure9a.eps (50mm,50mm)Figure9b.eps (50mm,50mm)Figure9c.eps
Figure 9: Combining a pair of target spectra taken in the CBS observation
mode. Left panel: Final $A-B$ image in which the positions of the targets
observed only at position $B$ are masked. Center panel: Inverted and
rearranged negative spectrum that will be combined with the corresponding
positive spectrum. Right panel: Combined spectrum. The S/N ratios of the
positive spectra are 1.4 times higher than the ratio of the corresponding
negative spectra in the combined image.
### 3.7 Object mask
In the process of fitting and subtract the residual background, some fraction
of the object flux can be subtracted. The best way to retain the object flux
is to mask the objects during the fit. Faint objects to be masked are
therefore selected by a human eye on the combined image and the reduction
process is repeated with these masks applied. The mask density should not be
too high to make the residual background subtraction work well: the maximum
density allowed is roughly 75% of the image.
### 3.8 Square-noise map
Before the mask edge correction process, the distribution of noise in an image
is truncated Poisson distribution caused by the bad pixel rejection process in
three subsequent stages. Although the noise level in an image is almost
homogeneous at this stage, the noise level map becomes complicated during the
throughput correction and CBS combine. To estimate the noise level of each
pixel, a frame is defined that consists of the squared noise level measured
along columns just before the mask edge correction. Here, the noise level is
measured by a 3$\sigma$ clipping algorithm iterated ten times for each column,
while some extra contribution proportional to the count is added to the
clipped pixels. This square-noise frame is divided by the squared image of the
relative dome flat, and reduced to half wherever a pair of spectra is averaged
during the CBS combine. Note that the noise level of the bright objects is
probably underestimated because the systematic uncertainty (i.e. sub-pixel
shifts of the spectra in raw images and tiny variation of the throughput
pattern caused by the instrument status such as temperature) is the major
contribution for them.
## 4 Flux calibration and check process
### 4.1 Outline
To calibrate the flux of the scientific targets, at least one bright
($J(AB)=15$–$18$mag) star in each science frames is needed as a spectral
reference, because it must have almost the same atmospheric absorption feature
as that for the scientific targets in the same field of view. All the spectra
are divided by the reference spectrum, and then multiplied by the expected
spectrum of the reference star whose flux and spectral type are known or can
be determined by the observed values in two different wavelength regions. If
the cataloged or estimated spectrum of the reference star is correct, all the
observed spectra will then be calibrated accurately. In the next subsection,
the method to estimate the reference star spectrum is described.
### 4.2 Template stellar spectra
Since the flux and slope of the spectrum of a reference star are determined
from the measured counts in an image, one needs the slope-removed template
spectra including the intrinsic absorption features of the star. We analyzed
128 stellar spectra in the IRTF Spectral
Library111http://irtfweb.ifa.hawaii.edu/spex/IRTF_Spectral_Library/References.html
(Rayner et al., 2009) split into seven groups: F0-F9IV/V (19 objects),
G0-G8IV/V (14 objects), K0-K7IV/V (10 objects), F0-F9I/II/III (21 objects),
G0-G9I/II/III (31 objects), K0-K3I/II/III (17 objects), and K4-K7I/II/III (16
objects). In each group, the spectra were averaged to improve the S/N ratio
and divided by the results of the linear fitting, so that slope-removed
template spectra of seven different types were prepared (cf. figure 10). Here,
only F, G, and K type stars were used to make the slope-removed template
spectra, because 1) their slope-removed spectra become roughly straight in
$\lambda$ \- $F_{\nu}$ plot, 2) these stars are quite popular and easy to
select in the target field, and 3) A and earlier type spectra are neither
available in the IRTF Spectral Library, nor in other database with similar
qualities.
Next, the correlations between the spectral types and the slopes of the
spectra have to be established. Figure 11 shows the distribution of measured
value of 128 stellar spectra from figure 10. The stellar types are numbered
from 0 (F0) to 29 (K9) while the slopes are defined by
$slope=\frac{F_{\nu}(1.55)-F_{\nu}(1.21)}{(1.55-1.21)F_{\nu}(1.31)}$ (1)
in $\mu{\rm m}^{-1}$. The correlations between the spectral types and the
slopes are determined by second-order polynomial fits to the distributions of
stellar types III and V:
$slope_{\rm III}=0.00132type^{2}+0.0229type-0.666,$ $slope_{\rm
V}=0.000456type^{2}+0.0253type-0.654.$ (2)
As a result, any spectrum from F0 to K9 can be synthesized by multiplying a
linear spectrum having a defined slope with the intrinsic absorption spectrum
from the interpolation of the nearest two slope-removed template spectra.
(80mm,80mm)Figure10a.ps (80mm,80mm)Figure10b.ps
Figure 10: Stellar spectra in the IRTF Spectral Library. The slope-removed
averaged spectrum of each stellar type (thick black lines) is used as the
intrinsic absorption template.
(80mm,80mm)Figure11.ps
Figure 11: Correlation between a stellar type and the slope of its spectrum.
The stellar types are numbered from 0 (F0) to 29 (K9), and the slope is
defined by equation (1). The open symbols are measured from figure 10: type I
(triangles), type II (diamonds), type III (squares), type IV (pentagons), and
type V (circles). The distributions of types III and V are fit to second-order
polynomials represented by the thick solid and dashed lines, respectively. The
thick dotted line indicate the intermediate correlation between types III and
V used when the stellar type of a reference star is unknown.
### 4.3 Calibration process
The first step in the calibration process is the relative throughput
correction of fibers in which the averaged throughput of position $A$ and $B$
is used for the merged spectra in the CBS combine process. The next step is to
remap the pixels with an increment of 5 Å/pixel (1.25 Å/pixel in the high-
resolution mode) based on the correlation between the observed wavelength and
the pixels determined in the preparatory reduction process. In this remapping
process, the observed wavelengths are multiplied by the count of each pixel in
order to convert the value from photon count to $F_{\nu}$. Next, the converted
count is divided by the atmospheric transmittance function of the airmass
present during the observation, so as to roughly correct for the effects of
atmospheric absorption. After this correction, the flux of the reference star
is estimated from the converted count at around 1.31$\mu$m, assuming that the
total system efficiency under good seeing condition is 2.5% including losses
at the entrance of the fibers. The slope of the spectrum is then measured with
a fixed efficiency ratio between 1.21 and 1.55 $\mu$m. (The value of this
ratio will be confirmed in the check process, along with the total system
efficiency of 2.5%.) Finally, the observed scientific spectra are divided by
the reference spectrum and multiplied by the stellar spectra from the measured
flux and slope. Here, the type V and III absorption templates are adopted as
the reference spectra, respectively bluer than G5 and redder than K1 (cf.
figure 11). As a consequence, two of the intrinsic absorption templates of FV,
GV, K1III, and K4III (in figure 10) are used to interpolate the absorption of
the reference spectrum. If the stellar type of the reference star is known,
the corresponding synthesized spectrum is used instead of this predicted
spectrum.
A resulting wavelength- and flux-calibrated image is shown in figure 12. The
size of this image is 1800$\times$1800 pixels, the wavelength in $\mu$m is
$\lambda=0.9+0.0005\times(ColumnNumber-1)\ $, the $n$-th spectrum is located
between $y=9\times(LineNumber)-8$ and $y=9\times(LineNumber)$, and the count
is in $\mu$Jy. The square-noise frame is converted in a similar way except
that all factors are multiplied twice. The 1D spectrum of each object is
extracted from this image with a user-defined mask of 9 pixels, together with
the square-noise frame. An example of the final 1D spectrum is shown in figure
13.
(80mm,80mm)Figure12a.eps (80mm,80mm)Figure12b.eps
Figure 12: Calibrated spectral image (at left) and a square-noise frame (at
right) with a size of 1800$\times$1800 pixels. The wavelength range is 0.9-1.8
$\mu$m with an increment of 5 Å/pixel. Each of the 200 spectra has a width of
9 pixels.
(80mm,80mm)Figure13.ps
Figure 13: Example of a processed 1D spectrum. The spectrum and the
corresponding 1$\sigma$ noise level are represented by the thick continuous
line and the thin dotted line, respectively. The wavelength has been converted
to the rest wavelength using the cataloged redshift value ($z=0.844$).
### 4.4 Check process
The results are checked by comparing the resulting flux in the $J$ and $H$
bands with the photometric data in the catalog. If all the factors
contributing to the system efficiency are normal, the flux should be
consistent with the catalog value. The accuracy of the efficiency ratio
between 1.21 and 1.55 $\mu$m can also be confirmed by these diagrams for most
cases. Figure 14 shows an example of such a comparison. However, weather
conditions, telescope and instrument focus, extended (non-stellar) morphology
of targets, and position accuracy of the catalog may cause the lower observed
flux of some targets than that expected from the catalog. Under sufficiently
good conditions, the observed flux should match the catalog magnitude as shown
in the thick line in figure 14 with some downward scatter due to small
allocation error of the fibers. Also, if the estimated efficiency ratio
between 1.21 and 1.55 is not correct, the points of $J$ and $H$ will have
small offset in this figure.
(80mm,80mm)Figure14.ps
Figure 14: Comparison of the observed flux with the catalog magnitude. The
triangles and circles represent the values in the $J$ and $H$ band,
respectively. The thick solid lines indicate the ideal values without loss of
flux.
## 5 Summary and Conclusions
The reduction process for FMOS images includes two special processing steps.
One is the segmented processing of the spectra to handle a given part of the
PSF as a unit, while the other is the automatic modeling of the reference
spectrum to calibrate the scientific targets. The segmented processing enable
to keep the original 2-dimensional information which has a large effect on bad
pixel filtering and detection of faint emission-lines. Most of the processes
are carried out automatically, but the object mask preparation and the
reference star selection require user judgement. The FIBRE-pac is available
from the FMOS instrument
page222http://www.naoj.org/Observing/Instruments/FMOS/ of the Subaru web site,
together with the sample dataset presented in this paper. The base reduction
platform is IRAF and the reduction scripts and sources are free and open, so
that users can check what is happening in each step by sending the commands
one at a time.
This work was supported by a Grant-in-Aid for Scientific Research (B) of Japan
(22340044) and by a Grant-in-Aid for the Global COE Program ”The Next
Generation of Physics, Spun from Universality and Emergence” from the Ministry
of Education, Culture, Sports, Science, and Technology (MEXT) of Japan.
IRAF is distributed by the National Optical Astronomy Observatories, which are
operated by the Association of Universities for Research in Astronomy, Inc.,
under cooperative agreement with the National Science Foundation.
## References
* Fixsen et al. (2000) Fixsen, D. J., Offenberg, J. D., Hanisch, R. J., et al. 2000, PASP, 112, 1350
* Iwamuro et al. (2001) Iwamuro, F., Motohara, K., Maihara, T., Hata, R., & Harashima, T. 2001, PASJ, 53, 355
* Kimura et al. (2010) Kimura, M., et al. 2010, PASJ, 62, 1135
* Pence (1999) Pence, W. 1999, Astronomical Data Analysis Software and Systems VIII, 172, 487
* Rayner et al. (2009) Rayner, J. T., Cushing, M. C., & Vacca, W. D. 2009, ApJS, 185, 289
|
arxiv-papers
| 2011-11-29T10:11:12 |
2024-09-04T02:49:24.715880
|
{
"license": "Public Domain",
"authors": "F. Iwamuro, Y. Moritani, K. Yabe, M. Sumiyoshi, K. Kawate, N. Tamura,\n M. Akiyama, M. Kimura, N. Takato, P. Tait, K. Ohta, T. Totani, Y. Suzuki, and\n M. Tonegawa",
"submitter": "Fumihide Iwamuro",
"url": "https://arxiv.org/abs/1111.6746"
}
|
1111.6805
|
# UCAC3 Proper Motion Survey. II.
DISCOVERY OF NEW PROPER MOTION STARS IN UCAC3
WITH 0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1 BETWEEN DECLINATIONS
$-$47$\arcdeg$ and 00$\arcdeg$
Charlie T. Finch, Norbert Zacharias finch@usno.navy.mil U.S. Naval
Observatory, Washington DC 20392–5420 Mark R. Boyd, Todd J. Henry Georgia
State University, Atlanta, GA 30302–4106 Nigel C. Hambly Scottish
Universities Physics Alliance, Institute for Astronomy, University of
Edinburgh, Royal Observatory, Blackford Hill, Edinburgh EH9 3HJ, Scotland, UK
###### Abstract
We present 474 new proper motion stellar systems in the southern sky having no
previously known components, with 0$\farcs$40 yr-1 $>$ $\mu$ $\geq$
0$\farcs$18 yr-1 between declinations $-$47$\arcdeg$ and 00$\arcdeg$. In this
second paper utilizing the U.S. Naval Observatory third CCD Astrograph Catalog
(UCAC3) we complete our sweep of the southern sky for objects in the proper
motion range targeted by this survey with R magnitudes ranging from 9.80 to
19.61. The new systems contribute a $\sim$16% increase in the number of new
stellar systems for the same region of sky reported in previous SuperCOSMOS
RECONS (SCR) surveys. Among the newly discovered stellar systems are 16
multiples, plus an additional 10 components that are new common proper motion
companions to previously known objects. A comparison of UCAC3 proper motions
to those from Hipparcos, Tycho-2, Southern Proper Motion (SPM4), and
SuperCOSMOS indicates that all proper motions are consistent to $\sim$10
mas/yr, with the exception of SuperCOSMOS. Distance estimates are derived for
all stellar systems having SuperCOSMOS Sky Survey (SSS) $B_{J}$, $R_{59F}$,
and $I_{IVN}$ plate magnitudes and Two-Micron All Sky Survey (2MASS) infrared
photometry. We find five new red dwarf systems estimated to be within 25 pc.
These discoveries support results from previous proper motion surveys
suggesting that more nearby stellar systems are to be found, particularly in
the fainter, slower moving samples.
In this second paper utilizing the U.S. Naval Observatory third CCD Astrograph
Catalog (UCAC3) we complete our sweep of the southern sky for objects in the
proper motion range targeted by this survey with R magnitudes ranging from
9.80 to 19.61.
solar neighborhood — stars: distances — stars: statistics — surveys —
astrometry
## 1 INTRODUCTION
The third U.S. Naval Observatory (USNO) CCD Astrograph Catalog (UCAC3)
(Zacharias et al., 2010) proper motion survey, addresses the possibility that
proper motion surveys using digitized scans of photographic plates might
overlook some proper motion systems. The UCAC3 obtained accurate proper
motions by combining CCD observations with early epoch photographic data. This
survey utilizes the UCAC3 proper motions to discover new systems that have
been missed in previous efforts. The first paper in this series (Finch et al.
2010b, ) (hereafter, U3PM1), confirmed this suspicion by revealing an
additional 25.3% stellar systems having a proper motion of 0$\farcs$40 yr-1
$>$ $\mu$ $\geq$ 0$\farcs$18 yr-1 between declinations $-$90$\arcdeg$ and
$-$47$\arcdeg$ over those found by the Research Consortium On Nearby Stars
(RECONS)111www.recons.org group using SuperCOSMOS Sky Survey (SSS) data. These
new discoveries provided the impetus for this second paper of the series,
which completes the sweep of the southern sky for systems with 0$\farcs$40
yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1 found in the UCAC3.
The data obtained from proper motion surveys aid astronomers in determining
accurate stellar luminosity and mass functions, thereby revealing how the
Galaxy’s stellar mass is divided among different types of stars. Our main goal
— identifying the Sun’s nearest neighbors — provides a vast sample of red
dwarf, subdwarf, and white dwarf stellar systems for studies of multiplicity,
activity, ages, and exoplanet searches. Because of their proximity, the nearby
stars offer the most accessible measurements of each of these characteristics.
Our UCAC3 proper motion survey is currently focused on the southern
hemisphere, which has not been surveyed as systematically as the northern sky,
where the pioneering surveys of Giclas (Giclas et al., 1971, 1978) and Luyten
(Luyten, 1979, 1980) were primarily carried out. Historically, proper motion
studies have been focused on blinking photographic plates taken at different
epochs to determine source motions. Recent surveys that complement the classic
efforts utilize various techniques, plate sets, modern computers, and
carefully tailored algorithms to effectively blink digitized images of
photographic plates. In the southern sky, such surveys include (Wroblewski &
Torres, 1994), (Wroblewski & Costa, 1999), (Scholz et al., 2000, 2002),
(Oppenheimer et al., 2001), (Pokorny et al., 2003), (Lépine, 2005, 2008),
(Deacon et al., 2005; Deacon & Hambly, 2007), and (Deacon et al., 2009).
In an effort to understand the stellar population of the solar neighborhood,
the RECONS group has been targeting the neglected southern sky to reveal new
stellar proper motion systems. To date, these discoveries have been reported
in six papers in The Solar Neighborhood (TSN) series (Hambly et al., 2004),
(Henry et al., 2004), (Subasavage et al. 2005a, ), (Subasavage et al. 2005b,
), (Finch et al., 2007), (Boyd et al., 2011). These new systems are discovered
using the SuperCOSMOS Sky Survey (SSS) data (Hambly et al. 2001a, ) and given
the name SCR (SuperCOSMOS-RECONS). Followup observations of intriguing systems
are performed at the Cerro Tololo Inter-American Observatory (CTIO) 0.9m
telescope, where RECONS operates a trigonometric parallax program focusing on
stars within 25 pc.
Our UCAC3 survey uses an approach fundamentally different from plate blinking
to reveal proper motion systems. We take advantage of observations reported in
many catalogs ranging in epochs from the early nineteenth to the early twenty-
first centuries, rather than directly using specific sets of digitized images
from photographic plates. In this investigation we focus on stellar systems in
the UCAC3 found between declinations $-$47$\arcdeg$ and 00$\arcdeg$ that have
0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1, completing a sweep of the
southern sky. The search region and proper motion range matches that in (Boyd
et al., 2011), hereafter TSN25, in which the lower proper motion cutoff was
chosen to match that of the NLTT catalog. TSN25 reports 2817 new SCR systems,
substantially adding to the number of new SCR systems found previously. In
Table 1, we summarize the number of new stellar systems discovered,
highlighting those estimated to be within 25 pc, for both the RECONS and UCAC3
surveys. In this paper we will focus in particular on the two SCR searches
(TSN18 and TSN25) that correspond to the same proper motion and declination
ranges as this UCAC3 survey (U3PM1 and this paper). New stellar objects from
this search are given USNO Proper Motion (UPM) names.
## 2 Method
### 2.1 UCAC3
The USNO CCD Astrograph Catalog (UCAC) project finished observations in late
2004 and has been producing astrometric catalogs since October 2000\. This
astrometric survey was conceived to densify the optical reference frame to
high accuracy beyond the Hipparcos and Tycho magnitudes. UCAC is the first
all-sky survey performed with a CCD detector utilizing the high level of
precision achievable with this technology. The first release, UCAC1,
(Zacharias et al., 2000), was a partial catalog covering 80% of the southern
sky. The second catalog, UCAC2, (Zacharias et al., 2004), contains roughly 80%
of the entire sky and includes improved proper motions from the use of early
epoch plates paired with the Astrograph CCD data. UCAC3 (Zacharias et al.,
2010), released in August 2009, is the first in the series to contain coverage
of the entire sky. UCAC3 also includes double star fitting and has a slightly
deeper limiting magnitude than UCAC2 due to a complete re-reduction of the
pixel data (Zacharias, 2010). In addition, data from the Two-Micron All Sky
Survey (2MASS) were used in UCAC3 to probe for and reduce systematic errors in
UCAC observations, providing a greater number of reference stars to stack up
residuals as a function of many parameters, such as observing site and
exposure time. A detailed description of the astrometric reductions of UCAC3
can be found in (Finch et al. 2010a, ). A detailed introduction to the UCAC3
can be found in the release paper (Zacharias et al., 2010) and the README file
of the data distribution. A new edition, UCAC4, (Zacharias et al., 2011) is
scheduled to be released later this year.
### 2.2 PROPER MOTIONS
The UCAC3 contains roughly 95 million calculated absolute proper motions. The
majority of these are derived proper motions from the use of early epoch
catalogs paired with the Astrograph CCD data. Earlier epoch data are all
reduced to the International Celestial Reference Frame (ICRF). UCAC3 mean
positions and proper motions are calculated using a weighted, least-squares
adjustment procedure.
Bright stars with R$\sim$8–12 in UCAC3 are combined with ground-based
photographic and transit circle catalogs. These include all catalogs used for
the production of the Tycho-2 project (Høg et al., 2000), unpublished measures
of over 5000 astrograph plates digitized on the StarScan machine (Zacharias et
al., 2008), new reductions of Southern Proper Motion (SPM) (Girard et al.,
2011) data, and data from the SuperCOSMOS project (Hambly et al. 2001a, ).
About 1.2 million star positions to about B$=$ 12 entered UCAC from digitizing
the AGK2 plates (epoch about 1930). The Hamburg Zone Astrograph and USNO Black
Birch Astrographs contributed another 7.3 million star positions, mainly in
the V$=$12 – 14 magnitude range, in fields covering about 30% of the sky, and
the Lick Astrograph plates taken around 1990 yielded over one million star
positions to V$=$16 in selected fields.
For all catalogs used to derive UCAC3 proper motions a systematic error
estimate was added to the root mean square (RMS) of the individual stars
random errors. The largest error floor added was 100 mas for the SuperCOSMOS
data due to zonal systematic errors ranging from 50–200 mas when compared to
2MASS data.
To identify previously known high proper motion (HPM) stars in the UCAC3, a
source list was compiled using the VizieR on-line data tool, along with
targeted supplements from published literature. In the north we used the LSPM-
North catalog (Lépine, 2005) containing 61977 new and previously found stars
having proper motions greater than 0$\farcs$15 yr-1. For the south we utilized
many surveys, notably including the Revised NLTT Catalog (Salim & Gould,
2003), which produced 17730 stars with proper motions greater than 0$\farcs$15
yr-1, and the RECONS efforts (SCR stars). For a full list of catalogs used,
see the UCAC3 README file. While this list is not comprehensive, this effort
tagged roughly 51000 known HPM stars in UCAC3 over the entire sky. These
previously identified HPM stars were given a mean position (MPOS) number
greater than 140 million and do not have derived UCAC3 proper motions. We
instead used the proper motion data from the catalogs themselves (see
$\S$4.5).
Proper motion errors in the UCAC3 catalog for stars brighter than R$\sim$12
are only $\sim$1–3 mas/yr in part because of the large epoch spread of roughly
100 years in some cases. The errors of the fainter stars range from $\sim$2–3
mas/yr if found in SPM4 and $\sim$6–8 mas/yr if SuperCOSMOS data are used in
lieu of SPM4 data.
### 2.3 SEARCH CRITERIA
In this second paper we survey the southern sky between declinations
$-$47$\arcdeg$ and 0$\arcdeg$ using the same proper motion range as in U3PM1,
0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1. In this area of the sky we
identify an initial sample of 212356 proper motion candidates. We utilize the
same search criteria as in U3PM1, using UCAC3 flags with values indicative of
real proper motion objects. A visual check from a sample of stars confirmed
that these flags still hold true in the region of the sky being surveyed. All
stars must (1) be in the 2MASS catalog with an e2mpho (2MASS photometry error)
less than or equal to 0.05 magnitudes in all three bands, (2) have a UCAC fit
model magnitude between 7 and 17 mag, (3) have a double star flag (dsf) equal
to 0, 1, 5 or 6, meaning a single star or fitted double, (4) have an object
flag (objt) between $-$2 and 2 to exclude positions that used only overexposed
images in the fit, (5) have an MPOS number less than 140 million, to exclude
already known high proper motion stars, and (6) have a LEDA galaxy flag of
zero, meaning that the source is not in the LEDA galaxy catalog. After all
these cuts, there remain 17516 “good” candidate list, fewer than expected for
this region of the sky, when compared to 9248 in U3PM1. A total of 7641
candidates were excluded from the “good” candidates due to being marked as
previously known in the UCAC3 catalog (MPOS number greater than 140 million).
These candidates were then cross-checked via VizieR and SIMBAD to determine if
they were previously known. All cross-checks are performed using a 90$\arcsec$
search radius, with one exception (the NLTT catalog). A larger search radius
of 180$\arcsec$ was used when comparing UPM candidates to the NLTT and LHS
catalog, which have been found to have inaccurate positions as reported in
(Bakos et al., 2002). Thus, UCAC3 proper motion candidates with positions
differing from Luyten’s or any other known object by less than 90$\arcsec$ are
considered known. Those differing from Luyten’s by 90–180$\arcsec$ are
considered new discoveries but are noted as possible NLTT stars in the tables.
Those differing by more than 180$\arcsec$ from Luyten are considered new
discoveries. All candidates matched to known stars had a final check to
determine if the proper motion and magnitudes matched — those that match are
considered known and not reported in this sample. As in U3PM1, it is not a
goal of this paper to revise the NLTT catalog and assign proper
identifications and accurate positions to NLTT entries; rather, the goal is to
identify new high proper motion stars.
After this, in effect, second cross-check for previously known stars, the list
was reduced to a manageable 3736 candidate proper motion objects. The 13780
known objects found during this cross-check shows how incomplete the UCAC3
catalog can be in identifying previously known high proper motion objects with
the given search criteria. Each of these candidates was then visually
inspected to confirm proper motion by blinking the $B_{J}$ and $R_{59F}$
SuperCOSMOS digitized plate images. During blinking, we noticed that for
declinations between roughly $-$33$\arcdeg$ and 0$\arcdeg$ the epoch spread
was insufficient ($\sim$3–5 years) to visually verify proper motion for all
candidates. For those candidates, a second sweep was done by blinking the
$POSS-IR$ and $R_{59F}$ SuperCOSMOS digitized plate images. Nearly 87% of the
candidates were found to have no verifiable proper motions and were discarded.
The final counts of new discoveries are 500 proper motion objects in 474
systems. Among these are 25 multiple systems (24 doubles and one triple), of
which ten were found to have CPM to previously known primaries.
For this search we find a successful hit rate — defined as the number of new
and known proper motion stars (21921) divided by the total “good” candidates
extracted (25157, including stars with an MPOS number $>$ 140 million) — of
87.1%, which is higher than the 81.4% hit rate found in TSN25. After looking
into the calculation used in U3PM1 to determine the successful hit rate a
counting error was found. The number of real objects excluded the known proper
motion objects tagged in the UCAC3 catalog (stars with an MPOS number $>$ 140
million). If we add these stars in the total for the U3PM1 count, we get a
total of 7975 real objects giving a new successful hit rate of 86.2%, which is
comparable to this paper. At least three factors mentioned in U3PM1 have been
identified that can lead to false detections in the UCAC3 proper motion
survey. First, some real objects are discarded during the sifting mentioned
above, particularly because of the 2MASS criterion which states that $JHK_{s}$
photometry errors must be less than 0.05 mag. Second, the UCAC3 contains many
phantom proper motion objects due to incorrect matches during proper motion
calculations. Third, other misidentifications arise from blended images, where
a single source in an earlier epoch catalog can be matched with two stars in
the UCAC3 data.
## 3 RESULTS
In Table 2, we list the 474 new proper motion stellar systems (500 objects)
discovered during this search. We highlight the five red dwarf systems
estimated to be within 25 pc in Table 3. In both tables we list names,
coordinates, proper motions, 1$\sigma$ errors in the proper motions, plate
magnitudes from SuperCOSMOS, near-IR photometry from 2MASS, the computed
$R_{59F}-J$ color, a distance estimate, and notes.
### 3.1 Positions and Proper Motions
All positions on the ICRF system, proper motions, and errors are taken
directly from UCAC3, unless otherwise noted. For a few stars that were found
during visual inspection without any UCAC3 data, information has been obtained
from alternate sources (see $\S$3.4). For this sample, the average positional
errors reported in the UCAC3 catalog are 51 mas in RA and 50 mas in Dec. For
proper motions, the average errors reported in the UCAC3 for this sample are
8.0 mas/yr in $\mu_{\alpha}\cos\delta$ and 7.7 mas/yr in $\mu_{\delta}$.
### 3.2 Photometry
In Tables 2 and 3, we give photographic magnitudes from the SuperCOSMOS and
2MASS surveys. From SuperCOSMOS, magnitudes are given from three plate
emulsions, $B_{J}$, $R_{59F}$, and $I_{IVN}$. Magnitude errors are typically
less than 0.3 mag for stars fainter than $\sim$15, with errors increasing for
brighter sources. From 2MASS, $JHK_{s}$ infrared photometry is given, with
errors typically 0.05 mag or less due to the search criteria. Additional
objects found during visual inspection are typically fainter with larger
photometric errors. The $R_{59F}-J$ color has been computed to indicate the
star’s color.
While SuperCOSMOS magnitudes are reported in the UCAC3, this sample was
checked against the SuperCOSMOS catalog to rectify some mismatches found in
the UCAC3 catalog. In some cases, SuperCOSMOS magnitudes are not given in the
tables, due to blending, no source detection, high chi-square or other
problems where no reliable magnitude is available. 2MASS magnitudes are given
for all but one object which was found visually that is not present in the
2MASS catalog, as indicated in the notes.
### 3.3 Distances
Plate photometric distance estimates are computed using the same method as in
U3PM1 and previous SCR searches. Using the relations in (Hambly et al., 2004),
11 distance estimates are generated based on colors computed from the six-band
photometry. This method assumes all objects are main sequence stars, and
provides distances accurate to 26%, determined from the mean differences
between the true distances for stars with accurate (errors less than 10 mas)
trigonometric parallaxes and distances estimated from the relations. Errors
are higher for stars with missing photometry, resulting in fewer than 11
relations, and stars that are not single, main sequence red dwarfs, e.g. cool
subdwarfs and white dwarfs. It is possible to produce a distance with only one
relation; however, six are needed to be considered “reliable” because that
allows for one magnitude dropout. Stars having fewer than six relations are
identified in the notes to Tables 2 and 3. If a star is identified as a
possible subdwarf, the distance estimate is expected to be too large and is
given in brackets.
### 3.4 Additional Objects
In Table 2 we include 17 additional proper motion objects found during visual
inspection of the candidate fields. These objects are CPM companion candidates
that either have fainter limiting magnitudes than implemented for this search,
were eliminated from the candidate list by the search criteria, or have UCAC3
proper motions less than 0$\farcs$18 yr-1. These new visual discoveries have
all been cross-checked with VizieR and SIMBAD using the same methods described
above for the main search. Proper motions have been obtained from UCAC3, SPM4,
PPMXL (Roeser et al., 2010), or SuperCOSMOS, in that order. For stars that
were not found in the UCAC3 data, positions were computed using the epoch,
coordinates, and proper motion obtained from the corresponding catalog.
Magnitudes are obtained using the 2MASS and SuperCOSMOS catalogs to compute
distance estimates.
## 4 ANALYSIS
### 4.1 Color-Magnitude Diagram
In Figure 1 we show a color-magnitude diagram of the 334 new UPM proper motion
objects (solid circles) and seven known objects (open triangles, companions to
UPM objects) from this search having $R_{59F}-J$ colors. Symbols that fall
below $R_{59F}\sim 17$ are CPM companion candidates noticed during visual
inspection. The brightest new object, UPM 0747-2537A, has $R_{59F}$ = 9.80 and
is estimated to be at a distance of 40.6 pc. The reddest object found in this
search is UPM 1848-0252 with $R_{59F}-J$ = 5.06, $R_{59F}$ = 16.57, at an
estimated distance of 26.9 pc.
The subdwarf population is not as well defined as in TSN18 and TSN25 because
there are far fewer new objects. Nonetheless, a separation can be seen below
the concentration of main sequence stars.
### 4.2 Reduced Proper Motion Diagram
In Figure 2, we show the reduced proper motion (RPM) diagram for all objects
also plotted in Figure 1, with similar symbols for new and known objects. The
RPM diagram is a good method to help separate white dwarfs and subdwarfs from
main-sequence stars, under the assumption that objects with larger distances
tend to have smaller proper motions. Using the same method as in U3PM1 and
TSN25 we obtain $H_{R_{59F}}$ via a modified distance modulus equation, in
which $\mu$ is substituted for distance:
$H_{R_{59F}}=R_{59F}+5+5\log\mu.$
The solid line seen in Figure 2 is used to separate white dwarfs from
subdwarfs. This is the same empirical line used in U3PM1 and previous TSN
papers. No white dwarf candidates have been found during this latest search.
Subdwarf candidates have been selected using the same method as in U3PM1 and
TSN25 — stars with $R_{59F}-J>$ 1.0 and within 4.0 mag in $H_{R}$ of the
empirical line separating the white dwarfs are considered subdwarfs. From this
survey there are 17 subdwarf candidates, all with distance estimates greater
than 122 pc, with the exception of one, UPM 1712-4432, with an estimated
distance of 33.9 (see $\S$4.4). Because the relations used to estimate
distances assume that stars are on the main sequence, underluminous cool
subdwarfs and white dwarfs have large distances, which can, in fact, be used
to identify such objects. The distance estimates for these stars are
presumably erroneous and are given in brackets in Tables 2, 3 and 4. Follow-up
spectroscopic observations will be needed to confirm all subdwarf candidates.
### 4.3 New Common Proper Motion Systems
In this search, we find 25 CPM candidate systems consisting of 24 binaries and
one triple. Included in these CPM systems are 16 new systems and nine known
systems with newly discovered components.
One binary system, UPM 0800-0617AB is a possible subdwarf binary system. The
lone triple is an SCR system with two newly discovered components. In Table 4,
we list the CPM system primaries and companions, their proper motions, and the
companions’ separations and position angles relative to the primaries (defined
to be the brightest star in each system using the UCAC bandpass, or an
alternate bandpass if a UCAC value is not available). We also provide distance
estimates for each component, where possible. Components were determined to be
potentially physically associated using distance estimates in conjunction with
the proper motions and visual inspections. However, most of the companions
were found during visual inspection, meaning that proper motions, 2MASS and/or
SuperCOSMOS magnitudes may be missing or suspect, as identified in the notes.
For systems with data missing in Table 4, the physical connection of the
system components should be considered tentative.
In Figure 3, we show comparisons of the proper motions in each coordinate for
the 19 CPM systems for which both components have a listed proper motion. CPM
candidates having proper motions from the UCAC3 are represented by solid
circles while those with proper motions from other sources are represented by
open circles. If a proper motion was not present in the UCAC3, data were
obtained manually from the SPM4, PPMXL or SuperCOSMOS databases, in that
order.
### 4.4 Notes on Specific Stars
UPM 0443-4129AB is a possible CPM binary. However, UPM 0443-4129A has a
suspect proper motion and the companion’s distance estimate uses fewer than 6
relations. It is possible that this pair is a case of a chance alignment. See
Table 4 for more details.
BD-04 2807AB is a possible CPM binary. However, the primary has a suspect
proper motion, a distance estimate that uses fewer than 6 relations, and there
is no distance estimate for the secondary. It is possible that this pair is a
case of a chance alignment. See Table 4 for more details.
UPM 0747-2537A is the brightest new discovery from this search with $R_{59F}=$
9.80 and an estimated distance of 40.6 pc. However, only one relation was
viable, making the distance estimate unreliable.
UPM 0800-0617AB is a possible candidate for a binary subdwarf system. The
primary is a possible subdwarf at an estimated distance of 175.5 pc. The
secondary is at a separation of 5.8$\arcsec$ at position angle 297.2∘ from the
primary. Color information is insufficient for a reliable distance estimate.
UPM 1226-3516B and C are in a candidate triple system with SCR 1226-3515A. The
A and B components are separated by 49.8$\arcsec$ at a position angle of
191.3∘. The C component has a separation of 97.0$\arcsec$ at a position angle
of 146.9∘ from the primary. The C component has a suspect proper motion and
the distance estimates for all there components are inconsistent. In
particular, the C component may not be a part of the system. See Table 4 for
more details.
UPM 1712-4432 is a subdwarf candidate with $R_{59F}=$ 13.04 and $R_{59F}-J=$
1.01 at a distance of 33.9 pc. However, only three relations were viable,
making the distance estimate unreliable. SuperCOSMOS magnitudes are indicative
of a blended image, meaning this is likely not what it seems.
UPM 1718-2245B has an estimated distance of only 13.2 pc based on 7 relations,
making it the nearest candidate in the sample. However, the primary has a
distance estimate of 25.4 pc based on 10 relations so we favor the larger
distance for the system.
UPM 1848-0252 is the reddest new discovery from this search, with $R_{59F}-J=$
5.06 and an estimated distance of 26.9 pc.
### 4.5 Comparison to Previous Proper-Motion Surveys
During production of the UCAC3 catalog, we made an effort to tag previously
known HPM stars. For these stars, proper motions were taken from their
respective catalogs rather than calculated using UCAC3 methodology, which made
comparisons to other catalogs/surveys difficult. However, during the present
search we have found 104 stars in both the Hipparcos and Tycho-2 catalogs that
are not tagged as HPM stars in the UCAC3 catalog — these stars are proper
motion candidates that were found to be in Tycho-2 during cross-checking. A
2.5$\arcsec$ radius was used to match these stars to sources in the Hipparcos
catalog so that we can compare the bright end of the UCAC3 proper motion stars
($R$ $\sim$ 7.13-13.66) to stars in both the Tycho-2 and Hipparcos catalogs.
In Figure 4, we compare proper motions in RA and Dec for these stars as given
in UCAC3, Hipparcos, and Tycho-2. These plots show that the differences in
proper motions are small, in general less than 10 mas/yr, and no significant
systematic errors as a function of declination are seen. The RMS differences
between UCAC3 proper motions in $\Delta\mu_{\alpha}\cos\delta$ and
$\Delta\mu_{\delta}$ and those from Hipparcos are 5.7 and 9.1 mas/yr,
respectively. Comparisons to Tycho-2 yield RMS differences of 5.2 and 8.3
mas/yr, respectively. Lower RMS differences of 3.0 mas/yr in
$\Delta\mu_{\alpha}\cos\delta$ and 3.2 mas/yr in $\Delta\mu_{\delta}$ are seen
when comparing the Hipparcos to Tycho-2 proper motions.
To investigate the fainter end of UCAC3, we compare results for 77 stars ($R$
$\sim$ 10.88-16.69) that are in both the SPM4 and SuperCOSMOS catalogs that
were not tagged as HPM stars in the UCAC3 catalog — these stars are proper
motion candidates that were found to be SCR stars during cross-checking. A
2.5$\arcsec$ radius was used to match these stars to sources in the SPM4
catalog. The SPM4 catalog only covers Dec = $-$90 to $-$20 sky area, limiting
the area included for this comparison. In Figure 5, we compare proper motions
in RA and Dec for these stars as given in UCAC3, SuperCOSMOS, and SPM4. These
plots show that differences in proper motions are similar to those found for
brighter stars when comparing UCAC3 and SPM4, but the differences are much
larger for the SuperCOSMOS results. The RMS differences between UCAC3 proper
motions in $\Delta\mu_{\alpha}\cos\delta$ and $\Delta\mu_{\delta}$ and those
in SPM4 are 6.0 and 5.7 mas/yr respectively. Comparisons to SuperCOSMOS yield
RMS differences of 16.5 and 14.1 mas/yr, respectively. In Figure 5, we also
see that proper motions in Dec appear to be systematically shifted in the
SuperCOSMOS data. These high RMS results and the systematic shift are also
seen in the comparison of the SPM4 to the SuperCOSMOS proper motions, yielding
RMS differences of 15.6 and 15.2 mas/yr in $\Delta\mu_{\alpha}\cos\delta$ and
$\Delta\mu_{\delta}$, respectively. The higher RMS differences for the
SuperCOSMOS proper motions are in agreement with the findings of TSN18 and
U3PM1 where SCR proper motions were found to have higher RMS differences when
compared to other external catalogs. It is worth noting that the SuperCOSMOS
proper motion RMS reported here are not representative of the entire catalog.
Objects having an R$\sim$16–19 with $\mu>$ 0$\farcs$10 yr-1 in the SuperCOSMOS
catalog should have an RMS no greater than 10 mas/yr, and considerably better
for fields with decades between the epochs (See Tables 1 and 3 from (Hambly et
al. 2001b, )).
Random and systematic differences of order 10 mas/yr in proper motions between
the various catalogs, particularly at the faint end, are expected because of
different data quality, measurements, reductions and epoch differences.
SuperCOSMOS for example uses Schmidt plates for both early and recent epoch
which typically show large errors. The proper motions of faint stars in UCAC3
are based on early epoch Schmidt plates for the sky area north of $-$20 deg
Dec and CCD observations for recent epoch data. A combination of CCD data and
early astrograph data (SPM plates) is used south of $-$20 deg, with
significantly smaller errors. The SPM4 proper motions are derived entirely on
SPM astrograph plates from 2 epochs. At the bright end proper motions are more
reliable due to higher quality of Hipparcos and Tycho data as well as
availability of many other star catalogs, most of which have been used in
common between Tycho-2 and UCAC3. However, there can be large differences
between Hipparcos and Tycho-2 for some stars because the Hipparcos PMs are
based on only about 3.5 years of observing (although with high quality), while
Tycho-2 PMs are based on typically 100 years epoch difference. Multiplicity
and residual orbital motions sometimes render Hipparcos PMs inferior in spite
of their small formal astrometric errors.
In TSN25 a total of 3073 objects were reported, all of which fit within the
proper motion and declination constraints of this paper. During this UCAC3
search, only 770 of the 3073 objects reported in TSN25 were recovered, or a
low 25.1% recovery rate. This is primarily due to the UCAC3 catalog having no
proper motion or a reported proper motion not meeting the criteria of this
paper (0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1) for $\sim$70% of
the new discoveries listed in TSN25.
The Hipparcos catalog contains 118218 total objects, of which 1690 meet the
proper motion and declination constraints of this paper. Tycho-2 contains
2539913 total objects in the main catalog, of which 3187 meet similar limits.
We recover 1316 Hipparcos stars and 2543 Tycho-2 stars using the search
criteria of this paper, yielding recovery rates of 77.9% and 79.8%
respectively. Objects missed in this UCAC3 survey are primarily due to UCAC3
lacking a source detection for $\sim$15% of the Tycho-2 objects. The
relatively high recovery rates of UCAC3, when compared to the Hipparcos and
Tycho-2 catalogs, implies the UCAC3 can be used as a reliable source to search
for new proper motion stars with $\mu$ = 0.18–0.40$\arcsec$ yr-1 for other
portions of the sky.
## 5 DISCUSSION
We have completed a sweep of the southern sky for new proper motion systems
using the UCAC3 catalog. So far, we have uncovered 916 new proper motion
systems, of which 474 are described in this paper. These systems constitute an
increase of 19.4% over the total number of SCR systems discovered in the
southern sky and an increase of 20.7% over SCR systems in the southern sky
with 0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1. This UCAC3 proper
motion survey has added 3.8% to the list of entries in the NLTT catalog south
of Dec $=$ 0$\arcdeg$ with 974 new proper motion objects from U3PM1 and this
paper.
In Figure 6, we show the sky distribution of systems found to date during the
UCAC3 proper motion survey. Plus signs represent objects from U3PM1 and solid
circles represent objects described in this paper. Overall, the distribution
of new objects is similar to that seen in Figure 6 of TSN25, including the
discovery of many new proper motion systems along the Galactic plane.
In Figure 7, we show a histogram of the number of proper motion systems
discovered to date during the UCAC3 proper motion survey, in 0$\farcs$01 yr-1
bins, and highlighting the number of those having distance estimates within 50
pc. Predictably, this plot shows that the slowest proper motion bins have the
most new systems. This confirms the trend reported in TSN18, TSN25 and U3PM1,
and suggests once again that more nearby stars are yet to be found at slower
proper motions.
We have found a total of 57 CPM candidate systems during this UCAC3 proper
motion survey, including 55 binaries and two triples. These systems have
separations of 1–359$\arcsec$ and will need further investigation to confirm
which of the systems are, in fact, gravitationally linked. In addition, we
have revealed a total of 48 subdwarf candidates, each of which is worthy of
followup observations, given the scarcity of nearby subdwarfs. Finally, we
have found 20 red dwarf systems likely to be within 25 pc. We plan to obtain
CCD photometry through $VRI$ filters for stars having estimated distances
within 25 pc in order to make more reliable distance estimates using the
$VRIJHK$ relations presented in (Henry et al., 2004). Stars estimated to be
within 10 pc will then be put on the CTIO parallax program, potentially to
join the ranks of the few hundred systems known to be so close to the Sun
(Henry et al., 2006).
We thank the entire UCAC team for making this proper motion survey possible,
and the USNO summer students, who helped with tagging HPM stars in the UCAC3
catalog. Special thanks go to members of the RECONS team at Georgia State
University for their support, and John Subasavage in particular for assistance
with the SCR searches. This work has made use of the SIMBAD, VizieR, and
Aladin databases operated at the CDS in Strasbourg, France. We have also made
use of data from the Two-Micron All Sky Survey, SuperCOSMOS Science Archive
and the Southern Proper Motion catalog.
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Figure 1: Color-apparent magnitude diagram for all proper motion systems in
the sample having an $R_{59F}-J$ color. New proper motion objects are
represented by solid circles while known objects (CPM companions to new
objects) are represented with open triangles. Data below $R_{59F}=$ 17 are CPM
candidates noticed during visual inspection. Figure 2: RPM diagram for all
proper motion systems in this sample having an $R_{59F}-J$ color. New proper
motion objects are represented by solid circles while known objects (CPM
companions to new objects) are represented with open triangles. The empirical
line separates the subdwarfs from where white dwarf candidates would be found.
No white dwarf candidates were found in the current search. Figure 3:
Comparisons of proper motions in each coordinate, $\mu_{\alpha}\cos\delta$
(top) and $\mu_{\delta}$ (bottom), for components in CPM systems. Proper
motions from the UCAC3 catalog are represented by solid circles while proper
motions manually obtained through other means are denoted by open circles. The
solid line indicates perfect agreement. Information on the outliers can be
found in $\S$4.4
Figure 4: Comparisons of UCAC3, Hipparcos and Tycho-2 proper motions per
coordinate, $\Delta\mu_{\alpha}\cos\delta$ (left column) and
$\Delta\mu_{\delta}$ (right column). Figure 5: Comparisons of UCAC3,
SuperCOSMOS and SPM4 proper motions per coordinate,
$\Delta\mu_{\alpha}\cos\delta$ (left column) and $\Delta\mu_{\delta}$ (right
column).
Figure 6: Sky distribution of all UCAC3 proper motion survey objects reported
in U3PM1 (plus signs) and this paper (solid circles), i.e. those between
declinations $-$90$\arcdeg$ and 0$\arcdeg$ having 0$\farcs$40 yr-1 $>$ $\mu$
$\geq$ 0$\farcs$18 yr-1. The curve represents the Galactic plane. Figure 7:
Histogram showing the number of proper motion objects in 0$\farcs$01 yr-1 bins
for the entire UCAC3 proper motion sample (empty bars) and the number of those
objects having distance estimates within 50 pc (filled bars).
Table 1: New Proper Motion Systems from the UCAC3 and SCR proper motion surveys Paper | New Systems | New Systems | References
---|---|---|---
| total | $\leq$ 25 pc |
U3PM1 | 442 | 15 | (Finch et al. 2010b, )
U3PM2 | 474 | 4 | this paper
TSN08 | 5 | 2 | (Hambly et al., 2004)
TSN10 | 4 | 4 | (Henry et al., 2004)
TSN12 | 141 | 12 | (Subasavage et al. 2005a, )
TSN15 | 152 | 25 | (Subasavage et al. 2005b, )
TSN18 | 1605 | 30 | (Finch et al., 2007)
TSN25 | 2817 | 79 | (Boyd et al., 2011)
totals | 5640 | 171 |
Table 2: New UCAC3 High Proper Motion Systems between Declinations $-$47$\arcdeg$ and 0$\arcdeg$ with 0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1 Name | RA J2000.0 | DEC J2000.0 | $\mu_{\alpha}\cos\delta$ | $\mu_{\delta}$ | sig$\mu_{\alpha}$ | sig$\mu_{\delta}$ | $B_{J}$ | $R_{59F}$ | $I_{IVN}$ | $J$ | $H$ | $K_{s}$ | $R_{59F}$ $-$ $J$ | Est Dist | Notes
---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---
| (deg) | (deg) | (mas/yr) | (mas/yr) | (mas/yr) | (mas/yr) | | | | | | | | (pc) |
UPM 0004-0833 | 1.1472383 | -8.5664958 | 182.2 | -26.3 | 8.6 | 8.6 | 14.680 | 12.649 | 11.643 | 10.874 | 10.255 | 10.072 | 1.775 | 63.8 |
UPM 0004-1258 | 1.2061531 | -12.9791222 | 102.0 | -156.8 | 10.2 | 6.6 | 14.504 | 12.478 | 11.385 | 10.774 | 10.160 | 9.973 | 1.704 | 62.3 |
UPM 0009-1539 | 2.4914975 | -15.6590597 | 184.0 | -15.7 | 8.7 | 8.7 | 17.100 | 15.050 | 13.144 | 12.387 | 11.793 | 11.562 | 2.663 | 93.0 |
UPM 0011-1448 | 2.9555050 | -14.8079581 | 179.8 | -34.9 | 9.0 | 9.0 | 14.785 | 13.131 | 12.454 | 12.005 | 11.353 | 11.246 | 1.126 | 115.6 |
UPM 0014-0029 | 3.6589169 | -0.4939803 | 161.0 | -104.7 | 13.7 | 14.0 | 17.819 | 15.829 | 13.947 | 12.292 | 11.708 | 11.457 | 3.537 | 58.7 |
UPM 0014-1219 | 3.7410256 | -12.3317842 | 183.9 | -21.6 | 7.0 | 3.9 | 17.189 | 15.068 | 13.388 | 12.808 | 12.230 | 11.974 | 2.260 | 130.3 |
UPM 0025-2547 | 6.3606617 | -25.7849942 | 170.4 | 65.2 | 9.4 | 9.3 | 20.968 | 18.826 | 16.962 | 15.167 | 14.514 | 14.163 | 3.659 | 179.8 | aaProper motions suspect
UPM 0044-1647 | 11.1954375 | -16.7984897 | 187.8 | 29.7 | 14.6 | 14.1 | $\cdots$ | $\cdots$ | 13.280 | 12.405 | 11.777 | 11.549 | $\cdots$ | 91.5 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0045-3602 | 11.2618550 | -36.0381975 | 118.0 | -136.9 | 3.2 | 2.4 | 15.862 | 13.863 | 11.982 | 11.398 | 10.888 | 10.610 | 2.465 | 67.0 |
UPM 0048-0217 | 12.1068206 | -2.2840133 | 163.9 | -75.0 | 8.9 | 8.3 | 16.715 | 14.628 | 12.563 | 11.042 | 10.488 | 10.194 | 3.586 | 31.8 |
UPM 0058-0158 | 14.6628417 | -1.9776981 | -48.2 | -181.1 | 11.2 | 11.2 | 18.275 | 16.265 | 14.745 | 13.405 | 12.870 | 12.662 | 2.860 | 141.3 |
UPM 0106-1342 | 16.7276508 | -13.7017203 | 173.8 | 95.7 | 7.4 | 7.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.372 | 11.775 | 11.689 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0122-0003 | 20.6797608 | -0.0530356 | 151.7 | 162.8 | 9.6 | 9.1 | $\cdots$ | $\cdots$ | 13.816 | 12.612 | 11.930 | 11.729 | $\cdots$ | 76.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0126-0000 | 21.7387347 | -0.0141164 | 188.1 | 43.8 | 12.2 | 12.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.074 | 11.374 | 11.219 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0137-0537 | 24.3814961 | -5.6331831 | 236.7 | -75.4 | 8.4 | 8.7 | $\cdots$ | $\cdots$ | $\cdots$ | 12.573 | 12.028 | 11.810 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0155-4344 | 28.9240619 | -43.7373506 | 148.1 | 106.2 | 2.0 | 3.5 | 16.195 | 13.912 | 12.128 | 11.368 | 10.829 | 10.575 | 2.544 | 59.8 |
UPM 0202-4056 | 30.6998675 | -40.9468992 | 180.1 | -67.9 | 2.6 | 2.6 | $\cdots$ | $\cdots$ | $\cdots$ | 10.454 | 9.754 | 9.572 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 0209-3339A | 32.4333397 | -33.6586869 | -86.1 | -166.9 | 4.2 | 2.0 | 16.526 | 14.533 | 12.628 | 11.437 | 10.885 | 10.605 | 3.096 | 49.5 | ddCommon proper motion companion; see Table 4
UPM 0209-3339B | 32.4371381 | -33.6580217 | -112.9 | -170.2 | 7.6 | 8.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.025 | 11.489 | 11.131 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 0218-0120 | 34.5873781 | -1.3488447 | 166.5 | -88.5 | 9.9 | 10.5 | 13.640 | 12.620 | 12.263 | 12.468 | 12.109 | 12.034 | 0.152 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0238-1348 | 39.6759683 | -13.8016906 | 173.3 | -114.7 | 9.9 | 9.6 | $\cdots$ | $\cdots$ | $\cdots$ | 11.998 | 11.375 | 11.263 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0241-1647 | 40.2660111 | -16.7848081 | 16.5 | -182.4 | 11.2 | 11.0 | 15.606 | 14.528 | 14.078 | 13.429 | 12.940 | 12.857 | 1.099 | [216.7] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 0302-1307 | 45.5532964 | -13.1308544 | 139.0 | -117.1 | 16.5 | 14.2 | 17.954 | 16.054 | 14.408 | 13.364 | 12.816 | 12.536 | 2.690 | 146.8 |
UPM 0308-0532 | 47.0978864 | -5.5470081 | 82.1 | -175.0 | 3.4 | 4.6 | 11.813 | 10.766 | 10.286 | 10.222 | 9.783 | 9.698 | 0.544 | 53.7 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0314-0150 | 48.5090014 | -1.8370214 | 114.2 | -186.2 | 8.8 | 8.6 | $\cdots$ | $\cdots$ | $\cdots$ | 11.672 | 11.088 | 10.843 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0314-0415 | 48.6207794 | -4.2554964 | 145.6 | -111.8 | 9.0 | 10.1 | 14.625 | 12.900 | 12.119 | 11.176 | 10.525 | 10.401 | 1.724 | 71.5 |
UPM 0330-0047 | 52.7083053 | -0.7902283 | 193.7 | -15.3 | 3.3 | 6.2 | 15.120 | 14.071 | 13.434 | 13.066 | 12.604 | 12.561 | 1.005 | [190.6] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 0356-0828 | 59.0545181 | -8.4810458 | 3.7 | -184.8 | 10.8 | 12.5 | $\cdots$ | 11.147 | 10.560 | 10.892 | 10.507 | 10.400 | 0.255 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0358-1617 | 59.5358836 | -16.2881436 | 184.6 | -37.0 | 12.3 | 12.5 | 17.926 | 15.806 | 13.993 | 12.933 | 12.442 | 12.212 | 2.873 | 115.0 |
UPM 0359-2449 | 59.8231094 | -24.8298217 | -105.0 | -149.1 | 4.6 | 4.6 | 17.990 | 15.911 | 13.990 | 12.757 | 12.211 | 11.956 | 3.154 | 88.6 |
UPM 0403-0635 | 60.7777225 | -6.5860825 | 205.2 | 13.7 | 11.2 | 10.6 | $\cdots$ | $\cdots$ | $\cdots$ | 12.140 | 11.563 | 11.276 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0405-1951 | 61.2543614 | -19.8521608 | 34.8 | -179.9 | 2.4 | 2.4 | $\cdots$ | 13.753 | 13.046 | 12.967 | 12.410 | 12.355 | 0.786 | 199.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 0410-0742 | 62.6567331 | -7.7134942 | 139.1 | -175.2 | 11.8 | 10.8 | 15.766 | 13.474 | 11.071 | 10.642 | 10.054 | 9.770 | 2.832 | 37.0 |
UPM 0413-4212 | 63.2980333 | -42.2011461 | 145.8 | 124.2 | 3.3 | 3.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.229 | 11.632 | 11.368 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0417-0431 | 64.3450542 | -4.5310789 | 244.9 | 78.0 | 10.3 | 10.6 | $\cdots$ | $\cdots$ | $\cdots$ | 11.392 | 10.819 | 10.545 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0424-0307 | 66.1372464 | -3.1186994 | 133.7 | -176.5 | 12.7 | 15.8 | 17.311 | 14.694 | 12.752 | 12.087 | 11.504 | 11.214 | 2.607 | 72.7 |
UPM 0443-4129A | 70.7825969 | -41.4844111 | 186.1 | 4.3 | 3.6 | 3.8 | 14.948 | 12.775 | 10.462 | 10.422 | 9.846 | 9.594 | 2.353 | 39.3 | aaProper motions suspect ddCommon proper motion companion; see Table 4
UPM 0446-4337 | 71.5897783 | -43.6211606 | 59.2 | 171.3 | 3.6 | 1.7 | 12.496 | 11.113 | 10.479 | 11.096 | 10.602 | 10.500 | 0.017 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0454-4217 | 73.5640600 | -42.2886506 | 159.7 | 83.2 | 10.9 | 10.5 | 16.433 | 14.437 | 12.943 | 12.095 | 11.431 | 11.275 | 2.342 | 91.8 |
UPM 0456-1138 | 74.1655053 | -11.6350681 | 30.5 | -250.8 | 11.5 | 11.3 | 17.224 | 15.471 | 13.518 | 12.269 | 11.716 | 11.471 | 3.202 | 74.7 |
UPM 0502-1938 | 75.6838831 | -19.6478917 | 54.4 | 172.2 | 4.0 | 2.8 | 16.562 | 14.733 | 12.874 | 11.574 | 10.917 | 10.641 | 3.159 | 48.8 |
UPM 0507-1302 | 76.8394261 | -13.0453781 | 184.1 | -119.4 | 9.8 | 9.7 | $\cdots$ | $\cdots$ | $\cdots$ | 12.601 | 12.004 | 11.786 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0508-0617 | 77.0132517 | -6.2872236 | 170.4 | -67.8 | 7.5 | 7.6 | $\cdots$ | 14.170 | 12.642 | 11.679 | 11.078 | 10.854 | 2.491 | 66.8 |
UPM 0521-0448 | 80.2826481 | -4.8134942 | -20.5 | -181.1 | 4.5 | 4.1 | 14.255 | 13.177 | 12.882 | 12.238 | 11.770 | 11.736 | 0.939 | 122.1 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0527-2006 | 81.8761928 | -20.1036386 | 25.6 | -183.7 | 3.4 | 3.4 | $\cdots$ | 10.973 | 9.523 | 11.168 | 10.591 | 10.335 | -0.195 | 60.1 | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0528-4313A | 82.0375869 | -43.2255300 | -75.6 | 164.7 | 4.0 | 4.0 | 16.795 | 14.566 | 12.807 | 11.889 | 11.312 | 11.067 | 2.677 | 70.4 | ddCommon proper motion companion; see Table 4
UPM 0528-4313B | 82.0298050 | -43.2357586 | -86.3 | 163.2 | 4.8 | 5.5 | 19.746 | 17.653 | 15.348 | 13.943 | 13.296 | 13.015 | 3.710 | 109.1 | ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 0532-1458 | 83.2114242 | -14.9684436 | 183.3 | -28.1 | 10.2 | 10.2 | 16.626 | 14.662 | 12.946 | 12.145 | 11.623 | 11.410 | 2.517 | 95.2 |
UPM 0534-1510 | 83.6021194 | -15.1696369 | 177.5 | -39.2 | 8.5 | 8.5 | 14.461 | 12.575 | 11.551 | 11.040 | 10.429 | 10.249 | 1.535 | 70.9 |
UPM 0540-2757 | 85.0146869 | -27.9660519 | 33.1 | -120.7 | 3.6 | 3.7 | 18.317 | 16.303 | 14.451 | 13.494 | 12.946 | 12.753 | 2.809 | 154.1 | aaProper motions suspect eeNot detected during automated search but noticed by eye during the blinking process
UPM 0542-4544 | 85.6736511 | -45.7491522 | 64.1 | 183.2 | 4.3 | 7.9 | 15.664 | 13.488 | 11.376 | 10.453 | 9.854 | 9.586 | 3.035 | 30.9 |
UPM 0545-0222 | 86.3577581 | -2.3675183 | -54.4 | -189.2 | 10.0 | 10.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.424 | 11.861 | 11.617 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0556-3937 | 89.0922544 | -39.6272753 | 84.1 | -206.8 | 4.7 | 5.6 | $\cdots$ | $\cdots$ | $\cdots$ | 13.607 | 13.063 | 12.885 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0557-1351 | 89.2654756 | -13.8644606 | 25.1 | -181.4 | 5.6 | 3.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.038 | 10.475 | 10.205 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0558-3745 | 89.7238911 | -37.7586756 | 140.6 | -115.6 | 3.9 | 5.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.662 | 11.035 | 10.785 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0601-1512 | 90.4218800 | -15.2024047 | 105.0 | -157.1 | 7.5 | 5.2 | 14.492 | 12.684 | 11.569 | 11.484 | 10.841 | 10.682 | 1.200 | 91.1 |
UPM 0602-1917 | 90.6749681 | -19.2984797 | -132.2 | -172.6 | 10.4 | 10.4 | 15.906 | 13.961 | 12.074 | 11.323 | 10.858 | 10.519 | 2.638 | 60.7 |
UPM 0603-1417 | 90.8013183 | -14.2981833 | -184.4 | -12.4 | 4.5 | 10.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.711 | 11.092 | 10.805 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0603-1433 | 90.8642300 | -14.5603703 | 49.8 | -195.9 | 3.6 | 3.6 | $\cdots$ | $\cdots$ | $\cdots$ | 12.090 | 11.491 | 11.327 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0614-3350 | 93.6930136 | -33.8383233 | 235.6 | 92.1 | 10.9 | 10.9 | $\cdots$ | 15.462 | 14.028 | 13.052 | 12.565 | 12.354 | 2.410 | 141.9 |
UPM 0615-0735 | 93.9365817 | -7.5865625 | 96.7 | -156.9 | 11.5 | 13.3 | $\cdots$ | 16.025 | 14.794 | 13.709 | 13.130 | 12.892 | 2.316 | 178.1 |
UPM 0620-0312 | 95.1532386 | -3.2076081 | -19.4 | -204.3 | 10.5 | 10.9 | 16.634 | 14.530 | $\cdots$ | 12.965 | 12.383 | 12.177 | 1.565 | [189.4] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 0629-1737 | 97.4615906 | -17.6326997 | 143.0 | -130.9 | 6.4 | 9.5 | $\cdots$ | $\cdots$ | $\cdots$ | 12.895 | 12.363 | 12.198 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0639-0451 | 99.9698181 | -4.8652011 | -229.8 | 121.5 | 7.5 | 7.6 | 16.020 | 13.918 | 12.144 | 11.205 | 10.624 | 10.385 | 2.713 | 52.0 |
UPM 0641-0255 | 100.2841003 | -2.9306503 | -72.7 | -167.6 | 18.7 | 11.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.551 | 10.989 | 10.716 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0645-0045 | 101.2847811 | -0.7658139 | -168.9 | -95.2 | 7.2 | 6.6 | 17.124 | 15.014 | 13.172 | 12.186 | 11.639 | 11.379 | 2.828 | 78.6 |
UPM 0652-0150 | 103.0572631 | -1.8487686 | 64.1 | -174.4 | 1.4 | 2.2 | 13.647 | 11.537 | 10.484 | 11.121 | 10.480 | 10.334 | 0.416 | 76.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0652-2243 | 103.2459989 | -22.7292186 | -83.2 | 169.2 | 2.8 | 3.2 | $\cdots$ | 14.509 | 12.303 | 11.303 | 10.719 | 10.419 | 3.206 | 43.4 |
UPM 0655-0715 | 103.8760131 | -7.2643236 | 61.7 | -240.9 | 7.8 | 7.9 | $\cdots$ | $\cdots$ | $\cdots$ | 10.392 | 9.779 | 9.563 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0659-0052A | 104.8421303 | -0.8801367 | -58.3 | -184.1 | 2.7 | 1.8 | 15.062 | 12.642 | 10.985 | 11.253 | 10.672 | 10.466 | 1.389 | 78.2 | ddCommon proper motion companion; see Table 4
UPM 0659-0052B | 104.8439419 | -0.8835767 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 13.969 | 13.436 | 13.172 | $\cdots$ | $\cdots$ | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 0702-2053 | 105.5694072 | -20.8837058 | 82.0 | -169.4 | 8.2 | 9.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.524 | 11.899 | 11.691 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0704-0602A | 106.1883350 | -6.0386081 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 12.488 | 11.140 | 12.557 | 11.932 | 11.758 | -0.069 | 123.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable aaProper motions suspect ccSuperCOSMOS plate magnitudes suspect ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 0704-0602B | 106.1882444 | -6.0352122 | 99.5 | -153.0 | 5.0 | 5.0 | $\cdots$ | 14.488 | 13.710 | 11.284 | 10.723 | 10.499 | 3.204 | 37.8 | ddCommon proper motion companion; see Table 4
UPM 0704-2033 | 106.1582275 | -20.5532164 | 90.1 | -158.2 | 10.3 | 10.1 | 17.435 | 15.252 | 13.485 | 12.546 | 12.008 | 11.769 | 2.706 | 98.2 |
UPM 0705-2830 | 106.3123578 | -28.5031700 | -23.5 | 182.2 | 4.7 | 19.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.831 | 11.247 | 10.980 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0706-2301 | 106.5726647 | -23.0297381 | -164.1 | 105.6 | 8.5 | 7.7 | 17.987 | 15.928 | 14.413 | 13.449 | 12.936 | 12.646 | 2.479 | 164.7 |
UPM 0708-2539 | 107.2233878 | -25.6601586 | 115.2 | -144.9 | 13.0 | 17.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.311 | 11.694 | 11.461 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0716-2342 | 109.2296050 | -23.7101622 | 139.8 | -113.7 | 9.7 | 9.5 | $\cdots$ | 15.025 | 12.960 | 12.396 | 11.860 | 11.611 | 2.629 | 98.8 |
UPM 0724-0949 | 111.1639358 | -9.8236900 | -192.9 | 15.5 | 5.0 | 5.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.127 | 10.519 | 10.353 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0725-0207 | 111.4631742 | -2.1191361 | -185.1 | -14.9 | 5.0 | 5.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.251 | 10.700 | 10.483 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0730-4042 | 112.7328597 | -40.7005344 | -237.6 | 113.1 | 7.4 | 7.0 | $\cdots$ | $\cdots$ | $\cdots$ | 10.452 | 9.848 | 9.609 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0736-4256 | 114.0765108 | -42.9425744 | 171.4 | -62.3 | 14.4 | 12.4 | $\cdots$ | 15.719 | 14.674 | 13.476 | 12.860 | 12.698 | 2.243 | 164.2 |
UPM 0740-2114 | 115.1509197 | -21.2383561 | -192.2 | 8.5 | 14.0 | 13.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.420 | 11.875 | 11.621 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0740-3055 | 115.1015369 | -30.9323542 | -235.9 | 96.7 | 7.7 | 7.8 | $\cdots$ | $\cdots$ | $\cdots$ | 12.146 | 11.616 | 11.364 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0742-2501 | 115.5580681 | -25.0244169 | 176.9 | 49.8 | 9.7 | 9.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.848 | 11.231 | 10.983 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0745-4149 | 116.2558594 | -41.8251314 | -89.8 | 159.3 | 2.2 | 2.4 | $\cdots$ | 14.674 | 13.178 | 12.175 | 11.583 | 11.358 | 2.499 | 83.7 |
UPM 0746-3729 | 116.5312928 | -37.4876147 | 11.2 | -182.4 | 4.4 | 4.2 | $\cdots$ | $\cdots$ | $\cdots$ | 12.527 | 11.994 | 11.710 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 0747-0320 | 116.8618994 | -3.3466481 | -158.2 | 89.5 | 7.0 | 7.6 | $\cdots$ | 14.709 | 12.737 | 11.930 | 11.345 | 11.126 | 2.779 | 71.9 |
UPM 0747-2537A | 116.9783592 | -25.6193264 | -148.5 | 101.9 | 4.2 | 2.1 | $\cdots$ | 9.801 | 8.849 | 9.593 | 8.958 | 8.814 | 0.208 | 40.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 0747-2537B | 116.9752431 | -25.6211242 | -151.3 | 102.3 | 3.5 | 3.5 | $\cdots$ | 13.947 | 13.441 | 11.500 | 10.965 | 10.741 | 2.447 | 47.3 | ddCommon proper motion companion; see Table 4
UPM 0748-0619 | 117.1063369 | -6.3225081 | 109.8 | -155.9 | 6.1 | 6.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.356 | 9.749 | 9.579 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0750-1807 | 117.5875447 | -18.1294489 | -182.7 | 45.5 | 7.0 | 6.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.248 | 11.756 | 11.471 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0751-0214 | 117.7670006 | -2.2466175 | 9.8 | -211.2 | 13.9 | 11.7 | $\cdots$ | $\cdots$ | $\cdots$ | 12.186 | 11.708 | 11.441 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0752-0751 | 118.0258133 | -7.8544356 | 12.8 | -195.6 | 5.9 | 5.9 | $\cdots$ | $\cdots$ | $\cdots$ | 10.377 | 9.794 | 9.516 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0752-1602 | 118.1642550 | -16.0351689 | -86.0 | -193.7 | 7.3 | 7.5 | $\cdots$ | $\cdots$ | $\cdots$ | 11.770 | 11.214 | 10.983 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0800-0617A | 120.2021822 | -6.2902914 | 135.2 | -233.8 | 9.0 | 9.4 | $\cdots$ | 14.669 | $\cdots$ | 12.944 | 12.321 | 12.176 | 1.725 | [175.5] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 0800-0617B | 120.2007742 | -6.2896208 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 14.720 | 14.193 | 13.983 | $\cdots$ | $\cdots$ | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 0801-1005 | 120.4445300 | -10.0905344 | 177.8 | -90.1 | 3.8 | 6.2 | $\cdots$ | $\cdots$ | $\cdots$ | 9.969 | 9.388 | 9.154 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0801-4347 | 120.4321706 | -43.7879044 | -100.9 | 158.0 | 3.6 | 3.4 | 14.186 | 11.870 | 10.471 | 10.146 | 9.546 | 9.327 | 1.724 | 44.4 |
UPM 0802-2010 | 120.6659497 | -20.1753642 | 191.5 | -210.6 | 7.7 | 7.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.925 | 11.310 | 11.141 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0803-4518 | 120.7579836 | -45.3141208 | 49.4 | 199.4 | 3.3 | 3.3 | $\cdots$ | $\cdots$ | $\cdots$ | 11.645 | 11.109 | 10.813 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0805-3827 | 121.3922181 | -38.4596086 | 153.3 | -151.8 | 6.6 | 6.2 | $\cdots$ | $\cdots$ | $\cdots$ | 12.108 | 11.584 | 11.251 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0807-0121 | 121.8875364 | -1.3593131 | -178.8 | 94.3 | 10.2 | 10.5 | $\cdots$ | $\cdots$ | $\cdots$ | 11.071 | 10.473 | 10.277 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0807-0930 | 121.8525969 | -9.5078494 | 197.5 | 32.8 | 3.9 | 12.3 | 13.602 | 11.228 | 9.767 | 9.758 | 9.148 | 8.921 | 1.470 | 38.5 |
UPM 0807-2025 | 121.8779764 | -20.4218667 | 110.5 | -144.7 | 3.0 | 3.0 | $\cdots$ | $\cdots$ | $\cdots$ | 10.948 | 10.335 | 10.137 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0808-0943 | 122.0687272 | -9.7311883 | 98.7 | -173.6 | 4.5 | 3.6 | 14.522 | 12.217 | 10.408 | 10.195 | 9.679 | 9.442 | 2.022 | 43.2 |
UPM 0809-0943 | 122.2784283 | -9.7207692 | 147.0 | -157.5 | 8.1 | 7.0 | 15.935 | 13.769 | 12.187 | 11.895 | 11.357 | 11.154 | 1.874 | 101.9 |
UPM 0809-4519 | 122.4626147 | -45.3295575 | 204.4 | -10.1 | 12.6 | 12.6 | $\cdots$ | $\cdots$ | $\cdots$ | 10.236 | 9.686 | 9.446 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0813-0429 | 123.3408675 | -4.4854875 | -32.9 | 225.8 | 8.2 | 9.0 | $\cdots$ | 15.922 | $\cdots$ | 13.271 | 12.743 | 12.499 | 2.651 | 137.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0813-1048 | 123.3918000 | -10.8166544 | -147.7 | -139.7 | 10.7 | 10.8 | 17.723 | 15.527 | 13.202 | 12.094 | 11.610 | 11.292 | 3.433 | 58.5 |
UPM 0813-4604 | 123.2772747 | -46.0709556 | -76.9 | 168.7 | 3.0 | 3.1 | 15.131 | 13.421 | 11.180 | 10.955 | 10.281 | 10.073 | 2.466 | 53.7 |
UPM 0814-0835 | 123.6868333 | -8.5917794 | -221.5 | 113.0 | 7.4 | 6.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.087 | 10.481 | 10.297 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0814-3645 | 123.7166631 | -36.7522964 | -80.9 | -167.1 | 10.2 | 4.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.660 | 11.150 | 10.963 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0815-3058 | 123.7880314 | -30.9746633 | -170.6 | 205.9 | 6.1 | 6.2 | $\cdots$ | $\cdots$ | 11.245 | 11.346 | 10.782 | 10.473 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0821-1452 | 125.3709208 | -14.8695183 | -22.6 | -183.9 | 3.6 | 5.5 | 12.697 | $\cdots$ | $\cdots$ | 10.073 | 9.481 | 9.376 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0821-4626 | 125.4251800 | -46.4334139 | -73.3 | 166.4 | 3.3 | 3.3 | 15.017 | 13.307 | 11.553 | 11.528 | 10.980 | 10.748 | 1.779 | 82.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0826-3942 | 126.5266758 | -39.7096897 | -98.3 | 196.1 | 6.2 | 6.3 | $\cdots$ | $\cdots$ | $\cdots$ | 10.745 | 10.153 | 9.911 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0828-1438 | 127.0004200 | -14.6352944 | -166.9 | 87.9 | 8.4 | 5.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.204 | 11.541 | 11.285 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0832-3942 | 128.0769122 | -39.7043753 | 224.2 | 4.0 | 14.0 | 13.9 | 18.269 | 16.870 | 14.829 | 13.047 | 12.517 | 12.256 | 3.823 | 87.7 |
UPM 0836-0059 | 129.0747936 | -0.9920019 | 109.8 | -147.3 | 6.8 | 6.7 | 15.555 | 13.318 | 11.750 | 10.932 | 10.342 | 10.051 | 2.386 | 49.2 |
UPM 0838-3247 | 129.6943083 | -32.7966136 | -106.7 | 147.5 | 4.6 | 3.9 | $\cdots$ | 14.149 | 12.693 | 12.312 | 11.637 | 11.492 | 1.837 | 115.5 |
UPM 0840-3641 | 130.0511394 | -36.6993433 | -173.6 | 49.8 | 4.9 | 3.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.154 | 10.611 | 10.320 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 0840-4437 | 130.1662658 | -44.6211017 | -84.9 | 161.1 | 2.2 | 2.2 | $\cdots$ | 12.483 | 11.812 | 11.471 | 11.056 | 11.013 | 1.012 | 106.9 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0842-0302 | 130.6869978 | -3.0436728 | 51.0 | -176.8 | 7.3 | 7.3 | $\cdots$ | 15.299 | 13.728 | 12.879 | 12.353 | 12.126 | 2.420 | 128.1 |
UPM 0842-0907 | 130.5108458 | -9.1258817 | 130.2 | -126.5 | 7.0 | 7.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.276 | 9.675 | 9.407 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0842-4532 | 130.7227961 | -45.5387481 | -151.7 | 98.1 | 6.7 | 4.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.610 | 12.124 | 11.903 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0843-3209 | 130.9895922 | -32.1573656 | -165.4 | 76.0 | 3.2 | 2.6 | 15.446 | 13.483 | 12.412 | 11.449 | 10.860 | 10.650 | 2.034 | 77.5 |
UPM 0846-2639 | 131.5501986 | -26.6632953 | -136.2 | 122.8 | 2.5 | 2.6 | $\cdots$ | $\cdots$ | 11.530 | 11.086 | 10.434 | 10.236 | $\cdots$ | 62.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0848-2542 | 132.1816267 | -25.7036467 | -181.8 | 7.0 | 1.8 | 6.2 | $\cdots$ | 13.801 | 12.021 | 11.335 | 10.765 | 10.529 | 2.466 | 61.3 |
UPM 0850-3052 | 132.6042867 | -30.8716572 | 73.5 | -167.7 | 5.9 | 3.8 | 17.121 | 14.799 | 13.286 | 12.239 | 11.726 | 11.439 | 2.560 | 86.4 |
UPM 0856-1741 | 134.0761817 | -17.6836933 | -124.7 | 134.3 | 1.7 | 2.9 | 17.414 | 15.472 | 13.932 | 12.293 | 11.719 | 11.460 | 3.179 | 69.5 |
UPM 0856-2909 | 134.0409614 | -29.1542111 | 255.3 | 2.9 | 9.1 | 8.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.988 | 11.453 | 11.215 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0900-4308 | 135.2269369 | -43.1379197 | -177.1 | 38.5 | 3.9 | 2.1 | $\cdots$ | 15.782 | 13.872 | 12.083 | 11.549 | 11.323 | 3.699 | 47.7 |
UPM 0913-4303 | 138.2865189 | -43.0630097 | -167.0 | 78.6 | 3.8 | 4.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.673 | 12.118 | 11.868 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0919-3821 | 139.7869492 | -38.3544064 | 181.0 | -27.6 | 10.7 | 5.2 | 18.201 | 16.059 | 13.594 | 12.398 | 11.839 | 11.577 | 3.661 | 60.3 |
UPM 0923-2518 | 140.8532544 | -25.3069511 | -174.6 | 66.3 | 5.1 | 5.0 | 18.212 | 16.305 | 14.672 | 13.104 | 12.527 | 12.241 | 3.201 | 99.2 |
UPM 0934-4355 | 143.5423842 | -43.9186631 | -178.6 | -36.3 | 6.1 | 4.5 | $\cdots$ | 15.225 | 12.743 | 11.673 | 11.138 | 10.879 | 3.552 | 49.2 |
UPM 0936-4557 | 144.2119747 | -45.9659628 | -175.8 | 51.4 | 3.4 | 3.3 | 12.959 | 11.119 | 9.990 | 9.830 | 9.194 | 9.028 | 1.289 | 41.8 |
UPM 0937-0014 | 144.4922225 | -0.2380892 | -186.6 | -52.5 | 6.9 | 7.1 | 17.167 | 15.231 | 14.180 | 13.184 | 12.723 | 12.483 | 2.047 | 184.6 |
UPM 0937-3214 | 144.4728728 | -32.2365067 | -92.6 | 155.7 | 3.6 | 3.4 | $\cdots$ | 15.513 | 14.370 | 13.263 | 12.703 | 12.492 | 2.250 | 152.5 |
UPM 0940-3918 | 145.0291622 | -39.3121189 | -75.3 | 164.2 | 6.9 | 6.9 | $\cdots$ | 13.681 | 11.832 | 11.519 | 10.957 | 10.738 | 2.162 | 76.4 |
UPM 0941-4439 | 145.3561000 | -44.6615028 | -203.9 | -71.1 | 11.3 | 11.0 | $\cdots$ | 15.616 | 13.837 | 12.855 | 12.327 | 12.069 | 2.761 | 109.0 |
UPM 0941-4518 | 145.3343672 | -45.3073242 | -174.8 | 89.0 | 10.9 | 6.3 | $\cdots$ | $\cdots$ | $\cdots$ | 11.682 | 11.132 | 10.884 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0948-4147 | 147.2015628 | -41.7863742 | -245.0 | -69.8 | 8.4 | 8.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.646 | 10.069 | 9.803 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 0951-4429 | 147.7806564 | -44.4948356 | -143.8 | 109.8 | 3.5 | 2.1 | 16.883 | 14.919 | 12.584 | 11.698 | 11.158 | 10.966 | 3.221 | 58.8 |
UPM 0956-0940 | 149.2293078 | -9.6741689 | -192.7 | -24.4 | 2.5 | 2.0 | 14.329 | $\cdots$ | 10.369 | 10.191 | 9.584 | 9.352 | $\cdots$ | 43.6 |
UPM 1003-2717 | 150.9001517 | -27.2959828 | -274.6 | -39.5 | 10.7 | 10.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.565 | 10.007 | 9.712 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1006-1144 | 151.5145706 | -11.7484167 | 66.9 | -170.5 | 2.0 | 2.0 | 13.619 | 11.759 | 10.923 | 11.193 | 10.586 | 10.439 | 0.566 | 90.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1009-0501B | 152.4402375 | -5.0217161 | -190.5 | 92.1 | 7.4 | 7.2 | $\cdots$ | $\cdots$ | $\cdots$ | 12.645 | 12.105 | 11.889 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 1011-4235 | 152.9185411 | -42.5901381 | -178.6 | -25.9 | 2.8 | 2.8 | 15.694 | 13.413 | $\cdots$ | 11.073 | 10.461 | 10.227 | 2.340 | 54.9 |
UPM 1015-0859 | 153.8675097 | -8.9998406 | -151.3 | 101.4 | 2.0 | 2.0 | 14.108 | 12.279 | 11.435 | 10.688 | 10.044 | 9.917 | 1.591 | 59.7 |
UPM 1020-0633A | 155.2034769 | -6.5554489 | -179.8 | -27.8 | 7.3 | 7.2 | 16.023 | 13.625 | 11.335 | 10.670 | 10.073 | 9.809 | 2.955 | 34.8 | ddCommon proper motion companion; see Table 4
UPM 1020-0642 | 155.0026631 | -6.7037208 | -211.0 | 70.6 | 7.6 | 7.6 | 16.635 | 14.687 | 12.883 | 11.920 | 11.350 | 11.148 | 2.767 | 75.2 |
UPM 1024-0317 | 156.1723939 | -3.2861519 | -143.2 | -126.3 | 7.9 | 7.9 | 16.780 | 14.686 | 12.708 | 11.846 | 11.276 | 11.025 | 2.840 | 67.0 |
UPM 1030-2400 | 157.5368944 | -24.0092022 | -164.3 | 76.2 | 4.0 | 2.0 | 16.458 | 14.603 | 12.580 | 11.291 | 10.690 | 10.449 | 3.312 | 42.7 |
UPM 1031-0024A | 157.7905397 | -0.4118389 | -207.4 | -105.6 | 11.4 | 10.3 | 14.016 | 11.944 | 10.371 | 10.561 | 10.048 | 9.747 | 1.383 | 55.1 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 1031-0024B | 157.7925844 | -0.4118883 | -142.5 | -96.9 | 6.3 | 6.3 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1041-3913 | 160.2807536 | -39.2238969 | 38.1 | -182.1 | 1.6 | 14.5 | 14.487 | 12.444 | 11.094 | 10.148 | 9.600 | 9.347 | 2.296 | 38.6 |
UPM 1044-2414 | 161.2433072 | -24.2419103 | -178.9 | 57.3 | 11.1 | 9.2 | 16.647 | 14.667 | 13.149 | 12.003 | 11.441 | 11.199 | 2.664 | 78.4 |
UPM 1046-3046 | 161.5409683 | -30.7693019 | -84.2 | -159.9 | 3.6 | 3.6 | $\cdots$ | $\cdots$ | $\cdots$ | 12.983 | 12.405 | 12.164 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1048-1538 | 162.1992581 | -15.6490136 | -177.7 | 42.4 | 3.3 | 4.4 | 17.396 | 15.382 | 13.938 | 13.043 | 12.412 | 12.221 | 2.339 | 141.8 |
UPM 1056-0542A | 164.1791192 | -5.7066733 | -98.1 | -173.8 | 8.2 | 8.3 | 15.970 | 14.281 | 12.816 | 11.682 | 11.141 | 10.939 | 2.599 | 76.5 | ddCommon proper motion companion; see Table 4
UPM 1056-0542B | 164.1816494 | -5.7064783 | -63.9 | -173.3 | 4.4 | 4.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.485 | 11.963 | 11.709 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1058-4441 | 164.5727686 | -44.6861050 | -168.3 | 66.0 | 9.8 | 4.2 | $\cdots$ | 16.249 | 15.218 | 13.872 | 13.176 | 13.037 | 2.377 | 175.5 |
UPM 1059-0020 | 164.9792050 | -0.3418581 | 102.1 | -152.1 | 2.9 | 5.0 | $\cdots$ | 11.771 | $\cdots$ | 10.356 | 9.748 | 9.586 | 1.415 | 57.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1059-3022 | 164.7671950 | -30.3801764 | -174.6 | 59.9 | 2.1 | 2.1 | 12.886 | 11.393 | 10.988 | 10.299 | 9.826 | 9.725 | 1.094 | 50.7 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1108-0644 | 167.1859397 | -6.7383533 | -205.9 | -1.5 | 8.7 | 8.5 | 14.710 | 12.648 | 11.392 | 10.927 | 10.290 | 10.088 | 1.721 | 65.0 |
UPM 1109-1032 | 167.3681956 | -10.5479519 | -162.3 | 89.0 | 10.1 | 10.1 | 17.490 | 15.334 | 13.246 | 12.484 | 11.943 | 11.663 | 2.850 | 89.5 |
UPM 1113-0113 | 168.3352431 | -1.2199069 | -184.6 | 23.4 | 10.5 | 10.0 | 17.881 | 15.813 | 14.194 | 12.982 | 12.530 | 12.253 | 2.831 | 119.8 |
UPM 1113-0148 | 168.3518608 | -1.8163381 | -174.3 | -64.9 | 7.4 | 11.6 | 17.348 | 15.208 | 13.359 | 12.395 | 11.858 | 11.551 | 2.813 | 84.9 |
UPM 1122-4530 | 170.6087600 | -45.5091333 | -179.3 | -29.6 | 1.9 | 1.9 | 17.065 | 15.177 | 12.936 | 11.666 | 11.061 | 10.816 | 3.511 | 46.4 |
UPM 1130-1622 | 172.5409028 | -16.3801233 | -200.6 | 57.3 | 6.6 | 7.0 | 17.842 | 15.713 | 14.155 | 13.372 | 12.854 | 12.562 | 2.341 | 166.1 | aaProper motions suspect
UPM 1131-0725 | 172.9437592 | -7.4266169 | -195.6 | -5.9 | 9.9 | 9.5 | 14.975 | 13.881 | 13.264 | 13.077 | 12.616 | 12.538 | 0.804 | 194.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1133-2302 | 173.3741581 | -23.0458669 | -191.9 | 0.6 | 1.7 | 1.7 | 16.127 | 14.100 | 12.923 | 12.523 | 11.928 | 11.745 | 1.577 | 145.4 |
UPM 1136-0525 | 174.1784361 | -5.4211494 | 158.0 | -140.2 | 13.6 | 13.5 | 17.294 | 15.429 | 13.996 | 12.813 | 12.249 | 11.989 | 2.616 | 117.1 |
UPM 1138-4553 | 174.6189286 | -45.8979317 | -191.9 | 1.9 | 3.8 | 2.5 | 16.204 | 14.087 | 12.336 | 11.305 | 10.616 | 10.411 | 2.782 | 49.2 |
UPM 1142-2055A | 175.5808186 | -20.9279653 | -186.7 | 44.2 | 2.4 | 2.5 | 14.178 | 11.408 | 10.020 | 10.028 | 9.348 | 9.172 | 1.380 | 41.2 | ddCommon proper motion companion; see Table 4
UPM 1142-2055B | 175.5814128 | -20.9302367 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 12.223 | 11.647 | 11.370 | $\cdots$ | $\cdots$ | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1143-4443 | 175.8747928 | -44.7271675 | -191.0 | 2.6 | 5.1 | 4.0 | 18.119 | 16.055 | 14.341 | 12.943 | 12.399 | 12.169 | 3.112 | 99.0 |
UPM 1149-0019B | 177.2613256 | -0.3233508 | -201.5 | 2.2 | 6.8 | 3.2 | $\cdots$ | $\cdots$ | $\cdots$ | 9.963 | 9.345 | 9.145 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 1153-1747 | 178.3206078 | -17.7875069 | -186.5 | -58.7 | 1.3 | 1.3 | 15.381 | 13.160 | 11.557 | 10.957 | 10.352 | 10.110 | 2.203 | 55.0 |
UPM 1158-0912 | 179.6672408 | -9.2024986 | -189.9 | 52.2 | 9.2 | 9.4 | 16.060 | 14.074 | 12.342 | 11.466 | 10.893 | 10.655 | 2.608 | 63.3 |
UPM 1159-0339 | 179.9315281 | -3.6635244 | -36.4 | -193.0 | 12.2 | 13.0 | $\cdots$ | 15.447 | 13.581 | 12.770 | 12.177 | 11.932 | 2.677 | 106.6 |
UPM 1159-3623A | 179.9150078 | -36.3841397 | -182.1 | -101.4 | 12.7 | 13.1 | $\cdots$ | 15.642 | 14.060 | 12.934 | 12.402 | 12.169 | 2.708 | 113.1 | ddCommon proper motion companion; see Table 4
UPM 1159-3623B | 179.9110964 | -36.3820372 | -172.6 | -92.4 | 8.4 | 6.5 | $\cdots$ | 18.260 | 16.344 | 14.447 | 14.007 | 13.646 | 3.813 | 132.4 | ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1202-0625 | 180.5348450 | -6.4258914 | 87.6 | -161.6 | 5.7 | 7.2 | 14.976 | 12.697 | 11.068 | 10.489 | 9.857 | 9.595 | 2.208 | 42.5 |
UPM 1203-0053 | 180.8805247 | -0.8924633 | -208.2 | -4.3 | 7.3 | 3.9 | $\cdots$ | 14.348 | 13.979 | 13.036 | 12.618 | 12.547 | 1.312 | [183.7] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1206-1851 | 181.5521064 | -18.8579533 | -190.4 | -26.2 | 5.6 | 5.9 | $\cdots$ | 15.193 | 14.577 | 13.870 | 13.261 | 13.128 | 1.323 | [255.4] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1208-0904 | 182.0832711 | -9.0831478 | -202.5 | 129.8 | 9.2 | 9.0 | 16.076 | 13.897 | $\cdots$ | 12.038 | 11.475 | 11.253 | 1.859 | 113.1 |
UPM 1209-0721 | 182.4658233 | -7.3607258 | -185.8 | 7.3 | 9.5 | 9.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.706 | 11.107 | 10.926 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1214-4016 | 183.6045086 | -40.2807997 | -180.5 | 3.8 | 14.4 | 2.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.455 | 10.789 | 10.516 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1217-1032 | 184.2883256 | -10.5362844 | -141.8 | -116.8 | 6.0 | 4.4 | 17.627 | 15.433 | 13.371 | 12.421 | 11.834 | 11.602 | 3.012 | 79.2 |
UPM 1219-0238 | 184.8517200 | -2.6367969 | -63.6 | -186.1 | 5.0 | 3.5 | 15.593 | 14.171 | 13.591 | 13.297 | 12.765 | 12.675 | 0.874 | 204.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1223-1757 | 185.8103789 | -17.9516678 | -183.7 | 56.5 | 3.3 | 2.4 | $\cdots$ | 14.157 | 12.414 | 11.911 | 11.341 | 11.106 | 2.246 | 87.1 |
UPM 1223-2947 | 185.9150278 | -29.7938564 | -158.1 | 97.2 | 2.5 | 2.5 | 16.417 | 14.610 | 12.933 | 11.803 | 11.181 | 10.935 | 2.807 | 66.8 |
UPM 1226-2020A | 186.6684756 | -20.3442347 | -137.6 | -119.8 | 8.0 | 14.1 | 13.216 | 11.819 | 11.024 | 10.896 | 10.327 | 10.198 | 0.923 | 72.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 1226-2020B | 186.6675500 | -20.3425083 | -146.3 | -117.3 | 5.6 | 5.7 | $\cdots$ | $\cdots$ | $\cdots$ | 13.457 | 12.901 | 12.690 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1226-3516B | 186.6989514 | -35.2778583 | -200.5 | 38.3 | 7.0 | 5.2 | 19.238 | 17.467 | 15.558 | 13.823 | 13.303 | 13.082 | 3.644 | 127.5 | ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1226-3516C | 186.7202739 | -35.2868633 | -115.2 | 8.5 | 7.7 | 6.1 | 19.891 | 18.036 | 16.387 | 14.792 | 14.321 | 14.107 | 3.244 | 243.4 | aaProper motions suspect ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1229-3351 | 187.4657011 | -33.8588425 | -276.0 | -141.9 | 15.2 | 15.2 | 17.955 | 16.382 | $\cdots$ | 13.235 | 12.747 | 12.476 | 3.147 | 129.1 |
UPM 1230-0436 | 187.7377119 | -4.6000967 | 56.9 | -175.8 | 3.4 | 3.2 | 15.575 | 13.562 | 11.944 | 11.098 | 10.568 | 10.359 | 2.464 | 59.1 |
UPM 1230-0439 | 187.5189581 | -4.6604111 | -155.7 | 105.0 | 4.8 | 4.5 | 16.484 | 14.340 | 12.884 | 11.915 | 11.353 | 11.065 | 2.425 | 78.8 |
UPM 1230-1444 | 187.6268586 | -14.7465403 | -192.9 | 28.8 | 2.1 | 2.3 | 15.721 | 14.376 | 13.613 | 13.446 | 12.910 | 12.826 | 0.930 | 244.1 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1232-0322 | 188.0907006 | -3.3744908 | -126.2 | -141.1 | 14.5 | 7.9 | 17.217 | 14.989 | 13.389 | 12.279 | 11.703 | 11.451 | 2.710 | 81.8 |
UPM 1232-4612 | 188.1934892 | -46.2156881 | -181.0 | -67.6 | 2.7 | 2.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.123 | 10.578 | 10.303 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1234-0023 | 188.5658922 | -0.3942053 | -159.8 | -94.0 | 1.9 | 3.6 | $\cdots$ | $\cdots$ | 12.901 | 12.589 | 12.062 | 11.998 | $\cdots$ | 148.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1235-0331 | 188.8826600 | -3.5177714 | -132.6 | -124.9 | 8.4 | 8.1 | $\cdots$ | $\cdots$ | 13.529 | 13.101 | 12.480 | 12.423 | $\cdots$ | 174.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1237-2708 | 189.4098558 | -27.1349803 | -165.6 | -89.2 | 18.4 | 5.4 | $\cdots$ | $\cdots$ | $\cdots$ | 10.255 | 9.641 | 9.414 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1238-0038 | 189.7220656 | -0.6494094 | -50.5 | -210.9 | 7.5 | 7.4 | 16.055 | 14.536 | 13.893 | 13.547 | 13.009 | 12.833 | 0.989 | 217.0 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1239-2521 | 189.9172039 | -25.3574425 | -131.8 | -125.4 | 1.0 | 0.9 | 12.470 | 11.395 | 10.838 | 10.671 | 10.251 | 10.189 | 0.724 | 66.7 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1240-1232 | 190.1730197 | -12.5345394 | -172.3 | -54.7 | 3.4 | 3.5 | 17.309 | 14.763 | $\cdots$ | 11.828 | 11.189 | 10.977 | 2.935 | 52.3 |
UPM 1243-0333 | 190.8523322 | -3.5658900 | -55.5 | -171.6 | 13.7 | 13.0 | $\cdots$ | 16.023 | 14.247 | 13.040 | 12.509 | 12.241 | 2.983 | 104.4 |
UPM 1246-3423 | 191.7417556 | -34.3920336 | -182.1 | 10.0 | 6.5 | 9.4 | 17.930 | 15.820 | 14.086 | 12.864 | 12.270 | 12.097 | 2.956 | 101.8 |
UPM 1250-3103 | 192.7028144 | -31.0625358 | -182.0 | -16.8 | 1.5 | 1.4 | 16.080 | 14.009 | 12.308 | 11.528 | 10.821 | 10.633 | 2.481 | 63.0 |
UPM 1254-0633 | 193.6115931 | -6.5503981 | -183.5 | 50.8 | 7.4 | 7.2 | 16.847 | 14.677 | 13.142 | 12.012 | 11.460 | 11.202 | 2.665 | 75.6 |
UPM 1255-0123 | 193.9155961 | -1.3989186 | -162.8 | 77.8 | 7.5 | 7.6 | 15.737 | 13.545 | 12.313 | 11.503 | 10.919 | 10.741 | 2.042 | 78.7 |
UPM 1255-0201 | 193.9661833 | -2.0191072 | -178.6 | -59.1 | 7.7 | 7.9 | 14.210 | 12.182 | 11.223 | 10.337 | 9.752 | 9.539 | 1.845 | 48.9 |
UPM 1301-2002 | 195.2599392 | -20.0495206 | -186.6 | -1.2 | 2.0 | 2.3 | 16.066 | 13.927 | 12.307 | 11.773 | 11.085 | 10.872 | 2.154 | 79.5 |
UPM 1302-1739 | 195.5729917 | -17.6554925 | 66.6 | -168.0 | 4.3 | 4.3 | $\cdots$ | 14.986 | 13.042 | 12.419 | 11.829 | 11.567 | 2.567 | 96.0 |
UPM 1303-0529 | 195.7612992 | -5.4846467 | -182.4 | 32.7 | 13.2 | 12.6 | 17.727 | 15.728 | 13.900 | 12.816 | 12.321 | 12.085 | 2.912 | 109.1 |
UPM 1305-0509 | 196.4285619 | -5.1588822 | -184.2 | -24.3 | 5.9 | 4.9 | 14.155 | 13.390 | 12.892 | 12.725 | 12.392 | 12.321 | 0.665 | 178.0 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1305-1023 | 196.3293772 | -10.3901519 | -146.7 | 104.8 | 3.9 | 4.8 | 13.943 | 11.997 | 11.040 | 10.876 | 10.238 | 10.101 | 1.121 | 68.8 |
UPM 1311-4557 | 197.8996369 | -45.9601217 | -170.7 | -71.8 | 14.0 | 7.3 | $\cdots$ | 15.818 | 13.636 | 12.716 | 12.123 | 11.900 | 3.102 | 91.1 |
UPM 1313-4112 | 198.3192592 | -41.2098522 | -20.5 | 208.9 | 11.3 | 11.2 | $\cdots$ | $\cdots$ | $\cdots$ | 11.363 | 10.788 | 10.487 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1314-0453 | 198.6887014 | -4.8900147 | 94.2 | -153.5 | 9.5 | 9.2 | 17.647 | 15.646 | 13.948 | 12.364 | 11.732 | 11.523 | 3.282 | 67.0 |
UPM 1315-0157 | 198.9190489 | -1.9517944 | -180.4 | -21.2 | 9.9 | 9.9 | 17.456 | 15.511 | 13.763 | 12.417 | 11.920 | 11.672 | 3.094 | 82.8 |
UPM 1315-2904A | 198.9894239 | -29.0780553 | -190.6 | -27.5 | 16.7 | 23.5 | 18.025 | 15.871 | 14.281 | 12.758 | 12.274 | 11.988 | 3.113 | 89.6 | ddCommon proper motion companion; see Table 4
UPM 1315-2904B | 198.9885206 | -29.0765469 | -209.7 | -1.5 | 12.8 | 11.1 | 20.271 | 17.824 | 15.971 | 14.325 | 13.751 | 13.613 | 3.499 | 149.3 | ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1318-3910 | 199.5230214 | -39.1766283 | -148.1 | -115.5 | 4.5 | 4.5 | 18.164 | 16.073 | 14.089 | 12.549 | 11.947 | 11.661 | 3.524 | 63.1 |
UPM 1323-3009 | 200.8448642 | -30.1521306 | -134.2 | -125.7 | 3.5 | 3.8 | 17.704 | 15.816 | 14.100 | 13.061 | 12.557 | 12.291 | 2.755 | 130.2 |
UPM 1323-4541 | 200.7692875 | -45.6888108 | -179.7 | 46.3 | 6.1 | 4.9 | 17.679 | 15.619 | 13.966 | 12.586 | 11.936 | 11.747 | 3.033 | 82.2 |
UPM 1324-2631 | 201.2206433 | -26.5186122 | -120.6 | -140.9 | 4.0 | 2.2 | 16.723 | 14.650 | 13.048 | 11.813 | 11.228 | 10.997 | 2.837 | 64.5 |
UPM 1327-4606 | 201.8371664 | -46.1093275 | -179.2 | -25.0 | 6.3 | 6.6 | 18.233 | 16.153 | 14.141 | 12.751 | 12.173 | 11.906 | 3.402 | 76.2 |
UPM 1329-3729 | 202.4278942 | -37.4868347 | -174.5 | -67.1 | 2.8 | 6.8 | 17.777 | 15.655 | 13.751 | 12.475 | 11.882 | 11.606 | 3.180 | 72.3 |
UPM 1331-1706 | 202.8177619 | -17.1053169 | -152.4 | -103.8 | 8.8 | 8.8 | 13.616 | 11.668 | 10.400 | 10.293 | 9.600 | 9.397 | 1.375 | 48.2 |
UPM 1335-1706 | 203.9090931 | -17.1108244 | -170.3 | 83.7 | 3.9 | 6.3 | 17.345 | 15.345 | 13.960 | 12.976 | 12.380 | 12.125 | 2.369 | 134.0 |
UPM 1337-0155 | 204.2700961 | -1.9167561 | -204.1 | -165.2 | 11.0 | 11.3 | 16.263 | 14.421 | 12.192 | 10.983 | 10.497 | 10.223 | 3.438 | 38.3 |
UPM 1338-1459 | 204.7328261 | -14.9929269 | -180.8 | 6.1 | 8.5 | 8.7 | 14.357 | 12.632 | 11.951 | 11.520 | 10.912 | 10.801 | 1.112 | 94.7 |
UPM 1343-0220 | 205.9086011 | -2.3335592 | -144.7 | 128.3 | 5.1 | 4.4 | 14.009 | 11.642 | 10.754 | 10.064 | 9.431 | 9.219 | 1.578 | 43.4 |
UPM 1343-3728 | 205.9051978 | -37.4830681 | -132.3 | -132.2 | 3.7 | 3.4 | 18.097 | 15.807 | 14.010 | 12.689 | 12.118 | 11.856 | 3.118 | 80.8 |
UPM 1344-0757 | 206.1694717 | -7.9526650 | -21.2 | -186.8 | 13.0 | 12.9 | 17.328 | 15.072 | 12.915 | 12.111 | 11.585 | 11.300 | 2.961 | 71.2 |
UPM 1346-2111B | 206.7059456 | -21.1849233 | -112.1 | -60.0 | 6.6 | 5.7 | 18.343 | 16.407 | 14.513 | 13.004 | 12.414 | 12.148 | 3.403 | 86.8 | aaProper motions suspect ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1347-0042 | 206.8132836 | -0.7141892 | -209.1 | 50.0 | 11.3 | 11.1 | 16.141 | 14.211 | 12.183 | 10.993 | 10.425 | 10.156 | 3.218 | 38.7 |
UPM 1349-4228 | 207.2552625 | -42.4784189 | -161.9 | -84.6 | 1.3 | 3.0 | 14.468 | 11.703 | $\cdots$ | 9.449 | 8.863 | 8.622 | 2.254 | 24.4 |
UPM 1349-4603 | 207.4748908 | -46.0662628 | -171.7 | -66.6 | 2.4 | 3.5 | 16.545 | 14.507 | 12.621 | 11.706 | 11.099 | 10.819 | 2.801 | 61.3 |
UPM 1350-2538 | 207.5852117 | -25.6338006 | -126.3 | -130.5 | 2.3 | 2.9 | 17.261 | 15.147 | 13.177 | 11.576 | 11.033 | 10.753 | 3.571 | 41.2 |
UPM 1355-2724 | 208.7766217 | -27.4151128 | -180.6 | -12.6 | 1.4 | 1.2 | 15.733 | 13.991 | 13.089 | 12.058 | 11.411 | 11.310 | 1.933 | 110.2 |
UPM 1355-3547 | 208.9479931 | -35.7875306 | -148.0 | -118.6 | 2.6 | 19.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.689 | 11.076 | 10.853 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1412-3518 | 213.1385153 | -35.3079467 | -177.6 | -39.0 | 2.1 | 3.8 | 17.023 | 14.997 | 13.834 | 12.895 | 12.301 | 12.085 | 2.102 | 145.3 |
UPM 1413-0615 | 213.4542447 | -6.2657081 | -205.7 | -9.9 | 18.1 | 6.5 | 14.141 | 12.262 | 11.387 | 11.476 | 10.845 | 10.750 | 0.786 | 103.0 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1413-2727 | 213.4797344 | -27.4651342 | -173.9 | -51.0 | 1.2 | 1.5 | $\cdots$ | 14.874 | 14.198 | 13.268 | 12.728 | 12.647 | 1.606 | [205.3] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1426-3807 | 216.7439033 | -38.1297039 | -173.8 | 59.9 | 3.3 | 3.3 | 17.549 | 15.605 | 13.932 | 12.591 | 12.052 | 11.855 | 3.014 | 92.7 |
UPM 1433-1006 | 218.2923342 | -10.1017839 | -178.4 | -34.4 | 3.9 | 3.7 | 17.319 | 15.500 | 14.052 | 13.008 | 12.426 | 12.204 | 2.492 | 138.0 |
UPM 1440-1216 | 220.0069028 | -12.2808661 | -95.2 | -154.7 | 2.1 | 2.1 | $\cdots$ | $\cdots$ | $\cdots$ | 10.770 | 10.118 | 9.951 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1440-2346 | 220.0484783 | -23.7699506 | -166.3 | -81.9 | 6.1 | 5.5 | 18.074 | 15.948 | 14.523 | 12.990 | 12.405 | 12.171 | 2.958 | 102.3 |
UPM 1443-1350 | 220.7918803 | -13.8418344 | -132.3 | -123.7 | 8.4 | 8.2 | 16.306 | 14.202 | 12.724 | 12.002 | 11.453 | 11.244 | 2.200 | 96.2 |
UPM 1443-3318 | 220.8361869 | -33.3092458 | -169.7 | -65.4 | 2.6 | 1.8 | $\cdots$ | 13.706 | 11.765 | 11.365 | 10.778 | 10.538 | 2.341 | 65.6 |
UPM 1444-1414 | 221.1479233 | -14.2457728 | -110.5 | -155.5 | 7.7 | 5.7 | $\cdots$ | $\cdots$ | 11.465 | 10.983 | 10.424 | 10.194 | $\cdots$ | 61.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1445-4228 | 221.4660325 | -42.4748283 | -160.6 | -90.7 | 7.4 | 7.5 | 16.627 | 14.542 | 13.149 | 12.083 | 11.525 | 11.310 | 2.459 | 88.6 |
UPM 1447-2307 | 221.8288436 | -23.1260342 | -169.5 | 68.6 | 3.9 | 3.3 | 17.332 | 15.364 | 13.932 | 12.472 | 11.945 | 11.733 | 2.892 | 91.3 |
UPM 1449-2906 | 222.4226139 | -29.1078767 | -184.3 | -32.6 | 1.9 | 1.7 | 16.131 | 13.917 | 12.176 | 11.357 | 10.700 | 10.475 | 2.560 | 55.6 |
UPM 1451-2451 | 222.9401378 | -24.8589333 | 126.6 | -131.3 | 1.6 | 1.6 | 16.739 | 14.847 | 13.100 | 11.702 | 11.110 | 10.884 | 3.145 | 55.3 |
UPM 1453-4446 | 223.4601097 | -44.7760186 | -204.1 | -63.0 | 19.9 | 25.5 | 14.117 | 12.089 | 10.891 | 10.938 | 10.262 | 10.120 | 1.151 | 69.6 | ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1454-3423 | 223.5418183 | -34.3933181 | -227.7 | -119.4 | 14.0 | 13.6 | $\cdots$ | $\cdots$ | $\cdots$ | 13.049 | 12.497 | 12.336 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1456-3806 | 224.1946236 | -38.1156314 | -174.1 | 52.7 | 3.1 | 2.1 | 13.065 | 11.263 | 10.578 | 10.002 | 9.379 | 9.216 | 1.261 | 46.1 |
UPM 1456-4218 | 224.1659378 | -42.3046658 | -220.1 | -37.0 | 7.1 | 7.8 | 15.332 | 13.396 | 11.931 | 10.866 | 10.258 | 10.035 | 2.530 | 48.5 |
UPM 1457-0555 | 224.4482581 | -5.9233250 | -115.2 | -179.5 | 9.0 | 8.7 | 17.059 | 15.107 | 13.937 | 13.015 | 12.418 | 12.184 | 2.092 | 154.0 |
UPM 1504-0235 | 226.2061903 | -2.5983772 | -22.5 | -194.4 | 9.3 | 9.5 | 16.692 | 14.908 | 14.139 | 13.546 | 12.938 | 12.775 | 1.362 | [234.4] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1506-4504 | 226.5962444 | -45.0685394 | -151.5 | -119.5 | 8.6 | 10.8 | 16.934 | 15.119 | 13.136 | 12.013 | 11.368 | 11.151 | 3.106 | 65.1 |
UPM 1509-1356 | 227.3431372 | -13.9350375 | -200.2 | 28.8 | 9.2 | 8.8 | $\cdots$ | 14.921 | 12.843 | 11.865 | 11.280 | 11.034 | 3.056 | 60.9 |
UPM 1514-0519 | 228.6674336 | -5.3324128 | -148.8 | -106.8 | 9.7 | 9.5 | 13.534 | 12.904 | 12.538 | 12.494 | 12.211 | 12.190 | 0.410 | 169.1 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1514-0712 | 228.7084150 | -7.2062508 | -168.0 | -75.5 | 6.4 | 7.2 | 13.339 | 11.338 | 10.441 | 9.950 | 9.295 | 9.178 | 1.388 | 44.3 |
UPM 1516-2731 | 229.0559556 | -27.5278219 | -131.4 | -129.1 | 4.8 | 4.7 | 17.417 | 15.647 | 13.761 | 12.139 | 11.608 | 11.364 | 3.508 | 61.2 |
UPM 1525-4622 | 231.3542089 | -46.3688022 | -181.0 | -151.9 | 6.3 | 5.8 | $\cdots$ | $\cdots$ | $\cdots$ | 12.271 | 11.915 | 11.824 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1529-0707 | 232.3161231 | -7.1222567 | -181.5 | -61.1 | 9.7 | 9.5 | 13.142 | 12.257 | 11.859 | 11.709 | 11.387 | 11.321 | 0.548 | 112.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1532-2833 | 233.0121631 | -28.5547736 | -20.5 | -132.9 | 7.1 | 6.7 | 19.359 | 17.550 | 15.606 | 13.724 | 13.168 | 12.897 | 3.826 | 104.9 | aaProper motions suspect eeNot detected during automated search but noticed by eye during the blinking process
UPM 1532-2834 | 233.0302039 | -28.5679217 | -120.7 | -24.6 | 7.7 | 7.3 | 19.707 | 17.755 | 15.689 | 13.749 | 13.193 | 12.916 | 4.006 | 94.7 | aaProper motions suspect eeNot detected during automated search but noticed by eye during the blinking process
UPM 1533-0251 | 233.4913458 | -2.8522467 | 56.7 | -179.0 | 10.2 | 10.9 | 13.565 | 12.658 | 11.989 | 11.098 | 10.676 | 10.588 | 1.560 | 77.8 |
UPM 1533-2126 | 233.3307392 | -21.4486364 | -171.6 | -68.5 | 2.2 | 2.1 | 14.692 | 12.723 | 11.510 | 11.830 | 11.200 | 10.980 | 0.893 | 97.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1534-3744 | 233.7133039 | -37.7359719 | -31.5 | -260.7 | 16.1 | 14.4 | 17.894 | 15.895 | 14.857 | 13.171 | 12.605 | 12.408 | 2.724 | 130.8 |
UPM 1535-1652 | 233.8295736 | -16.8676697 | -5.0 | -239.3 | 14.7 | 14.3 | 17.677 | 15.648 | 13.758 | 12.655 | 12.135 | 11.875 | 2.993 | 94.0 |
UPM 1536-2307 | 234.0091364 | -23.1279594 | -162.8 | 78.1 | 9.0 | 9.2 | 17.168 | 15.250 | 13.539 | 11.960 | 11.442 | 11.192 | 3.290 | 60.2 |
UPM 1542-4520 | 235.5850419 | -45.3483308 | -130.8 | -130.5 | 5.0 | 5.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.572 | 10.999 | 10.746 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1546-2553 | 236.6958903 | -25.8843933 | -159.9 | -135.9 | 8.8 | 9.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.478 | 10.909 | 10.700 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1548-2228 | 237.1036078 | -22.4715383 | -103.0 | -151.7 | 14.2 | 4.1 | 16.985 | 15.801 | 14.246 | 12.667 | 12.144 | 11.938 | 3.134 | 107.0 |
UPM 1551-0438 | 237.8945833 | -4.6492825 | -226.5 | -3.9 | 8.7 | 8.9 | 16.552 | 15.226 | 13.429 | 11.657 | 11.074 | 10.847 | 3.569 | 51.3 |
UPM 1551-1335 | 237.9823981 | -13.5938706 | -269.4 | -96.6 | 12.3 | 14.4 | 17.386 | 15.241 | 13.040 | 12.180 | 11.646 | 11.376 | 3.061 | 72.1 |
UPM 1552-1033 | 238.0100178 | -10.5601664 | -13.0 | -184.1 | 8.8 | 9.6 | 16.537 | 14.730 | 13.454 | 12.042 | 11.431 | 11.245 | 2.688 | 80.5 |
UPM 1552-1511 | 238.1625247 | -15.1866033 | -272.3 | -209.5 | 11.8 | 12.2 | 14.647 | 13.024 | 13.305 | 12.638 | 11.919 | 11.698 | 0.386 | 107.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1552-3825 | 238.2234847 | -38.4229861 | -58.0 | -177.4 | 4.8 | 4.8 | 14.200 | 12.733 | 12.214 | 11.750 | 11.166 | 11.067 | 0.983 | 92.9 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1553-0244 | 238.2634961 | -2.7434919 | -36.1 | -176.6 | 5.8 | 17.9 | 15.437 | 13.462 | 11.798 | 10.683 | 10.149 | 9.892 | 2.779 | 41.4 |
UPM 1600-0137 | 240.0560244 | -1.6199883 | 67.4 | -178.8 | 19.8 | 4.9 | $\cdots$ | $\cdots$ | $\cdots$ | 9.820 | 9.247 | 8.981 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1605-0010 | 241.3375458 | -0.1716083 | -197.4 | -64.0 | 8.3 | 8.5 | 17.541 | 15.676 | 14.006 | 12.518 | 11.979 | 11.736 | 3.158 | 82.7 |
UPM 1606-3534 | 241.5239356 | -35.5674578 | -175.1 | -51.7 | 2.4 | 2.2 | 15.487 | 13.696 | 12.562 | 10.799 | 10.175 | 9.954 | 2.897 | 40.1 |
UPM 1607-2307 | 241.8580489 | -23.1169094 | -145.2 | -133.7 | 20.8 | 5.3 | 16.383 | 14.608 | 12.686 | 11.766 | 11.127 | 10.854 | 2.842 | 64.3 |
UPM 1609-4639 | 242.2884686 | -46.6565003 | -140.9 | -196.7 | 6.6 | 6.6 | $\cdots$ | 16.743 | 15.550 | 13.013 | 12.434 | 12.241 | 3.730 | 67.3 |
UPM 1610-0227 | 242.7059622 | -2.4608150 | -9.3 | -186.0 | 8.9 | 8.2 | 17.288 | 15.364 | 14.344 | 12.667 | 12.074 | 11.838 | 2.697 | 102.0 |
UPM 1614-4033 | 243.6298703 | -40.5574253 | -128.0 | -138.4 | 5.2 | 4.3 | 17.591 | 16.426 | 14.743 | 12.402 | 11.821 | 11.543 | 4.024 | 60.5 |
UPM 1619-2602 | 244.8460106 | -26.0425961 | -133.9 | -154.6 | 4.3 | 3.1 | 15.425 | 14.502 | 14.568 | 11.879 | 11.310 | 11.109 | 2.623 | 55.3 |
UPM 1621-0031 | 245.3020250 | -0.5209803 | -115.7 | -183.5 | 7.6 | 7.7 | 14.421 | 13.350 | 12.756 | 12.074 | 11.596 | 11.532 | 1.276 | [122.1] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1623-0641 | 245.7802481 | -6.6980106 | -186.4 | 40.7 | 6.5 | 6.5 | $\cdots$ | $\cdots$ | $\cdots$ | 11.051 | 10.424 | 10.149 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1624-3133 | 246.2206158 | -31.5640650 | -161.0 | 87.9 | 9.8 | 9.8 | $\cdots$ | $\cdots$ | $\cdots$ | 12.711 | 11.997 | 11.845 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1634-1540 | 248.5736717 | -15.6766797 | -72.1 | -197.0 | 11.3 | 10.9 | $\cdots$ | $\cdots$ | $\cdots$ | 12.908 | 12.156 | 11.902 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1638-1033 | 249.6061678 | -10.5661297 | -217.3 | -41.9 | 10.2 | 9.8 | 15.959 | 14.378 | 13.123 | 11.491 | 10.891 | 10.645 | 2.887 | 58.2 |
UPM 1638-2541 | 249.7123400 | -25.6911056 | -136.2 | -139.7 | 3.6 | 3.6 | 16.105 | 13.846 | 13.824 | 12.311 | 11.583 | 11.458 | 1.535 | 105.3 |
UPM 1638-3439 | 249.6381294 | -34.6511661 | -96.1 | -156.2 | 6.4 | 14.4 | $\cdots$ | $\cdots$ | 13.568 | 12.165 | 11.622 | 11.370 | $\cdots$ | 58.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1639-2331 | 249.7811856 | -23.5252967 | -131.3 | -134.4 | 9.2 | 3.7 | 16.818 | 15.295 | 13.685 | 11.451 | 10.828 | 10.559 | 3.844 | 37.5 |
UPM 1642-2833 | 250.6688403 | -28.5619450 | -166.5 | -100.9 | 6.4 | 22.1 | $\cdots$ | $\cdots$ | $\cdots$ | 13.377 | 12.849 | 12.652 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1646-3709 | 251.7445328 | -37.1657722 | 193.7 | -2.2 | 7.7 | 7.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.056 | 10.530 | 10.293 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1647-3213 | 251.8940731 | -32.2208019 | -188.2 | -37.3 | 7.8 | 8.3 | $\cdots$ | 15.988 | 15.028 | 12.602 | 12.081 | 11.870 | 3.386 | 66.4 |
UPM 1648-3459 | 252.1200667 | -34.9967942 | 178.6 | 142.5 | 4.4 | 4.4 | $\cdots$ | 14.631 | $\cdots$ | 10.687 | 10.161 | 9.907 | 3.994 | 22.1 |
UPM 1648-3538 | 252.0849597 | -35.6427183 | -102.0 | -156.6 | 7.6 | 5.0 | 15.970 | 14.748 | $\cdots$ | 11.008 | 10.503 | 10.266 | 3.740 | 38.5 |
UPM 1648-3539 | 252.0724997 | -35.6630103 | 76.1 | -169.7 | 2.3 | 2.3 | 15.961 | 14.549 | 13.571 | 10.890 | 10.316 | 10.056 | 3.659 | 29.8 |
UPM 1650-2440 | 252.7400719 | -24.6745247 | -37.8 | -177.8 | 35.2 | 30.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.982 | 11.355 | 11.127 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1654-3105 | 253.6846164 | -31.0961000 | -32.4 | -215.9 | 7.4 | 7.2 | 15.122 | 13.553 | 11.977 | 10.072 | 9.482 | 9.237 | 3.481 | 23.8 | ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1658-2811 | 254.7136367 | -28.1977644 | 190.7 | 179.0 | 10.3 | 10.3 | 17.979 | 16.677 | 14.129 | 13.417 | 12.904 | 12.677 | 3.260 | 149.9 |
UPM 1658-3931 | 254.5610944 | -39.5308172 | -140.3 | -191.4 | 8.8 | 8.9 | $\cdots$ | 16.787 | 14.502 | 12.161 | 11.556 | 11.266 | 4.626 | 30.2 |
UPM 1700-0857 | 255.0531956 | -8.9576528 | -19.6 | -199.7 | 13.2 | 13.2 | 17.447 | 15.705 | 14.262 | 12.535 | 11.870 | 11.639 | 3.170 | 77.4 |
UPM 1700-2913 | 255.1488828 | -29.2189083 | -158.8 | 96.5 | 3.5 | 6.1 | $\cdots$ | 14.142 | 12.927 | 11.843 | 11.326 | 11.094 | 2.299 | 79.9 |
UPM 1701-0657 | 255.3763442 | -6.9526483 | 105.6 | -160.9 | 13.9 | 14.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.695 | 12.064 | 11.808 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1701-2502 | 255.4740292 | -25.0337903 | -78.8 | -190.1 | 4.6 | 4.6 | $\cdots$ | $\cdots$ | $\cdots$ | 13.016 | 12.491 | 12.276 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1701-3150 | 255.3130186 | -31.8452472 | -165.1 | -154.8 | 9.2 | 9.2 | $\cdots$ | $\cdots$ | $\cdots$ | 11.344 | 10.841 | 10.566 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1701-3345 | 255.2964308 | -33.7500906 | -157.2 | -99.6 | 5.6 | 5.6 | $\cdots$ | 14.898 | 13.449 | 11.496 | 10.928 | 10.692 | 3.402 | 38.2 |
UPM 1704-1459 | 256.2319803 | -14.9893325 | -114.1 | -148.4 | 4.4 | 4.5 | 17.336 | $\cdots$ | $\cdots$ | 12.347 | 11.791 | 11.571 | $\cdots$ | 85.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1704-3141 | 256.2325125 | -31.6892650 | -175.2 | -169.6 | 7.3 | 7.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.655 | 12.006 | 11.834 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1706-0254 | 256.6927197 | -2.9013867 | -122.8 | -153.3 | 12.3 | 11.9 | $\cdots$ | 16.016 | 15.214 | 13.721 | 13.074 | 12.860 | 2.295 | 164.7 |
UPM 1707-0345 | 256.8042544 | -3.7595239 | 6.0 | -191.4 | 7.8 | 7.8 | $\cdots$ | $\cdots$ | $\cdots$ | 10.865 | 10.218 | 10.073 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1707-1438 | 256.7945428 | -14.6436947 | -94.6 | -164.4 | 4.7 | 4.8 | 17.411 | 15.843 | 15.096 | 13.158 | 12.519 | 12.288 | 2.685 | 123.9 |
UPM 1709-1715 | 257.4969336 | -17.2529044 | -196.8 | -4.0 | 6.0 | 6.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.313 | 11.758 | 11.530 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1709-2858 | 257.4467800 | -28.9713886 | -173.0 | -126.4 | 10.8 | 10.3 | $\cdots$ | $\cdots$ | $\cdots$ | 11.234 | 10.636 | 10.270 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1710-1055 | 257.6343403 | -10.9268356 | -89.1 | -246.0 | 11.3 | 11.3 | 17.190 | 15.900 | 13.995 | 12.060 | 11.508 | 11.226 | 3.840 | 54.9 |
UPM 1711-3942 | 257.8640486 | -39.7006067 | -168.1 | -100.2 | 11.1 | 8.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.759 | 11.166 | 10.948 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1712-4432 | 258.1053478 | -44.5375533 | -149.3 | -158.6 | 7.2 | 6.7 | 14.386 | 13.041 | 14.109 | 12.028 | 11.541 | 11.404 | 1.013 | [ 33.9] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1712-4657 | 258.0571681 | -46.9649431 | -154.8 | -109.7 | 4.3 | 4.3 | $\cdots$ | $\cdots$ | $\cdots$ | 11.992 | 11.409 | 11.147 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1713-3512 | 258.3514606 | -35.2062569 | -182.8 | -132.3 | 7.5 | 7.3 | 17.994 | 16.680 | 16.283 | 12.680 | 11.851 | 11.656 | 4.000 | 40.7 |
UPM 1714-2118 | 258.5668300 | -21.3145236 | -183.9 | -129.4 | 3.1 | 3.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.710 | 11.032 | 10.878 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1718-2245A | 259.6129147 | -22.7616683 | -160.7 | -160.8 | 7.2 | 6.8 | 15.469 | 13.836 | 13.155 | 10.385 | 9.806 | 9.572 | 3.451 | 25.4 | ddCommon proper motion companion; see Table 4
UPM 1718-2245B | 259.6213031 | -22.7746183 | -161.1 | -154.8 | 10.9 | 9.6 | $\cdots$ | 14.787 | 13.289 | 10.207 | 9.608 | 9.375 | 4.580 | 13.2 | ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1720-1252 | 260.1161614 | -12.8698767 | -41.6 | -184.2 | 17.1 | 12.3 | $\cdots$ | 13.214 | $\cdots$ | 11.034 | 10.467 | 10.254 | 2.180 | 61.0 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1721-4500 | 260.3965597 | -45.0118219 | 111.2 | -181.3 | 7.0 | 6.6 | 14.978 | 12.715 | 11.539 | 10.365 | 9.776 | 9.537 | 2.350 | 38.8 |
UPM 1722-4136 | 260.5690697 | -41.6059956 | -170.4 | -72.7 | 5.3 | 20.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.924 | 11.366 | 11.112 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1724-0318A | 261.1610800 | -3.3011636 | 143.7 | -126.7 | 8.7 | 8.9 | 15.536 | 13.860 | 13.019 | 11.833 | 11.201 | 10.994 | 2.027 | 92.0 | ddCommon proper motion companion; see Table 4
UPM 1724-0318B | 261.1602147 | -3.2999100 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 15.536 | 13.860 | 13.019 | 12.925 | 12.291 | 12.042 | 0.935 | 169.5 | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 1724-2217 | 261.2276967 | -22.2935397 | -251.6 | -93.6 | 7.0 | 7.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.042 | 11.492 | 11.260 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1725-1703 | 261.4728428 | -17.0606703 | -80.5 | -163.3 | 2.8 | 2.8 | $\cdots$ | 12.446 | 12.710 | 10.927 | 10.559 | 10.513 | 1.519 | 61.0 |
UPM 1725-1749 | 261.4613711 | -17.8228856 | -163.8 | -123.4 | 6.4 | 7.3 | $\cdots$ | $\cdots$ | $\cdots$ | 10.639 | 10.250 | 10.035 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1731-1316 | 262.9469089 | -13.2742833 | 73.1 | -199.8 | 19.3 | 11.7 | 17.463 | 16.004 | 14.528 | 12.690 | 12.191 | 11.941 | 3.314 | 93.3 |
UPM 1732-1639 | 263.0863078 | -16.6663169 | -170.8 | -147.8 | 11.0 | 10.7 | $\cdots$ | 13.326 | 12.859 | 11.800 | 11.231 | 11.145 | 1.526 | 93.7 |
UPM 1733-2051 | 263.2673667 | -20.8637517 | 41.8 | -211.8 | 14.5 | 5.7 | 14.941 | 13.518 | 13.304 | 10.787 | 10.182 | 9.974 | 2.731 | 37.3 |
UPM 1737-2324 | 264.4041736 | -23.4064956 | -244.7 | -127.3 | 8.3 | 8.3 | $\cdots$ | $\cdots$ | $\cdots$ | 11.322 | 10.711 | 10.416 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1738-1917 | 264.7263142 | -19.2951633 | 21.8 | -202.1 | 4.6 | 5.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.358 | 10.832 | 10.622 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1740-4110 | 265.2443694 | -41.1801894 | -68.5 | -166.7 | 6.8 | 7.3 | $\cdots$ | 14.331 | 12.963 | 11.261 | 10.678 | 10.382 | 3.070 | 38.2 |
UPM 1741-4536 | 265.4011836 | -45.6043703 | 168.4 | 194.1 | 10.4 | 10.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.021 | 10.512 | 10.236 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1743-3957 | 265.9756108 | -39.9611264 | -191.0 | 84.3 | 8.3 | 8.0 | $\cdots$ | 15.433 | 13.811 | 12.184 | 11.692 | 11.431 | 3.249 | 61.1 |
UPM 1745-1322 | 266.4679789 | -13.3765100 | -183.1 | -35.8 | 6.2 | 6.1 | 17.039 | 15.378 | 14.462 | 12.724 | 12.077 | 11.902 | 2.654 | 111.3 |
UPM 1745-4336 | 266.3273608 | -43.6112189 | 23.4 | -220.2 | 8.6 | 7.9 | $\cdots$ | $\cdots$ | $\cdots$ | 10.374 | 9.790 | 9.503 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1746-1246 | 266.5727764 | -12.7717856 | 44.4 | -186.4 | 6.8 | 6.6 | 17.622 | 16.100 | 14.686 | 12.556 | 12.011 | 11.767 | 3.544 | 76.0 |
UPM 1749-3138 | 267.4376153 | -31.6363372 | -139.1 | -122.8 | 2.0 | 2.0 | 14.024 | 12.310 | 11.420 | 10.334 | 9.724 | 9.567 | 1.976 | 48.9 |
UPM 1749-4135 | 267.4756842 | -41.5986675 | -107.8 | -184.7 | 4.3 | 4.3 | $\cdots$ | $\cdots$ | $\cdots$ | 10.349 | 9.736 | 9.544 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1749-4313 | 267.4106522 | -43.2245928 | 64.1 | -172.6 | 2.7 | 2.7 | $\cdots$ | 13.347 | 11.682 | 11.380 | 10.832 | 10.596 | 1.967 | 75.3 |
UPM 1749-4404B | 267.4638828 | -44.0790978 | 0.0 | -204.9 | 4.3 | 3.9 | $\cdots$ | $\cdots$ | $\cdots$ | 11.688 | 11.147 | 10.870 | $\cdots$ | $\cdots$ | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect ddCommon proper motion companion; see Table 4
UPM 1750-0406 | 267.6260231 | -4.1000422 | -108.5 | -145.0 | 7.6 | 8.5 | 16.733 | 15.701 | 14.428 | 11.896 | 11.381 | 11.146 | 3.805 | 51.0 |
UPM 1750-1456 | 267.7385731 | -14.9423142 | -61.0 | -172.0 | 22.7 | 4.5 | $\cdots$ | $\cdots$ | $\cdots$ | 8.874 | 8.281 | 8.030 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1751-2404 | 267.7996450 | -24.0717064 | -114.2 | -184.3 | 15.7 | 13.6 | $\cdots$ | $\cdots$ | $\cdots$ | 11.631 | 10.728 | 10.473 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1754-3805 | 268.6284406 | -38.0912072 | 57.7 | -206.3 | 7.1 | 6.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.641 | 11.077 | 10.839 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1756-2126 | 269.2477247 | -21.4437422 | -161.3 | 92.0 | 5.7 | 7.0 | 17.869 | 16.628 | 16.074 | 12.485 | 11.797 | 11.548 | 4.143 | 38.5 |
UPM 1756-4052 | 269.1487281 | -40.8780769 | -77.5 | -196.1 | 19.7 | 19.7 | $\cdots$ | 14.372 | 13.688 | 13.299 | 12.667 | 12.607 | 1.073 | [215.7] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1757-3936 | 269.4195217 | -39.6043317 | -266.3 | -158.6 | 7.3 | 7.0 | $\cdots$ | 13.759 | $\cdots$ | 11.284 | 10.734 | 10.545 | 2.475 | 61.1 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1757-4013 | 269.2794617 | -40.2270539 | -75.0 | -194.2 | 6.5 | 6.1 | $\cdots$ | 12.349 | $\cdots$ | 11.883 | 11.302 | 11.200 | 0.466 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1757-4632B | 269.3754317 | -46.5427658 | 59.7 | -185.1 | 10.6 | 6.5 | $\cdots$ | $\cdots$ | $\cdots$ | 10.847 | 10.264 | 10.007 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 1758-4615 | 269.5915203 | -46.2631747 | -171.7 | -82.6 | 25.0 | 7.2 | $\cdots$ | $\cdots$ | $\cdots$ | 12.842 | 12.323 | 12.077 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1759-4528 | 269.7754394 | -45.4816806 | -42.4 | -188.1 | 10.4 | 10.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.429 | 10.891 | 10.601 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1800-4642 | 270.1561303 | -46.7136864 | -44.7 | -206.8 | 25.0 | 19.7 | $\cdots$ | $\cdots$ | 11.532 | 10.872 | 10.267 | 10.009 | $\cdots$ | 51.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1804-0902 | 271.0980389 | -9.0451769 | 164.6 | -158.9 | 5.8 | 6.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.081 | 11.395 | 11.162 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1805-4112 | 271.4910153 | -41.2111136 | 14.8 | -185.5 | 7.9 | 7.9 | $\cdots$ | 14.013 | 12.634 | 12.757 | 12.060 | 11.870 | 1.256 | [155.9] | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1806-1700 | 271.7052058 | -17.0047383 | -23.7 | -181.0 | 4.8 | 4.9 | $\cdots$ | $\cdots$ | $\cdots$ | 10.873 | 10.214 | 9.957 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1810-4412 | 272.7016906 | -44.2066736 | 133.6 | -139.9 | 4.8 | 3.9 | $\cdots$ | $\cdots$ | $\cdots$ | 11.616 | 11.037 | 10.793 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1811-0139 | 272.9795842 | -1.6625278 | -13.3 | -182.6 | 6.5 | 6.5 | $\cdots$ | $\cdots$ | $\cdots$ | 11.236 | 10.637 | 10.563 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1811-3907 | 272.9450564 | -39.1204408 | 198.7 | 207.8 | 5.3 | 5.4 | $\cdots$ | 15.619 | 15.210 | 13.780 | 13.230 | 12.961 | 1.839 | [180.2] | ccSuperCOSMOS plate magnitudes suspect ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1812-0445 | 273.0286425 | -4.7643792 | 239.4 | -61.1 | 7.8 | 7.8 | 18.072 | 16.984 | 15.462 | 12.316 | 11.828 | 11.509 | 4.668 | 50.4 |
UPM 1812-3958 | 273.1333542 | -39.9774706 | -125.7 | -150.8 | 8.0 | 8.2 | $\cdots$ | $\cdots$ | 12.478 | 12.721 | 12.126 | 11.969 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1817-0833 | 274.4449242 | -8.5561069 | -65.6 | -189.9 | 12.2 | 9.1 | 17.222 | 16.383 | 16.030 | 12.504 | 11.794 | 11.599 | 3.879 | 51.0 |
UPM 1818-2854 | 274.5728911 | -28.9023606 | -179.0 | 56.8 | 13.3 | 12.7 | $\cdots$ | $\cdots$ | $\cdots$ | 12.613 | 12.037 | 11.831 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1818-3727 | 274.6185211 | -37.4610811 | -161.2 | -187.0 | 11.1 | 15.9 | $\cdots$ | $\cdots$ | $\cdots$ | 11.970 | 11.454 | 11.210 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1818-3931 | 274.6280806 | -39.5268339 | 134.9 | -221.7 | 8.6 | 8.6 | $\cdots$ | $\cdots$ | $\cdots$ | 11.548 | 10.966 | 10.776 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1819-0734 | 274.8175769 | -7.5811236 | -173.8 | -209.9 | 8.7 | 8.8 | 15.950 | 15.187 | 15.621 | 11.121 | 10.363 | 10.143 | 4.066 | 33.4 |
UPM 1822-3206 | 275.5753822 | -32.1135317 | -56.7 | -174.1 | 18.7 | 16.7 | $\cdots$ | $\cdots$ | $\cdots$ | 14.030 | 13.493 | 13.280 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1822-3810 | 275.7248661 | -38.1719125 | -152.6 | -111.2 | 3.9 | 4.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.978 | 11.444 | 11.237 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1823-3237 | 275.9087906 | -32.6205678 | -206.0 | 37.0 | 7.1 | 6.9 | $\cdots$ | 12.757 | 11.493 | 10.897 | 10.321 | 10.060 | 1.860 | 58.9 |
UPM 1823-4055 | 275.8996567 | -40.9293583 | -196.9 | -178.2 | 8.0 | 8.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.567 | 11.025 | 10.817 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1824-1843 | 276.0411947 | -18.7243439 | -181.5 | -26.4 | 4.6 | 11.4 | $\cdots$ | 11.434 | $\cdots$ | 10.554 | 10.150 | 10.057 | 0.880 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1827-3955 | 276.8721647 | -39.9218242 | -45.5 | -184.0 | 7.4 | 7.4 | $\cdots$ | 12.337 | 10.815 | 10.803 | 10.208 | 9.942 | 1.534 | 61.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1828-3011 | 277.1039697 | -30.1852867 | -227.6 | 9.7 | 34.3 | 7.7 | $\cdots$ | $\cdots$ | $\cdots$ | 12.077 | 11.530 | 11.244 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1828-3826 | 277.1754025 | -38.4387294 | -107.1 | -226.0 | 7.2 | 7.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.835 | 11.212 | 10.950 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1828-3918 | 277.1296564 | -39.3128139 | 39.5 | -217.8 | 7.2 | 7.0 | 16.517 | 14.283 | 12.758 | 11.920 | 11.401 | 11.185 | 2.363 | 86.3 |
UPM 1831-1919 | 277.8215872 | -19.3332856 | -121.3 | -139.6 | 4.1 | 4.8 | 16.419 | 14.678 | 13.446 | 11.001 | 10.419 | 10.211 | 3.677 | 33.7 |
UPM 1831-2005 | 277.8505681 | -20.0885764 | 2.2 | -198.3 | 22.2 | 10.8 | $\cdots$ | $\cdots$ | $\cdots$ | 11.451 | 10.748 | 10.566 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1831-3017 | 277.9810053 | -30.2859400 | -95.5 | -162.4 | 7.4 | 7.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.022 | 11.433 | 11.198 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1832-3853 | 278.1483619 | -38.8879883 | 24.3 | -185.5 | 4.8 | 4.9 | $\cdots$ | $\cdots$ | $\cdots$ | 12.239 | 11.703 | 11.441 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1834-2656 | 278.5639625 | -26.9406214 | 169.4 | 75.0 | 11.8 | 11.8 | $\cdots$ | $\cdots$ | $\cdots$ | 13.405 | 12.820 | 12.658 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1834-3308 | 278.5452358 | -33.1453625 | 46.8 | -198.0 | 44.6 | 23.6 | 15.886 | 13.239 | 11.394 | 11.229 | 10.667 | 10.424 | 2.010 | 63.6 |
UPM 1836-0915 | 279.0383847 | -9.2573003 | -154.3 | 153.0 | 6.0 | 6.1 | 17.858 | 16.610 | 15.380 | 12.580 | 12.055 | 11.792 | 4.030 | 58.7 |
UPM 1839-1913 | 279.9685653 | -19.2177800 | 151.3 | -104.3 | 8.4 | 8.5 | $\cdots$ | $\cdots$ | $\cdots$ | 13.262 | 12.550 | 12.334 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1840-1934 | 280.1299600 | -19.5831464 | 201.2 | -89.8 | 8.8 | 8.1 | $\cdots$ | 13.733 | 12.400 | 10.264 | 9.706 | 9.450 | 3.469 | 20.7 |
UPM 1840-2334 | 280.0762169 | -23.5667811 | -128.2 | -143.8 | 6.8 | 3.0 | 12.605 | 11.900 | 11.719 | 11.549 | 11.300 | 11.228 | 0.351 | 107.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1840-2614 | 280.2257300 | -26.2364583 | -176.7 | -77.1 | 9.7 | 22.6 | 13.515 | 12.155 | 11.377 | 10.890 | 10.258 | 10.105 | 1.265 | 68.3 |
UPM 1841-1841 | 280.4988975 | -18.6836078 | 72.4 | -210.7 | 11.8 | 32.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.705 | 11.141 | 10.844 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1841-1902 | 280.4059922 | -19.0437531 | -14.8 | -252.1 | 10.0 | 10.1 | $\cdots$ | $\cdots$ | $\cdots$ | 11.847 | 11.254 | 10.958 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1841-2852 | 280.3630056 | -28.8708058 | -113.3 | -182.5 | 10.6 | 10.6 | $\cdots$ | $\cdots$ | $\cdots$ | 12.277 | 11.757 | 11.445 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1841-4049 | 280.3558444 | -40.8271589 | 110.8 | -160.4 | 6.5 | 5.9 | $\cdots$ | $\cdots$ | $\cdots$ | 13.108 | 12.509 | 12.313 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1842-2736 | 280.7361064 | -27.6091083 | 77.9 | -163.2 | 9.3 | 9.3 | 13.426 | $\cdots$ | 10.951 | 10.023 | 9.356 | 9.184 | $\cdots$ | 29.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1846-0947 | 281.5706144 | -9.7963978 | -150.4 | -115.9 | 6.7 | 7.3 | 17.167 | $\cdots$ | $\cdots$ | 12.344 | 11.764 | 11.497 | $\cdots$ | 88.7 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1848-0252 | 282.0143783 | -2.8768681 | 32.0 | -183.3 | 6.0 | 8.4 | 17.550 | 16.574 | 15.273 | 11.510 | 10.973 | 10.640 | 5.064 | 26.9 |
UPM 1850-1011 | 282.7145042 | -10.1934358 | -102.2 | -165.6 | 7.5 | 7.5 | $\cdots$ | $\cdots$ | $\cdots$ | 10.893 | 10.296 | 10.068 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1851-1431 | 282.9501331 | -14.5258975 | -164.4 | -92.0 | 4.5 | 4.5 | 15.938 | 14.662 | 13.257 | 11.551 | 11.005 | 10.746 | 3.111 | 60.0 |
UPM 1851-2232 | 282.9338092 | -22.5417239 | 209.9 | 10.0 | 7.3 | 9.5 | 16.115 | 13.559 | 12.603 | 11.955 | 11.379 | 11.194 | 1.604 | 105.8 |
UPM 1851-3840 | 282.8074922 | -38.6741419 | -161.2 | 82.0 | 6.2 | 6.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.575 | 9.983 | 9.704 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1852-2751 | 283.0031728 | -27.8602172 | 82.6 | -174.6 | 7.9 | 8.8 | $\cdots$ | $\cdots$ | $\cdots$ | 12.419 | 11.857 | 11.593 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1858-3548 | 284.6541453 | -35.8062356 | -102.9 | 178.4 | 6.7 | 6.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.275 | 10.871 | 10.760 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1859-1701 | 284.8129344 | -17.0311458 | -127.7 | -164.7 | 8.3 | 8.3 | $\cdots$ | $\cdots$ | $\cdots$ | 10.923 | 10.345 | 10.034 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1900-0511 | 285.2324742 | -5.1836367 | -90.8 | -168.3 | 7.2 | 7.0 | 13.538 | 11.997 | 11.474 | 10.426 | 9.765 | 9.570 | 1.571 | 43.5 |
UPM 1902-3731 | 285.5237094 | -37.5289453 | 163.2 | -82.3 | 12.1 | 12.2 | 18.665 | 16.574 | 15.884 | 13.833 | 13.187 | 12.935 | 2.741 | 158.4 |
UPM 1907-0221 | 286.7786636 | -2.3570300 | -23.4 | -192.6 | 7.2 | 7.1 | 13.359 | 12.590 | 12.326 | 11.287 | 10.812 | 10.740 | 1.303 | 76.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1917-2949 | 289.4523744 | -29.8228444 | -70.6 | -223.4 | 10.8 | 10.4 | $\cdots$ | $\cdots$ | $\cdots$ | 11.927 | 11.374 | 11.094 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1920-1606 | 290.0811403 | -16.1009153 | -175.4 | 68.9 | 10.9 | 11.0 | $\cdots$ | $\cdots$ | 17.198 | 16.119 | 15.579 | 15.625 | $\cdots$ | 574.5 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1920-2206 | 290.1098181 | -22.1062475 | 122.9 | -161.4 | 14.4 | 9.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.085 | 11.542 | 11.280 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1923-4026 | 290.7523894 | -40.4484822 | -154.6 | -96.7 | 7.7 | 9.5 | $\cdots$ | 16.086 | 14.192 | 12.847 | 12.337 | 12.100 | 3.239 | 88.5 |
UPM 1925-0916 | 291.4896428 | -9.2686953 | -58.8 | -197.4 | 6.4 | 6.2 | $\cdots$ | 14.429 | $\cdots$ | 12.428 | 11.800 | 11.617 | 2.001 | 121.9 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1925-3712 | 291.4461181 | -37.2007175 | 138.7 | -165.4 | 15.2 | 15.1 | $\cdots$ | $\cdots$ | 13.683 | 12.134 | 11.536 | 11.289 | $\cdots$ | 47.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1931-4001 | 292.9411467 | -40.0227744 | -3.1 | -185.0 | 1.9 | 1.9 | 14.966 | 12.936 | 11.356 | 10.714 | 10.092 | 9.861 | 2.222 | 50.2 |
UPM 1933-2916 | 293.3358858 | -29.2750011 | -90.4 | -159.5 | 2.9 | 5.1 | 15.935 | 13.822 | 11.914 | 10.767 | 10.151 | 9.882 | 3.055 | 34.6 |
UPM 1940-0508 | 295.1044458 | -5.1410269 | 22.0 | -181.3 | 8.1 | 8.0 | $\cdots$ | $\cdots$ | $\cdots$ | 11.602 | 11.074 | 10.825 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1943-0035 | 295.9131372 | -0.5899383 | -256.9 | 22.2 | 11.4 | 10.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.195 | 11.669 | 11.372 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1947-3542 | 296.7613983 | -35.7157425 | 24.0 | -180.4 | 2.5 | 2.6 | 16.879 | 14.954 | 13.555 | 12.269 | 11.659 | 11.436 | 2.685 | 86.0 |
UPM 1950-0135 | 297.6302547 | -1.5941806 | 223.2 | 2.3 | 10.2 | 10.1 | $\cdots$ | $\cdots$ | $\cdots$ | 12.355 | 11.803 | 11.575 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1952-0813 | 298.1828144 | -8.2235433 | 103.0 | -220.6 | 7.8 | 9.1 | 15.568 | 13.655 | 11.782 | 10.585 | 9.968 | 9.683 | 3.070 | 32.6 |
UPM 1953-0508 | 298.3124231 | -5.1386344 | -45.1 | -196.0 | 7.3 | 7.1 | $\cdots$ | 14.484 | 13.942 | 12.870 | 12.395 | 12.296 | 1.614 | 160.4 |
UPM 1954-0139 | 298.6503828 | -1.6506567 | -93.6 | -164.6 | 9.0 | 9.1 | 16.468 | 15.236 | 14.497 | 13.541 | 12.951 | 12.815 | 1.695 | [213.2] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 1955-3426 | 298.9125511 | -34.4337736 | -92.3 | -161.9 | 4.8 | 4.7 | 15.620 | 15.545 | $\cdots$ | 12.908 | 12.406 | 12.128 | 2.637 | 116.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ccSuperCOSMOS plate magnitudes suspect
UPM 1957-3801 | 299.4516389 | -38.0268108 | 119.4 | -160.3 | 6.6 | 17.2 | $\cdots$ | $\cdots$ | $\cdots$ | 12.317 | 11.721 | 11.521 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 1958-3348 | 299.6281386 | -33.8102706 | 122.5 | -132.4 | 5.7 | 5.0 | $\cdots$ | $\cdots$ | $\cdots$ | 12.145 | 11.624 | 11.377 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2007-0132 | 301.7682706 | -1.5366978 | 6.6 | -184.1 | 16.4 | 16.1 | $\cdots$ | 15.506 | $\cdots$ | 12.448 | 11.871 | 11.591 | 3.058 | 71.5 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2007-0312 | 301.9236269 | -3.2085131 | -159.3 | -86.5 | 7.3 | 7.2 | 16.388 | 14.715 | 12.952 | 11.235 | 10.702 | 10.462 | 3.480 | 41.6 |
UPM 2009-0041 | 302.4692369 | -0.6924347 | 115.0 | -191.9 | 14.0 | 14.7 | 17.493 | 15.536 | 13.355 | 12.277 | 11.749 | 11.497 | 3.259 | 72.2 |
UPM 2009-3305 | 302.2734844 | -33.0964475 | 60.3 | -179.0 | 22.0 | 16.9 | 17.666 | 15.906 | 14.344 | 13.039 | 12.533 | 12.329 | 2.867 | 129.1 |
UPM 2011-0002 | 302.9155206 | -0.0494872 | 82.3 | -262.2 | 10.4 | 10.7 | $\cdots$ | $\cdots$ | $\cdots$ | 11.420 | 10.819 | 10.684 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2011-0010 | 302.7694014 | -0.1825764 | -141.8 | -130.9 | 11.1 | 8.4 | $\cdots$ | $\cdots$ | 12.726 | 11.881 | 11.339 | 11.119 | $\cdots$ | 79.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2011-0139 | 302.8264014 | -1.6649031 | 129.2 | -127.5 | 15.9 | 16.0 | 17.948 | 15.746 | 13.609 | 12.433 | 11.860 | 11.584 | 3.313 | 67.8 |
UPM 2012-0133 | 303.1161578 | -1.5506017 | 236.9 | -65.7 | 14.0 | 13.5 | $\cdots$ | $\cdots$ | $\cdots$ | 11.962 | 11.370 | 11.068 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2014-0124 | 303.6403675 | -1.4061056 | 144.1 | -231.7 | 7.8 | 7.9 | $\cdots$ | $\cdots$ | 12.139 | 11.284 | 10.729 | 10.508 | $\cdots$ | 59.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2014-0624 | 303.5245169 | -6.4119417 | 3.3 | -220.4 | 17.6 | 17.2 | 17.444 | 15.620 | 14.056 | 12.989 | 12.373 | 12.223 | 2.631 | 131.8 |
UPM 2014-1634 | 303.6036331 | -16.5819139 | -167.7 | -192.2 | 12.1 | 13.4 | $\cdots$ | 13.477 | $\cdots$ | 12.772 | 12.137 | 12.020 | 0.705 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2019-0754 | 304.8205911 | -7.9011086 | 127.5 | -157.6 | 6.9 | 7.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.462 | 12.101 | 12.066 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2020-0826 | 305.0236797 | -8.4449589 | 174.3 | -48.6 | 11.2 | 8.1 | $\cdots$ | 14.934 | 12.955 | 11.903 | 11.347 | 11.099 | 3.031 | 62.7 |
UPM 2024-0638 | 306.0952239 | -6.6395264 | 75.9 | -247.6 | 10.9 | 11.1 | 16.931 | 14.871 | 13.039 | 11.878 | 11.357 | 11.085 | 2.993 | 64.4 |
UPM 2045-0612 | 311.4779600 | -6.2101506 | -143.8 | -120.3 | 10.0 | 10.1 | $\cdots$ | 15.646 | 13.667 | 12.937 | 12.409 | 12.167 | 2.709 | 122.0 |
UPM 2047-0429 | 311.8002206 | -4.4950453 | 143.0 | -117.9 | 6.4 | 6.6 | 12.644 | 11.975 | 11.578 | 11.659 | 11.348 | 11.398 | 0.316 | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2047-2232 | 311.9131269 | -22.5408683 | -136.3 | -128.9 | 9.8 | 10.0 | $\cdots$ | $\cdots$ | $\cdots$ | 13.261 | 12.733 | 12.634 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2047-3246 | 311.9777267 | -32.7689467 | -242.1 | 197.9 | 7.2 | 7.4 | $\cdots$ | $\cdots$ | $\cdots$ | 10.917 | 10.373 | 10.143 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2049-0304A | 312.4151003 | -3.0785519 | 83.1 | -185.7 | 7.0 | 7.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.889 | 12.288 | 12.072 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 2049-0304B | 312.4151678 | -3.0771406 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 14.444 | 13.943 | 13.673 | $\cdots$ | $\cdots$ | aaProper motions suspect bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4 eeNot detected during automated search but noticed by eye during the blinking process
UPM 2050-1541 | 312.5709881 | -15.6922022 | 84.2 | -178.4 | 3.2 | 4.7 | 15.656 | 13.702 | 12.131 | 11.456 | 10.862 | 10.580 | 2.246 | 70.2 |
UPM 2050-4535 | 312.5937853 | -45.5929939 | 17.8 | -207.4 | 9.5 | 9.6 | 17.265 | 15.170 | 14.302 | 13.405 | 12.863 | 12.688 | 1.765 | [214.1] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 2057-0911 | 314.4559047 | -9.1875750 | 129.1 | -168.5 | 8.7 | 8.4 | $\cdots$ | 12.472 | 11.512 | 11.501 | 10.806 | 10.656 | 0.971 | 87.5 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2058-0332 | 314.6363908 | -3.5394206 | -38.0 | -198.2 | 6.6 | 8.1 | 13.785 | 11.439 | 10.095 | 10.436 | 9.823 | 9.618 | 1.003 | 44.8 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2105-1703 | 316.4376317 | -17.0536781 | 81.7 | -161.3 | 2.1 | 2.3 | 15.955 | 14.016 | 12.396 | 11.385 | 10.749 | 10.506 | 2.631 | 57.4 |
UPM 2115-0631 | 318.8872478 | -6.5276206 | 191.1 | 28.7 | 7.6 | 8.4 | $\cdots$ | $\cdots$ | $\cdots$ | 12.633 | 12.090 | 11.891 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2118-2101 | 319.5260700 | -21.0311750 | 188.4 | 30.4 | 15.9 | 17.3 | 18.188 | 15.914 | 14.188 | 13.103 | 12.484 | 12.212 | 2.811 | 108.6 |
UPM 2131-3027 | 322.7657644 | -30.4592144 | 47.4 | -173.7 | 3.7 | 3.7 | 17.497 | 15.514 | 13.704 | 12.264 | 11.687 | 11.415 | 3.250 | 65.9 |
UPM 2140-0613 | 325.0265956 | -6.2270500 | -56.1 | -210.2 | 7.9 | 8.2 | 13.765 | 12.664 | 13.526 | 12.762 | 12.356 | 12.341 | -0.098 | 156.2 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2145-4145 | 326.4178564 | -41.7660619 | 137.1 | -125.2 | 1.7 | 2.1 | 16.569 | 13.895 | 12.271 | 11.407 | 10.885 | 10.560 | 2.488 | 55.7 | ccSuperCOSMOS plate magnitudes suspect
UPM 2152-1147 | 328.2452208 | -11.7888700 | 150.3 | -123.1 | 1.9 | 4.1 | 16.766 | 14.875 | 12.787 | 11.466 | 10.888 | 10.632 | 3.409 | 44.4 |
UPM 2154-0143 | 328.6042781 | -1.7273786 | 259.8 | -7.5 | 11.2 | 11.7 | 13.464 | 11.889 | 10.739 | 11.170 | 10.542 | 10.353 | 0.719 | 79.4 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2157-0251 | 329.4304583 | -2.8516900 | 139.8 | -132.6 | 18.1 | 17.9 | 18.051 | $\cdots$ | 14.875 | 13.480 | 12.902 | 12.672 | $\cdots$ | 152.1 |
UPM 2222-3528 | 335.5426575 | -35.4820467 | -116.0 | -155.8 | 5.9 | 5.7 | $\cdots$ | $\cdots$ | $\cdots$ | 13.313 | 12.672 | 12.461 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2229-0432B | 337.4473828 | -4.5360572 | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | $\cdots$ | 12.140 | 11.556 | 11.319 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliableeeNot detected during automated search but noticed by eye during the blinking process ddCommon proper motion companion; see Table 4 aaProper motions suspect
UPM 2231-1642 | 337.7917883 | -16.7142497 | -123.3 | -139.2 | 1.8 | 1.9 | 16.421 | 14.446 | 12.668 | 11.850 | 11.283 | 11.036 | 2.596 | 76.2 |
UPM 2232-0225 | 338.1181078 | -2.4226789 | 181.2 | -31.2 | 11.1 | 11.5 | 16.226 | 14.169 | 12.605 | 11.514 | 10.938 | 10.684 | 2.655 | 60.9 |
UPM 2233-0003 | 338.2930958 | -0.0525181 | 247.8 | 99.7 | 8.7 | 8.8 | $\cdots$ | 15.656 | $\cdots$ | 12.763 | 12.209 | 11.978 | 2.893 | 94.6 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2237-0118 | 339.3505139 | -1.3070036 | -147.7 | -110.1 | 7.0 | 7.0 | 14.926 | 13.000 | 11.919 | 11.674 | 11.029 | 10.877 | 1.326 | 97.5 |
UPM 2244-4333 | 341.0835264 | -43.5574258 | 183.0 | -51.4 | 14.8 | 15.1 | 17.708 | 15.718 | 14.089 | 13.114 | 12.529 | 12.353 | 2.604 | 138.3 |
UPM 2246-0017 | 341.6494342 | -0.2863092 | 193.6 | -28.5 | 10.3 | 10.3 | 17.545 | 15.654 | 13.935 | 12.665 | 12.130 | 11.873 | 2.989 | 95.3 |
UPM 2248-3255 | 342.1992344 | -32.9243678 | -126.6 | -174.9 | 3.8 | 3.9 | $\cdots$ | $\cdots$ | $\cdots$ | 12.662 | 12.120 | 11.932 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2250-2908A | 342.5898331 | -29.1365486 | 186.9 | -9.2 | 5.0 | 3.0 | 15.217 | 13.358 | 11.795 | 12.126 | 11.517 | 11.293 | 1.232 | 115.3 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable ddCommon proper motion companion; see Table 4
UPM 2250-2908B | 342.5883269 | -29.1372669 | 182.3 | -17.0 | 3.7 | 3.9 | 15.217 | 13.358 | 11.795 | 13.184 | 12.599 | 12.327 | 0.174 | 155.7 | aaProper motions suspect ddCommon proper motion companion; see Table 4 bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable eeNot detected during automated search but noticed by eye during the blinking process
UPM 2251-0201 | 342.7661222 | -2.0260436 | 50.1 | -185.5 | 7.6 | 7.6 | 15.605 | 13.598 | 12.627 | 11.995 | 11.372 | 11.176 | 1.603 | 110.8 |
UPM 2305-4612 | 346.2905092 | -46.2043653 | 143.5 | -110.4 | 7.8 | 2.5 | 16.914 | 15.068 | 13.649 | 12.773 | 12.240 | 11.993 | 2.295 | 136.2 |
UPM 2306-0315 | 346.6112989 | -3.2637836 | 181.1 | -55.6 | 10.2 | 10.2 | 17.564 | 15.549 | 13.841 | 12.506 | 11.947 | 11.661 | 3.043 | 80.8 |
UPM 2308-1954 | 347.1959358 | -19.9085208 | 183.6 | -0.3 | 13.0 | 12.1 | 13.693 | 12.059 | 11.387 | 11.479 | 11.092 | 11.006 | 0.580 | 131.5 | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable aaProper motions suspect
UPM 2310-0410 | 347.5503997 | -4.1822806 | 155.6 | -93.0 | 11.0 | 11.1 | $\cdots$ | 15.438 | 14.685 | 13.920 | 13.313 | 13.160 | 1.518 | [263.6] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 2320-0006 | 350.0481022 | -0.1140686 | 78.7 | -165.2 | 9.3 | 2.5 | 13.636 | 11.776 | 10.884 | 10.299 | 9.637 | 9.479 | 1.477 | 50.0 |
UPM 2325-1715 | 351.4973083 | -17.2589567 | 180.0 | 254.5 | 11.5 | 10.5 | 16.160 | 13.342 | 11.790 | 11.113 | 10.486 | 10.218 | 2.229 | 50.7 |
UPM 2328-0546 | 352.0857586 | -5.7781261 | 183.3 | 62.8 | 6.6 | 7.2 | $\cdots$ | $\cdots$ | $\cdots$ | 10.995 | 10.415 | 10.162 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2329-0510 | 352.4821203 | -5.1685144 | 159.2 | 88.3 | 12.7 | 9.7 | 17.693 | 15.621 | 13.802 | 12.583 | 12.002 | 11.781 | 3.038 | 85.5 |
UPM 2331-0233 | 352.8670611 | -2.5567644 | -189.0 | -114.3 | 8.6 | 11.3 | 17.038 | 14.938 | 13.029 | 11.864 | 11.303 | 11.000 | 3.074 | 58.3 |
UPM 2331-0617 | 352.7987228 | -6.2871728 | 180.6 | 36.1 | 10.3 | 10.3 | $\cdots$ | $\cdots$ | $\cdots$ | 12.722 | 12.079 | 11.904 | $\cdots$ | $\cdots$ | bbNumber of relations used for distance estimate $<$ 6: plate distance less reliable
UPM 2334-4145 | 353.6053231 | -41.7568239 | -20.3 | -187.7 | 2.4 | 3.2 | 17.676 | 15.973 | 15.112 | 14.059 | 13.421 | 13.279 | 1.914 | [273.8] | ffSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
UPM 2349-1023 | 357.2962428 | -10.3857086 | 169.5 | 71.8 | 13.1 | 14.1 | 17.811 | 15.720 | 13.845 | 12.954 | 12.352 | 12.065 | 2.766 | 109.4 |
UPM 2356-0453 | 359.1178706 | -4.8861358 | 189.6 | -44.4 | 19.9 | 18.6 | 18.052 | 16.087 | 14.310 | 12.885 | 12.374 | 12.113 | 3.202 | 95.4 |
Table 3: New UCAC3 High Proper Motion Systems estimated to be within 25 pc between Declinations $-$47$\arcdeg$ and 0$\arcdeg$ with 0$\farcs$40 yr-1 $>$ $\mu$ $\geq$ 0$\farcs$18 yr-1 Name | RA J2000.0 | DEC J2000.0 | $\mu_{\alpha}\cos\delta$ | $\mu_{\delta}$ | sig $\mu_{\alpha}$ | sig $\mu_{\delta}$ | $B_{J}$ | $R_{59F}$ | $I_{IVN}$ | $J$ | $H$ | $K_{s}$ | $R_{59F}$ $-$ $J$ | Est Dist | Notes
---|---|---|---|---|---|---|---|---|---|---|---|---|---|---|---
| (deg) | (deg) | (mas/yr) | (mas/yr) | (mas/yr) | (mas/yr) | | | | | | | | (pc) |
UPM 1349-4228 | 207.2552625 | -42.4784189 | -161.9 | -84.6 | 1.3 | 3.0 | 14.468 | 11.703 | $\cdots$ | 9.449 | 8.863 | 8.622 | 2.254 | 24.4 |
UPM 1648-3459 | 252.1200667 | -34.9967942 | 178.6 | 142.5 | 4.4 | 4.4 | $\cdots$ | 14.631 | $\cdots$ | 10.687 | 10.161 | 9.907 | 3.944 | 22.1 |
UPM 1654-3105 | 253.6846164 | -31.0961000 | -32.4 | -215.9 | 7.4 | 7.2 | 15.122 | 13.553 | 11.977 | 10.072 | 9.482 | 9.237 | 3.481 | 23.8 | aaPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1718-2245A | 259.6129147 | -22.7616683 | -160.7 | -160.8 | 7.2 | 6.8 | 15.469 | 13.836 | 13.155 | 10.385 | 9.806 | 9.572 | 3.451 | 25.4 | bbCommon proper motion companion; see Table 4
UPM 1718-2245B | 259.6213031 | -22.7746183 | -161.1 | -154.8 | 10.9 | 9.6 | $\cdots$ | 14.787 | 13.289 | 10.207 | 9.608 | 9.375 | 4.580 | 13.2 | bbCommon proper motion companion; see Table 4 ccNot detected during automated search but noticed by eye during the blinking process
UPM 1840-1934 | 280.1299600 | -19.5831464 | 201.2 | -89.8 | 8.8 | 8.1 | $\cdots$ | 13.733 | 12.400 | 10.264 | 9.706 | 9.450 | 3.469 | 20.7 |
Table 4: Common Proper Motion Candidate Systems Primary | $\mu_{\alpha}\cos\delta$ | $\mu_{\delta}$ | Distance | Secondary/Tertiary | $\mu_{\alpha}\cos\delta$ | $\mu_{\delta}$ | Distance | Separation | $\theta$ | notes
---|---|---|---|---|---|---|---|---|---|---
| (mas/yr) | (mas/yr) | (pc) | | (mas/yr) | (mas/yr) | (pc) | ($\arcsec$) | ($\arcdeg$) |
UPM 0209-3339A | -86.1 | -166.9 | 49.5 | UPM 0209-3339B | -112.9 | -170.2 | $\cdots$ | 11.6 | 78.1 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 0443-4129A | 186.1 | 4.3 | 39.3 | 2MASS J04430760-4128575B | -107.5 | -53.1 | $\cdots$ | 6.8 | 339.2 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
UPM 0528-4313A | -75.6 | 164.7 | 70.4 | UPM 0528-4313B | -86.3 | 163.2 | 109.1 | 42.1 | 209.0 | aaNot detected during automated search but noticed by eye during the blinking process
UPM 0659-0052A | -58.3 | -184.1 | 78.2 | UPM 0659-0052B | $\cdots$ | $\cdots$ | $\cdots$ | 13.8 | 151.6 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
UPM 0704-0602A | $\cdots$ | $\cdots$ | 123.4 | UPM 0704-0602B | 99.5 | -153.0 | 37.8 | 12.2 | 359.1 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect ddSuperCOSMOS plate magnitudes suspect
UPM 0747-2537A | -148.5 | 101.9 | 40.6 | UPM 0747-2537B | -151.3 | 102.3 | 47.3 | 12.0 | 237.4 | bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 0800-0617A | 135.2 | -233.8 | [175.5] | UPM 0800-0617B | $\cdots$ | $\cdots$ | $\cdots$ | 5.8 | 297.2 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect eeSubdwarf candidate selected from RPM diagram; plate distance [in bracket] is incorrect
BD-04 2807A | -142.1 | -37.1 | 19.5 | UPM 1009-0501B | -190.5 | 92.1 | $\cdots$ | 20.9 | 338.5 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 1020-0633A | -179.8 | -27.8 | 34.8 | SCR 1020-0634B | -181.5 | -24.2 | 37.5 | 87.3 | 157.2 |
UPM 1031-0024A | -207.4 | -105.6 | 55.1 | UPM 1031-0024B | -142.5 | -96.9 | $\cdots$ | 7.4 | 91.4 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ffSource not in 2MASS
UPM 1056-0542A | -98.1 | -173.8 | 76.5 | UPM 1056-0542B | -63.9 | -173.3 | $\cdots$ | 9.0 | 86.2 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 1142-2055A | -186.7 | 44.2 | 41.2 | UPM 1142-2055B | $\cdots$ | $\cdots$ | $\cdots$ | 8.3 | 167.6 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
NLTT 28641A | -201.8 | 6.2 | $\cdots$ | UPM 1149-0019B | -201.5 | 2.2 | $\cdots$ | 27.3 | 128.2 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 1159-3623A | -182.1 | -101.4 | 113.1 | UPM 1159-3623B | -172.6 | -92.4 | 132.4 | 13.6 | 303.7 | aaNot detected during automated search but noticed by eye during the blinking process
UPM 1226-2020A | -137.6 | -119.8 | 72.3 | UPM 1226-2020B | -146.3 | -117.3 | $\cdots$ | 7.0 | 333.3 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
SCR 1226-3515A | -192.3 | 41.0 | 56.5 | UPM 1226-3516B | -200.5 | 38.3 | 127.5 | 49.8 | 191.3 | aaNot detected during automated search but noticed by eye during the blinking process ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
| | | | UPM 1226-3516C | -115.2 | 8.5 | 243.3 | 97.0 | 146.9 | aaNot detected during automated search but noticed by eye during the blinking process ccProper motions suspect ggPossible NLTT star with a position difference $>$ 90$\arcsec$ when compared to UCAC3 position
UPM 1315-2904A | -190.6 | -27.5 | 89.6 | UPM 1315-2904B | -209.7 | -1.5 | 149.3 | 5.9 | 332.6 | aaNot detected during automated search but noticed by eye during the blinking process
2MASS J13465039-2112266A | -186.8 | -44.0 | 46.8 | UPM 1346-2111B | -112.1 | -60.0 | 86.8 | 82.2 | 350.8 | aaNot detected during automated search but noticed by eye during the blinking process ccProper motions suspect
UPM 1718-2245A | -160.7 | -160.8 | 25.4 | UPM 1718-2245B | -161.1 | -154.8 | 13.2 | 54.3 | 149.2 | aaNot detected during automated search but noticed by eye during the blinking process
UPM 1724-0318A | 143.7 | -126.7 | 92.0 | UPM 1724-0318B | $\cdots$ | $\cdots$ | 169.5 | 5.5 | 325.4 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
TYC 7897 997 1A | -10.4 | -160.2 | 40.3 | UPM 1749-4404B | 0.0 | -204.9 | $\cdots$ | 19.8 | 254.0 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect ddSuperCOSMOS plate magnitudes suspect
TYC 8344 154 1A | 51.6 | -179.5 | $\cdots$ | UPM 1757-4632B | 59.7 | -185.1 | $\cdots$ | 30.7 | 296.9 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable
UPM 2049-0304A | 83.1 | -185.7 | $\cdots$ | UPM 2049-0304B | $\cdots$ | $\cdots$ | $\cdots$ | 5.3 | 1.5 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
2MASS J22294694-0432036A | 160.9 | -86.5 | 50.8 | UPM 2229-0432B | $\cdots$ | $\cdots$ | $\cdots$ | 8.7 | 134.7 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
UPM 2250-2908A | 186.9 | -9.2 | 115.3 | UPM 2250-2908B | 182.3 | -17.0 | 155.7 | 5.4 | 241.4 | aaNot detected during automated search but noticed by eye during the blinking process bbNumber of relations used for distance estimate $<$ 6; plate distance less reliable ccProper motions suspect
|
arxiv-papers
| 2011-11-29T13:19:39 |
2024-09-04T02:49:24.729963
|
{
"license": "Public Domain",
"authors": "C. T. Finch, N. Zacharias, M. R. Boyd, T. J. Henry, N. C. Hambly",
"submitter": "Charlie Finch",
"url": "https://arxiv.org/abs/1111.6805"
}
|
1111.6829
|
# Heavy Flavour Results at the LHC
P. Koppenburg Nikhef, Amsterdam, The Netherlands On behalf of the LHCb
collaboration, including Atlas, CMS and Alice results.
###### Abstract
We present a brief overview of the first flavour physics results at the LHC.
Cross-section for charm and beauty production have been measured by several
experiments and the first competitive results on $D$ and $B$ decays are
presented.
## I Introduction
Precise measurements of CP violation and searches for rare decays have a high
potential for discovering effects of New Physics. They are a specific task of
the LHCb experiment and are complementary to direct searches performed by
general purpose experiments. CP violation and rare decays are sensitive to new
particles and couplings in an indirect way via interferences with Standard
Model (SM) processes. This not only probes a potentially higher mass scale
than direct searches for new particles, but also gives access to amplitudes
and phases of the new couplings. The $b$ and $c$ quark decays are the best
laboratory for this programme.
## II The LHC as a the new Flavour Factory
After a very successful decade dominated by the $B$ factories, Belle and
Babar, the LHC is taking over as the new flavour factory. While PEP-2 and
KEK-B have produced around $10^{9}$ $b\bar{b}$ pairs during their lifetime,
the LHC has produced close to $4\cdot 10^{12}$ just in 2011 thanks to the very
large $b$ cross-section in high energy proton collisions. Of course the major
challenge at the LHC is to efficiently collect the most interesting among all
these events.
### II.1 The LHCb Experiment
The LHCb experiment [1] is dedicated to precision measurements of CP violation
and rare decays in beauty and charm decays. Its forward geometry covering the
range $2<\eta<5$ exploits the dominant heavy flavour production mechanism at
the LHC and covers about 40% of the differential cross-section.
Among the features unique to LHCb are its high precision vertex detector,
which is retracted away from the beam at injection and moved as close as
$8\>$mm from the beam during data taking. It is followed by a tracking system
located in a dipole magnet of which the polarity can be reversed, allowing to
cancel detector asymmetries in CP-violation measurements. A system of two Ring
Imaging Cherenkov detectors (RICH) allows a very good pion/kaon/proton
separation over the momentum range $1$–$100\>{\rm GeV}/c$.
Due to the very large total cross section an effective on-line event selection
is required where the rate is reduced by a first level hardware trigger and
further by two levels of software triggers [2]. LHCb uses hadrons, muons,
electrons and photons throughout the trigger chain, thus maximising the
trigger efficiency on all heavy quark decays.
In order not to saturate the trigger and keep the event multiplicity low, the
LHCb experiment is not operating at the maximum LHC luminosity. During most of
2011 LHCb has kept their luminosity at the constant value of $3.5\cdot
10^{32}\>\rm cm^{-2}s^{-1}$ by displacing the proton beams laterally in real
time.
### II.2 Flavour Physics at Atlas, CMS and Alice
The Atlas [3] and CMS [4] detectors are multi-purpose central detectors
optimised for searches of heavy objects. At high luminosity, their potential
for flavour physics is limited by their triggering capabilities and focus
mainly on $b$ and charmonium decays involving dimuons. But at the lower
luminosities at which the LHC operated during 2010 and part of 2011, a more
open trigger allowed an interesting favour physics programme that is
complementary to LHCb’s in particular for cross-section measurements in the
central region.
The Alice [5] detector is optimised for heavy ion collisions and, while
covering mostly the central region, has similar particle ID and tracking
capabilities as LHCb. They are thus a key player for cross-section
measurements (which are often a necessary normalisation or their QGP
programme) but are not competitive for CP violation and rare decays searches
due to the lower luminosity at which they operate.
### II.3 Data Samples
In 2011, Atlas and CMS have been delivered around $5.7\>\rm fb^{-1}$ each,
LHCb $1.2\>\rm fb^{-1}$ and Alice $5\>\rm pb^{-1}$, of which more than 90%
have been recorded and are useful for physics. The measurements reported below
use only a fraction of these data, in most cases the $1.1\>\rm fb^{-1}$
(Atlas, CMS) and $370\>\rm pb^{-1}$ (LHCb) collected until end of June 2011.
Most cross-section measurements use the lower-luminosity data sample of about
$40\>\rm pb^{-1}$ collected during 2010.
## III Heavy Flavour Production
The four LHC experiments offer a vast coverage of rapidity: Atlas and CMS
cover the central region up to $2.5$, Alice $0$–$1$ and $2.5$–$4$ and LHCb
$2$–$5.5$. A combination of differential cross-section measurements would
allow to cover the range $0<|\eta|<5.5$, but in most cases such combinations
have not yet been performed. Yet, many measurements from the various
experiments are available and give a good picture of heavy flavour production
in $pp$ collisions at $\sqrt{s}=7\>\rm TeV$.
### III.1 Charmonium
The prompt $J/\psi$ production has been measured by all four experiments [6,
7, 8, 9] with 2010 data in bins of rapidity and transverse momentum (See Fig.
1) and compared with theoretical models. No large discrepancies are seen with
the present level of uncertainties. The unknown polarisation is the main
uncertainty in all measurements, and more data is needed to be able to resolve
it. LHCb also reports a cross-section for double $J/\psi$ production, which is
very sensitive to the production mechanism [10].
Figure 1: Differential $J/\psi$ cross section measurement at CMS [6].
### III.2 Open Charm
The open charm cross-section has been measured using $D^{0}$, $D^{+}$, $D_{s}$
and $D^{\ast+}$ modes by Atlas and LHCb [11, 12] using the first few $\rm
nb^{-1}$ delivered by the LHC in 2010. The low luminosity during this period
allowed to profit from a very open trigger which helps keeping systematic
errors low. The unexpected high total $c\bar{c}$ cross section of $6.10\pm
0.93\rm\>mb$, extrapolated from these measurements is very encouraging for
charm physics in $7\>$TeV collisions. This is about 10% of the total inelastic
cross-section.
### III.3 Beauty
The inclusive production of beauty and charm hadrons in pp collisions has been
measured by LHCb. In particular using semi-leptonic decays $b\to
D^{0}(K\pi)\mu\bar{\nu}X$ [13] the cross section $\sigma(pp\to
b\bar{b}X)=284\pm 20\pm 49\mu b$ is obtained [14], extrapolating to the full
phase space. All species of beauty hadrons can be produced in pp collisions,
including $b$ baryons [15, 16] and $B_{c}^{+}$ [17].
The knowledge of the relative fractions of the various $b$ hadron species is
of crucial importance for all measurements of branching rations, most
prominently for $B_{s}\to\mu\mu$. LHCb have measured the $B_{s}$ to $B_{d}$
production ratio using semileptonic $B$ meson decays [18] and $SU(3)$ partner
decays $B_{d}\to DK$ and $B_{s}\to D_{s}\pi$ [19]. Both measurements are
consistent and get an average of
$f_{s}/f_{d}=0.267{\>}^{+{\>}0.021}_{-{\>}0.020}$ [20].
Figure 2: Invariant mass relative to threshold of $B^{0}\pi$ system
($m_{B^{0}\pi}-m_{B}-m_{\pi}$, top) and normalised fit residuals (bottom)
[21].
LHCb also studies orbitally excited $B$ mesons, notably observing for the
first time states decaying to $B^{0}\pi^{+}$ (Fig 2) [21].
## IV Flavour Physics
The LHCb experiment has been optimised for flavour physics at the LHC and
therefore all the results presented in this Section have been obtained by this
experiment, with the notable exception of the $B_{s}\to\mu\mu$ result from
CMS.
### IV.1 Charm Mixing
The most interesting topic of charm physics is the characterisation of neutral
$D$ meson mixing and the hunt for $CP$ violation in $D$ meson decays. As for
any other neutral long-lived meson (e.g. $K^{0}$, $B^{0}$, $B_{s}^{0}$), the
neutral $D$ system can be described in terms of two flavour eigenstates
$D^{0}$, $\overline{D}^{0}$ or two mass eigenstates:
$D_{1,2}=p\left|D^{0}\right>\pm q\left|\overline{D}^{0}\right>$
of masses $m_{1,2}$ and decay widths $\Gamma_{1,2}$. This allows to define the
quantities
$x=\frac{m_{2}-m_{1}}{2\Gamma}\qquad\text{and}\qquad
y=\frac{\Gamma_{2}-\Gamma_{1}}{2\Gamma}.$
The HFAG averages for these quantities [22] differ from the no mixing
hypothesis $x=0,y=0$ by $10.2\sigma$ but no single measurement excludes this
hypothesis at $5\sigma$.
Using two-body $D$ decays selected in $26\>\rm pb^{-1}$ of 2010 data the LHCb
experiment measures a linear combination of these quantities as [23]
$\displaystyle y_{CP}$ $\displaystyle=$
$\displaystyle\frac{\hat{\Gamma}(D^{0}\to K^{-}K^{+})}{\hat{\Gamma}(D^{0}\to
K^{-}\pi^{+})}-1$ $\displaystyle=$ $\displaystyle
y\cos\phi-x\sin\phi\left(\frac{A_{m}}{2}+A_{\text{prod}}\right)$
$\displaystyle=$ $\displaystyle\left(-0.55\pm 0.63\pm 0.41\right)\%,$
where $1+A_{m}=|q/p|$ and the production asymmetry $A_{\text{prod}}$ is
measured to be very small. In the limit of vanishing $CP$ violation
$y_{CP}=y$. Using the same data sample, LHCb also measure the lifetime
difference [24]
$\displaystyle A_{\Gamma}$ $\displaystyle=$
$\displaystyle\frac{\tau(\overline{D}^{0}\to K^{+}{}K^{-})-\tau(D^{0}\to
K^{+}{}K^{-})}{\tau(\overline{D}^{0}\to K^{+}{}K^{-})+\tau(D^{0}\to
K^{+}{}K^{-})}$ $\displaystyle=$ $\displaystyle\left(-0.59\pm 0.59\pm
0.21\right).$
Figure 3: $D^{0}$ impact parameter (left) and proper time [24]
The key to the lifetime measurement is a good separation of prompt and
secondary charm (from $b$ decays), illustrated in Fig 3.
### IV.2 $CP$ Violation in Charm
$CP$ violation is expected to be vanishingly small in the charm sector in the
Standard Model. A non-zero $CP$ asymmetry above a few per-mille in $D^{0}\to
h^{+}h^{-}$ ($h=\pi,K$) decays would be strong sign of new physics.
Experimentally the flavour of the $D$ meson is tagged using the decay
$D^{\ast}\to D^{0}\pi^{+}$. The raw $CP$ asymmetry of tagged $D^{0}\to f$ and
$\overline{D}^{0}\to f$ can be factorised as
$A_{\text{RAW}}(f)=A_{CP}(f)+A_{\text{D}}(f)+A_{\text{D}}(\pi_{s})+A_{\text{P}}(D^{\ast})$
where $A_{D}$ are detector asymmetries related to the final state $f$ and the
bachelor pion $\pi_{s}$ and $A_{\text{P}}(D^{\ast})$ is the production
asymmetry in $pp$ collisions. The detection asymmetry for $f$ vanishes when
one uses decays to $CP$ eigenstates, e.g. $\pi^{+}\pi^{-}$ or $K^{+}K^{-}$.
All other asymmetries can be cancelled at first order by measuring the
difference of the two $CP$ asymmetries in these two channels. While writing
these proceedings the following interesting result has become available. Using
$580\>\rm pb^{-1}$ of 2011 data the LHCb collaboration gets very clean samples
of $1.44\cdot 10^{6}$ tagged $D\to K^{+}K^{-}$ and $0.38\cdot 10^{6}$
$D\to\pi^{+}\pi^{-}$ (Fig 4) [25].
Figure 4: Fits to the mass difference spectrum in tagged $D\to K^{+}K^{-}$
(left) and $D\to\pi^{+}\pi^{-}$ (right) decays [25]
Due to the different lifetime acceptance of the two channels a small
contribution from mixing induced $CP$ violation does not cancel out in the
measurement but its magnitude can be extracted from data:
$\displaystyle\Delta A_{CP}$ $\displaystyle\equiv$ $\displaystyle A_{CP}^{\rm
raw}(K^{+}K^{-})-A_{CP}^{\rm raw}(\pi^{+}\pi^{-})$ $\displaystyle=$
$\displaystyle A_{CP}^{\rm dir}(K^{+}K^{-})-A_{CP}^{\rm
dir}(\pi^{+}\pi^{-})+0.098\>A_{CP}^{\rm ind}$ $\displaystyle=$
$\displaystyle\left(\\-0.82\pm 0.21\pm 0.11\right)\%.$
The measured difference of $CP$ asymmetries the first ($3.5\sigma$) evidence
of $CP$ violation in the charm sector.
### IV.3 Rare $b$ Decays
The SM prediction for the Branching Ratios (BR) of the decays
$B_{q}\to\mu^{+}\mu^{-}$ have been computed to be ${\rm
BR}(B_{s}\to\mu^{+}\mu^{-})=(3.2\pm 0.2)\cdot 10^{-9}$ and ${\rm
BR}(B_{d}\to\mu^{+}\mu^{-})=(0.10\pm 0.01)\cdot 10^{-9}$ [26]. However, many
extensions of the SM predict large enhancements to these BR. The first search
for this at the LHC decay was reported by the LHCb collaboration with 2010
data [27]. Recently LHCb and CMS collaborations have presented new searches
based on $0.3$ and $1.1\>\rm fb^{-1}$ samples collected in 2011, respectively
[28, 29].
Figure 5: $B_{s}\to\mu\mu$ search windows at LHCb [28] (top) and CMS [29]
(bottom).
While CMS uses a cut-based approach, LHCb uses a boosted decision tree
calibrated on $B\to hh$ ($h=\pi,K$) decays which have the same topology as the
signal. The estimated yield is then normalised using $B_{d}\to K\pi$ (LHCb
only), $B^{+}\to J/\psi K$ and $B_{s}\to J/\psi\phi$. No excess of signal is
observed at neither of the two experiments (Fig. 5) and upper limits are set
by LHCb as ${\rm BR}(B_{s}\to\mu^{+}\mu^{-})<1.6\cdot 10^{-8}$ (95% C.L.) and
${\rm BR}(B_{d}\to\mu^{+}\mu^{-})<5.1\cdot 10^{-9}$, and CMS as ${\rm
BR}(B_{s}\to\mu^{+}\mu^{-})<1.8\cdot 10^{-8}$. The combined LHC result is
${\rm BR}(B_{s}\to\mu^{+}\mu^{-})<1.1\cdot 10^{-8}$ [30] thus not confirming
the excess reported by CDF [31]. The limits set by the LHC strongly constrain
the allowed SUSY parameter space, especially at large $\tan\beta$ [32].
The rare decay $B_{d}\to\mu\mu K^{\ast}$ is a $b\to s$ flavour changing
neutral current decay which is in the SM mediated by electroweak box and
penguin diagrams. It can be a highly sensitive probe for new right handed
currents and new scalar and pseudoscalar couplings. These New Physics
contributions can be probed by its contribution to the angular distributions
of the $B^{0}$ daughter particles. The most prominent observable is the
forward-backward asymmetry of the muon system ($A_{\rm FB}$). $A_{\rm FB}$
varies with the invariant mass-squared of the dimuon pair ($q^{2}$) and in the
SM changes sign at a well defined point, where the leading hadronic
uncertainties cancel. In many NP models the shape of $A_{\rm FB}$ as a
function of $q^{2}$ can be dramatically altered.
Figure 6: $B_{d}\to\mu\mu K^{\ast}$ forward-backward asymmetry $A_{\rm FB}$
(top) and $K^{\ast}$ polarisation fraction $F_{L}$ in bins of dimuon mass
$q^{2}$ [33].
The latest LHCb analysis [33] uses $309\rm\>pb^{-1}$ of data collected during
2011 to measure $A_{\rm FB}$, the fraction of longitudinal polarisation of the
$K^{\ast}$, $F_{L}$, and the differential branching fraction, $dB/dq^{2}$, as
a function of the dimuon invariant mass squared, $q^{2}$. There is good
agreement between recent Standard Model predictions and the LHCb measurement
of $A_{\rm FB}$, $F_{L}$ and $dB=dq^{2}$ in the six $q^{2}$ bins (Fig. 6). In
a $1<q^{2}<6\rm\>GeV^{2}$ bin, LHCb measures $A_{\rm FB}=0.10\pm 0.14\pm
0.05$, to be compared with theoretical predictions of $A_{\rm FB}=0.04\pm
0.03$. The experimental uncertainties are presently statistically dominated,
and will improve with a larger data set. Such a data set would also enable
LHCb to explore a wide range of new observables.
Using a very similar selection, LHCb also searched for Majorana Neutrinos [34]
giving raise to $B^{+}\to K^{-}\mu^{+}\mu^{+}$ and
$B^{+}\to\pi^{-}\mu^{+}\mu^{+}$ decays. No excess was found and 95% C.L.
limits have been set at $5.4\cdot 10^{-8}$ and $5.8\cdot 10^{-8}$,
respectively.
Figure 7: $B\to K^{\ast}\gamma$ and $B_{s}\to\phi\gamma$ mass peaks [35].
Other $b\to s$ transitions of interest are radiative decays $B\to
K^{\ast}\gamma$ and $B_{s}\to\phi\gamma$. LHCb reports the first measurement
of the ratio of branching fractions of these two decays as $1.52\pm 0.15\pm
0.10\pm 0.12$ where the last error comes from ($f_{d}/f_{s}$) [35]. The mass
resolution (Fig 7) is dominated by the photon energy resolution. LHCb already
has the largest sample of $B_{s}\to\phi\gamma$, which will become to measure
or constrain non-standard right-handed currents.
### IV.4 CP violation in $B$ decays
Decays of neutral $B$ mesons provide a unique laboratory to study CP-violation
originating from a non-trivial complex phase in the CKM matrix. The relative
phase between the direct decay amplitude and the amplitude of decay via mixing
gives rise to time-dependent CP-violation, a difference in the proper decay
time distribution of $B$-meson and anti-$B$-meson decays. The decay $B_{s}\to
J/\psi\phi$ is considered the golden modes for measuring this type of CP-
violation, In the Standard Model the CP-violating phase in this decay is
predicted to be $\phi_{s}\simeq-2\beta_{s}$ where
$\beta_{s}=\arg(-V_{ts}V_{tb}^{\ast}/V_{cs}V_{cb}^{\ast})$. The indirect
determination via global fits to experimental data gives
$2\beta_{s}=(0.0363{\>}^{+{\>}0.0016}_{-{\>}0.0015})\>\rm rad$. New Physics
contributions could significantly alter this phase.
The channel $B_{s}\to J/\psi f_{0}(980)$ is also sensitive to the same phase.
It has been first observed by the LHCb collaboration [36] using 2010 data and
quickly confirmed by Belle [37] and CDF [38].
LHCb report measurements of the phase $\phi_{s}$ for each of these channels
using $338\>\rm fb^{-1}$, and also performing a simultaneous fit to both
channels. In both cases, flavour-tagged and untagged events are used, and the
tagging efficiency is calibrated to control channels. The trigger and
selection bias, in particular with respect to lifetime, is also extracted from
the data itself. Due to the vector nature of the $\phi$ meson, the $B_{s}\to
J/\psi\phi$ needs an angular analysis to disentangle the CP-even and CP-odd
final states (Figs. 8 and 9). This is not necessary in the $f_{0}(980)$ case.
Figure 8: Definitions of the decay angles in $B_{s}\to J/\psi\phi$ decay.
Figure 9: Fits projections to proper time (top left), and the three angles in
the $B_{s}\to J/\psi\phi$ decay [39].
The results are [39, 40, 41]
$\displaystyle\phi_{s}^{J/\psi f_{0}}$ $\displaystyle=$ $\displaystyle-0.44\pm
0.44\pm 0.02~{}\mathrm{rad}$ $\displaystyle\phi_{s}^{J/\psi\phi}$
$\displaystyle=$ $\displaystyle+0.13\pm 0.18\pm 0.07~{}\mathrm{rad}$
$\displaystyle\phi_{s}^{\rm Comb}$ $\displaystyle=$ $\displaystyle+0.03\pm
0.16\pm 0.07~{}\mathrm{rad}$
which are consistent with the SM prediction. A sign ambiguity remains under
the sign reversal of ($\phi_{s}$) and ($\Delta\Gamma_{s}$). The allowed
regions are shown in Fig 10.
Figure 10: Constraints of LHCb’s measurement on the
$\phi_{s}$–$\Delta\Gamma_{s}$ plane from $B_{s}\to J/\psi\phi$ [39].
With increasing precision on $CP$ violating phases in $b\to c\bar{c}s$
transitions, assumed to be dominated by tree-level topologies, it will become
crucial in the future to understand contributions from penguin topologies [42,
43]. These can be studied using Cabibbo-suppressed decays that are related by
$U$-spin symmetry. One example is the $B_{s}\to J/\psi K_{S}^{0}$ decay, which
is the partner of the golden mode $B_{d}\to J/\psi K_{S}^{0}$. CDF [44] and
LHCb [45] have recently reported on the branching ratio of this channel and
more precision will become available when more data is collected.
## V Conclusions
With its large $b$ and $c$ cross-sections, the LHC is the new flavour factory.
Many flavour physics results, mostly from LHCb, are becoming available
yielding an unexpected and interesting pattern of measurements. The long
awaited $B_{s}\to\mu\mu$ decay has not yet been observed, thus excluding large
regions of the SUSY parameter space. Similarly the $CP$ violating phase in
$B_{s}$ decays is compatible with the SM expectation, as well as the angular
distributions in $B\to K^{\ast}\mu\mu$. Yet all these measurements are still
affected by statistical errors much larger than the theoretical errors, which
leaves a lot of room for observations of new physics.
The biggest surprise comes from the measurement of $\Delta A_{CP}$ which
exhibits a $3.5\sigma$ evidence for new physics. This will all have to be
followed up very closely with increasing statistics. As high-precision beauty
and charm physics is sensitive to energy scales much beyond the LHC centre-of-
mass energy it is likely that flavour physics is paving the way for direct
observations of new particles by the general purpose detectors at the LHC.
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|
arxiv-papers
| 2011-11-29T14:50:12 |
2024-09-04T02:49:24.767495
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Patrick Koppenburg",
"submitter": "Patrick Koppenburg",
"url": "https://arxiv.org/abs/1111.6829"
}
|
1111.7004
|
# Arcsecond resolution mapping of Sulfur Dioxide emission
in the circumstellar envelope of VY Canis Majoris
Roger Fu11affiliation: Harvard-Smithsonian Center for Astrophysics, 60 Garden
Street, Cambridge, MA 22affiliation: Department of Earth, Atmospheric, and
Planetary Sciences, Massachusetts Institute of Technology, Cambridge, MA ,
Arielle Moullet11affiliation: Harvard-Smithsonian Center for Astrophysics, 60
Garden Street, Cambridge, MA , Nimesh A. Patel11affiliation: Harvard-
Smithsonian Center for Astrophysics, 60 Garden Street, Cambridge, MA
33affiliation: Author for correspondence: npatel@cfa.harvard.edu ,
John Biersteker11affiliation: Harvard-Smithsonian Center for Astrophysics, 60
Garden Street, Cambridge, MA , Kimberly L. Derose11affiliation: Harvard-
Smithsonian Center for Astrophysics, 60 Garden Street, Cambridge, MA , Kenneth
H. Young11affiliation: Harvard-Smithsonian Center for Astrophysics, 60 Garden
Street, Cambridge, MA
###### Abstract
We report Submillimeter Array observations of SO2 emission in the
circumstellar envelope of the red supergiant VY CMa, with an angular
resolution of $\approx 1"$. SO2 emission appears in three distinct outflow
regions surrounding the central continuum peak emission that is spatially
unresolved. No bipolar structure is noted in the sources. A fourth source of
SO2 is identified as a spherical wind centered at the systemic velocity. We
estimate the SO2 column density and rotational temperature assuming local
thermal equilibrium (LTE) as well as perform non-LTE radiative transfer
analysis using RADEX. Column densities of SO2 are found to be $\sim 10^{16}$
cm-2 in the outflows and in the spherical wind. Comparison with existing maps
of the two parent species OH and SO shows the SO2 distribution to be
consistent with that of OH. The abundance ratio $f_{SO_{2}}/f_{SO}$ is greater
than unity for all radii greater than at least $3\times 10^{16}$ cm. SO2 is
distributed in fragmented clumps compared to SO, PN, and SiS molecules. These
observations lend support to specific models of circumstellar chemistry that
predict $f_{SO_{2}}/f_{SO}>1$ and may suggest the role of localized effects
such as shocks in the production of SO2 in the circumstellar envelope.
stars: individual (VY CMa (catalog )) — stars: late-type — circumstellar
matter — submillimeter — radio lines: stars
††slugcomment: Accepted on 22 November 2011 for publication in ApJ.
## 1 Introduction
Massive stars ($M_{*}\gtrsim 8\,M_{\odot}$), enter a red supergiant (RSG)
phase during which the star experiences mass-loss at rates of
$\dot{M}\sim\,10^{-5}-10^{-3}M_{\odot}yr^{-1}$ (van Loon et al., 2005). The
time variation of this mass-loss rate is not well-constrained by theoretical
studies (Yoon and Cantiello, 2010). As a result, the total amount of mass lost
over the course of the RSG phase remains uncertain for a given initial mass
(Smith et al., 2009). Observations of mass-loss events have shown them to be
sporadic and spatially anisotropic (de Wit et al., 2008).
VY Canis Majoris (VY CMa) is an oxygen-rich red supergiant with an estimated
mass of $M_{*}\approx 25\,M_{\odot}$ and a mass-loss rate estimated to be
$\dot{M}\sim\,2-4\times 10^{-4}M_{\odot}yr^{-1}$ (Danchi et al., 1994; Smith
et al., 2009). Optical images of this source show multiple discrete and
asymmetric mass-loss events, ranging in age from 1700 to 157 years ago, that
are distinct from the general flow of diffuse material (Humphreys et al.,
2007). Detailed studies of mm/sub-mm molecular spectra of VY CMa have been
carried out, revealing the chemical complexity in the envelope (Ziurys et al.,
2007; Royer et al., 2010; Tenenbaum et al., 2010). Spatial structures in the
visible and IR bands have also been obtained (Smith et al., 2001). High
angular resolution Submillimeter Array (SMA)111The Submillimeter Array is a
joint project between the Smithsonian Astrophysical Observatory and the
Academia Sinica Institute of Astronomy and Astrophysics, and is funded by the
Smithsonian Institution and the Academia Sinica. observations of VY CMa
produced maps of the spatial distribution of CO and SO (Muller et al., 2007).
Millimeter wavelength observations have long shown evidence of SO2 in the
circumstellar envelopes (CSE) of oxygen-rich red supergiants (e.g., Guilloteau
et al., 1986) and SO2 emission in VY CMa was first detected by Sahai & Wannier
(1992). Later infrared observations have also shown the production of SO2
within several radii of O-rich AGB stars (in the ”inner wind”) Yamamura et al.
(1999). From a theoretical perspective, the presence of SO2 in circumstellar
envelopes has been predicted by isotropic, non-equilibrium models of stellar
chemistry (Scalo & Slavsky, 1980; Nejad & Millar, 1988; Willacy & Millar,
1997). This class of models assume an isotropic geometry of the CSE with
specified values for the mass loss rate and expansion velocity. The resulting
model CSEs consist of an inner region where assumed parent species in the
outflow are broken down, a intermediate region where chemical reactions lead
to high abundances of daughter species, and an outer region where the daughter
species are destroyed by photodissociation. SO2 in these models has been
assigned as a daughter species, although it has since been shown to exist as a
possible parent species created in the inner wind (Decin et al., 2010).
Cherchneff (2006) studied the role of shocks in SO2 formation in the inner
wind within a few stellar radii of the photosphere.
Observational works (e.g., Jackson & Nguyen, 1988; Yamamura et al., 1999;
Decin et al., 2010) have repeatedly shown SO2 abundance to be higher both in
the extended CSE and the inner wind than expected from CSE and inner wind
chemistry models (Willacy & Millar, 1997; Cherchneff, 2006). Local abundance
of SO2 may be enhanced in the ISM by the passage of shocks (Hartquist et al.,
1980), and shock chemistry has been proposed to explain the high observed SO2
relative abundances in CSEs (Jackson & Nguyen, 1988). Alternatively, low model
abundances of SO2 presented in Nejad & Millar (1988) compared to that of Scalo
& Slavsky (1980) may be due to the latter’s lower assumed value for the
photodissociation of SO2, which is the main mode of destruction of this
molecule at large radii (Sahai & Wannier, 1992).
Interferometric mapping of SO2 distribution in the CSE of an evolved star may
contribute to the improved understanding of sulfur chemistry by directly
constraining the relative abundances of sulfur-bearing species as a function
of radius. Furthermore, associations of SO2 enhancement with discrete outflow
features may favor the idea of localized production and provide evidence for
non-isotropic processes in CSEs (Jackson & Nguyen, 1988). In this paper we
present results of high spatial resolution maps obtained from SMA observations
of SO2 around VY CMa. We describe four spatially discrete sources of SO2
emission and use multiple transitions of SO2 and derive rotational
temperatures and column densities (assuming local thermodynamical
equilibrium). Furthermore, we perform non-LTE radiative transfer analysis
using the RADEX package (van der Tak et al., 2007) to derive local kinetic
temperatures, H2 densities, and SO2 abundances of the SO2 emitting regions. We
then compare our maps to published visible, IR, and submillimeter observations
to evaluate their consistency with the above cited models of circumstellar
chemistry.
## 2 Observations
VY CMa was observed with the SMA on 2009 February 18, with the array in
extended configuration, offering baselines from 44.2 to 225.9 m. The projected
baseline lengths were from 14 to 188 m. The frequency coverage was 234.36 to
236.34 GHz in the lower sideband and 244.36 to 246.34 GHz in the upper
sideband. The phase center was at
$\alpha(2000)=07^{h}22^{m}58^{s}.27,\delta(2000)=-25^{\circ}46^{\prime}03.4^{\prime\prime}$.
The quasars 0730-116 and 0538-440 were observed every 20 minutes for gain
calibration, and the spectral band-pass was calibrated using quasar 3c273.
Flux calibration was done using recent SMA measurements of 0730-116 (2.75 Jy
at 1 mm) and 0538-440 (3.96 Jy at 1 mm). Nominal flux calibration accuracy is
15 to 20%, depending on the phase stability. In our observations, the
uncertainty appears to be better than 10%, based on the good agreement with
the ARO spectra (see section 3.2 and Figure 3). The on-source integration time
on VY CMa was 5.66 hours. Tsys (SSB) varied approximately from 200 to 400 K
during the track with an atmospheric zenith optical depth of $\sim$0.1 at the
standard reporting frequency of 225 GHz, measured at the nearby Caltech
Submillimeter Observatory. The conversion factor between Kelvin and Jansky for
our observations is $\sim 14$ K/Jy.
The visibility data were calibrated using the MIR package in IDL and imaged
with the MIRIAD software222http://www.cfa.harvard.edu/$\sim$cqi/mircook.html
(Sault, Teuben & Wright, 1995). The field of view (FWHM primary beam) varies
from 53′′ to 56′′ whereas the largest angular extent of the source is expected
to be about 5′′. The synthesized beam size, representing the obtained spatial
resolution, is $1.^{\prime\prime}48\times 1.03^{\prime\prime}$. The adopted
weighting mode for imaging is “natural” weighting (specified in the Miriad
task INVERT). The final rms noise level is $\sim$60 mJy/beam per channel in
the spectral line images and 3.9 mJy/beam in the continuum map. Line emission
was subtracted b yvisually examining the spectra in visibility amplitudes, and
using the line-free channels specification to the Miriad task UVLIN.
## 3 Results
### 3.1 Continuum emission
We detected continuum emission in both the upper and lower sidebands, which is
well described by a point source centered at RA=$07^{h}22^{m}58.336^{s}$,
Dec=$-25^{\circ}46^{\prime}03.063^{\prime\prime}$. A 2D Gaussian fit in the
image plane yields an integrated flux of $335.0\pm 66$ mJy at $235.4$ GHz and
$359.3\pm 72$ mJy at $245.4$ GHz (maximum errors are of $20\%$, mainly
resulting from the uncertainty in flux calibration). Previous SMA observations
of VY CMa at $215$ and $225$ GHz have yielded continuum fluxes of $270\pm 40$
mJy and $288\pm 25$ mJy respectively (Shinnaga et al., 2004; Muller et al.,
2007). These four values are consistent under the Rayleigh-Jeans approximation
with black-body radiation with a brightness temperature higher than 12 K.
### 3.2 Line emission
In the two sidebands of $\sim$4 GHz bandwidth, we detected a total of 14
lines, which show emissions from CS, H2O, PN, SiS and SO2. In Figure 1 we
present the spectra and maps of the SiS J=13-12 and the PN J=5-4 lines. These
lines are representative of a relatively compact emission that is centered on
the peak of the continuum emission, hence they allow us to estimate the
systemic velocity of the star. The SiS spectrum shows the distinctive triple
peak morphology as described by Ziurys et al. (2007); Tenenbaum et al. (2010)
that is interpreted as the signature of a slowly expanding shell and a pair of
faster outflows nearly collimated with the line of sight. We use the fitted
velocity of the central peak of these lines to find the systemic velocity of
the star $v_{LSR}=19.5$ km s-1. This value is consistent with other results
from mm and sub-mm observations (Ziurys et al., 2007; Muller et al., 2007) and
is adopted as the star rest-frame velocity in the following sections.
For the remainder of this work we focus only on the SO2 lines, which are
summarized in Table 1. The excitation energies of the four identified SO2
lines vary from 19 to 130 K. All four of these lines were detected in the
spectral line survey of Tenenbaum et al. (2010). Our interferometric
observations reveal the spatial distribution of the SO2 emission for the first
time. The spectra and integrated maps of all four lines are shown in Figure 2.
Figure 3 compares the spectra of two SO2 lines with single dish data from the
Submillimeter Telescope of the Arizona Radio Observatory, while channel maps
of all four lines are presented in Figures 4, 5, 6 and 7.
Maps of SO2 lines maps show well resolved spatially extended emission over an
area of $3^{\prime\prime}\times 5^{\prime\prime}$, in which we distinguish
four distinct components, hereafter designated as sources A, B, C, and D.
Figure 8 shows all four sources in a position-velocity space diagram of the
245563.4 MHz line. Source positions, orientations, and intensities were fitted
to 2D Gaussians using the MIRIAD routine IMFIT. Source A is a very strong
blue-shifted source with $v_{lsr}$ between $-18$ and $+10$ km s-1. It is
spatially offset to the East of the stellar position by $\sim 0.7"$. Of
similar intensity to source A is source B, an elongated source offset by $\sim
0.8"$ to the West. Source B has a broader spectral profile than source A and
is heavily red-shifted with $v_{lsr}$ between $+26$ and $+66$ km s-1. Source C
is a strong source with a similar radial velocity as source B. In contrast to
sources A and B, however, it is heavily offset to the West by $\sim 2.5"$.
Source D is significantly weaker than sources A, B, and C. In velocity space,
it is centered on the star’s frame of motion ($12$ km s-1 $<v_{LSR}$$<26$ km
s-1) with a distinctly narrower spectral profile and similar spatial
dimensions as the other sources. It is also spatially centered at the star’s
location.
Assuming local thermal equilibrium, we can use the integrated intensities
measured in the detected SO2 lines in each of the four identified SO2 sources
to constrain their respective column density $N$ and rotational temperature
$T_{rot}$, using the expression from Friedel et al. (2005):
$\ln[\frac{3c^{2}I_{beam}}{16\pi^{3}B\theta_{a}\theta_{b}\nu^{3}S\mu^{2}}]=\ln[\frac{N}{Z}]-\frac{E_{u}}{T_{rot}}$
(1)
where $I_{beam}$ is the integrated flux in Jy km s-1 beam-1,
$\theta_{a}\theta_{b}$ is the beam size, and $B$ the beam filling factor
varying from $0.8$ to $0.95$ for the four distinct sources, as derived from
Snyder et al. (2005): $\Theta_{source}/(\Theta_{beam}+\Theta_{source})$. $Z$
is the rotational partition function (Müller et al., 2005), $S$ is the line
strength, $\mu$ the dipole moment (1.63 Debye for SO2), $E_{u}$ the line’s
upper rotational state energy in K taken from the JPL spectroscopic database
(Pickett, 1991), and $\Theta$ the solid angle.
By performing linear fits of the rotational temperature diagram shown in
Figure 9, we find that the three faster outflows (sources A, B, and C) show
variations in rotational temperature. Sources A, B, and C have rotational
temperature of $61_{-17}^{+39}$ K (one $\sigma$ error), $110_{-29}^{+63}$ K,
and $69_{-17}^{+34}$ K, respectively. The column densities (N) of all three
sources are similar (A: $9.2_{-4.4}^{+7.8}\times 10^{15}$ cm-2; B:
$1.9_{-0.6}^{+0.8}\times 10^{16}$ cm-2; C: $1.1_{-0.4}^{+0.7}\times 10^{16}$
cm-2). The rotational temperature of the close-in source D is much higher at
$240_{-52}^{+91}$ K, while its column density is likely lower than that of the
others at $7.6\pm 1.0\times 10^{15}$ cm-2.
The derived parameters for source D are uncertain since a significant portion
of flux in its velocity range was missed by our interferometric observations
due to the lack of short spacing data (our shortest projected baseline 14 m).
Comparison with single dish data from the Arizona Radio Observatory shows that
features as large as 4” are fully recovered (Figure 3). Features much larger
than 4” may be subject to an underestimation of its flux, although 4” is an
highly conservative estimate. To address this deficiency, we repeated the
rotational temperature diagram analysis for source D using a single dish
observation made by the Arizona Radio Observatory 10 meter Submillimeter
Telescope (Tenenbaum et al., 2010). Assuming a source size of 4” by 4” we
derive a lower rotational temperature of $97_{-31}^{+84}$ and a column density
of $1.9^{+2.9}_{-1.8}\times 10^{15}$ cm-2. Because this column density
represents an average over our assumed source size, it cannot be directly
compared to our derived column densities, which represent the values at the
center of each source. This averaged column density also represents an upper
limit due to our assumption of the smallest possible angular size.
The validity of the LTE approximation can be assessed by checking the
linearity of the data points in the rotational temperature diagram. Deviation
from a linear trend suggests non-LTE conditions in the source or line
misidentification. In the case of source D, deviation from linearity is very
small compared to the error. We are therefore confident that LTE is a valid
assumption for this source. For the outflow sources A through C, a systematic
concave up shape is noted in the rotation temperature diagram, suggesting
departure from LTE conditions.
These findings justify the need to consider a non-LTE radiative transfer model
to interpret the measured fluxes. We perform this analysis using the RADEX
package (van der Tak et al., 2007), which calculates expected line intensities
for a given column density of SO2, kinetic temperature (Tkin), and volume
density of H2. For each source, we assume spatially homogeneous kinetic
temperature and volume density. The SO2 column density is fixed for each
source and its value is drawn from the LTE analysis, which gives $10^{16}$
cm-2 for sources A, B, and C and $10^{15}$ cm-2 for source D. Varying column
density one order of magnitude in each direction does not significantly affect
the results. We perform calculations for a wide range of values for the
kinetic temperature and the H2 density. We then evaluated the resulting ratios
between the flux of each line and that of the $235151.7$ MHz transition for
the outflow sources A, B, and C. Due to the lack of emission from the
$235151.7$ MHz line in the vicinity of source D, the 236216.7 MHz line was
used instead to calculate line ratios.
Fits to our data are presented as the $\chi^{2}$ statistic p-value for each
assumed value of Tkin, and H2 density (Figure 10) and best fit results are
tabulated in Table 3. We see that the non-LTE RADEX results for $T_{kin}$ show
general agreement with $T_{rot}$ derived with LTE assumption. A reliable
estimate of H2 density is elusive. Because the flux of each line becomes more
similar at greater values of H2 density, we are able only to constrain a lower
bound on this value.
The constraints on the temperature and H2 density in source D are much weaker.
However, the best fit values of Tkin span the 130 K to 330 K range, which
brackets the value of T${}_{rot}=240$ K obtained from the LTE analysis. This
agreement between Trot and Tkin is consistent with LTE conditions in source D.
Our RADEX analysis also allows us to evaluate our assumption of optically thin
SO2 lines necessary for the LTE analysis performed above. Although for lower
H2 densities ($\apprle 10^{7}$ cm-3), the 235.151 GHz line comes close to
being optically thick, for the inferred H2 densities, all four SO2 have
optical thickness below $10^{-2}$.
Assuming that the extent of the sources in the line of sight direction is
similar to that in the plane of the sky, and that the density is homogenous,
we estimate the fractional abundance of SO2, $f_{SO_{2}}$ in each source.
These are tabulated in Table 3.
## 4 Discussion
### 4.1 Source properties
We begin by discussing the positions of the four identified SO2 sources.
Source C, offset $2.5"$ to the West, is clearly an isolated body with no
antipodal companion to the East. Its existence was suspected in Ziurys et al.
(2007), but was not treated as a distinct source. On the other hand sources A
and B are found to be offset in opposite directions relative to the star
position with similar blue and redshift velocities relative to the stellar
frame. SMA observations of the CO and SO analogs of sources A and B were
interpreted to be antipodal companions oriented $15^{\circ}$ from the line of
sight by Muller et al. (2007). On the other hand, Ziurys et al. (2009) found
that single dish observations of CO and other molecules are best explained by
a blueshifted and a redshifted source at 20∘ and 45∘ from the line of sight.
Similarly, visible and IR HST observations (Humphreys et al., 2007; Smith et
al., 2001) have found no evidence of antipodal structure around VY CMa. A
careful inspection of sources A and B in our observation shows that the faster
(in radial velocity relative to the star’s reference frame) sections of both
bodies are offset towards the southwest, while the slower sections are offset
towards the northeast. This observation argues against a bipolar geometry, in
which antipodal subsections of the two outflows should have similar
velocities. We therefore treat sources A, B, and C as three distinct outflows
probably unrelated to the symmetry axis of the star. These sources do not seem
to correspond to any visible or IR features. Our maps show that all SO2
outflow sources are too close to the star to be identified as the ”curved
nebulous tail” or the numbered arcs presented in Smith et al. (2001), which
are located $\sim 3.5"$ from the star.
We identify sources A and B with the blue and redshifted outflows of the
previous authors (Ziurys et al., 2007; Muller et al., 2007; Ziurys et al.,
2009). The P-V space morphology of these sources match closely with outflows
of both CO and SO lines mapped by Muller et al. (2007) (Figure 8). These
bodies are hypothesized to have originated in an episode of anomalously high
mass loss at uncorrelated locations on the stellar surface (Smith et al.,
2001). Assuming that the ages of the outflows are similar and adopting the
$\sim 500$ year age found by Muller et al. (2007), we can attempt to derive
their locations in three-dimensions. We find that the deprojected radii of
sources A and B are $4.2\times 10^{16}$ and $4.8\times 10^{16}$ cm and that
both are situated at $22^{\circ}$ from the line of sight. Their deprojected
star-frame velocities of $28$ and $30$ km s-1 fall within the range of
measured outflow velocities from the multi-epoch observations of Humphreys et
al. (2007). If we assume the same age as for sources A and B, then source C is
found at a similar radius from the star ($6.9\times 10^{16}$ cm), but it is
much faster at $44$ km s-1 and is situated $54^{\circ}$ from the line of
sight. SO2 abundance at large radii is expected to be controlled by a balance
between the rates of production, expansion, and photodissociation (Scalo &
Slavsky, 1980). Outflows with faster expansion velocity are expected to
maintain high SO2 abundance out to greater radii, as in the case of our source
C.
Source D is elongated and centered at the stellar position, with a resolved
minor and major radii of $1"$ and $1.6"$, corresponding to $2.2\times 10^{16}$
and $3.5\times 10^{16}$ cm. The minor axis radius corresponds to between $180$
and $540$ stellar radii, depending on the adopted stellar radius (Monnier et
al., 1999; Massey et al., 2006). A significant proportion of source D flux is
missing from our observation when compared to single dish results (Figure 3).
The source we observe therefore appears to represent the warm core of a larger
extended envelope found at the systemic velocity with lower average column
density. This extended source is analogous to the spherical wind described in
previous millimeter wavelength observations of VY CMa by Ziurys et al. (2007),
Ziurys et al. (2009), Muller et al. (2007), and Tenenbaum et al. (2010);
however, only the first and last of these works observed SO2 and did not
identify it in the spherical wind. These previous works have found a
relatively low expansion velocity of between $15$ and $20$ km s-1, which is
high given the narrow velocity range of source in our channel maps (Figures 4
\- 7). However, this discrepancy may be due to missing source D flux in our
observations. The rotational temperatures derived from SMA and ARO data
indicate the existence of a thermalized compact region with elevated
temperatures with diameter $\approx 3\times 10^{16}$ cm. In the sections
below, we refer to this inner region as the core of source D.
Our derived values of $T_{rot}$ and $T_{kin}$ distinguish between the lower
temperatures of sources A, B, and C and a much hotter core of source D. In
comparison to previous works, our temperatures for source A, B, and C bracket
the range of temperatures derived by Muller et al. (2007) and Ziurys et al.
(2009) ($57$K and $85$K, respectively). This may be expected, as the preceding
authors adopted the same best fit power-law temperature profile for both
outflows, which are assumed to be at the same radius. As such, these previous
results may represent an average of the temperatures of the outflows.
Our inferred H2 densities from RADEX radiative transfer analysis are higher
than the values found in previous studies. Ziurys et al. (2009) adopted an
isotropic H2 density profile based on an assumed mass-loss rate that gives a
value of $\sim 1\times 10^{6}$ cm-3 at $10^{16}$ cm radius. Muller et al.
(2007) use a similar procedure to arrive at a lower value of $4.5\times
10^{5}$ cm-3 at the same location. A higher density of $3\times 10^{6}$ cm-3
is assumed for the inner wind region in Tenenbaum et al. (2007).
Part of this discrepancy between our values for H2 density and that of
previous authors may be explained by the latter’s assumption of isotropic mass
flow, which does not account for density concentration in the spatially
confined outflow regions. However, this reason alone may not be able to
account for the more than two orders of magnitude difference. More
significantly, our derived densities of between $5\times 10^{6}$ and $2\times
10^{8}$ cm-3 may be due to the presence of SO2 in regions of local density
enhancement above the expected values from an isotropic model. We speculate
that such high density regions are the result of shocks, and their existence
around VY CMa can be inferred from the observations of OH masers, which
overlap with sources A and B of SO2 emission (Bowers et al., 1983). The
activation of the observed 1612 MHz maser line requires H2 densities of
between $1\times 10^{6}$ and $3\times 10^{7}$ cm-3 at 100 K (Pavlakis &
Kylafis, 1996). The upper end of this range is similar to our inferred value
of minimum H2 density for sources B and C, while that of source A is much
higher. However, OH masers may still be active in our source A despite its
high density since its temperature of $\sim 55$K is cooler than that assumed
in Pavlakis & Kylafis (1996).
### 4.2 Sulfur chemistry in the CSE
For the remainder of the Discussion, we address the implications of our
observations for circumstellar chemistry by comparing SO2 distribution with
those of the OH and SO molecules. SO2 in CSEs is formed via the following
reaction (Scalo & Slavsky, 1980; Nejad & Millar, 1988; Willacy & Millar, 1997;
Cherchneff, 2006):
$SO+OH\longrightarrow SO_{2}+H$
The radial abundance of SO2 is therefore expected to reflect that of the two
reactant molecules. A striking similarity between the spectral profiles of SO2
and OH masers has already been noted by Ziurys et al. (2007). Maps of the 1612
MHz OH maser line show that it coincides with SO2 in sources A and B, but it
is weak or undetectable in the outflow source C or the spherical wind source
D. The lack of OH maser detection in source C (or perhaps a very weak
detection; see Bowers et al., 1983) may be explained by anisotropic nature of
maser radiation. Assuming that the velocity of outflows is oriented radially
away from the star, relative velocities between different clumps of gas along
a photon’s line of travel are smallest when the path is parallel or
antiparallel to the gas expansion velocity. The most efficient pumping of a
masering state is then achieved when photons travel radially inward or outward
from the star. Therefore, the strongest maser emissions are observed from
sources along our line of sight (Elitzur et al., 1976). This condition is
nearly met for sources A and B ($15^{\circ}$ to $22^{\circ}$ from the line of
sight). On the other hand source C, found at a line of sight angle of
$54^{\circ}$, may also produce the 1612 MHz OH maser, but its signal is weak
along the line of sight.
The weak OH maser emission in the stellar velocity frame may be attributed to
the high kinetic temperature of gas in this region. For a given volume density
of gas, temperatures above a certain threshold tend to induce thermal
equilibrium in the emitting body, undoing the population inversions
responsible for maser emissions (Pavlakis & Kylafis, 1996). For temperatures
of $\sim 200$K, the maximum allowable H2 density for the 1612 MHz OH maser is
$3\times 10^{6}$ cm-3, which is more than an order of magnitude lower than our
inferred density for source D. The highly linear trend of Source D data in the
rotational temperature diagram (Figure 9) corroborates the prevalence of LTE
conditions in source D.
We therefore find that OH and SO2 distributions are generally similar and that
their differences are reconcilable.
Finally, we address the discrepancies between the SO2 and SO distributions of
Muller et al. (2007), who have mapped the distribution of SO around VY CMa at
a similar resolution to ours using a single rotational line of SO (J = 65 \-
54; $E_{u}=35$K). SO was found in all four source regions described in this
work, although the red-shifted SO emitter is not fragmented into two discrete
sources as in the case of SO2. Because of the strongly contrasting values of
H2 density adopted in Muller et al. (2007) and this work, direct comparisons
between column densities are more instructive than comparisons between
fractional abundances. For column density to act as a valid proxy for
abundance, we must assume that the line-of-sight dimension of the
corresponding SO and SO2 sources are similar and that the distribution of the
each molecules within each source is homogeneous. While we cannot be certain
that the line-of-sight thickness of the SO and SO2 sources are equal, their
plane of sky dimensions are similar ($\sim$30% discrepancies). The combined
column density of both SO outflows and the inner wind where they overlap along
the line of sight was found to be $\sim 10^{16}$ cm-2. This value reflected
the combined column density from both outflows and the spherical wind
(corresponding to our sources A, B, and D). Given that the three regions
contribute approximately equal amounts to the total SO column density, the
value of $N_{SO}$ in each region is on the order of a few $10^{15}$ cm-2. In
contrast, we find column density of SO2 to be $1$ to $2\times 10^{16}$ cm-2 in
$\it{each}$ outflow source. The column density of SO2 in the outflow sources A
and C is therefore $\apprge 3$ times greater than that of SO in the same
regions. This statement likely applies as well to outflow source C, given its
similar SO2 column density compared to the other outflows and the lack of a
discrete SO source at its location. Furthermore, even with missing flux, the
column density of SO2 in the compact core of source D is within error as that
of the outflow sources, implying that $N_{SO_{2}}>N_{SO}$ in the core of the
spherical wind $<3\times 10^{16}$ cm from the star.
Non-equilibrium CSE chemistry models (Scalo & Slavsky, 1980; Nejad & Millar,
1988; Willacy & Millar, 1997) make differing predictions about ratio of SO2 to
SO abundance ($f_{SO_{2}}/f_{SO}$). The Scalo & Slavsky (1980) (SS) model
adopts a reaction rate for SO2 formation that is fast compared to the rate of
SO2 destruction via photodissociation. Under these conditions, SO is quickly
converted into SO2, which becomes the primary S-bearing species. SO2 therefore
reaches maximum abundance at a greater radius than SO. SS has predicted the
radius of this transition where $f_{SO_{2}}/f_{SO}>1$ to be around several
times $10^{15}$ cm.
Our observations suggest that $f_{SO_{2}}/f_{SO}>1$ for the outflow regions
($>4\times 10^{16}$ cm from star) and in the compact region within $3\times
10^{16}$ cm of the star. If there does exist a transitional radius at which
$f_{SO_{2}}/f_{SO}=1$, it must be less than the latter radius, making our
study consistent with the Scalo & Slavsky (1980) model. On the other hand,
models that predict $f_{SO_{2}}/f_{SO}<1$ at all radii are inconsistent with
our data (Nejad & Millar, 1988; Willacy & Millar, 1997). In addition, the SS
model predicts a steady increase in the value of $f_{SO_{2}}/f_{SO}$ outward
from the $f_{SO_{2}}/f_{SO}=1$ transition radius at $\sim 2\times 10^{15}$ cm.
The value of $f_{SO_{2}}/f_{SO}$ increases by more than an order of magnitude
by radius 1016 cm. Observations of the outflow sources A, B, and C do not
support this view, as $f_{SO_{2}}/f_{SO}$ in on the order of 3 for all four
sources. More precise comparisons between SO2 and SO abundance in each source
is hindered by the lack of more detailed interpretations of SO maps (the
Muller et al. (2007) observation included only one SO line). Therefore,
quantitative comparison between model radial distribution of the two species
and our results is elusive and we regard our support of the SS model as only
qualitative.
Other uncertainties remain in our understanding of sulfur chemistry in the CSE
of VY CMa. Our observations cannot be used to constrain whether the SO2
originates in the CSE, as assumed by the cited models, or in shocks in an
”inner wind” within a few stellar radii of the photosphere as modeled by
Cherchneff (2006) and observed by Decin et al. (2010) and Yamamura et al.
(1999). Furthermore, as proposed earlier by Jackson & Nguyen (1988) and
Willacy & Millar (1997), shock events similar to those that occur in the inner
wind may be the source of SO2 enrichment in the outflow lobes. Indeed,
comparison of SO2 and maps of SO, PN, and SiS (Muller et al., 2007; Figure 1,
this work) shows that SO2 exhibits a more fragmented distribution with
discrete sources in the red-shifted outflow. Unlike these other molecules, the
red-shifted SO2 outflow is partitioned into two regions of high local
abundance (sources B and C), suggesting that local effects may participate in
the formation of SO2. However, CSE chemistry models that include the effect of
shocks (similar to Cherchneff, 2006) remain to be done.
A further uncertainty involves the very high mass loss rate of VY CMa, which
is, for example, about two orders of magnitude faster than the asymptotic
giant branch star IK Tau (Olofsson et al., 1998). Varying the mass loss rate
in CSE chemistry models does not significantly affect the radial abundance
profiles of chemicals species while it does result in globally larger
envelopes for higher mass loss rates Willacy & Millar (1997). Therefore the
discussions of the $f_{SO_{2}}/f_{SO}$ profile in this work are also relevant
to stars with lower mass loss rates, while the actual measured radii of peak
SO2 abundance around VY CMa are unique this object.
### 4.3 Summary
1. 1.
Four rotational lines of the SO2 molecule are mapped with $\sim 1"$ resolution
around VY CMa. SO2 is found in four discrete sources, three of which are fast
($28$ to $44$ km s-1) outflows far from the star and one is a slower spherical
wind near the star. No symmetrical relationship among the faster outflows or
visible and IR features are found.
2. 2.
The three fast outflows are found at similar distances from the star and
probably originated around 500 years ago.
3. 3.
Comparison between our SO2 maps and those of the 1612 MHz OH maser line
suggests that the two species are strongly correlated and that the OH maser
detection may be limited by high temperature and density in the spherical
wind.
4. 4.
SO2 is more abundant than SO in all three outflow sources, supporting the non-
equilibrium chemistry model of Scalo & Slavsky (1980). It is inconsistent with
models that predict $f_{SO_{2}}/f_{SO}<1$ for all radii (e.g., Nejad & Millar,
1988; Willacy & Millar, 1997). The distribution of SO2 in discrete clumps when
compared to other molecules may point to the role of localized effects, such
as shocks, in the enhancement of SO2 abundance.
We are grateful to Raymond Blundell, Thomas Dame and Patrick Thaddeus for the
opportunity to carry out this research using the SMA as part of the Laboratory
Astrophysics (Astro 191) course at Harvard University. We also thank Carl
Gottlieb for helpful comments on the manuscript.
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Table 1: Summary of observed SO2 lines. SO2 | Frequency | Eu11Eu is the upper rotational state energy. | $S$22$S$ is the line strength.
---|---|---|---
transition | (MHz) | (K) |
$4_{2,2}-3_{1,3}$ | 235151.7 | 19.0 | 1.71
$10_{3,7}-10_{2,8}$ | 245563.4 | 72.7 | 5.44
$15_{2,14}-15_{1,15}$ | 248057.4 | 119.3 | 5.27
$16_{1,15}-15_{2,14}$ | 236216.7 | 130.7 | 6.05
Table 2: Description of SO2 line emission sources**Elliptical gaussian fits of the identified SO2 emission sources. Offsets are measured in arcseconds from the center of the continuum emission (RA=$07^{h}22^{m}58.336^{s}$, Dec=$-25^{\circ}46^{\prime}03.063^{\prime\prime}$). Source size represents the major and minor axis of the ellipse, and its orientation angle is that of the major axis. Angles are given as rotations from north through east. Fits were performed on maps of the 235151.7 MHz line. Source | Offset | Position | Size | Orientation
---|---|---|---|---
(vlsr range in km s-1) | (“) | Angle (∘) | (“) |
A: ($-18,\,+10$) | $0.66\pm 0.10$ | $33^{\circ}$ | $2.5\times 2.4$ | $126^{\circ}$
B: ($+26,\,+66$) | $0.83\pm 0.10$ | $316^{\circ}$ | $2.9\times 1.3$ | $41^{\circ}$
C: ($+26,\,+66$) | $2.45\pm 0.09$ | $272^{\circ}$ | $2.9\times 1.8$ | $127^{\circ}$
D: ($+12,\,+26$) | $\sim 0$ | $N/A$ | $3.2\times 1.9$ | $\sim-40^{\circ}$
Table 3: Summary of LTE and RADEX modeling results Source | $T_{rot}$ | SO2 Column Density | $T_{kin}$ | H2 Density | SO2 Abundance ( $f_{SO_{2}}$)
---|---|---|---|---|---
| (K) | (cm-2) | (K) | (cm-3) |
A | $61_{-17}^{+39}$ | $9.2_{-4.4}^{+7.8}\times 10^{15}$ cm-2 | 50 | $\apprge 2$$\times 10^{8}$ | $7\times 10^{-10}$
B | $110_{-29}^{+63}$ | $1.9_{-0.6}^{+0.8}\times 10^{16}$ | $\sim$ 75 | $\apprge 3$$\times 10^{7}$ | $2\times 10^{-8}$
C | $69_{-17}^{+34}$ | $1.1_{-0.4}^{+0.7}\times 10^{16}$ | $\sim$ 100 | $\apprge 3$$\times 10^{7}$ | $8\times 10^{-9}$
D | $240_{-52}^{+91}$ | $7.6\pm 1.0\times 10^{15}$**Values for source D suffer from missing extended flux. Using single dish data, average $T_{rot}=97$ K and $N=1.8\times 10^{15}$ cm-2 assuming that the emitter is an uniform 4” by 4” region. | $\sim$ 220 | $\apprge 4$$\times 10^{7}$ | $5\times 10^{-10}$
Figure 1: Left: PN J=5–4 and SiS J=13–12 spectra showing the peak emission at
19.49 km s-1 (from the line rest frequencies). We adopt this velocity as that
of the stellar frame. ${\it Right:}$ Integrated intensity maps of PN (top) and
SiS (bottom) emission, which show maxima close to the peak of the continuum
emission. Figure 2: Left: SO2 spectra averaged in a region of 3”x3” around
the continuum peak. Right: Integrated intensity maps corresponding SO2 lines
over velocity intervals $-20$ and $10$ km s-1 are shown in blue contours. The
red contours show the integrated intensity over velocity intervals 20 to 60 km
s-1. The starting value and interval of contour levels are set to 35 mJy
beam-1 km s-1. The continuum emission is shown in grey scale (repeated in all
panels). The location of the three identified sources A, B and C is shown in
the bottom two maps. Figure 3: Left: SO2 spectra from our (black) and 10
meter single dish observations by the Arizona Radio Observatory’s
Submillimeter Telescope (red) for the two lowest temperature transitions
(235151.7 and 245563.4 MHz). The SMA data was smoothed at the appropriate
scale to approximate the field of view of the single dish telescope ( 30”).
Figure 4: Channel maps of SO2 4(2,2)–3(1,3) emission at 235151.7 MHz toward VY
CMa. In each panel, emission is integrated over a 10 km s-1 wide velocity
range centered on the velocity indicated on the top right corner. Contour
levels are same as in Fig. 2. Figure 5: Same as in Fig.4 for SO2
10(3,7)–10(2,8) emission at 245563.4 MHz. Figure 6: Same as in Fig.4 for SO2
15(2,14)–15(1,15) emission at 248057.4 MHz. Figure 7: Same as in Fig.4 for
SO2 15(1,15)–15(2,14) emission at 236216.6 MHz. Figure 8: Left: A position
velocity cut along east-west at Declination offset of 0′′ in the SO2 line
emission at 245563.4 MHz, at position angle 105∘. The contours correspond to
10% of the peak emission. Sources A, B, C and D are clearly separated. Figure
9: Rotational temperature diagram of the four SO2 emission sources. Red, blue,
and black curves and points indicate sources A, B, and C respectively. Orange
represents the spherical wind (source D). Variable “L” on the vertical axis
represents the left-hand side of Equation 1. The dashed orange line represents
the spherical wind based on ARO single dish data Tenenbaum et al. (2010). The
errors plotted are set to three sigma. Horizontal coordinates of the plotted
points have been staggered for clarity. Figure 10: Fits between our measured
line ratios and predictions with Radex radiative transfer. Solid, dashed, and
dotted contours enclose regions of $\chi^{2}$ test p-values greater than 5%,
30%, and 95%. Owing to lower S/N and weak dependence of line ratios to these
quantities, the derived temperature and density of source D are much less
well-constrained than the other three.
|
arxiv-papers
| 2011-11-29T22:01:38 |
2024-09-04T02:49:24.778084
|
{
"license": "Public Domain",
"authors": "Roger Fu, Arielle Moullet, Nimesh A. Patel, John Biersteker, Kimberly\n L. DeRose, Kenneth H. Young",
"submitter": "Nimesh Patel",
"url": "https://arxiv.org/abs/1111.7004"
}
|
1111.7068
|
# Unveiling the super-orbital modulation of LS I $+$61∘303 in X-rays
Jian Li11affiliation: Laboratory for Particle Astrophysics, Institute of High
Energy Physics, Beijing 100049, China. Email: jianli@ihep.ac.cn , Diego F.
Torres22affiliation: Institut de Ciències de l’Espai (IEEC-CSIC), Campus UAB,
Torre C5, 2a planta, 08193 Barcelona, Spain 33affiliation: Institució Catalana
de Recerca i Estudis Avançats (ICREA). , Shu Zhang11affiliation: Laboratory
for Particle Astrophysics, Institute of High Energy Physics, Beijing 100049,
China. Email: jianli@ihep.ac.cn , Daniela Hadasch22affiliation: Institut de
Ciències de l’Espai (IEEC-CSIC), Campus UAB, Torre C5, 2a planta, 08193
Barcelona, Spain , Nanda Rea22affiliation: Institut de Ciències de l’Espai
(IEEC-CSIC), Campus UAB, Torre C5, 2a planta, 08193 Barcelona, Spain , G.
Andrea Caliandro22affiliation: Institut de Ciències de l’Espai (IEEC-CSIC),
Campus UAB, Torre C5, 2a planta, 08193 Barcelona, Spain , Yupeng
Chen11affiliation: Laboratory for Particle Astrophysics, Institute of High
Energy Physics, Beijing 100049, China. Email: jianli@ihep.ac.cn , Jianmin
Wang11affiliation: Laboratory for Particle Astrophysics, Institute of High
Energy Physics, Beijing 100049, China. Email: jianli@ihep.ac.cn
###### Abstract
From the longest monitoring of LS I $+$61∘303 done to date by the Rossi X-ray
Timing Explorer (RXTE) we found evidence for the long-sought, years-long
modulation in the X-ray emission of the source. The time evolution of the
modulated fraction in the orbital lightcurves can be well fitted with a
sinusoidal function having a super-orbital period of 1667 days, the same as
the one reported in non-contemporaneous radio measurements. However, we have
found a 281.8 $\pm$ 44.6 days shift between the super-orbital variability
found at radio frequencies extrapolated to the observation time of our
campaign and what we found in the super-orbital modulation of the modulated
fraction of our X-ray data. We also find a super-orbital modulation in the
maximum count rate of the orbital lightcurves, compatible with the former
results, including the shift.
X-rays: binaries, X-rays: individual (LS I $+$61∘303)
## 1 Introduction
LS I $+$61∘303 is one of a handful high-mass X-ray binaries that have been
detected at all frequencies, including TeV and GeV energies. Its nature is
still under debate, with rotationally-powered pulsar-composed systems (see
Maraschi & Treves 1981, Dubus 2006) and microquasar jets (see Bosch-Ramon &
Khangulyan 2009 for a review) being discussed. Recently, evidence favoring LS
I $+$61∘303 as the source of a very short X-ray burst led to the analysis of a
third alternative, in which LS I $+$61∘303 is a magnetar binary (see Torres et
al. 2011, also Bednarek 2009, and Dubus 2010). Long-term monitoring of the
source, ideally across all wavelengths, is a key ingredient to disentangle
differences in behavior which could point to the underlying source nature.
Here, we report on the analysis of _RXTE_ -Proportional Counter Array (PCA)
monitoring observations of the $\gamma$-ray binary system LS I $+$61∘303\. The
dataset we consider covers the period between 2007 September and 2011
September (2007-08-28 – 2011-09-15), and it is the longest monitoring campaign
done for this source. Smaller datasets included in this campaign have been
analyzed and results have been presented in our previous papers I and II
(Torres et al. 2010, and Li et al. 2011, respectively). In them, we focused on
establishing the orbit-to-orbit X-ray variability and on studying the spectral
properties and the flares. The current data enhances in more than one
additional year the reported coverage on the source, and for the first time,
it is sufficient to consider the possible super-orbital modulation of the
X-ray emission from LS I $+$61∘303.
## 2 Observations, data analysis, and results
Our current analysis includes 473 _RXTE_ -PCA pointed observations identified
by proposal numbers 93100, 93101, 93102, 93017, 94102, 95102 and 96102. The
analysis of PCA data was performed using HEASoft 6.9. We filtered the data
using the standard _RXTE_ -PCA criteria. Only PCU2 (in the 0-4 numbering
scheme) has been used for the analysis because it was the only Proportional
Counter Unit (PCU) that was 100% on during the observations. We select time
intervals where the source elevation is $>$10∘ and the pointing offset is
$<0.02^{\circ}$. The background file used in the analysis of PCA data is the
most recent available one from the HEASARC Web site for faint sources, and
detector breakdown events have been removed.111The background file is
pca_bkgd_cmfaintl7_eMv20051128.mdl. The data have been barycentered using the
FTOOLS routine faxbary using the JPL DE405 solar system ephemeris. Our flux
and count rate values are given for an energy range of 3–30 keV.
At first, we consider the complete X-ray lightcurve of our campaign, and
fitted it with a constant to obtain an average flux, resulting in
(1.616$\pm$0.006) counts s-1. In order to remove the influence of the several
ks-long flares detected from the source, we cut all observations that
presented a count rate larger than three times this average. The remaining on-
source time amounts to 684.3 ks (460 observations), and it is uniformly
distributed in the system’s orbital (between 60 and 80 ks per each 0.1 of
orbital phase bin) as well as in the system’s super-orbital phase (around 60
ks per each 0.1 of super-orbital phase bin, except at super-orbital phase 0.8,
where an intensive campaign increased the exposure to 120 ks) as defined by
the radio observations (Gregory 2002).
Given a 6-months time bin (or approximately 6.8 orbits of the system), we take
the peak X-ray flux in that period and compute the modulated flux fraction.
The latter is defined as ($c_{max}-c_{min})/(c_{max}+c_{min}$ ), where
$c_{max}$ and $c_{min}$ are the maximum and minimum count rates in the 3–30
keV orbital profile of that period, after background subtraction. The minimum
count rates are roughly constant around 1 count s-1 all along the observation
time. Results are shown in Figure 1. Table 1 presents the values of the
reduced $\chi^{2}$ for fitting different models to the modulation fraction and
the peak flux in X-rays. It compares the results of fitting a horizontal line,
a linear fit, and two sinusoidal functions. One of the latter has the same
period and phase of the radio modulation (from Gregory 2002, labelled as
‘Radio’ in Table 1, dotted line in Figure 1). The other sine function has the
same period as in radio but allowing for a phase shift from it (a solid line
in Figure 1, labelled as ‘Shifted’ in Table 1). It is clear that there is
variability in the data (and thus that a constant fit is unacceptable), as
well as that the sinusoidal description with a phase-shift is better than the
linear one, which is obvious by visual inspection of Figure 1. The phase shift
derived by fitting the modulated fraction is 281.8 $\pm$ 44.6 days,
corresponding in phase to $\sim 0.2$ of the 1667$\pm$8 days super-orbital
period. The phase shift derived by fitting instead the maximum flux is 300.1
$\pm$ 39.1 days, and results compatible with the former.
Table 1: Reduced $\chi^{2}$ for fitting different models to the modulation fraction and the peak flux in X-rays. | Constant | Linear | Radio | Shifted
---|---|---|---|---
Modulation Fraction | 88.2 / 7 | 38.0 / 6 | 42.1 / 6 | 1.1 / 5
Peak Flux | 212.8 / 7 | 114.8 / 6 | 91.8 / 6 | 4.9 / 5
Figure 2 analyzes whether there is anything special concerning the 6 months
time bin just chosen to present the former results. It presents the peak flux
and modulated fraction in different time bins, from 4 to 12 months. The black
lines in Figure 2 shows the sine fitting with a 1667 days period fixed. The
black lines in all left (right) panels are the same as the one depicted in the
left (right) panel of Figure 1. It is interesting to note that the smaller the
time bin the larger it is the scattering around the sinusoidal fit, which can
be understood as an effect of the increasing similarity between the time bin
itself and the orbital period of the system (of 26.4960$\pm$0.0028, Gregory
2002). Indeed, orbit-to-orbit variability is known to exist in our data, and
can be similar or larger than the super-orbital induced variability at times
(e.g., see the variations found for the same phase-bin in contiguous orbits in
Figure 4 of Torres et al. 2010). Thus, the shorter the time bin, the less
likely it is that the super-orbital induced variation could be detected, which
may be sub-dominant to the local-in-time changes. On the contrary, as soon as
the integrated time bin is large enough in comparison with the orbital period
of LS I $+$61∘303, the super-orbital variability is consistently observable.
This same fact makes a direct comparison of the sinusoidal trend of the super-
orbital X-ray modulation of Figure 1 with the flux obtained from LS I
$+$61∘303 under short and isolated observations untrustable. For example the
observations at soft X-rays conducted by XMM-Newton (Neronov & Chernyakova
2006; Chernyakova et al. 2006, Sidoli et al. 2006), Chandra (Paredes et al.
2007, Rea et al. 2010), ASCA (Leahy et al. 1997), ROSAT (Goldoni & Mereghetti
1995; Taylor et al. 1996), and Einstein (Bignami et al. 1981) were all too
short to cover even a single full orbit. Similarly, earlier campaigns with PCA
during the month of March 1996 (Harrison et al. 2000, Greiner and Rau 2001)
are also too short. Comparing their obtained flux values with the super-
orbital trend would be meaningless since there is not enough time coverage to
average out the possible orbit-to-orbit variations.
Finally, we note that we have also analyzed 15 years of _RXTE_ -ASM data on LS
I $+$61∘303, but the larger error bars on the count rates preclude to obtain
any conclusion from that dataset.
Figure 1: Left: Peak count rate of the X-ray emission from LS I $+$61∘303 as a
function of time and the super-orbital phase. Right: modulated fraction, see
text for details. The dotted line shows the sine fitting to the modulated flux
fraction and peak flux with a period and phase fixed at the radio parameters
(from Gregory 2002). The solid curve stands for sinusoidal fit obtained by
fixing the period at the 1667 days value, but letting the phase vary. The time
bin corresponds to six months. The colored boxes represent the times of the
TeV observations that covered the broadly-defined apastron region. The boxes
in green denote the times when TeV observations are in low state while boxes
in yellow are TeV observations in high state.
Figure 2: Peak flux (left) and modulation fraction (right) at different time
scales. From top to bottom we plot the results obtained by binning the data in
4, 6, 8, and 12 months bins. The black line shows the sine fitting with a 1667
days period fixed, as in the corresponding panels of Figure 1
## 3 Discussion
We notice that the X-ray long-term monitoring (2007–2011) on LS I $+$61∘303
started about 7 years later from the end of the campaign used to determine the
super-orbital period in radio (1977–2000, see Gregory 2002). We will assume
then that the radio-determined super-orbital modulation of the source is,
although possibly variable, active today with similar features as the ones
claimed a decade earlier. This appears possible given that recent reports of
variation in the orbital radio maxima are of only $\sim$0.1 in phase (see
Trushkin & Nizhelskij 2010). Under such assumptions, we showed that there is a
phase shift between the radio and the X-ray super orbital modulation.
Interestingly, this shift is the same as the one hinted at between the radio
and the H$\alpha$ line (Zamanov et al. 1999, Zamanov & Martí 2000). Indeed,
the optical observations that covered the period 1989–1999 were fitted with a
period of $\sim$1584 days (Zamanov et al. 1999), a value reported prior to the
work by Gregory (2002), where the super-orbital period was refined to
1667$\pm$8 days. To investigate further the long-term modulation of LS I
$+$61∘303 at optical wavelengths and how it compares with the current findings
in X-rays, we took the H$\alpha$ data from Table 4 of Paredes et al. (1994),
Table 1 of Zamanov et al. (1999), and Figure 1 of Zamanov et al. (2000), and
as stated in Zamanov et al. (1999), considered an error of 10% for all the
equivalent widths. We performed an analysis similar to the one done for X-rays
in the previous section, and derived and optical phase lag of 290.1$\pm$16.7
days with respect to the radio phase. Thus, the optical phase lag is
coincident with the one derived at X-rays, although the observations at the
two bands are about 8 years apart. The stellar disk of Be stars are well known
to grow larger as the equivalent width of the H$\alpha$ emission line
increases (e.g., Hanushik et al. 1988; Grundstrom & Gies 2006). The optical
variability has been most likely attributed to the cyclic variation of Be
circumstellar disk. Thus, the possible coincidence with the X-ray phase lag
suggests that the stellar disk may play an important role also for the X-ray
emission, and probably for the higher-energy non-thermal emission of LS I
$+$61∘303 too.
The coincidence between the X-ray and optical shift with respect to the 1667
days radio modulation has to be taken with the necessary prudence prompted by
it being based on non-simultaneous observations. In particular, it seems that
the optical observations present the largest degree of variation in time. We
checked that in addition of the H$\alpha$ measurements mentioned above, there
are more recent ones in the works of e.g., Liu et al. (2005), Grundstrom et
al. (2007), Zamanov et al. (2007), and McSwain et al. (2010). However, the
latter span 0.51 (at best, being usually much shorter) of the super-orbital
period, and as such we can not directly use them for a comparison in long-
terms. Nevertheless, they seem to hint at that the H$\alpha$ variability is
not strictly periodic or at least at a changing amplitude.
Based (among other reasons discussed in Torres et al. 2011) on the analysis of
a the _Swift_ -BAT detection of a short, magnetar-like burst from the
direction of LS I $+$61∘303, we have proposed that the system’s compact object
is a high magnetic field, slow period pulsar. In that case, we proved that the
LS I $+$61∘303-system would most likely be subject to a flip-flop behavior,
from a rotationally powered regime (in apastron, also known as ejector), to a
propeller regime (in periastron) along each of the system’s eccentric orbits.
The multi-wavelength phenomenology can be put in the context of the former
model, and in particular, also the highest energy TeV emission, which has also
shown low and high states which are apparently modulated by the same super-
orbital period as well. Within this model, we notice that an increase in the
accreted mass onto the compact object (unavoidably linked to the mass-loss
rate of the star) by a factor of a few222Estimations of the cyclical
variations in the mass loss-rate from the Be star in LS I $+$61∘303 are given
as the ratio between maximal and minimal values obtained either from radio
emission (a factor of 4 was determined by Gregory & Neish 2002) or from
H$\alpha$ measurements, which span from a factor of 5.6 (Gregory et al. 1989)
to 1.5 (Zamanov et al. 1999). can put the system in a permanent propeller
stage along the orbit, including at the apastron region. This change of
behavior for such an small change in mass loss rate can be the reason behind
the evolution of the modulated fraction. Indeed, using the formulae in Torres
et al. (2011), and considering to simplify that the condition $R_{m}=R_{lc}$,
where $R_{m}$ stands for the magnetic radius and $R_{lc}$ for the light
cylinder, establishes both the out-of-ejector and into-ejector condition, one
sees that the period–mass-loss–magnetic field relation for the apastron of LS
I $+$61∘303 is $\left({P}/{\rm 1\;s}\right)\sim
a\\!\times\\!15({B}/{10^{14}\;{\rm
G}})^{4/7}({\dot{M}_{*}}/{10^{18}{\rm\;g\;s^{-1}}})^{-2/7}$ where $a$
represents a constant of order 1, and we have assumed an eccentricity of 0.6
and a semi-major axis of $6\times 10^{12}$ cm. For periods shorter than the
former, the system is in an ejector phase. For larger periods, it is in a
propeller stage (see Torres et al. 2011 for details). High values of magnetic
field and slow periods would make the transition possible: a cyclical change
by a factor of a few in $\dot{M}_{*}$ can make the system to abandon the
ejector phase in apastron. For instance, under a variation by a factor of 4 in
$\dot{M}_{*}$, a case that leads to a super-orbital induced transition is
given by a magnetic field of $5\times 10^{13}$ G, and a typical period of
magnetars ($\sim 7$ s). This may also happen for smaller values of the
magnetic field but only in the case of a relatively long period. For instance,
again under a variation by a factor of 4 in $\dot{M}_{*}$ and for $B=10^{12}$
G, the period should be between 700 ms and 1s in order for the system to flip-
flop in the super-orbital evolution, although no known pulsar in these
parameter ranges has a rotational energy in excess of $10^{36}$ erg s-1 (ATNF
Catalogue version: 1.43), which would be needed to account for the multi-
wavelength output of the system. Note in particular that the behavior of the
LS I $+$61∘303 system containing a pulsar with $B\sim 10^{12}$ G and $P<700$
ms would be unaffected by the cyclical variation of the mass-loss rate: it
would act as an ejector in apastron along the whole super-orbital period.
The flip-flop mechanism can then be used to qualitatively explain why LS I
$+$61∘303 has entered in a low TeV state (see, e.g., Acciari et al. 2010) when
at the maximum of the radio super-orbital variability, but perhaps also to
explain why the modulated X-ray flux fraction varies as we found in Figure 1.
When the mass-loss rate is low, the inter-wind shock formed at the collision
region between the pulsar and the stellar wind would be present at the broad
apastron region (and so will the TeV emission there), disappearing at
periastron. In this situation, there are two contributors to the X-ray
emission along the orbit, expected to be roughly at the same level (e.g.,
Zamanov et al. 1999); the shock at apastron and the propeller at periastron,
and the modulated fraction is consistently low. When at the maximum of the
mass-loss rate, the inter-wind shock may not form, and abundant TeV particles
would not be produced since shocks at the magnetic radius are unable to reach
TeV energies. Thus there is only one process generating the X-ray emission
along the system’s orbit, the propeller, and the modulated fraction is then
maximum. The exact position of the X-ray maximum along each of the orbits
would depend on the local-in-time conditions of the accreted mass onto the
compact object, which established the relative weight of the two X-ray
contributors, and it is thus expected to vary beyond the super-orbital trend
in short timescales, and not always be located at periastron. However, given
that the H$\alpha$ cycle represents the cyclical modulation of the mass loss
rate, it would be natural to expect that the X-ray emission be correlated with
it in long timescales (i.e., with how much mass is falling towards the compact
object, e.g., see Bednarek 2009 or Bednarek & Pabich 2011).
Zamanov et al. (2001) already discussed when the radio emission is expected to
peak in each of the system orbits: The switch on of the ejector phase will
activate the pulsar wind, creating a cavern around the neutron star which will
start to expand. This means that the radio outburst will peak with some delay
after the change of regimes, which is supposed to happen somewhere after the
periastron, when the accretion rate onto the compact object diminishes enough.
In a cyclical variability of the mass loss rate of the star, the ejector-
propeller transition moves in phase: at lower mass loss rates, the ejector
will switch on earlier, and the radio outburst will peak at earlier orbital
phases than at higher mass loss rates. A generic TeV and radio anti-
correlation is thus expected since the more mass fuels the propeller phase the
more violent the radio outburst will be, and the less effectively the inter-
wind shock will generate TeV particles.
Figure 3: Peak flux per orbit in TeV shown in red (all of them happening in
the 0.6–1.0 orbital phase range) as a function of superorbital phase, together
with radio (top panel) and H$\alpha$ (bottom panel) data(black) as described
in the text. The upper gray dashed line stands for the TeV flux level at
discovery of the source in 2006, whereas the lower dashed line stands for 1/3
of this flux value. Two super-orbital phases are shown for clarity. Whenever
there are both MAGIC and VERITAS compatible observations for the same orbit,
they are averaged.
The colored boxes in Figure 1 represent the times of the TeV observations that
covered the broadly-defined apastron region (from Albert et al. 2006, 2008,
Anderhub et al. 2009, Aleksic et al. 2011; Acciari et al. 2008, 2009, 2010).
Those boxes colored in green denote the times for which the observations led
only to imposing an upper-limit or to a detection with a flux that is about 3
times less than the one obtained at the discovery observations of 2006 (Albert
et al. 2006), which defines the low state. The yellow boxes stand for those
observations for which the level of the TeV emission was roughly compatible
with the original discovery. There is a trend for finding a low TeV state
towards the maximum of the super-orbital low-frequency cycles. This is perhaps
more clearly seen in Figure 3, where we plot the peak flux per orbit in TeV
(all of them happening in the 0.6–1.0 orbital phase range) as a function of
superorbital phase, together with radio and H$\alpha$ data. However, the
scarcity (and non-simultaneity) of the TeV coverage precludes reaching a
definite conclusion on whether there is an anti-correlation of the TeV
emission with the radio or with the H$\alpha$ curves. It would seem, however,
that the TeV emission is rather anti-correlated with the radio flux and not
with H$\alpha$, but this could not be quantitatively proven with the data at
hand, especially given the caveats of dealing with non-contemporaneous
observations. A simultaneous optical-TeV campaign is needed to establish the
nature of the anti-correlation. The latter would be particularly useful for
the forthcoming extrapolated radio maximum around October-November 2012.
We acknowledge support from the grants AYA2009-07391 and SGR2009-811, as well
as the Formosa Program TW2010005, by the National Natural Science Foundation
of China via NSFC-10325313, 10521001, 10733010,11073021, and 10821061, the CAS
key Project KJCX2-YW-T03, and 973 program 2009CB824800. YPC thanks the Natural
Science Foundation of China for support via NSFC-11103020 and 11133002. NR is
supported by a Ramon y Cajal Fellowship. We also acknowledge the use of the
High Energy Astrophysics Science Archive Research Center (HEASARC), provided
by NASA’s Goddard Space Flight Center. JL acknowledges the hospitality of
IEEC-CSIC, where this research was conducted.
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|
arxiv-papers
| 2011-11-30T07:10:06 |
2024-09-04T02:49:24.789950
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Jian Li, Diego F. Torres, Shu Zhang, Daniela Hadasch, Nanda Rea, G.\n Andrea Caliandro, Yupeng Chen, Jianmin Wang",
"submitter": "Jian Li",
"url": "https://arxiv.org/abs/1111.7068"
}
|
1111.7084
|
# Failed Gamma-Ray Bursts: Thermal UV/Soft X-ray Emission Accompanied by
Peculiar Afterglows
M. Xu11affiliation: Department of Astronomy, Nanjing University, Nanjing
210093, China; hyf@nju.edu.cn 22affiliation: Yukawa Institute for Theoretical
Physics, Oiwake-cho, Kitashirakawa, Sakyo-ku, Kyoto 606-8502, Japan
33affiliation: Department of Physics, Yunnan University, Kunming 650091, China
, S. Nagataki22affiliation: Yukawa Institute for Theoretical Physics, Oiwake-
cho, Kitashirakawa, Sakyo-ku, Kyoto 606-8502, Japan , Y. F.
Huang11affiliation: Department of Astronomy, Nanjing University, Nanjing
210093, China; hyf@nju.edu.cn , and S.-H. Lee22affiliation: Yukawa Institute
for Theoretical Physics, Oiwake-cho, Kitashirakawa, Sakyo-ku, Kyoto 606-8502,
Japan
###### Abstract
We show that the photospheres of “failed” Gamma-Ray Bursts (GRBs), whose bulk
Lorentz factors are much lower than 100, can be outside of internal shocks.
The resulting radiation from the photospheres is thermal and bright in UV/Soft
X-ray band. The photospheric emission lasts for about one thousand seconds
with luminosity about several times $10^{46}$ erg/s. These events can be
observed by current and future satellites. It is also shown that the
afterglows of failed GRBs are peculiar at the early stage, which makes it
possible to distinguish failed GRBs from ordinary GRBs and beaming-induced
orphan afterglows.
gamma-ray bursts: general — radiation mechanisms: thermal
## 1 Intronduction
Gamma-Ray Bursts (GRBs) are the most powerful explosion in the universe. The
origin of prompt emission remains unresolved, owing to the fact that the
prompt emission has a large explosion energy showing non-thermal spectrum with
rapid time-variabilities.
It is widely accepted that the prompt emission is coming from a highly-
relativistic flow, because it can reduce the optical depth of the flow which
makes the radiation spectrum non-thermal (Rees & Mészáros 1994). In fact, in
the internal shock scenario, which is one of the most promising scenarios, the
relativistic shells collide with each other after the system becomes optically
thin (e.g. Piran 1994 for a review). However, it was pointed out by Mészáros &
Rees (2000) that even such a relativistic flow should have a photosphere
inevitably and thermal radiation should be coming from there. Since then,
there have been many theoretical (Daigne & Mochkovitch 2002; Pe′er et al.
2006; Pe′er 2008; Pe′er & Ryde 2011) and observational (Ghirlanda et al. 2003;
Ryde 2004, 2005; Ryde et al. 2010; Guiriec et al. 2010; Ryde et al. 2011)
studies on how the thermal component contributes to the prompt emission (or
the precursor). Nowadays, the internal shock model with a photosphere is
frequently discussed (e.g. Toma et al. 2011; Wu & Zhang 2011). In this
picture, the radius of the photosphere, $R_{\rm PS}$, is usually smaller than
the radius $R_{\rm IS}$ where internal shocks are happening.
The above scenario is based on the assumption that the bulk Lorentz factor of
the jet is as large as 100-1000. But what happens if the bulk Lorentz factor
is not so high? Theoretically, it is natural to consider such a case, because
it is very hard to realize such a clean, highly-relativistic flow. Especially,
in case of long GRBs, some bursts are at least coming from the death of
massive stars where a lot of baryons should be surrounding the central engine
(MacFadyen & Woosley 1999; Proga et al. 2003; Nagataki et al. 2007; Nagataki
2009, 2010). Thus we can expect there are a lot of “failed GRBs” that have
dirty, not so highly-relativistic flows in the universe (e.g. Dermer et al.
1999; Huang et al. 2002; Paragi et al. 2010; Xu et al. 2011). Recently,
Bromberg et al. (2011a, 2011b) suggested the existence of a large population
of failed GRBs if the jets failed to break out of the progenitor stars in
collapsar model.
Qualitatively, $R_{\rm IS}$ becomes smaller if the bulk Lorentz factor of the
flow is smaller, while $R_{\rm PS}$ increases with the decreasing of bulk
Lorentz factor. Thus we can expect that the photosphere will become outside of
the internal shock region for some lower Lorentz factors (see Fig. 1a). In
such a case, $\gamma$-rays from the internal shocks cannot escape. Instead,
softer thermal radiation from the photosphere followed by an afterglow will be
seen.
We note that such a situation was considered in case of a “successful GRB”
(Rees & Mészáros 2005; Lazzati et al. 2009; Mizuta et al. 2011; Nagakura et
al. 2011; Ryde & Pe′er 2009; Ryde et al. 2010, 2011). In the photospheric
model, it is considered that the photospheric emission itself is the origin of
the prompt emission: its spectrum is modified to non-thermal due to the
heating by relativistic electrons that are produced at the internal shocks
inside the photosphere (Beloborodov 2010; Vurm et al. 2011). Thus the
photospheric model can explain normal successful GRBs.
Here, in this study, we consider the case of “failed GRBs”. We show that the
direct emission from the photosphere of a failed-GRB will be a UV/soft X-ray
burst, followed by an afterglow with a peculiar spectrum at early stage. It is
also shown that the afterglows of failed GRBs can be distinguished from
ordinary GRB afterglows and beaming-induced orphan afterglows.
Figure 1: (a) A sketch illustrating the emission regions of a failed-GRB.
Note that the internal shock radius is smaller than the photosphere radius.
(b) Schematic diagram showing the path of a photospheric photon escaping from
the ejecta in the stellar frame. A photon emitted at point E will escape from
the ejecta at point G.
## 2 Photospheric emission
In this study, we consider axisymmetric jet and assume that the observer is on
the axis $\overline{OD}$ in Fig. 1b. This is the extension of the 1-D
formulation derived by Daigne & Mochkovitch (2002). For a baryon-rich ejecta
in the stellar frame (i.e., burst source frame, and from now on all variables
are defined in this frame), we assume the ejecta has been accelerated at a
distance $r_{\rm acc}$ from the central engine. The mass flux of the ejecta is
written as $\dot{M}=\dot{E}/\Gamma c^{2}$, where $\Gamma$ is the Lorentz
factor of the ejecta and $\dot{E}$ is the energy injection rate. The energy
injection begins at $t_{\rm inj}$=0, and stops at $t_{\rm inj}=t_{\rm w}$,
i.e., the central engine activity lasts for a period of time $t_{\rm w}$.
The ejecta can be subdivided into a series of concentric layers, where each
layer has been injected at $r_{\rm acc}$ at a certain injection time $t_{\rm
j}$. Each ejecta layer becomes transparent when it has expanded to a distance
$r(t_{\rm inj}=t_{\rm j},t)$ at a specific time $t$, where $r(t_{\rm
inj}=t_{\rm j},t)=r_{\rm acc}+\beta c(t-t_{\rm j})$ and
$\beta=\sqrt{1-1/\Gamma^{2}}$. The photons emitted at $r(t_{\rm inj}=t_{\rm
j},t)$ will escape from the ejecta at a time $t_{\rm esc}$ at a distance
$r_{\rm esc}(t_{\rm inj}=t_{\rm j},t)$. Here $r_{\rm esc}$ and $t_{\rm esc}$
are defined where a photon emitted at time t by the shell ejected at time
$t_{\rm inj}=t_{\rm j}$ escapes the outflow, i.e. reaches the first shell
emitted at time $t_{\rm inj}=0$. In other words, the optical depth from
$r(t_{\rm inj}=t_{\rm j},t)$ to $r_{\rm esc}(t_{\rm inj}=t_{\rm j},t)$ is
unity.
A geometric sketch illustrating the escape path of photons inside the ejecta
is shown in Fig. 1b. A photon emitted at point E (at a distance $r(t_{\rm
inj}=t_{\rm j},t)$ and propagation angle $\varphi$) will escape from the
ejecta at point G, such that the optical depth from E to G is
$\tau(t_{\rm j},\varphi)=\int^{r_{\rm esc}(t_{\rm inj}=t_{\rm j},t)}_{r(t_{\rm
inj}=t_{\rm j},t)}d\tau(r),$ (1)
where $d\tau(r)$ can be estimated as (Abramowicz et al. 1991; Daigne &
Mochkovitch 2002; Pe′er 2008)
$d\tau(r)=\frac{\kappa\dot{M}(1-\beta{\rm cos}\varphi)}{4\pi r^{2}}dr.$ (2)
We define the photospheric radius $R_{\rm PS}(t_{\rm j})=r(t_{\rm j},t)$ at
which $\tau$ is equal to unity. In light of Fig. 1b, we choose a cylindrical
coordinate system (Pe′er 2008) with the central point $O$ being the stellar
center, and the observer located along the +z-direction (defined by the
direction of $\overline{OD}$). Photons are emitted at a perpendicular distance
$r_{\rm min}=r(t_{\rm j},t){\rm sin}\varphi$ from the z-axis and a distance
$z_{\rm min}=r(t_{\rm j},t){\rm cos}\varphi$ along the z-axis from point $O$.
The escape radius can be estimated from the triangle $OEG$, i.e., $r_{\rm
esc}(t_{\rm inj}=t_{\rm j},t)=r(t_{\rm j},t)+\beta c(t_{\rm j}-0)+\beta
c(t_{\rm esc}-t)=\\{r(t_{\rm j},t)^{2}+[c(t_{\rm esc}-t)]^{2}-2r(t_{\rm
j},t)c(t_{\rm esc}-t){\rm cos}(\pi-\varphi)\\}^{\frac{1}{2}}$,
The integration for the optical depth can be conveniently rewritten in
cylindrical coordinates as the following
$\tau(t_{\rm j},\varphi)=\int^{z_{\rm max}}_{z_{\rm
min}}\frac{\kappa\dot{M}(1-\beta{\rm cos}\varphi)}{4\pi
r^{2}}\frac{dr}{dz}dz,$ (3)
where $r=\sqrt{z^{2}+r_{\rm min}^{2}}$, $dz/dr=z/\sqrt{z^{2}+r_{\rm min}^{2}}$
and $z_{\rm max}=\sqrt{r_{\rm esc}^{2}-r_{\rm min}^{2}}$. The photospheric
radius ($R_{\rm PS}$), which depends on the propagation angle ($\varphi$), can
be found readily by defining $\tau=1$.
For a relativistic ejecta with Lorentz factor $\Gamma$, the arrival time of
photons emitted at the photosphere in the observer frame are delayed relative
to that measured in the stellar frame
$t_{\rm obs}=t-R_{\rm PS}{\rm cos}\varphi/c=t_{\rm j}+(1-\beta{\rm
cos}\varphi)R_{\rm PS}/\beta c,$ (4)
i.e., the observer time is a function of injection time and propagation angle.
In this equation, we have neglected the effect of the acceleration radius
($r_{\rm acc}$) because it is much smaller than the photosphere radius.
The evolution of photospheric radius with propagation angle is shown in Fig.
2b. The parameters are taken as $\Gamma=10$, $\dot{E}=10^{51}$erg/s and
$t_{\rm w}=2000$s. The solid curve shows its evolution at observer time
$t_{\rm obs}=100$ s and the dashed curve at $t_{\rm obs}=2000$ s. From this
panel, we can see that in the observer frame, the photospheric radius
decreases with propagation angle.
Figure 2: Evolution of the photospheric radius and temperature with respect
to the observer time and photon propagation angle. The parameters are taken as
$\Gamma=10$, $\dot{E}=10^{51}$erg/s and $t_{\rm w}=2000$s.
(a) The photospheric radius vs. the observer time. The solid curve is plot for
photons propagating along the expansion direction of the ejecta ($\varphi=0$).
The dashed curve corresponds to $\varphi=0.1$ rad. (b) The photospheric radius
vs. the photon propagation angle. The solid and dashed curves correspond to
the observer time of 100 s and 2000 s respectively. (c) The photospheric
temperature vs. observer time. The solid and dashed curves correspond to
photon propagation angles of 0 and 0.1 rad, respectively. (d) The photospheric
temperature vs. the photon propagation angle. The solid and dashed curves are
for $t_{\rm obs}$=100 s and 2000 s, respectively.
We also show the evolution of photospheric radius with observer time in Fig.
2a. The solid, dashed curves present the cases for $\varphi=0$ and
$\varphi=0.1$ rad respectively. As is shown in this panel, the photospheric
radius increases with time, and the duration of the photospheric emission is
prolonged at larger propagation angle. The end points of the two curves
indicate the observer time when the last layer of the ejecta becomes
transparent.
According to the fireball model, the temperature of a layer at its
photospheric radius is given by (Piran 1999)
$kT_{\rm PS}=\frac{D}{\Gamma}kT^{0}\big{(}\frac{R_{\rm PS}}{r_{\rm
acc}}\big{)}^{-2/3},$ (5)
where $D=[\Gamma(1-\beta{\rm cos}\varphi)]^{-1}$ is the Doppler factor,
$r_{\rm acc}$ is the saturation radius and $T^{0}$ is its blackbody
temperature. In Fig. 2c and Fig. 2d, we show the evolution of $T_{\rm PS}$
with respect to the propagation angle and observer time respectively. From the
two panels, we find that the photosperic temperature decreases with the
propagation angle and time in the observer frame.
The evolution of the injection time with respect to the observer time and the
propagation angle is shown in Fig. 3a and 3b, respectively. For the “standard”
parameter set ($\Gamma=10$, $\dot{E}=10^{51}$ erg/s, $t_{\rm w}$=2000 s),
$t_{\rm inj}(t_{\rm obs})$ is mildly smaller than $t_{\rm obs}$ for different
$\varphi$, and $t_{\rm inj}(\varphi)$ is almost independent of $\varphi$ for
different $t_{\rm obs}$. In Fig. 3c and Fig. 3d, we also show the evolution of
the escaping radius $r_{\rm esc}$ with respect to the observer time and the
propagation angle, respectively. The relations of $t_{\rm obs}$ and $\varphi$
with respect to $r_{\rm esc}$ are similar to that of the photospheric radius
$R_{\rm PS}$.
Figure 3: Evolution of the injection time and escaping radius with respect to
the observer time and photon propagation angle. The parameters are the same as
those in Fig. 2. (a) The injection time vs. the observer time. The solid curve
is plot for photons propagating along the expansion direction of the ejecta
($\varphi=0$). The dashed curve is the evolution of the injection time for
$\varphi=0.1$ rad. (b) The injection time vs. the photon propagation angle.
The solid and dashed curves are observed at 100s and 2000s respectively. (c)
The escaping radius vs. the observer time. The solid and dashed curves are
plot with a propagation angle of 0 and 0.1 rad, respectively. (d) The escaping
radius vs. the photon propagation angle. The solid and dashed curves are for
$t_{\rm obs}$=100 s and 2000 s, respectively.
As for a jet with half-opening angle $\theta=0.1$ rad, constant Lorentz factor
$\Gamma=10$, and energy injection from $t_{\rm inj}=0$ to $t_{\rm inj}=t_{\rm
w}=2000$ s with energy injection rate per solid angle
$\dot{E}/4\pi=10^{51}/4\pi$ erg/s, we can estimate that $r_{\rm acc}\simeq
9\times 10^{7}$cm and $kT^{0}\simeq 0.41$ MeV for a fireball model (Piran
1999; $\rm M\acute{e}sz\acute{a}ros~{}\&~{}Rees$ 2000; Daigne & Mochkovitch
2002). If the line-of-sight is along the jet central axis, we can find that
the photospheric radius is about $1.1\times 10^{14}$ cm when the last layer
becomes transparent, the observer’s time can be calculated from Eq. 4, which
is found to be about 2020 s and corresponds to the end point of the solid
curve in Fig. 2a.
The observed luminosity of photospheric emission can be determined by
integrating over the surface of the photosphere
$L=\int_{0}^{\theta}\sigma T_{\rm PS}^{4}dS{\rm cos}\vartheta$ (6)
where $\vartheta$ is the angle between the tangential direction of the
photosphere surface and the line-of-sight when the propagation angle is
$\varphi$. $dS{\rm cos}\vartheta=2\pi R_{\rm PS}(t_{\rm obs},\varphi){\rm
sin}\varphi[R_{\rm PS}(t_{\rm obs},\varphi+d\varphi){\rm
sin}(\varphi+d\varphi)-R_{\rm PS}(t_{\rm obs},\varphi){\rm sin}\varphi]$ is
the photospheric surface area from propagation angle $\varphi$ to
$\varphi+d\varphi$. The evolution of the photospheric luminosity with observer
time is shown by the solid curve in Fig. 4. There is a break in the light
curve at about $t_{\rm obs}\simeq$ 2020 s, which is attributed to the stop of
energy injection by the central engine and when the last layer of the ejecta
became transparent as the photons propagate along the line-of-sight.
Afterwards, only photospheric emission at high latitude (large propagation
angles) contributes to the observed luminosity. The photospheric emission
ceases when the last layer with propagation angle $\varphi=\theta=0.1$ become
transparent, which is about 2050 s in the observer frame.
Figure 4: Evolution of the photospheric luminosity (solid curve) and
effective temperature (dashed curve) with observer time for a jet with
parameters of $\theta=0.1$ rad, $\Gamma=10$, $\dot{E}=10^{51}$erg/s and
$t_{\rm w}$=2000 s.
We can also define an effective temperature for the photosphere
$T_{\rm eff}=\frac{\int_{0}^{\theta}TdL}{\int_{0}^{\theta}dL}.$ (7)
This effective temperature is shown as a dashed curve in Fig. 4. We can find
that the photospheric emission of a failed GRB is presented as a short soft
X-ray burst and then becomes a UV burst which lasts for about several thousand
seconds.
We also investigated the parameter effect on the photospheric emission, which
are shown in Fig. 5 and Fig. 6. All the curves are derived when the last layer
of ejecta along the line-of-sight became transparent , i.e., $\varphi=0$ and
$t_{\rm j}=t_{\rm w}$. As is shown in Fig. 5a, the photospheric radii are
decreasing with the increase of Lorentz factor for different sets of
parameters. The solid curve corresponds to the standard parameters
($\dot{E}=10^{51}$erg/s, $t_{\rm w}=2000$ s), while the parameters for the
dashed curve and the dotted curve are $\dot{E}=10^{49}$erg/s, $t_{\rm w}=2000$
s and $\dot{E}=10^{51}$erg/s, $t_{\rm w}=200$ s, respectively. A lower energy
injection rate and shorter injection time will decease the radius of the
photosphere. Fig. 5b shows the evolution of the observer time with Lorentz
factor. The parameters for each curve are the same as Fig. 5a. This time
period can be interpreted as the duration of the photospheric emission. From
this panel, we can find that the duration of photospheric emission decreases
with an increase of Lorentz factor. Lower energy injection rate and shorter
injection time will decease the duration of the photospheric emission.
Figure 5: Parameter dependence of the photospheric emission. All curves are
obtained when the last layer of the ejecta became transparent along the line-
of-sight, i.e., $\varphi=0$ and $t_{\rm j}=t_{\rm w}$. (a) Parameter
dependence of the photospheric radius. The solid curve is derived using the
standard parameters, i.e., $\dot{E}=10^{51}$erg/s and $t_{\rm w}$=2000 s; the
dashed curve is for $\dot{E}=10^{49}$erg/s and $t_{\rm w}$=2000 s; and the
dotted curve is for $\dot{E}=10^{51}$erg/s and $t_{\rm w}$=200 s. The
evolution of the internal shock radii $R_{\rm IS}$ for a variability timescale
of $\delta t$=1 s, 0.33 s and 0.1 s are shown and marked correspondingly. (b)
Parameter dependence of the observer time. The identities of the curves are
the same as in (a).
In Fig. 6, we show the parameter effect on the photospheric luminosity and
effective temprature. The parameters of each curve are the same as Fig. 5.
From Fig. 6a, we can find that the luminosity of the photospheric emission are
low in both high and low Lorentz factor. Lower energy injection rate results
in lower luminosity. As is shown in Fig. 6b, the effective temperatures are
increasing with the increase of Lorentz factor. Lower energy injection rate
and shorter injection time will decease the radii of the photosphere and hence
results in a higher effective temperature.
Figure 6: Parameter dependence of the photospheric luminosity (a) and the
effective temperature (b) in the Lorentz factor space. All curves are obtained
when the last layer of ejecta became transparent along the line-of-sight,
i.e., $\varphi=0$ and $t_{\rm j}=t_{\rm w}$. The identities of the curves are
the same as in Figure 5.
As for a jet with half-opening angle of about 0.1 rad, Lorentz factor
$\Gamma=2-20$, energy injection rate $\dot{E}=10^{49-51}$erg/s and injection
time $t_{\rm w}$=200-2000s, from Fig. 5 and 6, we can conclude that the prompt
emission for a failed GRB is thermal soft X-ray or UV photospheric emission,
there will be no significant non-thermal gamma-ray emission. The photospheric
luminosity is about $10^{46}$erg/s and last for about one thousand seconds.
Note that the photospheric luminosity is far lower than the energy injection
power, most of of the energy is re-converted into the ejecta’s kinetic energy.
From Fig. 5a, we find that the radius of the photosphere is about $10^{14}$
cm, which is larger then the prediction for the internal shock’s radius, i.e.,
$R_{\rm IS}\simeq\Gamma^{2}c\delta t\simeq 10^{12}$ cm (Mészáros 2006), where
$\delta t\sim 0.33$ s here is the variability timescale of the prompt
emission. The evolution of $R_{\rm IS}$ with respect to $\Gamma$ for $\delta
t$=1 s, 0.33 s and 0.1 s are shown and marked in Fig. 5a correspondingly. This
radius is consistent with Fig. 1a.
This thermal radiation will be in the UV or soft X-ray band. The lower band of
Swift-XRT ($0.2-10$ keV) may cover this energy range and it is sensitive
enough to detect such a photospheric emission component (Gehrels et al. 2004).
MAXI-SSC monitors all-sky in the energy range of $0.5-10$ keV and also has a
chance to detect such events (Matsuoka et al. 1997). Future UV satellites may
be also have the capability to detect these events, such as TAUVEX (wavelength
range 120nm-350nm) (Safonova et al. 2008).
## 3 Afterglow emission
As the outflow expands outward, it will collide with the surrounding medium
and afterglow will be produced. The dynamical evolution of a relativistic jet
in interstellar medium has been studied by Huang et al. (1999). Their codes
can be used in both ultra-relativistic and non-relativistic phases.
In our model, we consider a jet with the bulk Lorentz factor $\Gamma=10$, the
half-opening angle $\theta=0.1$ and the isotropic energy $E=10^{50}$ erg. The
jet expands laterally at the comoving sound speed and collides with a medium
whose number density is $n_{\rm ISM}=1$ cm-3. We also assume typical values
for some other parameters of the jet, i.e., the electron energy fraction
$\epsilon_{e}=0.1$, the magnetic energy fraction $\epsilon_{B}=0.01$ and the
power-law index of the energy distribution function of electrons $p=2.5$.
Multiband afterglow emission is expected from synchrotron radiation of
relativistic electrons. Using this exquisite model, we numerically calculated
the afterglow light curves and spectra with line of sight parallel to the jet
axis. We assume a redshift $z=1$ and a standard cosmology with
$\Omega_{M}=0.27$, $\Omega_{\Lambda}=0.73$ and with the Hubble constant of
$H_{0}=71$ km s-1 Mpc-1.
Our results for the afterglow spectra of failed GRBs are shown in Fig. 7. The
thick and thin solid curves are the spectra observed at $10^{3}$ s and
$10^{6}$ s respectively. In this figure, we can find that the spectrum becomes
softer with the elapse of the observational time. At the early stage ($10^{3}$
s), the peak flux appears at about $5\times 10^{12}$ Hz, i.e., in the IR band.
Both the peak flux and peak frequency decrease with time. At late time
($10^{6}$ s), the peak flux is more than one magnitude less than that in the
early stage. The peak frequency decreases to about $10^{9}$ Hz at late stage.
It is in the radio band and more than three magnitudes less than the peak
frequency at the early stage.
Figure 7: Evolution of the afterglow spectra for the three types of GRBs. The
solid, dashed and dotted curves are spectra of a failed GRB afterglow, an
ordinary GRB afterglow and a beaming-induced orphan afterglow respectively.
The thick and thin curves are the spectra observed at $10^{3}$ s and $10^{6}$
s respectively.
For comparison, we also show the afterglow spectra of an ordinary GRB in Fig.
7. Here we choose the same parameters as the failed GRB except for a much
larger bulk Lorentz factor ($\Gamma=300$). The dashed curves show the
afterglow spectra of this GRB with observing angle $\theta_{\rm obs}=0$ (the
line of sight is parallel to the jet axis). As is shown in Fig. 7, the spectra
of the failed GRB and the ordinary GRB are similar at the late stage (thin
solid and thin dashed curves) because their energies and Lorentz factors are
both similar at this moment. But at early stage, they are very different
(thick solid and thick dashed curves) due to their very different initial
Lorentz factors and the corresponding minimum Lorentz factors of electrons
(Sari et al. 1998). The peak frequency of the GRB afterglow is much larger
than that of the failed GRB afterglow.
In Fig. 7, we also show the spectra of a beaming-induced orphan afterglow
(afterglow from an ordinary highly collimated GRB outflow, but with the
observing angle larger than the jet half-opening angle so that no prompt
gamma-rays can be observed in the main burst phase, Rhoads 1997; Huang et al.
2002). Here we assume the same parameters as the ordinary GRB except
$\theta_{\rm obs}=0.125$. The early and late spectra of this orphan afterglow
are shown in Fig. 7 with thick dotted curve and thin dotted curve. From this
figure, we can find the spectra of beaming-induced orphan afterglow are
similar to that of the ordinary GRB. Although it is hard to distinguish a
beaming-induced orphan afterglow from a failed GRB afterglow through their
afterglow light curves (Huang et al. 2002), they can be potentially
distinguished from their spectra at the early stages. Their spectra of early
afterglows are very different: the peak frequency of a failed GRB afterglow is
far lower than that of a beaming-induced orphan afterglow. Another way to
distinguish them is through their early light curves. Early afterglow of a
beaming-induced orphan afterglow will show a rebrightening while failed GRB
will not (Huang et al. 1999, 2002; Xu & Huang 2010).
## 4 Conclusion and Discussions
The analysis in this paper shows that the emission of ejecta with low Lorentz
factors is very different from that expected from ejecta with high Lorentz
factors. Prompt emission of a GRB is non-thermal and bright in the gamma-ray
band. For a failed GRB, however, the emission originates from the photosphere
with a thermal spectrum, and is bright in the UV or soft X-ray band instead of
gamma-rays. This photospheric emission lasts for about a thousand seconds with
a luminosity about several times $10^{46}$ erg/s.
Since the photospheric emission manifests as a UV or soft X-ray transient, it
can be detected by some current and future satellites, such as Swift-XRT,
MAXI-SSC and TAUVEX etc. On 2008 January 9, Swift-XRT discovered a peculiar
X-ray transient 080109 in NGC 2770 (Berger & Soderberg 2008; Page et al.
2008). No gamma-ray emission was detected. This X-ray transient reached its
peak at about $60s$ and lasted for about $600$ s. Its spectrum can be fitted
with an absorbed double blackbody model with temperatures about $0.36$ keV and
$1.24$ keV respectively (Li 2008). This transient may be a candidate of
photospheric emission from a failed GRB. Meanwhile, some unidentified X-ray
transients have been detected by MAXI during its one-year monitoring (Nakajima
et al. 2009; Suzuki et al. 2010). These transients generally showed an
absorbed blackbody spectrum and lasted for tens of seconds. It is possible
that some of them are photospheric emission from failed GRBs.
If we extend the injection time to about $10^{5}$ s and the jet half-opening
to about 0.4 rad in our model, we find that the photospheric radius is about
$10^{15}$ cm and the effective temperature is deceased to lower than 1 eV,
i.e., there will be an optical burst. This kind of optical burst will last for
about several thousand seconds with a luminosity about $10^{42}$ erg/s, which
may be detected by the Hyper-Suprime Camera of the $Subaru$ telescope in the
future.
In this work, we have assumed that the prompt emission is thermal radiation
coming from the photosphere where the optical depth is unity. Due to the low
density of GRB jets, however, it has been pointed out that the last-scattering
positions of the observed photons may not simply coincide with the
photosphere, but instead possess a finite distribution around it (e.g. Pe′er
et al. 2006; Pe′er 2008; Beloborodov 2010; Pe′er & Ryde 2011). This stochastic
effect can lead to differentiation of the observed spectrum from a thermal one
of purely photospheric origin. Such mechanism can work even in failed GRBs,
and it is our future work to study how the spectrum will be reshaped using
Monte-Carlo calculations. We are planning to investigate this effect in the
context of failed GRBs as a next step of our study.
From the comparison of afterglow emissions from failed and ordinary GRBs,
while we find it challenging to distinguish them at their late stage of
evolution, their spectra at the early stage are profoundly different. We
conclude that it is possible to identify failed GRBs by observing their
afterglow emission in the early stage. The typical frequency at peak flux in
the afterglow phase for failed GRBs is much lower than that for ordinary GRBs
(or beaming-induced orphan GRBs). We can thus define a hardness ratio, for
instance, as the flux contrast between $10^{12}$ Hz and $10^{14}$ Hz at an
observed time of $1000$ s, i.e., $f_{\rm 1ks}\equiv F_{10^{12}{\rm
Hz}}/F_{10^{14}{\rm Hz}}$. If $f_{\rm 1ks}>1$, then it is quite likely that
the emission is coming from a failed GRB. If $f_{\rm 1ks}<1$, then it would be
more likely to come from an ordinary GRB afterglow or a beaming-induced orphan
afterglow. In addition, at the early afterglow stage, a rebrightening phase
will be present in the case of a beaming-induced orphan GRB, while it is not
expected for ordinary or failed GRBs. Therefore, the afterglows of failed GRBs
can be distinguished from both ordinary GRB afterglows and beaming-induced
orphan afterglows through observations at the early stages.
We thank the anonymous referee for many of the useful suggestions and
comments. We also would like to thank P. Mészáros, T. Piran and J. Aoi for
helpful discussions. This work was supported by the National Natural Science
Foundation of China (Grant No. 11033002), the National Basic Research Program
of China (973 Program, Grant No. 2009CB824800) and the Grant-in-Aid for the
’Global COE Bilateral International Exchange Program’ of Japan, Grant-in-Aid
for Scientific Research on Priority Areas No. 19047004 and Scientific Research
on Innovative Areas No. 21105509 by Ministry of Education, Culture, Sports,
Science and Technology (MEXT), Grant-in-Aid for Scientific Research (S) No.
19104006 and Scientific Research (C) No. 21540404 by Japan Society for the
Promotion of Science (JSPS).
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|
arxiv-papers
| 2011-11-30T08:54:53 |
2024-09-04T02:49:24.797762
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Xu, S. Nagataki, Y.-F. Huang, and S.-H. Lee",
"submitter": "Ming Xu",
"url": "https://arxiv.org/abs/1111.7084"
}
|
1111.7114
|
# Novel Fermi Liquid of 2D Polar Molecules
Zhen-Kai Lu1,2,3 and G. V. Shlyapnikov2,4 1Max-Planck-Institut für
Quantenoptik, Hans-Kopfermann-Straße 1, 85748 Garching, Germany
2Laboratoire de Physique Théorique et Modèles Statistiques, CNRS and
Université Paris Sud, UMR8626, 91405 Orsay, France
3 Département de Physique, École Normale Supérieure, 75005, Paris, France
4Van der Waals-Zeeman Institute, University of Amsterdam, Science Park 904,
1098 XH Amsterdam, The Netherlands
###### Abstract
We study Fermi liquid properties of a weakly interacting 2D gas of single-
component fermionic polar molecules with dipole moments $d$ oriented
perpendicularly to the plane of their translational motion. This geometry
allows the minimization of inelastic losses due to chemical reactions for
reactive molecules and, at the same time, provides a possibility of a clear
description of many-body (beyond mean field) effects. The long-range character
of the dipole-dipole repulsive interaction between the molecules, which scales
as $1/r^{3}$ at large distances $r$, makes the problem drastically different
from the well-known problem of the two-species Fermi gas with repulsive
contact interspecies interaction. We solve the low-energy scattering problem
and develop a many-body perturbation theory beyond the mean field. The theory
relies on the presence of a small parameter $k_{F}r_{*}$, where $k_{F}$ is the
Fermi momentum, and $r_{*}=md^{2}/\hbar^{2}$ is the dipole-dipole length, with
$m$ being the molecule mass. We obtain thermodynamic quantities as a series of
expansion up to the second order in $k_{F}r_{*}$ and argue that many-body
corrections to the ground-state energy can be identified in experiments with
ultracold molecules, like it has been recently done for ultracold fermionic
atoms. Moreover, we show that only many-body effects provide the existence of
zero sound and calculate the sound velocity.
## I Introduction
The recent breakthrough in creating ultracold diatomic polar molecules in the
ground ro-vibrational state Ni ; Deiglmayr2008 ; Inouye ; Nagerl and cooling
them towards quantum degeneracy Ni has opened fascinating prospects for the
observation of novel quantum phases Baranov2008 ; Lahaye2009 ; Pupillo2008 ;
Wang2006 ; Buchler2007 ; Taylor ; Cooper2009 ; Pikovski2010 ; Lutchyn ; Potter
; Barbara ; Sun ; Miyakawa ; Gora ; Ronen ; Parish ; Babadi ; Baranov . A
serious problem in this direction is related to ultracold chemical reactions,
such as KRb+KRb$\Rightarrow$K2+Rb2 observed in the JILA experiments with KRb
molecules Jin ; Jin2 , which places severe limitations on the achievable
density in three-dimensional samples. In order to suppress chemical reactions
and perform evaporative cooling, it has been proposed to induce a strong
dipole-dipole repulsion between the molecules by confining them to a
(quasi)two-dimensional (2D) geometry and orienting their dipole moments (by a
strong electric field) perpendicularly to the plane of the 2D translational
motion Bohn1 ; Baranov1 . The suppression of chemical reactions by nearly two
orders of magnitude in the quasi2D geometry has been demonstrated in the
recent JILA experiment Ye . At the same time, not all polar molecules of
alkali atoms, on which experimental efforts are presently focused, may undergo
these chemical reactions Jeremy . In particular, they are energetically
unfavorable for RbCs bosonic molecules obtained in Innsbruck Nagerl , or for
NaK and KCs molecules which are now being actively studied by several
experimental groups (see, e.g. MITZ ). It is thus expected that future
experimental studies of many-body physics will deal with non-reactive polar
molecules or with molecules strongly confined to the 2D regime.
Therefore, the 2D system of fermionic polar molecules attracts a great deal of
interest, in particular when they are in the same internal state. Various
aspects have been discussed regarding this system in literature, in particular
the emergence and beyond mean field description of the topological
$p_{x}+ip_{y}$ phase for microwave-dressed polar molecules Cooper2009 ; Gora ,
interlayer superfluids in bilayer and multilayer systems Pikovski2010 ; Ronen
; Potter ; Zinner , the emergence of density-wave phases for tilted dipoles
Sun ; Miyakawa ; Parish ; Babadi ; Baranov . The case of superfluid pairing
for tilted dipoles in the quasi2D geometry beyond the simple BCS approach has
been discussed in Ref. Baranov . The Fermi liquid behavior of this system has
been addressed by using the Fourier transform of the dipole-dipole interaction
potential Das1 ; Das2 ; Miyakawa ; Taylor ; Pu ; Das3 ; Baranov and then
employing various types of mean field approaches, such as the Hartree-Fock
approximation Miyakawa or variational approaches Taylor ; Pu . It should be
noted, however, that the short-range physics can become important for the
interaction between such polar molecules, since in combination with the long-
range behavior it introduces a peculiar momentum dependence of the scattering
amplitude Gora .
On the other hand, there is a subtle question of many-body (beyond mean field)
effects in the Fermi liquid behavior of 2D polar molecules, and it can be
examined in ultracold molecule experiments. For the case of atomic fermions, a
milestone in this direction is the recent result at ENS, where the experiment
demonstrated the many-body correction to the ground state energy of a short-
range interacting two-species fermionic dilute gas Salomon1 ; Salomon2 . This
correction was originally calculated by Huang, Lee, and Yang Huang ; Lee by
using a rather tedious procedure. Later, it was found by Abrikosov and
Khalatnikov Abr in an elegant way based on the Landau Fermi liquid theory
Landau .
In this paper, we study a weakly interacting 2D gas of fermionic polar
molecules which are all in the same internal state. It is assumed that each
molecule has an average dipole moment $d$ which is perpendicular to the plane
of the translational motion, so that the molecule-molecule interaction at
large separations $r$ is
$U(r)=\frac{d^{2}}{r^{3}}=\frac{\hbar^{2}r_{*}}{mr^{3}},$ (1)
where $r_{*}=md^{2}/\hbar^{2}$ is the characteristic dipole-dipole distance,
and $m$ is the molecule mass. The value of $d$ depends on the external
electric field. At ultralow temperatures that are much smaller than the Fermi
energy, characteristic momenta of particles are of the order of the Fermi
momentum $k_{F}$, and the criterion of the weakly interacting regime is:
$k_{F}r_{*}\ll 1.$ (2)
As a consequence, the Fermi liquid properties of this system, such as the
ground state energy, compressibility, effective mass, can be written as a
series of expansion in the small parameter $k_{F}r_{*}$. We obtain explicit
expressions of these quantities up to the second order in $k_{F}r_{*}$, which
requires us to reveal the role of the short-range physics in the scattering
properties and develop a theory beyond the mean field. Our analysis shows that
only many-body (beyond mean field) effects provide the existence of undamped
zero sound in the collisionless regime.
The paper is organized as follows. In Section II we analyze the low-energy 2D
scattering of the polar molecules due to the dipole-dipole interaction. We
obtain the scattering amplitude for all scattering channels with odd orbital
angular momenta. The leading part of the amplitude comes from the so-called
anomalous scattering, that is the scattering related to the interaction
between particles at distances of the order of their de Broglie wavelength.
This part of the amplitude corresponds to the first Born approximation and,
due to the long-range $1/r^{3}$ character of the dipole-dipole interaction, it
is proportional to the relative momentum $k$ of colliding particles for any
orbital angular momentum $l$. We then take into account the second Born
correction, which gives a contribution proportional to $k^{2}$. For the
$p$-wave scattering channel it is necessary to include the short-range
contribution, which together with the second Born correction leads to the term
behaving as $k^{2}\ln{k}$. In Section III, after reviewing the Landau Fermi
liquid theory for 2D systems, we specify two-body (mean field) and many-body
(beyond mean field) contributions to the ground state energy for 2D fermionic
polar molecules in the weakly interacting regime. We then calculate the
interaction function of quasiparticles on the Fermi surface and, following the
idea of Abrikosov-Khalatnikov Abr , obtain the compressibility, ground state
energy, and effective mass of quasiparticles. In Section IV we calculate the
zero sound velocity and stress that the many-body contribution to the
interaction function of quasiparticles is necessary for finding the undamped
zero sound. We conclude in Section V, emphasizing that the 2D gas of fermionic
polar molecules represents a novel Fermi liquid, which is promising for
revealing many-body effects. Moreover, we show that with present facilities it
is feasible to obtain this system in both collisionless and hydrodynamic
regimes.
## II Low-energy scattering of fermionic polar molecules in 2D
### II.1 General relations
We first discuss low-energy two-body scattering of identical fermionic polar
molecules undergoing the 2D translational motion and interacting with each
other at large separations via the potential $U(r)$ (1). The term low-energy
means that their momenta satisfy the inequality $kr_{*}\ll 1$. In order to
develop many-body theory for a weakly interacting gas of such molecules, we
need to know the off-shell scattering amplitude defined as
$f(\mathbf{k}^{\prime},\mathbf{k})=\int\exp(-i\mathbf{k}^{\prime}\mathbf{r})U(r)\tilde{\psi}_{\mathbf{k}}(\mathbf{r})d^{2}\mathbf{r},$
(3)
where $\tilde{\psi}_{\mathbf{k}}(\mathbf{r})$ is the true wavefunction of the
relative motion with momentum $\mathbf{k}$. It is governed by the Schrödinger
equation
$\left(-\frac{\hbar^{2}}{m}\Delta+U(r)\right)\tilde{\psi}_{\mathbf{k}}(\mathbf{r})=\frac{\hbar^{2}k^{2}}{m}\tilde{\psi}_{\mathbf{k}}(\mathbf{r}).$
(4)
For $|{\bf k}^{\prime}|=|{\bf k}|$ we have the on-shell amplitude which enters
an asymptotic expression for $\psi_{\mathbf{k}}(\mathbf{r})$ at
$r\rightarrow\infty$ Lan2 ; Gora :
$\tilde{\psi}_{\mathbf{k}}(\mathbf{r})=\exp(i{\bf
kr})-\frac{m}{\hbar^{2}}\sqrt{\frac{i}{8\pi kr}}f(k,\varphi)\exp(ikr),$ (5)
with $\varphi$ being the scattering angle, i.e. the angle between the vectors
${\bf k}^{\prime}$ and ${\bf k}$.
The wavefunction $\tilde{\psi}_{{\bf k}}({\bf r})$ can be represented as a sum
of partial waves $\tilde{\psi}_{l}(k,r)$ corresponding to the motion with a
given value of the orbital angular momentum $l$:
$\tilde{\psi}_{{\bf k}}({\bf
r})=\sum_{l=-\infty}^{\infty}\tilde{\psi}_{l}(k,r)i^{l}\exp(il\varphi).$ (6)
Using the relation
$\exp(i{\bf
kr})=\sum_{l=-\infty}^{\infty}i^{l}J_{l}(kr)\exp[il(\varphi_{k}-\varphi_{r})],$
(7)
where $J_{l}$ is the Bessel function, and $\varphi_{k}$ and $\varphi_{r}$ are
the angles of the vectors ${\bf k}$ and ${\bf r}$ with respect to the
quantization axis. Eqs. (6) and (7) allow one to express the scattering
amplitude as a sum of partial-wave contributions:
$f(\mathbf{k}^{\prime},\mathbf{k})=\sum_{l=-\infty}^{\infty}\exp(il\varphi)f_{l}(k^{\prime},k),$
(8)
with the off-shell $l$-wave amplitude given by
$f_{l}(k^{\prime},k)=\int_{0}^{\infty}J_{l}(k^{\prime}r)U(r)\tilde{\psi}_{l}(k,r)2\pi
rdr.$ (9)
Similar relations can be written for the on-shell scattering amplitude:
$\displaystyle f(k,\varphi)=\sum_{l=-\infty}^{\infty}\exp(il\varphi)f_{l}(k),$
(10) $\displaystyle
f_{l}(k)=\int_{0}^{\infty}J_{l}(k^{\prime}r)U(r)\tilde{\psi}_{l}(k,r)2\pi
rdr.$ (11)
The asymptotic form of the wavefunction of the $l$-wave relative motion at
$r\rightarrow\infty$ may be represented as
$\tilde{\psi}_{l}(k,r)\propto\frac{\cos(kr-\pi/4+\delta_{l}(k))}{\sqrt{kr}},$
(12)
where $\delta_{l}(k)$ is the scattering phase shift. This is obvious because
in the absence of scattering the $l$-wave part of the plane wave $\exp(i{\bf
kr})$ at $r\rightarrow\infty$ is $(kr)^{-1/2}\cos(kr-\pi/4)$. Comparing Eq.
(12) with the $l$-wave part of Eq. (5) we obtain a relation between the
partial on-shell amplitude and the phase shift:
$f_{l}(k)=-\frac{4\hbar^{2}}{m}\frac{\tan\delta_{l}(k)}{1-i\tan\delta_{l}(k)}.$
(13)
Note that away from resonances the scattering phase shift is small in the low-
momentum limit $kr_{*}\ll 1$.
For the solution of the scattering problem it is more convenient to normalize
the wavefunction of the radial relative motion with orbital angular momentum
$l$ in such a way that it is real and for $r\rightarrow\infty$ one has:
$\displaystyle\psi_{l}(k,r)$
$\displaystyle=\left[J_{l}(kr)-\tan\delta_{l}(k)N_{l}(kr)\right]$ (14)
$\displaystyle\propto\cos(kr-l\pi/2-\pi/4+\delta_{l}(k)),$
where $N_{l}$ is the Neumann function. One checks straightforwardly that
$\tilde{\psi}_{l}(k,r)=\frac{\psi_{l}(k,r)}{1-i\tan\delta_{l}(k)}.$
Using this relation the off-shell scattering amplitude (9) can be represented
as
$f_{l}(k^{\prime},k)=\frac{{\bar{f}}_{l}(k^{\prime},k)}{1-i\tan\delta_{l}(k)},$
(15)
where ${\bar{f}}_{l}(k^{\prime},k)$ is real and follows from Eq. (9) with
$\tilde{\psi}_{l}(k,r)$ replaced by $\psi_{l}(k,r)$. Setting $k^{\prime}=k$ we
then obtain the related on-shell scattering amplitude:
${\bar{f}}_{l}(k,k)\equiv{\bar{f}}_{l}(k)=-\frac{4\hbar^{2}}{m}\tan\delta_{l}(k).$
(16)
### II.2 Low-energy $p$-wave scattering
As we will see, the slow $1/r^{3}$ decay of the potential $U(r)$ at
sufficiently large distances makes the scattering drastically different from
that of short-range interacting atoms. For identical fermionic polar
molecules, only the scattering with odd orbital angular momenta $l$ is
possible. For finding the amplitude of the $p$-wave scattering in the
ultracold limit, $kr_{*}\ll 1$, we employ the method developed in Ref. Gora
and used there for the scattering potential containing an attractive $1/r^{3}$
dipole-dipole tail. We divide the range of distances into two parts: $r<r_{0}$
and $r>r_{0}$, where $r_{0}$ is in the interval $r_{*}\ll r_{0}\ll k^{-1}$. In
region I where $r<r_{0}$, the $p$-wave relative motion of two particles is
governed by the Schrödinger equation with zero kinetic energy:
$-\frac{\hbar^{2}}{m}\left(\frac{d^{2}\psi_{I}}{dr^{2}}+\frac{1}{r}\frac{d\psi_{I}}{dr}-\frac{\psi_{I}}{r^{2}}\right)+U(r)\psi_{I}=0.$
(17)
At distances where the potential $U(r)$ already acquires the form (1), the
solution of Eq. (17) can be expressed in terms of growing and decaying Bessel
functions:
$\psi_{I}(r)\propto\left[AI_{2}\left(2\sqrt{\frac{r_{*}}{r}}\right)+K_{2}\left(2\sqrt{\frac{r_{*}}{r}}\right)\right].$
(18)
The constant $A$ is determined by the behavior of $U(r)$ at short distances
where Eq. (1) is no longer valid. If the interaction potential $U(r)$ has the
form (1) up to very short distances, then $A=0$, so that for $r\rightarrow 0$
equation (18) gives an exponentially decaying wavefunction.
It should be noted here that for the quasi2D regime obtained by a tight
confinement of the translational motion in one direction, we can encounter the
situation where $r_{*}\lesssim l_{0}$, with $l_{0}$ being the confinement
length. However, we may always select $r_{0}\gg l_{0}$ if the condition
$kl_{0}\ll 1$ is satisfied. Therefore, our results for the 2D $p$-wave
scattering obtained below in this section remain applicable for the quasi2D
regime. The character of the relative motion of particles at distances
$r\lesssim l_{0}$ is only contained in the value of the coefficient $A$, and
the extra requirement is the inequality $kl_{0}\ll 1$.
At large distances, $r>r_{0}$, the relative motion is practically free and the
potential $U(r)$ can be considered as perturbation. To zero order, the
relative wavefunction is given by
$\psi_{II}^{(0)}(r)=J_{1}(kr)-\tan\delta_{I}(k)N_{1}(kr),$ (19)
where the phase shift $\delta_{I}(k)$ is due to the interaction between
particles in region I. Equalizing the logarithmic derivatives of $\psi_{I}(r)$
and $\psi_{II}^{(0)}$ at $r=r_{0}$ we obtain:
$\\!\\!\tan\delta_{I}\\!=\\!-\frac{\pi
k^{2}r_{0}r_{*}}{8}\left[1\\!+\\!\frac{r_{*}}{r_{0}}\left(2C\\!-\\!\frac{1}{2}\\!-\\!2A\\!+\\!\ln\frac{r_{*}}{r_{0}}\right)\right],\\!\\!$
(20)
with $C=0.5772$ being the Euler constant.
We now include perturbatively the contribution to the $p$-wave scattering
phase shift from distance $r>r_{0}$. In this region, to first order in $U(r)$,
the relative wavefunction is given by
$\\!\psi^{(1)}_{II}(r)=\psi^{(0)}_{II}(r)\\!-\\!\int_{r_{0}}^{\infty}G(r,r^{\prime})U(r^{\prime})\psi^{(0)}_{II}(r^{\prime})2\pi
r^{\prime}dr^{\prime}\\!,$ (21)
where the Green function for the free $p$-wave motion obeys the radial
equation:
$\displaystyle-\frac{\hbar^{2}}{m}\left(\frac{d^{2}}{dr^{2}}+\frac{1}{r}\frac{d}{dr}-\frac{1}{r^{2}}+k^{2}\right)G(r,r^{\prime})=\frac{\delta(r-r^{\prime})}{2\pi
r}.$
For the normalization of the relative wavefunction chosen in Eq. (14), we
have:
$G(r,r^{\prime})=-\frac{m}{4\hbar^{2}}\begin{cases}\psi^{(0)}_{II}(r^{\prime})N_{1}(kr),&r>r^{\prime}\\\
\\\ \psi^{(0)}_{II}(r)N_{1}(kr^{\prime}).&r<r^{\prime}\end{cases}$ (22)
Substituting this Green function into Eq. (21) and taking the limit
$r\rightarrow\infty$, for the first order contribution to the phase shift we
have:
$\\!\\!\tan\delta_{1}^{(1)}(k)\\!=\\!\tan\delta_{I}(k)\\!-\\!\frac{m}{4\hbar^{2}}\int_{r_{0}}^{\infty}\\!\\![\psi^{(0)}_{II}(r)]^{2}U(r)2\pi
rdr.$ (23)
Using Eqs. (19) and (20) we then obtain:
$\\!\\!\\!\tan\delta_{1}^{(1)}(k)\\!=\\!-\frac{2kr_{*}}{3}\\!-\\!\frac{\pi
k^{2}r_{*}^{2}}{8}\left(\\!\\!-2A\\!+\\!2C\\!+\\!\ln\frac{r_{*}}{r_{0}}\\!-\\!\frac{3}{2}\\!\right).\\!$
(24)
To second order in $U(r)$, we have the relative wavefunction:
$\displaystyle\psi^{(2)}_{II}(r)$
$\displaystyle=\psi^{(1)}_{II}(r)+\int_{r_{0}}^{\infty}G(r,r^{\prime})U(r^{\prime})2\pi
r^{\prime}dr^{\prime}$ (25)
$\displaystyle\times\int_{r_{0}}^{\infty}G(r^{\prime},r^{\prime\prime})U(r^{\prime\prime})\psi^{(0)}_{II}(r^{\prime\prime})2\pi
r^{\prime\prime}dr^{\prime\prime}.$
Taking the limit $r\rightarrow\infty$ in this equation we see that including
the second order contribution, the scattering phase shift becomes:
$\displaystyle\tan\delta_{1}(k)$
$\displaystyle=\tan\delta^{(1)}(k)-\frac{m^{2}}{8\hbar^{4}}\int_{r_{0}}^{\infty}\psi^{(0)}_{II}(r)^{2}U(r)2\pi
rdr$ (26)
$\displaystyle\times\int_{r}^{\infty}N_{1}(kr^{\prime})U(r^{\prime})\psi^{(0)}_{II}(r^{\prime})2\pi
r^{\prime}dr^{\prime}.$
As we are not interested in terms that are proportional to $k^{3}$ or higher
powers of $k$, we may omit the term $\tan\delta_{I}(k)N_{1}(kr)$ in the
expression for $\psi^{(0)}_{II}(r)$. Then the integration over $dr^{\prime}$
leads to:
$\displaystyle\tan\delta_{1}(k)=\tan\delta_{1}^{(1)}(k)-\frac{(\pi
kr_{*})^{2}}{2}\int_{kr_{0}}^{\infty}\frac{J^{2}_{1}(x)}{x^{2}}dx$
$\displaystyle\times\Big{[}\frac{2}{3}x\left(N_{0}(x)J_{2}(x)-N_{1}(x)J_{1}(x)\right)$
$\displaystyle\;\;\;\;-\frac{1}{2}N_{0}(x)J_{1}(x)+\frac{1}{6}N_{1}(x)J_{2}(x)-\frac{1}{\pi
x}\Big{]}.$ (27)
For the first four terms in the square brackets, we may put the lower limit of
integration equal to zero and use the following relations:
$\displaystyle\int_{0}^{\infty}J^{3}_{1}(x)N_{1}(x)\frac{dx}{x}=-\frac{1}{4\pi},$
$\displaystyle\int_{0}^{\infty}J^{2}_{1}(x)J_{2}(x)N_{0}(x)\frac{dx}{x}=\frac{1}{8\pi},$
$\displaystyle\int_{0}^{\infty}J^{3}_{1}(x)N_{0}(x)\frac{dx}{x^{2}}=\frac{1}{16\pi},$
$\displaystyle\int_{0}^{\infty}J^{2}_{1}(x)J_{2}(x)N_{1}(x)\frac{dx}{x^{2}}=-\frac{1}{16\pi}.$
For the last term in the square brackets we have:
$\int_{kr_{0}}^{\infty}J^{2}_{1}(x)\frac{dx}{x^{3}}\approx\frac{1}{16}-\frac{C}{4}+\frac{\ln
2}{4}-\frac{1}{4}\ln kr_{0}.$ (28)
We then obtain:
$\displaystyle\tan\delta_{1}(k)$
$\displaystyle=\tan\delta_{1}^{(1)}(k)-\frac{\pi(kr_{*})^{2}}{8}\left[\frac{7}{12}+C-\ln
2+\ln kr_{0}\right]$ $\displaystyle=-\frac{2kr_{*}}{3}-\frac{\pi
k^{2}r^{2}_{*}}{8}\ln\xi kr_{*},$ (29)
where:
$\xi=\exp\left(3C-\ln 2-\frac{11}{12}-2A\right).$ (30)
Using Eqs. (16) and (29) we represent the on-shell $p$-wave scattering
amplitude ${\bar{f}}_{1}(k)$ in the form:
${\bar{f}}_{1}(k)={\bar{f}}^{(1)}_{1}(k)+{\bar{f}}^{(2)}_{1}(k),$ (31)
with
${\bar{f}}^{(1)}_{1}(k)=\frac{8\hbar^{2}}{3m}kr_{*}$ (32)
and
${\bar{f}}^{(2)}_{1}(k)=\frac{\pi\hbar^{2}}{2m}(kr_{*})^{2}\ln\xi kr_{*}.$
(33)
The leading term is ${\bar{f}}^{(1)}_{1}(k)\propto k$. It appears to first
order in $U(r)$ and comes from the scattering at distances $r\sim 1/k$. This
term can be called “anomalous scattering” term (see Lan2 ). The term
$f^{(2)}_{1}(k)\propto k^{2}\ln\xi kr_{*}$ comes from both large distances
$\sim 1/k$ and short distances. Note that the behavior of the wavefunction at
short distances where $U(r)$ is no longer given by Eq. (1), is contained in
Eq. (29) only through the coefficient $\xi$ under logarithm.
### II.3 Scattering with $|l|>1$
The presence of strong anomalous $p$-wave scattering, i.e. the scattering from
interparticle distances $\sim 1/k$, originates from the slow $1/r^{3}$ decay
of the potential $U(r)$ at large $r$. The strong anomalous scattering is also
present for partial waves with higher $l$. In this section we follow the same
method as in the case of the $p$-wave scattering and calculate the amplitude
of the $l$-wave scattering with $|l|>1$. For simplicity we consider positive
$l$, having in mind that the scattering amplitude and phase shift depend only
on $|l|$.
To zero order in $U(r)$, the wavefunction of the $l$-wave relative motion at
large distances $r>r_{0}$ is written as:
$\psi^{(0)}_{l(II)}(k,r)=\left[J_{l}(kr)-\tan\delta_{l(I)}(k)N_{l}(kr)\right],$
(34)
where $\delta_{l(I)}(k)$ is the $l$-wave scattering phase shift coming from
the interaction at distances $r<r_{0}$. We then match
$\psi^{(0)}_{l(II)}(k,r)$ at $r=r_{0}$ with the short-distance wavefunction
$\psi_{l(I)}(r)$ which follows from the Schrödinger equation for the $l$-wave
relative motion in the potential $U(r)$ at $k=0$. This immediately gives a
relation:
$\tan\delta_{l(I)}(k)=\frac{kJ^{\prime}_{l}(kr_{0})-w_{l}J_{l}(kr_{0})}{kN^{\prime}_{l}(kr_{0})-w_{l}N_{l}(kr_{0})},$
(35)
where the momentum-independent quantity $w_{l}$ is the logarithmic derivative
of $\psi_{l(I)}(r)$ at $r=r_{0}$. Since we have the inequality $kr_{0}\ll 1$,
the arguments of the Bessel functions in Eq. (35) are small and they reduce to
$J_{l}(x)\sim x^{l}\;,\;N_{l}(x)\sim x^{-l}$. This leads to
$\tan\delta_{l(I)}(k)\sim(kr_{0})^{2l}$. Thus, the phase shift coming from the
interaction at short distances is of the order of $(kr_{0})^{2l}$. As we
confine ourselves to second order in $k$, we may put $\tan\delta_{l(I)}(k)=0$
for the scattering with $|l|>1$.
Then, like for the $p$-wave scattering, we calculate the contribution to the
phase shift from distances $r>r_{0}$ by considering the potential $U(r)$ as
perturbation. To first and second order in $U(r)$, at $r>r_{0}$ we have
similar expressions as Eq. (23), (25) for the relative wavefunction of the
$l$-wave motion. Following the same method as in the case of the $p$-wave
scattering and retaining only the terms up to $k^{2}$, for the first order
phase shift we have:
$\displaystyle\tan\delta_{l}^{(1)}(k)=-\frac{m}{4\hbar^{2}}\int_{r_{0}}^{\infty}[\psi^{(0)}_{l(II)}(r)]^{2}U(r)2\pi
rdr$ $\displaystyle\simeq-\frac{\pi
kr_{*}}{2}\int_{kr_{0}}^{\infty}J_{l}^{2}(x)\frac{1}{x^{2}}dx=-\frac{2kr_{*}}{4l^{2}-1}.$
(36)
The second order phase shift is:
$\displaystyle\tan\delta^{(2)}_{l}(k)=-\frac{m^{2}}{8\hbar^{4}}\int_{r_{0}}^{\infty}\psi^{(0)}_{l(II)}(r)^{2}U(r)2\pi
rdr$
$\displaystyle\;\;\times\int_{r}^{\infty}N_{l}(kr^{\prime})U(r^{\prime})\psi^{(0)}_{l(II)}(r^{\prime})2\pi
r^{\prime}dr^{\prime}$ $\displaystyle\simeq-\frac{(\pi
kr_{*})^{2}}{2}\int_{kr_{0}}^{\infty}\frac{J^{2}_{l}(x)}{x^{2}}dx\int_{x}^{\infty}\frac{N_{l}(y)J_{l}(y)}{y^{2}}dy,$
(37)
and we may put the lower limit of integration equal to zero. For the integral
over $dy$, we obtain :
$\displaystyle\int_{x}^{\infty}\frac{N_{l}(y)J_{l}(y)}{y^{2}}dy$
$\displaystyle=\frac{1}{2l(2l-1)}J_{l}(x)N_{l-1}(x)+\frac{1}{2l(2l+1)}J_{l+1}(x)N_{l}(x)$
$\displaystyle\\!\\!+\frac{2x}{4l^{2}-1}\big{[}N_{l-1}(x)J_{l+1}(x)\\!-\\!J_{l}(x)N_{l}(x)\big{]}-\frac{1}{\pi
lx}.$ (38)
Then, using the relations:
$\int_{0}^{\infty}\frac{J_{l}^{2}(x)}{x^{3}}dx=\frac{1}{4l(l^{2}-1)},$
$\int_{0}^{\infty}\frac{J_{l}^{2}(x)}{x}N_{l-1}(x)J_{l+1}(x)dx=\frac{1}{4l(l+1)\pi},$
$\int_{0}^{\infty}\frac{J_{l}^{3}(x)}{x}N_{l}(x)dx=-\frac{1}{4l^{2}\pi},$
$\int_{0}^{\infty}\frac{J_{l}^{2}(x)}{x^{2}}J_{l}(x)N_{l-1}(x)dx=\frac{1}{8l^{2}(l+1)\pi},$
$\int_{0}^{\infty}\frac{J_{l}^{2}(x)}{x^{2}}J_{l+1}(x)N_{l}(x)dx=-\frac{1}{8l^{2}(l+1)\pi},$
we find the following result for the second order phase shift:
$\displaystyle\tan\delta^{(2)}_{l}(k)=\frac{3\pi(kr_{*})^{2}}{8}\frac{1}{l(l^{2}-1)(4l^{2}-1)}.$
(39)
So, the total phase shift is given by
$\displaystyle\tan\delta_{l}(k)$
$\displaystyle=\tan\delta_{l}^{(1)}(k)+\tan\delta_{l}^{(2)}(k)$
$\displaystyle=-\frac{2kr_{*}}{4l^{2}-1}+\frac{3\pi(kr_{*})^{2}}{8l(l^{2}-1)(4l^{2}-1)}.$
(40)
Then, according to Eq. (16) the on-shell scattering amplitude
${\bar{f}}_{l}(k)$ is
${\bar{f}}_{l}(k)={\bar{f}}^{(1)}_{l}(k)+{\bar{f}}^{(2)}_{l}(k),$ (41)
where
${\bar{f}}^{(1)}_{l}(k)=\frac{8\hbar^{2}kr_{*}}{m}\frac{1}{4l^{2}-1},$ (42)
${\bar{f}}^{(2)}_{l}(k)=-\frac{3\pi\hbar^{2}}{2m}(kr_{*})^{2}\frac{1}{|l|(l^{2}-1)(4l^{2}-1)}.$
(43)
Note that Eqs. (42) and (43) do not contain short-range contributions as those
are proportional to $k^{2|l|}$ and can be omitted for $|l|>1$.
### II.4 First order Born approximation and the leading part of the
scattering amplitude
As we already said above, in the low-momentum limit for both $|l|=1$ and
$|l|>1$ the leading part of the on-shell scattering amplitude
${\bar{f}}_{l}(k)$ is ${\bar{f}}_{l}^{(1)}(k)$ and it is contained in the
first order contribution from distances $r>r_{0}$. For $|l|>1$ it is given by
Eq. (42) and follows from Eq. (II.3) with $\psi^{(0)}_{l(II)}=J_{l}(kr)$. In
the case of $|l|=1$ this leading part is given by Eq. (32) and follows from
the integral term of Eq. (23) in which one keeps only $J_{1}(kr)$ in the
expression for $\psi^{(0)}_{II}(r)$. This means that ${\bar{f}}_{l}^{(1)}(k)$
actually follows from the first order Born approximation.
The off-shell scattering amplitude can also be represented as
${\bar{f}}_{l}(k^{\prime},k)={\bar{f}}_{l}^{(1)}(k^{\prime},k)+{\bar{f}}_{l}^{(2)}(k^{\prime},k)$,
and the leading contribution ${\bar{f}}_{l}^{(1)}(k^{\prime},k)$ follows from
the first Born approximation. It is given by Eq. (9) in which one should
replace $\tilde{\psi}_{l}(k,r)$ by $J_{l}(kr)$:
${\bar{f}}_{l}^{(1)}(k^{\prime},k)=\int_{0}^{\infty}J_{l}(kr)J_{l}(kr^{\prime})U(r)2\pi
rdr.$ (44)
Note that it is not important that we put zero for the lower limit of the
integration, since this can only give a correction which behaves as $k^{2}$ or
a higher power of $k$. Then, putting $U(r)=\hbar^{2}r_{*}/mr^{3}$ in Eq. (44),
we obtain:
$\displaystyle{\bar{f}}^{(1)}_{l}(k^{\prime},k)=$
$\displaystyle\frac{\pi\hbar^{2}}{m}\frac{\Gamma(l-1/2)}{\sqrt{\pi}}\frac{k^{l}r_{*}}{(k^{\prime})^{l-1}}$
$\displaystyle\times
F\left(-\frac{1}{2},-\frac{1}{2}+l,1+l,\frac{k^{2}}{k^{\prime 2}}\right),$
(45)
where $F$ is the hypergeometric function. The result of Eq. (II.4) corresponds
to $k<k^{\prime}$, and for $k>k^{\prime}$ one should interchange $k$ and
$k^{\prime}$.
For identical fermions the full scattering amplitude contains only partial
amplitudes with odd $l$. Since the scattered waves with relative momenta ${\bf
k}^{\prime}$ and $-{\bf k}^{\prime}$ correspond to interchanging the identical
fermions, the scattering amplitude can be written as (see, e.g. Lan2 ):
$\tilde{f}({\bf k}^{\prime},{\bf k})=f({\bf k}^{\prime},{\bf k})-f(-{\bf
k}^{\prime},{\bf k}).$ (46)
Then, according to equation (10) one can write:
$\tilde{f}({\bf k}^{\prime},{\bf
k})=2\sum_{l\,odd}f_{l}(k^{\prime},k)\exp(il\varphi).$ (47)
In the first Born approximation there is no difference between
$f_{l}(k^{\prime},k)$ and ${\bar{f}}_{l}(k^{\prime},k)$ because
$\tan\delta_{l}(k)$ in the denominator of Eq. (15) is proportional to $k$ and
can be disregarded. Therefore, one may use ${\bar{f}}_{l}^{(1)}(k^{\prime},k)$
of Eq.(II.4) for $f_{l}(k^{\prime},k)$ in Eq. (47). One can represent
$\tilde{f}({\bf k}^{\prime},{\bf k})$ in a different form recalling that in
the first Born approximation we have:
$f({\bf k}^{\prime},{\bf k})=\int U(r)\exp[i({\bf k}-{\bf k}^{\prime}){\bf
r}]d^{2}r.$ (48)
Performing the integration in this equation, with $U(r)$ given by Eq. (1), and
using Eq. (46) we obtain:
$\tilde{f}({\bf k}^{\prime},{\bf k})=\frac{2\pi\hbar^{2}r_{*}}{m}\\{|{\bf
k}+{\bf k}^{\prime}|-|{\bf k}-{\bf k}^{\prime}|\\}.$ (49)
Equation (49) is also obtained by a direct summation over odd $l$ in Eq. (47),
with $f_{l}(k^{\prime},k)$ following from Eq. (II.4).
## III Thermodynamics of a weakly interacting 2D gas of fermionic polar
molecules at $T=0$
### III.1 General relations of Fermi liquid theory
Identical fermionic polar molecules undergoing a two-dimensional translational
motion and repulsively interacting with each other via the potential (1)
represent a 2D Fermi liquid. General relations of the Landau Fermi liquid
theory remain similar to those in 3D (see, e.g. Landau ). The number of
“dressed” particles, or quasiparticles, is the same as the total number of
particles $N$, and the (quasi)particle Fermi momentum is
$k_{F}=\sqrt{\frac{4\pi N}{S}},$ (50)
where $S$ is the surface area. At $T=0$ the momentum distribution of free
quasiparticles is the step function
$n({\bf k})=\theta(k_{F}-k),$ (51)
i.e. $n({\bf k})=1$ for $k<k_{F}$ and zero otherwise.The chemical potential is
equal to the boundary energy at the Fermi circle,
$\mu=\epsilon_{F}\equiv\epsilon(k_{F})$.
The quasiparticle energy $\epsilon({\bf k})$ is a variational derivative of
the total energy with respect to the distribution function $n({\bf k})$. Due
to the interaction between quasiparticles, the deviation $\delta n$ of this
distribution from the step function (51) results in the change of the
quasiparticle energy:
$\delta\epsilon(\mathbf{k})=\int F(\mathbf{k},\mathbf{k}^{\prime})\delta
n(\mathbf{k}^{\prime})\frac{d^{2}k^{\prime}}{(2\pi)^{2}}.$ (52)
The interaction function of quasiparticles $F(\mathbf{k},\mathbf{k}^{\prime})$
is thus the second variational derivative of the total energy with regard to
$n({\bf k})$. The quantity $\delta n({\bf k})$ is significantly different from
zero only near the Fermi surface, so that one may put ${\bf k}=k_{F}{\bf n}$
and ${\bf k}^{\prime}=k_{F}{\bf n}^{\prime}$ in the arguments of $F$ in Eq.
(52), where ${\bf n}$ and ${\bf n}^{\prime}$ are unit vectors in the
directions of ${\bf k}$ and ${\bf k}^{\prime}$. The quasiparticle energy near
the Fermi surface can be written as:
$\epsilon(\mathbf{k})=\epsilon_{F}+\hbar v_{F}(k-k_{F})+\int
F(\mathbf{k},\mathbf{k}^{\prime})\delta
n(\mathbf{k}^{\prime})\frac{d^{2}k^{\prime}}{(2\pi)^{2}}.$ (53)
The quantity $v_{F}=\partial\epsilon({\bf k})/\hbar\partial k|_{k=k_{F}}$ is
the Fermi velocity, and the effective mass of a quasiparticle is defined as
$m^{*}=\hbar k_{F}/v_{F}$. It can be obtained from the relation (see Landau ):
$\displaystyle\frac{1}{m}=\frac{1}{m^{*}}+\frac{1}{(2\pi\hbar)^{2}}\int_{0}^{2\pi}F(\theta)\cos\theta
d\theta,$ (54)
where $\theta$ is the angle between the vectors ${\bf n}$ and ${\bf
n}^{\prime}$, and $F(\theta)=F(k_{F}{\bf n},k_{F}{\bf n}^{\prime})$.
The compressibility $\kappa$ at $T=0$ is given by Landau :
$\kappa^{-1}=\frac{N^{2}}{S}\frac{\partial\mu}{\partial N}.$ (55)
The chemical potential is $\mu=\epsilon_{F}$, and the variation of $\mu$ due
to a change in the number of particles can be expressed as
$\delta\mu=\int F(k_{F}{\bf n},\mathbf{k}^{\prime})\delta
n(\mathbf{k^{\prime}})\frac{d^{2}k^{\prime}}{(2\pi)^{2}}+\frac{\partial\epsilon_{F}}{\partial
k_{F}}\delta k_{F}.$ (56)
The quantity $\delta n(\mathbf{k^{\prime}})$ is appreciably different from
zero only when $\mathbf{k^{\prime}}$ is near the Fermi surface, so that we can
replace the interaction function $F$ by its value on the Fermi surface. Then
the first term of Eq. (56) becomes
$\displaystyle\int F(\theta)\frac{d\theta}{2\pi}\int\delta
n(\mathbf{k^{\prime}})\frac{d^{2}k^{\prime}}{(2\pi)^{2}}=\frac{\delta N}{2\pi
S}\int F(\theta)d\theta.$
The second term of Eq. (56) reduces to
$\displaystyle\frac{\partial\epsilon_{F}}{\partial k_{F}}\delta
k_{F}=\frac{\hbar^{2}k_{F}}{m^{*}}\delta
k_{F}=\frac{2\pi\hbar^{2}}{m^{*}}\frac{\delta N}{S}.$ (57)
We thus have (see Landau ):
$\displaystyle\frac{\partial\mu}{\partial N}$ $\displaystyle=\frac{1}{2\pi
S}\int_{0}^{2\pi}F(\theta)d\theta+\frac{2\pi\hbar^{2}}{m^{*}S}$
$\displaystyle=\frac{2\pi\hbar^{2}}{mS}+\frac{1}{2\pi
S}\int_{0}^{2\pi}(1-\cos\theta)F(\theta)d\theta.$ (58)
Equation (III.1) shows that the knowledge of the interaction function of
quasiparticles on the Fermi surface, $F(\theta)$, allows one to calculate
$\partial\mu/\partial N$ and, hence, the chemical potential $\mu=\partial
E/\partial N$ and the ground state energy $E$. This elegant way of finding the
ground state energy has been proposed by Abrikosov and Khalatnikov Abr . It
was implemented in Ref. Abr for a two-component 3D Fermi gas with a weak
repulsive contact (short-range) interspecies interaction.
We develop a theory beyond the mean field for calculating the interaction
function of quasiparticles for a single-component 2D gas of fermionic polar
molecules in the weakly interacting regime. We obtain the ground state energy
as a series of expansion in the small parameter $k_{F}r_{*}$ and confine
ourselves to the second order. In this sense our work represents a sort of
Lee-Huang-Yang Huang ; Lee and Abrikosov-Khalatnikov Abr calculation for
this dipolar system. As we will see, the long-range character of the dipole-
dipole interaction makes the result quite different from that in the case of
short-range interactions.
### III.2 Two-body and many-body contributions to the ground state energy
We first write down the expression for the kinetic energy and specify two-body
(mean field) and many-body (beyond mean field) contributions to the
interaction energy. The Hamiltonian of the system reads:
$\\!\\!\hat{\cal
H}\\!=\\!\sum_{\mathbf{k}}\frac{\hbar^{2}k^{2}}{2m}\hat{a}_{\mathbf{k}}^{{\dagger}}\hat{a}_{\mathbf{k}}\\!+\\!\frac{1}{2S}\\!\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{q}}\\!\\!\\!U(\mathbf{q})\hat{a}_{\mathbf{k_{1}}+\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}-\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}}\hat{a}_{\mathbf{k_{1}}},\\!\\!$
(59)
where $\hat{a}^{\dagger}_{{\bf k}}$ and $\hat{a}_{{\bf k}}$ are creation and
annihilation operators of fermionic polar molecules, and $U(\mathbf{q})$ is
the Fourier transform of the interaction potential $U(r)$:
$U(\mathbf{q})=\int d^{2}\mathbf{r}U(r)e^{-i\mathbf{q}\cdot\mathbf{r}},$ (60)
The first term of Eq. (59) represents the kinetic energy and it gives the main
contribution to the total energy $E$ of the system. This term has only
diagonal matrix elements, and using the momentum distribution (51) at $T=0$ we
have:
$\displaystyle\frac{E_{kin}}{S}=\int_{0}^{k_{F}}\frac{\hbar^{2}k^{2}}{2m}\frac{2\pi
kdk}{(2\pi)^{2}}=\frac{\hbar^{2}k_{F}^{4}}{16m}.$ (61)
The interaction between the fermionic molecules is described by the second
term in Eq. (59) and compared to the kinetic energy it provides a correction
to the total energy $E$. The first order correction is given by the diagonal
matrix element of the interaction term of the Hamiltonian:
$\displaystyle E^{(1)}$
$\displaystyle=\frac{1}{2S}\sum_{{\bf{k_{1}}},{\bf{k_{2}}},{\bf
q}}U(\mathbf{q})\langle\hat{a}_{\mathbf{k_{1}}+\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}-\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}}\hat{a}_{\mathbf{k_{1}}}\rangle$
$\displaystyle=\frac{1}{2S}\sum_{\mathbf{k_{1}},\mathbf{k_{2}}}\left[U(0)-U(\mathbf{k_{2}}-\mathbf{k_{1}})\right]n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}.$
(62)
The second order correction to the energy of the state $\left|{j}\right>$ of a
non-interacting system can be expressed as:
$\displaystyle E_{j}^{(2)}=\sum_{m\neq j}\frac{V_{jm}V_{mj}}{E_{j}-E_{m}},$
(63)
where the summation is over eigenstates $\left|{m}\right>$ of the non-
interacting system, and $V_{jm}$ is the non-diagonal matrix element. In our
case the symbol $j$ corresponds to the ground state and the symbol $m$ to
excited states. The non-diagonal matrix element is
$\\!\\!V_{jm}\\!=\\!\frac{1}{2S}\left<\\!m\left|\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{q}}U(\mathbf{q})\hat{a}_{\mathbf{k_{1}}+\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}-\mathbf{q}}^{{\dagger}}\hat{a}_{\mathbf{k_{2}}}\hat{a}_{\mathbf{k_{1}}}\right|j\\!\right>.\\!$
(64)
This matrix element corresponds to the scattering of two particles from the
initial state $\mathbf{k_{1}}$, $\mathbf{k_{2}}$ to an intermediate state
$\mathbf{k^{\prime}_{1}}$, $\mathbf{k^{\prime}_{2}}$, and the matrix element
$V_{mj}$ describes the reversed process in which the two particles return from
the intermediate to initial state. Taking into account the momentum
conservation law ${\bf k}_{1}+{\bf k}_{2}={\bf k}^{\prime}_{1}+{\bf
k}^{\prime}_{2}$ the quantity $V_{jm}V_{mj}=|V_{jm}|^{2}$ is given by
$\displaystyle|V_{jm}|^{2}$
$\displaystyle=\frac{1}{(2S)^{2}}n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}(1-n_{\mathbf{k_{1}^{\prime}}})(1-n_{\mathbf{k_{2}^{\prime}}})$
$\displaystyle\times\left|U(\mathbf{k^{\prime}_{1}}-\mathbf{k_{1}})-U(\mathbf{k^{\prime}_{2}}-\mathbf{k_{1}})\right|^{2},$
(65)
and the second order correction to the ground state energy takes the form:
$\displaystyle E^{(2)}=$
$\displaystyle\frac{1}{(2S)^{2}}\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k_{1}^{\prime}}}\Bigg{[}\left|U({\bf
k}^{\prime}_{1}-{\bf k}_{1})-U({\bf k}^{\prime}_{2}-{\bf k}_{1})\right|^{2}$
$\displaystyle\times\frac{n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}(1-n_{\mathbf{k_{1}^{\prime}}})(1-n_{\mathbf{k_{2}^{\prime}}})}{\hbar^{2}(\mathbf{k^{2}_{1}}+\mathbf{k^{2}_{2}}-\mathbf{k^{\prime
2}_{1}}-\mathbf{k^{\prime 2}_{2}})/2m}\Bigg{]}.$ (66)
From Eq. (III.2) we see that the second order correction diverges because of
the term proportional to $n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}$, which is
divergent at large $k^{\prime}_{1}$. This artificial divergence is eliminated
by expressing the energy correction in terms of a real physical quantity, the
scattering amplitude. The relation between the Fourier component of the
interaction potential and the off-shell scattering amplitude is given by Lan2
:
$f(\mathbf{k^{\prime}},\mathbf{k})=U(\mathbf{k}^{\prime}-\mathbf{k})+\frac{1}{S}\sum_{\mathbf{k^{\prime\prime}}}\frac{U(\mathbf{k}^{\prime}-\mathbf{k^{\prime\prime}})f(\mathbf{k^{\prime\prime}},\mathbf{k})}{(E_{\mathbf{k}}-E_{\mathbf{k^{\prime\prime}}}-i0)},$
(67)
where $E_{\mathbf{k}}=\hbar^{2}{\bf k}^{2}/m$ and
$E_{\mathbf{k^{\prime\prime}}}=\hbar^{2}{\bf k}^{\prime\prime 2}/m$ are
relative collision energies. Obviously, we have:
$E_{\mathbf{k}}-E_{\mathbf{k^{\prime\prime}}}=\hbar^{2}(\mathbf{k_{1}}^{2}+\mathbf{k_{2}}^{2}-\mathbf{k^{\prime\prime}_{1}}^{2}-\mathbf{k^{\prime\prime}_{2}}^{2})/2m$,
with ${\bf k}_{1}$, ${\bf k}_{2}$ (${\bf k}^{\prime\prime}_{1}$, ${\bf
k}^{\prime\prime}_{2}$) being the momenta of colliding particles in the
initial (intermediate) state, as the relative momenta are given by
$\mathbf{k}=(\mathbf{k}_{1}-\mathbf{k}_{2})/2$,
$\mathbf{k^{\prime\prime}}=(\mathbf{k^{\prime\prime}_{1}}-\mathbf{k^{\prime\prime}_{2}})/2$.
We thus can write:
$\\!\\!\\!U(\mathbf{k}^{\prime}\\!-\mathbf{k})\\!=\\!f(\mathbf{k}^{\prime}\\!,\\!\mathbf{k})\\!-\\!\frac{2m}{\hbar^{2}S}\sum_{\mathbf{k_{1}^{\prime\prime}}}\frac{U(\mathbf{k}^{\prime}\\!-\\!\mathbf{k}^{\prime\prime})f(\mathbf{k^{\prime\prime}}\\!,\\!\mathbf{k})}{\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime\prime 2}_{2}}\\!-\\!i0}.\\!\\!$ (68)
Then, putting ${\bf k}^{\prime}={\bf k}$ we have
$\displaystyle U(0)=f({\bf k},{\bf k})-\frac{2m}{\hbar^{2}S}\sum_{{\bf
k}^{\prime\prime}_{1}}\frac{U({\bf k}-{\bf k}^{\prime\prime})f({\bf
k}^{\prime\prime},{\bf
k})}{\mathbf{k^{2}_{1}}+\mathbf{k^{2}_{2}}-\mathbf{k^{\prime\prime
2}_{1}}-\mathbf{k^{\prime\prime 2}_{2}}-i0},$
and setting ${\bf k}^{\prime}=-{\bf k}$ we obtain
$\displaystyle\\!\\!U({\bf k}_{2}\\!-\\!{\bf k}_{1})=f(-{\bf k},{\bf
k})\\!-\\!\frac{2m}{\hbar^{2}S}\sum_{{\bf k}^{\prime\prime}_{1}}\frac{U(-{\bf
k}\\!-\\!{\bf k}^{\prime\prime})f({\bf k}^{\prime\prime},{\bf
k})}{\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime\prime 2}_{2}}\\!-\\!i0},$
Using these relations the first order correction (III.2) takes the form:
$\displaystyle
E^{(1)}=\frac{1}{2S}\sum_{\mathbf{k_{1}},\mathbf{k_{2}}}\left[f(\mathbf{k},\mathbf{k})-f(\mathbf{-k},\mathbf{k})\right]n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}$
$\displaystyle\\!\\!-\frac{1}{2S^{2}}\\!\\!\\!\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k^{\prime}_{1}}}\\!\\!\frac{[U(\mathbf{k}\\!-\\!\mathbf{k^{\prime}})\\!-\\!U({\\!-\bf
k}\\!-\\!{\bf
k}^{\prime})]f(\mathbf{k^{\prime}},\mathbf{k})}{\hbar^{2}(\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime 2}_{2}}\\!-\\!i0)/2m}n_{{\bf k}_{1}}n_{{\bf
k}_{2}}.\\!\\!$ (69)
The quantity $[U({\bf k}-{\bf k}^{\prime})-U(-{\bf k}-{\bf k}^{\prime})]$ in
the second term of Eq. (III.2), being expanded in circular harmonics
$\exp(il\varphi)$ contains terms with odd $l$. Therefore, partial amplitudes
with even $l$ in the expansion of the multiple $f({\bf k}^{\prime},{\bf k})$
vanish after the integration over $d^{2}k^{\prime}$. Hence, this amplitude can
be replaced by $[f({\bf k}^{\prime},{\bf k})-f({\bf k}^{\prime},-{\bf k})]/2$.
As we are interested only in the terms that behave themselves as $\sim k$ or
$\sim k^{2}$, the amplitudes in the second term of Eq. (III.2) are the ones
that follow from the first Born approximation and are proportional to $k$.
Therefore, we may put $[U({\bf k}-{\bf k}^{\prime})-U(-{\bf k}-{\bf
k}^{\prime})]=[f({\bf k},{\bf k}^{\prime})-f(-{\bf k},{\bf k}^{\prime})]$ and
$f({\bf k}^{\prime},{\bf k})=f^{*}({\bf k},{\bf k}^{\prime})$. Then the first
order correction takes the form:
$\displaystyle
E^{(1)}=\frac{1}{2S}\sum_{\mathbf{k_{1}},\mathbf{k_{2}}}\left[f(\mathbf{k},\mathbf{k})-f(\mathbf{-k},\mathbf{k})\right]n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}-\frac{1}{(2S)^{2}}$
$\displaystyle\times\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k^{\prime}_{1}}}\frac{|f(\mathbf{k^{\prime}},\mathbf{k})\\!-\\!f({\bf
k}^{\prime},-{\bf
k})|^{2}}{\hbar^{2}(\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime 2}_{2}}\\!-\\!i0)/2m}n_{{\bf k}_{1}}n_{{\bf
k}_{2}}.$ (70)
Using the expansion of the full scattering amplitude in terms of partial
amplitudes as given by Eq. (47) we represent the first order correction as
$\displaystyle E^{(1)}=\frac{1}{S}\sum_{{\bf k}_{1},{\bf
k}_{2}}\sum_{l\,odd}f_{l}(k)n_{{\bf k}_{1}}n_{{\bf
k}_{2}}-\frac{1}{S^{2}}\sum_{{\bf k}_{1},{\bf k}_{2}}\sum_{l\,odd}$
$\displaystyle\\!\times\\!\\!\int\\!\frac{d^{2}k^{\prime}}{(2\pi)^{2}}\frac{f_{l}^{2}(k)}{\hbar^{2}(\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime 2}_{2}}\\!-\\!i0)/2m}n_{{\bf k}_{1}}n_{{\bf
k}_{2}}.\\!\\!$ (71)
The contribution of the pole in the integration over $d^{2}k^{\prime}$ in the
second term of Eq. (III.2) gives $imf_{l}^{2}(k)/4\hbar^{2}$ for each term in
the sum over ${\bf k}_{1}$, ${\bf k}_{2}$, and $l$, and we may use here the
amplitude ${\bar{f}}^{(1)}_{l}(k)$. In the first term of Eq. (III.2) we should
use $f_{l}(k)=f_{l}^{(1)}(k)+f_{l}^{(2)}(k)$. However, we may replace
$f_{l}^{(2)}$ by ${\bar{f}}_{l}^{(2)}$ because the account of $\tan\delta(k)$
in the denominator of Eq. (15) leads to $k^{3}$ terms and terms containing
higher powers of $k$. For the amplitude $f_{l}^{(1)}(k)$, we use the
expression:
$f_{l}^{(1)}(k)={\bar{f}}^{(1)}_{l}+i\tan\delta(k){\bar{f}}^{(1)}_{l}={\bar{f}}^{(1)}_{l}-im[{\bar{f}}^{(1)}_{l}]^{2}/4\hbar^{2},$
which assumes a small scattering phase shift. The second term of this
expression, being substituted into the first line of Eq. (III.2), exactly
cancels the contribution of the pole in the second term of (III.2). Thus, we
may use the amplitude ${\bar{f}}_{l}$ in the first term of equation (III.2)
and take the principal value of the integral in the second term. The resulting
expression for the first order correction reads:
$\displaystyle E^{(1)}$ $\displaystyle=\frac{1}{S}\sum_{{\bf k}_{1},{\bf
k}_{2}}{\bar{f}}({\bf k})n_{{\bf k}_{1}}n_{{\bf k}_{2}}-\frac{1}{(2S)^{2}}$
$\displaystyle\times\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k^{\prime}_{1}}}\frac{2m|f(\mathbf{k},\mathbf{k^{\prime}})-f({-\bf
k},{\bf
k}^{\prime})|^{2}}{\hbar^{2}(\mathbf{k^{2}_{1}}+\mathbf{k^{2}_{2}}-\mathbf{k^{\prime
2}_{1}}-\mathbf{k^{\prime 2}_{2}})}n_{{\bf k}_{1}}n_{{\bf k}_{2}},$ (72)
where ${\bar{f}}({\bf k})=\sum_{l\,\,odd}{\bar{f}}_{l}(k)$.
The second order correction (III.2) can also be expressed in terms of the
scattering amplitude by using Eq.(67). Replacing $U({\bf k}_{1}-{\bf
k}^{\prime}_{1})=U({\bf k}-{\bf k}^{\prime})$ and $U({\bf k}^{\prime}_{2}-{\bf
k}_{1})=U(-{\bf k}-{\bf k}^{\prime})$ by $f({\bf k}^{\prime},{\bf k})$ and
$f(-{\bf k},{\bf k}^{\prime})$, respectively, we have:
$\displaystyle E^{(2)}=\frac{1}{(2S)^{2}}$
$\displaystyle\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k_{1}^{\prime}}}\Big{[}\frac{\left|f(\mathbf{k^{\prime}},\mathbf{k})-f(\mathbf{k^{\prime}},-\mathbf{k})\right|^{2}}{\hbar^{2}(\mathbf{k^{2}_{1}}+\mathbf{k^{2}_{2}}-\mathbf{k^{\prime
2}_{1}}-\mathbf{k^{\prime 2}_{2}})/2m}$ $\displaystyle\times
n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}(1-n_{\mathbf{k_{1}^{\prime}}})(1-n_{\mathbf{k_{2}^{\prime}}})\Big{]}.$
(73)
Note that the divergent term proportional to $n_{{\bf k}_{1}}n_{{\bf k}_{2}}$
in Eq. (III.2) and the (divergent) second term of Eq. (III.2) exactly cancel
each other, and the sum of the first and second order corrections can be
represented as $E^{(1)}+E^{(2)}=\tilde{E}^{(1)}+\tilde{E}^{(2)},$ where
$\tilde{E}^{(1)}=\frac{1}{S}\sum_{{\bf k}_{1},{\bf k}_{2}}{\bar{f}}({\bf
k})n_{{\bf k}_{1}}n_{{\bf k}_{2}},$ (74)
and
$\displaystyle\tilde{E}^{(2)}=\frac{1}{(2S)^{2}}$
$\displaystyle\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k^{\prime}_{1}}}\Big{\\{}\frac{|f(\mathbf{k^{\prime}},\mathbf{k})-f({\bf
k^{\prime}},-{\bf
k})|^{2}}{\hbar^{2}(\mathbf{k^{2}_{1}}+\mathbf{k^{2}_{2}}-\mathbf{k^{\prime
2}_{1}}-\mathbf{k^{\prime 2}_{2}})/2m}$ $\displaystyle\times n_{{\bf
k}_{1}}n_{{\bf k}_{2}}[(1-n_{{\bf k}^{\prime}_{1}})(1-n_{{\bf
k}^{\prime}_{2}})-1]\Big{\\}}.$ (75)
The term $\tilde{E}^{(1)}$ originates from the two-body contributions to the
interaction energy and can be quoted as the mean field term. The term
$\tilde{E}^{(2)}$ is the many-body contribution, which is beyond mean field.
It is worth noting that the term proportional to the product of four
occupation numbers vanishes because its numerator is symmetrical and the
denominator is antisymmetrical with respect to an interchange of ${\bf
k}_{1},{\bf k}_{2}$ and ${\bf k}^{\prime}_{1},{\bf k}^{\prime}_{2}$. The terms
containing a product of three occupation numbers, $n_{{\bf k}_{1}}n_{{\bf
k}_{2}}n_{{\bf k}^{\prime}_{1}}$ and $n_{{\bf k}_{1}}n_{{\bf k}_{2}}n_{{\bf
k}^{\prime}_{2}}$ are equal to each other because the denominator is
symmetrical with respect to an interchange of ${\bf k}^{\prime}_{1}$ and ${\bf
k}^{\prime}_{2}$. We thus reduce Eq. (III.2) to
$\\!\\!\\!\\!\tilde{E}^{(2)}\\!\\!=\\!\\!-\frac{1}{2S^{2}}\\!\\!\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k^{\prime}_{1}}}\\!\\!\\!\frac{2m|f(\mathbf{k^{\prime}},\mathbf{k})\\!-\\!f({\bf
k^{\prime}},\\!-{\bf
k})|^{2}}{\hbar^{2}(\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime
2}_{1}}\\!\\!-\\!\mathbf{k^{\prime 2}_{2}})}n_{{\bf k}_{1}}\\!n_{{\bf
k}_{2}}\\!n_{{\bf k}^{\prime}_{1}}.\\!\\!\\!\\!$ (76)
Equations (74) and (76) allow a direct calculation of the ground state energy.
With respect to the mean field term $\tilde{E}^{(1)}$ this is done in Appendix
A. However, a direct calculation of the many-body correction $\tilde{E}^{(2)}$
is even a more tedious task than in the case of two-component fermions with a
contact interaction. We therefore turn to the Abrikosov-Khalatnikov idea of
calculating the ground state energy (and other thermodynamic quantities)
through the interaction function of quasiparticles on the Fermi surface.
### III.3 Interaction function of quasiparticles
The interaction function of quasiparticles $F({\bf k},{\bf k}^{\prime})$ is
the second variational derivative of the total energy with respect to the
distribution $n_{{\bf k}}$. The kinetic energy of our system is linear in
$n_{{\bf k}}$ (see Eq. (59)), and the second variational derivative is related
to the variation of the interaction energy $\tilde{E}$. We have Landau :
$\delta\tilde{E}=\frac{1}{2S}\sum_{{\bf k},{\bf k}^{\prime}}F({\bf k},{\bf
k}^{\prime})\delta n_{\mathbf{k}}\delta n_{\mathbf{k^{\prime}}},$ (77)
where $\tilde{E}=\tilde{E}^{(1)}+\tilde{E}^{(2)}$, and the quantities
$\tilde{E}^{(1)}$ and $\tilde{E}^{(2)}$ are given by equations (74) and (76).
On the Fermi surface we should put $|{\bf k}|=|{\bf k}^{\prime}|=k_{F}$, so
that the interaction function will depend only on the angle $\theta$ between
${\bf k}$ and $\mathbf{k^{\prime}}$. Hereinafter it will be denoted as
$\tilde{F}(\theta)$.
The contribution $\tilde{F}^{(1)}(\theta)=2S\delta\tilde{E}^{(1)}/\delta
n_{{\bf k}}\delta n_{{\bf k}^{\prime}}$ is given by
$\tilde{F}^{(1)}(\theta)=2f\left(\frac{|{\bf k}-{\bf
k}^{\prime}|}{2}\right)=2\sum_{l\,odd}{\bar{f}}_{l}\left(k_{F}|\sin{\frac{\theta}{2}}|\right),$
(78)
where ${\bar{f}}_{l}={\bar{f}}_{l}^{(1)}+{\bar{f}}_{l}^{(2)}$, and the
amplitudes ${\bar{f}}_{l}^{(1)}$ and ${\bar{f}}_{l}^{(2)}$ follow from Eqs.
(32) and (33) at $|l|=1$, and from Eqs. (42), (43) at $|l|>1$. We thus may
write equation (42),
${\bar{f}}^{(1)}_{l}(k)=\frac{8\hbar^{2}}{m}\frac{1}{4l^{2}-1}kr_{*},$
for any odd $l$, and
$\displaystyle{\bar{f}}^{(2)}_{l}(k)=\frac{\pi\hbar^{2}}{2m}(kr_{*})^{2}\times\begin{cases}\ln(\xi
kr_{*});&\text{$|l|=1$}\\\
-\frac{3}{|l|(l^{2}-1)(4l^{2}-1)};&\text{$|l|>1$}\end{cases}$
with $\xi$ from Eq. (30). Making a summation over all odd $l$ we obtain:
${\bar{f}}^{(1)}(k)=\sum_{l\,odd}f^{(1)}_{l}(k)=\frac{2\pi\hbar^{2}}{m}kr_{*},\\\
$ (79)
$\\!\\!\\!{\bar{f}}^{(2)}(k)\\!=\\!\\!\\!\sum_{l\,odd}f^{(2)}_{l}(k)\\!\\!=\\!\\!\frac{\pi\hbar^{2}}{m}(kr_{*}\\!)^{2}\\!\\!\left[\ln(\xi
kr_{*}\\!)\\!-\\!\frac{25}{12}\\!+\\!3\\!\ln 2\\!\right]\\!\\!.\\!\\!\\!$ (80)
Putting $k=k_{F}|\sin(\theta/2)|$ and substituting the results of equations
(79) and (80) into Eq. (78) we find:
$\displaystyle\tilde{F}^{(1)}(\theta)$
$\displaystyle=\frac{4\pi\hbar^{2}k_{F}r_{*}}{m}|\sin\frac{\theta}{2}|+\frac{2\hbar^{2}}{m}(k_{F}r_{*})^{2}$
$\displaystyle\times\pi\sin^{2}\frac{\theta}{2}\left[\ln|\xi
r_{*}k_{F}\sin\frac{\theta}{2}|-\frac{25}{12}+3\ln 2\right].$ (81)
The many-body correction (76) we represent as
$\tilde{E}^{(2)}=\tilde{E}_{1}^{(2)}+\tilde{E}_{2}^{(2)}$, where
$\displaystyle\tilde{E}_{1}^{(2)}\\!\\!\\!=\\!\\!-\frac{8(\pi\hbar
r_{*})^{2}}{mS^{2}}\\!\\!\\!\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k_{1}^{\prime}}}\\!\\!\\!\frac{|\mathbf{k^{\prime}_{1}}\\!-\\!\mathbf{k_{1}}|^{2}}{\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\\!\mathbf{k^{\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime
2}_{2}}}n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}n_{\mathbf{k_{1}^{\prime}}}\\!,\\!\\!\\!$
(82) $\displaystyle\tilde{E}_{2}^{(2)}\\!\\!\\!=\\!\frac{8(\pi\hbar
r_{*})^{2}}{mS^{2}}\\!\\!\\!\\!\\!\sum_{\mathbf{k_{1}},\mathbf{k_{2}},\mathbf{k_{1}^{\prime}}}\\!\\!\\!\frac{|\mathbf{k_{1}}\\!-\\!\mathbf{k^{\prime}_{1}}|\\!\cdot\\!|\mathbf{k_{2}}\\!-\\!\mathbf{k^{\prime}_{1}}|}{\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}_{2}}\\!-\mathbf{k^{\prime
2}_{1}}\\!-\\!\mathbf{k^{\prime
2}_{2}}}n_{\mathbf{k_{1}}}n_{\mathbf{k_{2}}}n_{\mathbf{k_{1}^{\prime}}}\\!,\\!\\!$
(83)
and we used Eqs (46) and (49) for the scattering amplitudes. The contribution
to the interaction function from $\tilde{E}_{1}^{(2)}$ is calculated in
Appendix B and it reads:
$\\!\\!\\!\tilde{F}_{1}^{(2)}\\!(\theta)\\!\\!=\\!\\!\frac{2\hbar^{2}(k_{F}r_{*}\\!)^{2}}{m}\\!\left[\\!3\pi\\!+\\!2\pi\sin^{2}\frac{\theta}{2}\left(\\!\frac{4}{3}\\!-\\!\ln\\!|\tan\frac{\theta}{2}|\right)\right]\\!.\\!\\!$
(84)
The contribution from $\tilde{E}_{2}^{(2)}$ is calculated in Appendix C. It is
given by
$\displaystyle\tilde{F}^{(2)}_{2}(\theta)=$
$\displaystyle\frac{2\hbar^{2}k^{2}_{F}r^{2}_{*}}{m}\Big{\\{}-\sin^{2}\frac{\theta}{2}\left(\pi\ln
2+\frac{\pi}{2}-\pi\ln|\sin\frac{\theta}{2}|+4\ln|\cos\frac{\theta}{2}|-4\ln(1+|\sin\frac{\theta}{2}|)+\mathcal{G}(\theta)+\frac{4\arcsin|\sin\frac{\theta}{2}|-2\pi}{|\cos\frac{\theta}{2}|}\right)$
$\displaystyle-\frac{1}{|\cos\frac{\theta}{2}|}\left(\pi-2\arcsin|\sin\frac{\theta}{2}|+|\sin\theta|\right)-4\left[\cos^{2}\frac{\theta}{2}\ln\frac{1+|\sin(\theta/2)|}{1-|\sin(\theta/2)|}+2|\sin\frac{\theta}{2}|\right]\Big{\\}},$
(85)
where
$\mathcal{G}(\theta)=\int_{0}^{\pi}2\sin^{2}\varphi\ln\left(\sin\varphi+\sqrt{\sin^{2}\frac{\theta}{2}+\cos^{2}\frac{\theta}{2}\sin^{2}\varphi}\right)d\varphi,$
(86)
so that
$\frac{d\mathcal{G}(\theta)}{d\theta}=\frac{\pi}{2}\cot\frac{\theta}{2}-\frac{|\sin(\theta/2)|}{\sin(\theta/2)}\frac{1}{\cos(\theta/2)}+\frac{\arcsin|\cos(\theta/2)|}{|\cos(\theta/2)|}\left(\tan\frac{\theta}{2}-\cot\frac{\theta}{2}\right).$
(87)
We thus have
$\tilde{F}(\theta)=\tilde{F}^{(1)}(\theta)+\tilde{F}^{(2)}_{1}(\theta)+\tilde{F}^{(2)}_{2}(\theta)$,
where $\tilde{F}^{(1)}$, $\tilde{F}^{(2)}_{1}$, $\tilde{F}^{(2)}_{2}$ follow
from Eqs. (81), (84), and (85). This allows us to proceed with the calculation
of thermodynamic quantities.
### III.4 Compressibility, ground state energy, and effective mass
We first calculate the compressibility at $T=0$. On the basis of Eq. (III.1)
we obtain:
$\displaystyle\frac{\partial\mu}{\partial
N}=\frac{2\pi\hbar^{2}}{mS}+\frac{1}{2\pi
S}\int(1-\cos\theta)\left[\tilde{F}^{(1)}(\theta)+\tilde{F}_{1}^{(2)}(\theta)+\tilde{F}_{2}^{(2)}(\theta)\right]d\theta$
$\displaystyle=\frac{2\pi\hbar^{2}}{mS}+\frac{32\hbar^{2}}{3mS}k_{F}r_{*}+\frac{3\pi\hbar^{2}}{2mS}(k_{F}r_{*})^{2}\left(\ln[4\xi
k_{F}r_{*}]-\frac{3}{2}\right)+\frac{6\pi\hbar^{2}}{mS}(k_{F}r_{*})^{2}-\frac{\hbar^{2}}{\pi
mS}(k_{F}r_{*})^{2}(30-8G+21\zeta(3)),$ (88)
where $G=0.915966$ is the Catalan constant, and $\zeta(3)=1.20206$ is the
Riemann zeta function. Calculating coefficients and recalling that
$k_{F}=\sqrt{4\pi N/S}$ we represent the inverse compressibility following
from Eq. (55) in a compact form:
$\\!\\!\kappa^{-1}\\!=\\!\frac{\hbar^{2}k_{F}^{2}}{2m}\frac{N}{S}\\!\left(\\!1\\!+\\!\frac{16}{3\pi}k_{F}r_{*}\\!+\\!\frac{3}{4}(k_{F}r_{*})^{2}\ln(\zeta_{1}k_{F}r_{*})\right)\\!,\\!\\!\\!\\!$
(89)
where we obtain the coefficient $\zeta_{1}=2.16\exp(-2A)$ by using Eq. (30)
for the coefficient $\xi$ which depends on the short-range behavior through
the constant $A$ (see Eq. (18)). For the chemical potential and ground state
energy we obtain:
$\displaystyle\mu$ $\displaystyle=$
$\displaystyle\frac{2\pi\hbar^{2}N}{mS}+\frac{64\hbar^{2}N}{9mS}k_{F}r_{*}+\frac{3\pi\hbar^{2}N}{4mS}(k_{F}r_{*})^{2}\left(\ln[4\xi
k_{F}r_{*}]-\frac{7}{4}\right)+\frac{3\pi\hbar^{2}N}{mS}(k_{F}r_{*})^{2}-\frac{\hbar^{2}N}{2\pi
mS}(k_{F}r_{*})^{2}(30-8G+21\zeta(3))$ (90) $\displaystyle=$
$\displaystyle\frac{\hbar^{2}k_{F}^{2}}{2m}\left(1+\frac{32}{9\pi}k_{F}r_{*}+\frac{3}{8}(k_{F}r_{*})^{2}\ln(\zeta_{2}k_{F}r_{*})\right).$
$\displaystyle\frac{E}{N}$ $\displaystyle=$
$\displaystyle\frac{\pi\hbar^{2}N}{mS}+\frac{128\hbar^{2}N}{45mS}k_{F}r_{*}+\frac{\pi\hbar^{2}N}{4mS}(k_{F}r_{*})^{2}\left(\ln[4\xi
k_{F}r_{*}]-\frac{23}{12}\right)+\frac{\pi\hbar^{2}N}{mS}-\frac{\hbar^{2}N}{6\pi
mS}(30-8G+21\zeta(3))$ (91) $\displaystyle=$
$\displaystyle\frac{\hbar^{2}k_{F}^{2}}{4m}\left(1+\frac{128}{45\pi}k_{F}r_{*}+\frac{1}{4}(k_{F}r_{*})^{2}\ln(\zeta_{3}k_{F}r_{*})\right),$
with numerical coefficients $\zeta_{2}=1.68\exp(-2A)$ and
$\zeta_{3}=1.43\exp(-2A)$. Note that the first term in the second line of Eq.
(III.4) and the first terms in the first lines of Eqs. (90) and (91) represent
the contributions of the kinetic energy, the second and third terms correspond
to the contributions of the mean field part of the interaction energy, and the
last two terms are the contributions of the many-body effects.
The effective mass is calculated in a similar way by using Eq. (54):
$\displaystyle\frac{1}{m^{*}}\\!=\\!\frac{1}{m}\\!-\\!\frac{1}{(2\pi\hbar)^{2}}\int_{0}^{2\pi}\\!(F^{(1)}(\theta)\\!+\\!F^{(2)}_{1}(\theta)\\!+\\!F^{(2)}_{2}(\theta))\cos\theta
d\theta\\!\\!$
$\displaystyle\\!=\\!\\!\frac{1}{m}\\!\\!\left[\\!1\\!\\!+\\!\frac{4k_{F}r_{\\!*}}{3\pi}\\!+\\!\frac{(k_{F}r_{\\!*})^{\\!2}}{4}\\!\\!\left(\\!\ln{\\![4k_{F}r_{\\!*}\xi]}\\!\\!-\\!\frac{8}{3}\\!+\\!\frac{48G\\!\\!-\\!\\!20\\!-\\!\\!14\zeta(3)}{\pi^{2}}\\!\right)\\!\right]\\!\\!$
$\displaystyle\\!=\frac{1}{m}\left[1+\frac{4}{3\pi}k_{F}r_{*}+\frac{1}{4}(k_{F}r_{*})^{2}\ln(\zeta_{4}k_{F}r_{*})\right],$
(92)
where the numerical coefficient $\zeta_{4}=0.65\exp(-2A)$. Note that if the
potential $U(r)$ has the dipole-dipole form (1) up to very short distances, we
have to put $A=0$ in the expressions for the coefficients
$\zeta_{1},\,\zeta_{2},\,\zeta_{3},\,\zeta_{4}$. Considering the quasi2D
regime, this will be the case for $r_{*}$ greatly exceeding the length of the
sample in the tightly confined direction, $l_{0}$. Then, as one can see from
equations (89), (90), (91), and (92), the terms proportional to
$(k_{F}r_{*})^{2}$ are always negative in the considered limit $k_{F}r_{*}\ll
1$. These terms may become significant for $k_{F}r_{*}>0.3$.
## IV Zero sound
In the collisionless regime of the Fermi liquid at very low temperatures,
where the frequency of variations of the momentum distribution function
greatly exceeds the relaxation rate of quasiparticles, one has zero sound
waves. For these waves, variations $\delta n({\bf q},{\bf r},t)$ of the
momentum distribution are related to deformations of the Fermi surface, which
remains a sharp boundary between filled and empty quasiparticle states. At
$T\rightarrow 0$ the equilibrium distribution $n_{\bf q}$ is the step function
(51), so that $\partial n_{\bf q}/\partial{\bf q}=-{\bf
n}\delta(q-k_{F})=-\hbar{\bf v}\delta(\epsilon_{q}-\epsilon_{F})$, where ${\bf
v}=v_{F}{\bf n}$, with ${\bf n}$ being a unit vector in the direction of ${\bf
q}$. Then, searching for the variations $\delta n$ in the form:
$\delta n({\bf q},{\bf r},t)=\delta(\epsilon_{q}-\epsilon_{F})\nu({\bf
n})\exp{i({\bf kr}-\omega t)}$
and using Eq. (52), from the kinetic equation in the collisionless regime:
$\displaystyle\frac{\partial\delta n}{\partial
t}+\mathbf{v}\cdot\frac{\partial\delta n}{\partial\mathbf{r}}-\frac{\partial
n_{\bf
q}}{\partial\mathbf{q}}\cdot\frac{\partial\delta\epsilon_{q}}{\hbar\partial\mathbf{r}}=0,$
one obtains an integral equation for the function $\nu({\bf n})$ representing
displacements of the Fermi surface in the direction of ${\bf n}$ Landau :
$\displaystyle(\omega-
v_{F}\mathbf{n}\cdot\mathbf{k})\nu(\mathbf{n})=\frac{k_{F}}{(2\pi)^{2}\hbar}\mathbf{n}\cdot\mathbf{k}\int
F(k_{F}\mathbf{n},k_{F}\mathbf{n^{\prime}})\nu(\mathbf{n^{\prime}})d{\bf
n}^{\prime}.$
Introducing the velocity of zero sound $u_{0}=\omega/k$ and dividing both
sides of this equation by $v_{F}k$ we have:
$(s-\cos\theta)\nu(\theta)=\frac{m^{*}\cos\theta}{(2\pi\hbar)^{2}}\int_{0}^{2\pi}\tilde{F}(\theta-\theta^{\prime})\nu(\theta^{\prime})d\theta^{\prime},$
(93)
where $s=u_{0}/v_{F}$, and $\theta,\,\theta^{\prime}$ are the angles between
${\bf k}$ and ${\bf n},\,{\bf n}^{\prime}$, so that $\theta-\theta^{\prime}$
is the angle between ${\bf n}$ and ${\bf n}^{\prime}$. The dependence of the
interaction function of quasiparticles
$\tilde{F}=\tilde{F}^{(1)}+\tilde{F}^{(2)}_{1}+\tilde{F}^{(2)}_{2}$ on
$(\theta-\theta^{\prime})$ follows from Eqs. (81), (84), and (85) in which one
has to replace $\theta$ by $(\theta-\theta^{\prime})$.
The solution of equation (93) gives the function $\nu(\theta)$ and the
velocity of zero sound $u_{0}$, and in principle one may obtain several types
of solutions. It is important to emphasize that undamped zero sound requires
the condition $s>1$, i.e. the sound velocity should exceed the Fermi velocity
Landau . We will discuss this issue below.
For solving Eq. (93) we represent the interaction function $\tilde{F}$ as a
sum of the part proportional to $k_{F}r_{*}$ and the part proportional to
$(k_{F}r_{*})^{2}$. As follows from Eqs. (81), (84), and (85), we have:
$\\!\\!\\!\\!\tilde{F}(\theta\\!-\\!\theta^{\prime})\\!=\\!\frac{4\pi\hbar^{2}}{m}k_{F}r_{*}\\!\left|\sin\frac{\theta\\!-\\!\theta^{\prime}}{2}\right|+\frac{2\hbar^{2}}{m}(k_{F}r_{*}\\!)^{2}\Phi(\theta\\!-\\!\theta^{\prime}),\\!\\!\\!\\!$
(94)
where the function $\Phi(\theta-\theta^{\prime})$ is given by the sum of three
terms. The first one is the term in the second line of Eq. (81), the second
term is the expression in the square brackets in Eq. (84), the third term is
the one in curly brackets in Eq. (85), and we should replace $\theta$ by
$(\theta-\theta^{\prime})$ in all these terms. It is important that the
function $\Phi(\theta-\theta^{\prime})$ does not have singularities and
$\Phi(0)=\Phi(\pm 2\pi)=2\pi$. Using Eq. (94) the integral equation (93) is
reduced to the form:
$\displaystyle(s-\cos\theta)\nu(\theta)$
$\displaystyle=\beta\cos\theta\int_{0}^{2\pi}\nu(\theta^{\prime})\left|\sin\frac{\theta-\theta^{\prime}}{2}\right|d\theta^{\prime}$
$\displaystyle+\frac{\beta^{2}m}{2m^{*}}\cos\theta\int_{0}^{2\pi}\nu(\theta^{\prime})\Phi(\theta-\theta^{\prime})d\theta^{\prime},$
(95)
where $\beta=(m^{*}/\pi m)k_{F}r_{*}\ll 1$.
We now represent the function $\nu(\theta)$ as
$\nu(\theta)=\sum_{p=0}^{\infty}C_{p}\cos p\theta.$ (96)
Then, integrating over $d\theta^{\prime}$ in Eq. (95), multiplying both sides
of this equation by $\cos{j\theta}$ and integrating over $d\theta$, we obtain
a system of linear equations for the coefficients $C_{j}$. We write this
system for the coefficients $\eta_{j}=C_{j}(1-\beta/(j^{2}-1/4))$, so that
$C_{j}=\eta_{j}(1+\beta U_{j})$, where $U_{j}=(j^{2}-1/4-\beta)^{-1}$. The
system reads:
$\displaystyle\\!\\!\\!\\!(s-1)(1+\beta
U_{0})\eta_{0}+[\eta_{0}-\frac{1}{2}\eta_{1}]+\beta
U_{0}\eta_{0}=\frac{\beta^{2}}{2}{\bar{\Phi}}_{0};\\!$ (97)
$\displaystyle\\!\\!\\!\\!(s\\!-\\!1)(1\\!+\\!\beta
U_{1})\eta_{1}\\!+\\![\eta_{1}-\eta_{0}-\frac{1}{2}\eta_{2}]+\beta
U_{1}\eta_{1}\\!\\!=\\!\frac{\beta^{2}}{2}{\bar{\Phi}}_{1};$ (98)
$\displaystyle\\!\\!\\!\\!(\\!s\\!\\!-\\!\\!1\\!)(\\!1\\!\\!+\\!\\!\beta
U_{j}\\!)\eta_{j}\\!\\!+\\!\\![\eta_{j}\\!\\!-\\!\\!\frac{1}{2}\\!(\\!\eta_{j\\!-\\!1}\\!\\!+\\!\eta_{j\\!+\\!1)\\!}\\!]\\!\\!+\\!\\!\beta
U_{j}\eta_{j}\\!\\!=\\!\\!\frac{\beta^{2}}{2}\\!{\bar{\Phi}}_{j}\\!;\,j\\!\\!\geq\\!\\!2,\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!\\!$
(99)
where
$\\!\\!{\bar{\Phi}}_{j}\\!=\\!\frac{\tilde{C}_{j}}{\pi}\\!\\!\int_{0}^{2\pi}\\!\\!\\!\\!\\!\\!\\!\\!\cos\theta\,\cos
j\theta d\theta\\!\\!\int_{0}^{2\pi}\\!\\!\sum_{p=0}^{\infty}C_{p}\cos
p\theta^{\prime}\Phi(\theta-\theta^{\prime})d\theta^{\prime}\\!,\\!\\!\\!\\!\\!\\!\\!$
(100)
with $\tilde{C}_{j}=1$ for $j\geq 1$ and $\tilde{C}_{0}=1/2$, and we put
$m^{*}=m$ in the terms proportional to $\beta^{2}$.
In the weakly interacting regime the velocity of zero sound is close to the
Fermi velocity and, hence, we have $(s-1)\ll 1$ (see, e.g. Landau ). Since
$\beta\ll 1$, we first find coefficients $\eta_{j}$ omitting the terms
proportional to $\beta$ and $\beta^{2}$ in Eqs. (97)-(99). For $j\gg 1$
equation (99) then becomes:
$(s-1)\eta_{j}-\frac{1}{2}\frac{d^{2}\eta_{j}}{dj^{2}}=0,$
and searching for $s>1$ we may write
$\eta_{j}\simeq\exp\\{-\sqrt{2(s-1)}j\\};\,\,\,j\gg 1.$ (101)
If $j\ll 1/\sqrt{s-1}$, then we may also omit the terms proportional to
$(s-1)$ in the system of linear equations for $\eta_{j}$ (97)-(99). The system
then takes the form:
$\displaystyle\eta_{0}-\frac{1}{2}\eta_{1}=0;$
$\displaystyle\eta_{1}-\eta_{0}-\frac{1}{2}\eta_{2}=0;$
$\displaystyle\eta_{j}-\frac{1}{2}[\eta_{j-1}+\eta_{j+1}]=0;\,\,\,\,j\geq 2.$
Without loss of generality we may put $\eta_{0}=1/2$. This immediately gives
$\eta_{j}=1$ for $j\geq 1$, which is consistent with Eq. (101) at $j\ll
1/\sqrt{s-1}$. We thus have the zero order solution:
$\begin{cases}\eta_{0}=1/2;\\\ \eta_{j}=1;\,\,\,1\leq j\ll
1/\sqrt{s-1}.\end{cases}$ (102)
In order to find the coefficients $\eta_{j}$ taking into account the terms
linear in $\beta$, we consider $j$ such that $\beta
U_{j}\sim\beta/j^{2}\gg(s-1)$, i.e. $j\ll\sqrt{\beta/(s-1)}$. Then we may omit
the terms proportional to $(s-1)$ in equations (97)-(99). Omitting also the
terms proportional to $\beta^{2}$ this system of equations becomes:
$\displaystyle\eta_{0}-\frac{1}{2}\eta_{1}+\beta U_{0}\eta_{0}=0;$ (103)
$\displaystyle\eta_{1}-\eta_{0}-\frac{1}{2}\eta_{2}+\beta U_{1}\eta_{1}=0;$
(104) $\displaystyle\eta_{j}-\frac{1}{2}[\eta_{j-1}+\eta_{j+1}]+\beta
U_{j}\eta_{j}=0;\,\,j\geq 2.$ (105)
Putting again $\eta_{0}=1/2$ the solution of these equations reads:
$\displaystyle\eta_{1}=1+\beta U_{0};$ $\displaystyle\eta_{j}=1+\beta
jU_{0}+2\beta\sum_{p=1}^{j-1}(j-p)U_{p};\,\,\,j\geq 2.$
Confining ourselves to terms linear in $\beta$ we put $U_{p}=1/(p^{2}-1/4)$
and, hence, $U_{0}=-4$. Then, using the relation
$\sum_{p=1}^{j-1}\frac{1}{p^{2}-1/4}=\frac{4(j-1)}{2j-1},$
which is valid for $j\geq 2$, we obtain:
$\displaystyle\begin{cases}&\eta_{0}=\frac{1}{2};\\\ \\\
&\eta_{1}=1-4\beta;\\\ \\\
&\eta_{j}=1-2\beta\left\\{\frac{2j}{2j-1}+\sum_{p=1}^{j-1}\frac{p}{p^{2}-1/4}\right\\};\,\,\,j\geq
2.\end{cases}$ (106)
For $j\gtrsim\sqrt{\beta/(s-1)}$ we should include the terms proportional to
$(s-1)$ in Eq. (105). This leads to the solution in the form of the decaying
Bessel function:
$\eta_{j}\simeq\sqrt{2(s-1)/\pi}K_{\sqrt{1/4+\beta}}(\sqrt{s-1}j)$, which for
small $\beta$ is practically equivalent to Eq. (101).
We now make a summation of equations (97)-(99) from $j=0$ to $j=j_{*}\ll
1/\sqrt{s-1}$. The summation of the second terms of these equations gives
$\sqrt{(s-1)/2}$, whereas the contribution of the terms proportional to
$(s-1)$ is much smaller and will be omitted. The sums
$\sum_{j=0}^{j_{*}}U_{j}\eta_{j}$ and $\sum_{j=0}^{j_{*}}{\bar{\Phi}}_{j}$
converge at $j\ll 1/\sqrt{s-1}$, and the upper limit of summation in these
terms can be formally replaced by infinity. We thus obtain a relation:
$\sqrt{\frac{s-1}{2}}+\sum_{j=0}^{\infty}\beta\eta_{j}U_{j}-\frac{\beta^{2}}{2}\sum_{j=0}^{\infty}{\bar{\Phi}}_{j}=0.$
(107)
Confining ourselves to contributions up to $\beta^{2}$, in the second term on
the left hand side of Eq. (107) we use coefficients $\eta_{j}$ given by Eqs.
(106), and write $U_{j}=1/(j^{2}-1/4)+\beta/(j^{2}-1/4)^{2}$. In the
expressions for ${\bar{\Phi}}_{j}$ we use $C_{0}=1/2$ and $C_{p}=1$ for $p\geq
1$. We then have:
$\\!\\!\sum_{j=0}^{\infty}\beta\eta_{j}U_{j}\\!=\\!-2\beta\\!+\\!\sum_{j=1}^{\infty}\frac{\beta}{\\!j^{2}\\!-\\!1/4}\\!+\\!\beta^{2}\\{\\!8\\!+\\!S_{1}\\!-\\!2S_{2}\\!-2S_{3}\\!\\}.\\!\\!\\!\\!$
(108)
The contribution linear in $\beta$ vanishes because
$\sum_{j=1}^{\infty}1/(j^{2}-1/4)=2$. The quantities $S_{1},\,S_{2}$, and
$S_{3}$ are given by
$\displaystyle\\!\\!\\!S_{1}=\sum_{j=1}^{\infty}\frac{1}{(j^{2}-1/4)^{2}}=\pi^{2}-8;$
$\displaystyle\\!\\!\\!S_{2}\\!=\\!\sum_{j=2}^{\infty}\frac{1}{j^{2}\\!-\\!1/4}\\!\sum_{p=1}^{j-1}\frac{p}{p^{2}\\!-\\!1/4}\\!=\\!\sum_{j=1}\frac{j}{(j^{2}\\!-\\!1/4)(j\\!+\\!1/2)};$
$\displaystyle\\!\\!\\!S_{3}=\sum_{j=1}^{\infty}\frac{j}{(j-1/2)(j^{2}-1/4)},$
so that
$S_{2}+S_{3}=\sum_{j=1}^{\infty}\frac{2j^{2}}{(j^{2}-1/4)^{2}}=\frac{\pi^{2}}{2}.$
We thus see that the contribution quadratic in $\beta$ also vanishes because
the term in the curly brackets in Eq. (108) is exactly equal to zero. Hence,
we have $\sum_{j=0}^{\infty}\beta\eta_{j}U_{j}=0$ up to terms proportional to
$\beta^{2}$.
The sum in the third term on the left hand side of Eq. (107), after putting
$C_{0}=1/2$ and $C_{p}=1$ for $p\geq 1$ in the relations for
${\bar{\Phi}}_{j}$, reduces to
$\displaystyle\sum_{j=0}^{\infty}{\bar{\Phi}}_{j}$
$\displaystyle=\frac{1}{4\pi}\\!\sum_{j=-\infty}^{\infty}\\!\int_{0}^{2\pi}\\!\\!\\!\\!\cos\theta\,\cos{j\theta}d\theta$
$\displaystyle\times\sum_{p=-\infty}^{\infty}\\!\int_{0}^{2\pi}\\!\\!\\!\\!\cos{p\theta^{\prime}}\Phi(\theta\\!-\\!\theta^{\prime})d\theta^{\prime}.\\!\\!\\!\\!$
(109)
For $\theta$ in the interval $0\leq\theta\leq 2\pi$ we have a relation:
$\sum_{j=-\infty}^{\infty}\cos{j\theta}=\pi[\delta(\theta)+\delta(\theta-2\pi)],$
which transforms Eq. (IV) to
$\sum_{j=0}^{\infty}{\bar{\Phi}}_{j}=\frac{\pi}{4}[2\Phi(0)+\Phi(2\pi)+\Phi(-2\pi)]=2\pi^{2},$
(110)
and equation (107) becomes:
$\displaystyle\sqrt{\frac{s-1}{2}}-\beta^{2}\pi^{2}=0.$
This gives $s=1+2(\beta\pi)^{4}$, and recalling that $\beta=k_{F}r_{*}/\pi$
(we put $m^{*}=m$) we obtain for the velocity of zero sound:
$u_{0}=v_{F}[1+2(k_{F}r_{*})^{4}].$ (111)
Note that in contrast to the 3D two-species Fermi gas with a weak repulsive
contact interaction (scattering length $a$), where the correction
$(u_{0}-v_{F})$ exponentially depends on $k_{F}a$, for our 2D dipolar gas we
obtained a power law dependence. This is a consequence of dimensionality of
the system.
It is important that confining ourselves to only the leading part of the
interaction function $\tilde{F}$, which is proportional to $k_{F}r_{*}$ and is
given by the first term of Eq. (81), we do not obtain undamped zero sound
($s>1$) comment . This corresponds to omitting the terms
$\beta^{2}{\bar{\Phi}}_{j}/2$ in equations (97)-(99) and is consistent with
numerical calculations Baranov . Only the many-body corrections to the
interaction function of quasiparticles, given by equations (84) and (85),
provide non-zero positive values of $\Phi(0)$ and $\Phi(\pm 2\pi)$, thus
leading to a positive value of $(u_{0}-v_{F})$. One then sees that many-body
effects are crucial for the propagation of zero sound.
In principle, we could obtain the result of Eq. (111) in a simpler way,
similar to that used for the two-species Fermi gas with a weak repulsive
interaction (see, e.g. Landau ). Representing the function $\nu(\theta)$ as
$\nu(\theta)=\cos{\theta}\tilde{\nu}(\theta)/(s-\cos{\theta})$ we transform
Eq. (93) to the form:
$\tilde{\nu}(\theta)=\frac{m^{*}}{(2\pi\hbar)^{2}}\int_{0}^{2\pi}\frac{\tilde{F}(\theta-\theta^{\prime})\tilde{\nu}(\theta^{\prime})\cos{\theta^{\prime}}}{s-\cos{\theta^{\prime}}}d\theta^{\prime}.$
(112)
Since $s$ is close to unity, it looks reasonable to assume that the main
contribution to the integral in Eq. (112) comes from $\theta^{\prime}$ close
to zero and to $2\pi$. Using the fact that
$\tilde{F}(\theta)=\tilde{F}(2\pi-\theta)$ we then obtain:
$\tilde{\nu}(\theta)=\frac{m^{*}\tilde{F}(\theta)\tilde{\nu}(0)}{4\pi\hbar^{2}}\sqrt{\frac{2}{s-1}}.$
(113)
We now take the limit $\theta\rightarrow 0$ and substitute
$\tilde{F}(0)=(4\pi\hbar^{2}/m)(k_{F}r_{*})^{2}$ as follows froms Eqs. (81),
(84), and (85). Putting $m^{*}=m$ we then obtain $s=1+2(k_{F}r_{*})^{4}$ and
arrive at Eq. (111).
Note, however, that for very small $\theta$ or $\theta$ very close to $2\pi$
the dependence $\tilde{F}(\theta)$ is very steep. For $\theta\rightarrow 0$
the leading part of the interaction function, which is linear in $k_{F}r_{*}$,
vanishes, and only the quadratic part contributes to $\tilde{F}(0)$.
Therefore, strictly speaking the employed procedure of calculating the
integral in Eq. (112) is questionable for very small $\theta$. This prompted
us to make the analysis based on representing $\nu(\theta)$ in the form (96)
and on solving the system of linear equations (97)-(99).
Equation (112) is useful for understanding why undamped zero sound requires
the condition $s>1$ so that $u_{0}>v_{F}$. For $s<1$ there is a pole in the
integrand of Eq. (112), which introduces an imaginary part of the integral. As
a result, the zero sound frequency $\omega$ will also have an imaginary part
at real momenta $k$, which means the presence of damping (see, e.g. Landau ).
We could also consider an odd function $\nu(\theta)$, namely such that
$\nu(2\pi-\theta)=-\nu(\theta)$ and $\nu(0)=\nu(2\pi)=0$. In this case,
however, we do not obtain an undamped zero sound.
## V Concluding remarks
We have shown that (single-component) fermionic polar molecules in two
dimensions constitute a novel Fermi liquid, where many-body effects play an
important role. For dipoles oriented perpendicularly to the plane of
translational motion, the many-body effects provide significant corrections to
thermodynamic functions. Revealing these effects is one of the interesting
goals of up-coming experimental studies. The investigation of the full
thermodynamics of 2D polar molecules, including many-body effects, can rely on
the in-situ imaging technique as it has been done for two-component atomic
Fermi gases Salomon1 ; Salomon2 . This method can also be extended to 2D
systems for studying thermodynamic quantities Zwierlein ; Dalibard2 . Direct
imaging of a 3D pancake-shaped dipolar molecular system has been recently
demonstrated at JILA Ye2 . For 2D polar molecules discussed in our paper,
according to equations (89)-(91), the contribution of many-body corrections
proportional to $(k_{F}r_{*})^{2}$ can be on the level of $10\%$ or $20\%$ for
$k_{F}r_{*}$ close to $0.5$. Thus, finding many-body effects in their
thermodynamic properties looks feasible.
It is even more important that the many-body effects are responsible for the
propagation of zero sound waves in the collisionless regime of the 2D Fermi
liquid of polar molecules with dipoles perpendicular to the plane of
translational motion. This is shown in Section IV of our paper, whereas mean-
field calculations do not find undamped zero sound Baranov . Both
collisionless and hydrodynamic regimes are achievable in on-going experiments.
This is seen from the dimensional estimate of the relaxation rate of
quasiparticles. At temperatures $T\ll\epsilon_{F}$ the relaxation of a non-
equilibrium distribution of quasiparticles occurs due to binary collisions of
quasiparticles with energies in a narrow interval near the Fermi surface. The
width of this interval is $\sim T$ and, hence, the relaxation rate contains a
small factor $(T/\epsilon_{F})^{2}$ (see, e.g. Landau ). Then, using the Fermi
Golden rule we may write the inverse relaxation time as
$\tau^{-1}\sim(g_{eff}^{2}/\hbar)(m/\hbar^{2})n(T/\epsilon)^{2}$, where $n$ is
the 2D particle density, the quantity $\sim m/\hbar^{2}$ represents the
density of states on the Fermi surface, and the quantity $g_{eff}$ is the
effective interaction strength. Confining ourselves to the leading part of
this quantity, from Eqs. (III.2) and (79) we have
$g_{eff}\sim\hbar^{2}k_{F}r_{*}/m$. We thus obtain:
$\frac{1}{\tau}\sim\frac{\hbar
n}{m}(k_{F}r_{*})^{2}\left(\frac{T}{\epsilon_{F}}\right)^{2}.$ (114)
Note that as $\epsilon_{F}\approx\hbar^{2}k_{F}^{2}/2m\approx
2\pi\hbar^{2}n/m$, for considered temperatures $T\ll\epsilon_{F}$ the
relaxation time $\tau$ is density independent. Excitations with frequencies
$\omega\ll 1/\tau$ are in the hydrodynamic regime, where on the length scale
smaller than the excitation wavelength and on the time scale smaller than
$1/\omega$ the system reaches a local equilibrium. On the other hand,
excitations with frequencies $\omega\gg 1/\tau$ are in the collisionless
regime. Assuming $T\sim 10$nK, for KRb molecules characterized by the dipole
moment $d\simeq 0.25$ D in the electric field of $5$kV/cm as obtaind in the
JILA experiments, we find $\tau$ on the level of tens of milliseconds. The
required condition $T\ll\epsilon_{F}$ is satisfied for $\epsilon_{F}\gtrsim
70$ nK, which corresponds to $n\gtrsim 2\cdot 10^{8}$ cm-2. In such conditions
excitations with frequencies of the order of a few Hertz or lower will be in
the hydrodynamic regime, and excitations with larger frequencies in the
collisionless regime.
The velocity of zero sound is practically equal to the Fermi velocity
$v_{F}=\hbar k_{F}/m^{*}$. This is clearly seen from Eq. (111) omitting a
small correction proportional to $(k_{F}r_{*})^{4}$. Then, using Eq. (92) for
the effective mass and retaining only corrections up to the first order in
$k_{F}r_{*}$, we have:
$u_{0}\simeq\frac{\hbar k_{F}}{m}\left(1+\frac{4}{3\pi}k_{F}r_{*}\right).$
(115)
In the hydrodynamic regime the sound velocity is:
$u=\sqrt{\frac{N}{m}\frac{\partial\mu}{\partial N}}\simeq\frac{\hbar
k_{F}}{m}\left(1+\frac{8}{3\pi}k_{F}r_{*}\right),$ (116)
where we used Eq. (III.4) for $\partial\mu/\partial N$ and retained
corrections up to the first order in $k_{F}r_{*}$. The hydrodynamic velocity
$u$ is slightly larger than the velocity of zero sound $u_{0}$, and the
difference is proportional to the interaction strength. This is in sharp
contrast with the 3D two-component Fermi gas, where $u_{0}\approx
v_{F}>u\approx v_{F}/\sqrt{3}$.
We thus see that it is not easy to distinguish between the hydrodynamic and
collisionless regimes from the measurement of the sound velocity. A promising
way to do so can be the observation of damping of driven excitations, which in
the hydrodynamic regime is expected to be slower. Another way is to achieve
the values of $k_{F}r_{*}$ approaching unity and still discriminate between
$u_{0}$ and $u$ in the measurement of the sound velocity. For example, in the
case of dipoles perpendicular to the plane of their translational motion the
two velocities are different from each other by about $20\%$ at
$k_{F}r_{*}\simeq 0.5$. These values of $k_{F}r_{*}$ are possible if the 2D
gas of dipoles still satisfies the Pomeranchuk criteria of stability. These
criteria require that the energy of the ground state corresponding to the
occupation of all quasiparticle states inside the Fermi sphere, remains the
minimum energy under an arbitrarily small deformation of the Fermi sphere. The
generalization of the Pomeranchuk stability criteria to the case of the 2D
single-component Fermi liquid with dipoles perpendicular to the plane of their
translational motion reads:
$1+\frac{m^{*}}{(2\pi\hbar)^{2}}\int_{0}^{2\pi}\tilde{F}(\theta)\cos{j\theta}\,d\theta>0,$
(117)
and this inequality should be satisfied for any integer $j$. As has been found
in Ref. Baranov , the Pomeranchuk stability criteria (117) are satisfied for
$k_{F}r_{*}$ approaching unity from below if the interaction function of
quasiparticles contains only the first term of Eq. (81), which is the leading
mean field term. We have checked that the situation with the Pomeranchuk
stability does not change when we include the full expression for the
interaction function,
$\tilde{F}(\theta)=\tilde{F}^{(1)}(\theta)+\tilde{F}^{(2)}_{1}(\theta)+\tilde{F}^{(2)}_{2}(\theta)$,
following from Eqs. (81), (84), and (85). Thus, achieving $k_{F}r_{*}$
approaching unity looks feasible. For KRb molecules with the (oriented) dipole
moment of $0.25$ D the value $k_{F}r_{*}\approx 0.5$ requires densities
$n\approx 2\cdot 10^{8}$ cm-2.
Finally, we would like to emphasize once more that our results are applicable
equally well for the quasi2D regime, where the dipole-dipole length $r_{*}$ is
of the order of or smaller than the confinement length
$l_{0}=(\hbar/m\omega_{0})^{1/2}$, with $\omega_{0}$ being the frequency of
the tight confinement. The behavior at distances $r\lesssim l_{0}$ is
contained in the coefficient $A$ defined in Eq. (18). Therefore, the results
for the velocity of zero sound which is independent of $A$, are universal in
the sense that they remain unchanged when going from $r_{*}\gg l_{0}$ to
$r_{*}\lesssim l_{0}$. The only requirement is the inequality $k_{F}l_{0}\ll
1$. It is, however, instructive to examine the ratio $r_{*}/l_{0}$ that can be
obtained in experiments with ultracold polar molecules. Already in the JILA
experiments using the tight confinement of KRb molecules with frequency
$\omega_{0}\approx 30$ kHz and achieving the average dipole moment $d\simeq
0.25$ D in electric fields of $5$ kV/cm, we have $r_{*}\simeq 100$ nm and
$l_{0}\simeq 50$ nm so that $r_{*}/l_{0}\simeq 2$. A decrease of the
confinement frequency to $5$ kHz and a simultaneous decrease of the dipole
moment by a factor of 2 leads to $r_{*}/l_{0}\sim 0.2$. On the other hand, for
$d$ close to $0.5$ (which is feasible to obtain for other molecules) one can
make the ratio $r_{*}/l_{0}$ close to $10$ at the same confinement length.
## Acknowledgements
We are grateful to M.A. Baranov and S.I. Matveenko for fruitful discussions.
We acknowledge support from EPSRC Grant No. EP/F032773/1, from the IFRAF
Institute, and from the Dutch Foundation FOM. This research has been supported
in part by the National Science Foundation under Grant No. NSF PHYS05-51164.
LPTMS is a mixed research unit No. 8626 of CNRS and Université Paris Sud.
## Appendix A Direct calculation of the first order contribution to the
interaction energy
For directly calculating the first order (mean field) contribution to the
interaction energy $\tilde{E}^{(1)}$ (74), we represent it as
$\tilde{E}^{(1)}=\tilde{E}^{(1)}_{1}+\tilde{E}^{(1)}_{2}$ where
$\displaystyle\tilde{E}^{(1)}_{1}=\int{\bar{f}}^{(1)}\left(\frac{|{\bf
k}_{1}-{\bf k}_{2}|}{2}\right)n_{{\bf k}_{1}}n_{{\bf
k}_{2}}\frac{d^{2}k_{1}d^{2}k_{2}}{(2\pi)^{4}},$ (118)
$\displaystyle\tilde{E}^{(1)}_{2}=\int{\bar{f}}^{(2)}\left(\frac{|{\bf
k}_{1}-{\bf k}_{2}|}{2}\right)n_{{\bf k}_{1}}n_{{\bf
k}_{2}}\frac{d^{2}k_{1}d^{2}k_{2}}{(2\pi)^{4}},$ (119)
and the amplitudes ${\bar{f}}^{(1)}$ and ${\bar{f}}^{(2)}$ are given by Eqs.
(79) and (80), respectively. In the calculation of the integrals for
$\tilde{E}^{(1)}_{1}$ and $\tilde{E}^{(1)}_{2}$ we turn to the variables ${\bf
x}=({\bf k}_{1}-{\bf k}_{2})/2k_{F}$ and ${\bf y}=({\bf k}_{1}+{\bf
k}_{2})/2k_{F}$, so that $d^{2}k_{1}d^{2}k_{2}=8\pi
k_{F}^{4}d^{2}xd^{2}yd\varphi$, where $\varphi$ is the angle between the
vectors ${\bf x}$ and ${\bf y}$, and the integration over $d\varphi$ should be
performed from $0$ to $2\pi$. The distribution functions $n_{{\bf k}_{1}}$ and
$n_{{\bf k}_{2}}$ are the step functions (51). The integration over $dk_{1}$
and $dk_{2}$ from $0$ to $k_{F}$ corresponds to the integration over $dy$ from
$0$ to $y_{0}(x,\varphi)=-x|\cos\varphi|+\sqrt{1-x^{2}\sin 2\varphi}$ and over
$dx$ from $0$ to $1$. Using Eq. (79) we reduce Eq. (118) to
$\tilde{E}^{(1)}_{1}=\frac{S\hbar^{2}k_{F}^{4}}{\pi^{2}m}k_{F}r_{*}I_{1},$
(120)
where
$\displaystyle
I_{1}=\int_{0}^{2\pi}d\varphi\int_{0}^{1}x^{2}dx\int_{0}^{y_{0}(x,\varphi)}ydy=\frac{1}{2}\int_{0}^{2\pi}d\varphi\int_{0}^{1}x^{2}dx$
$\displaystyle\times[1-2|\cos\varphi|\sqrt{1-x^{2}\sin^{2}\varphi}+x^{2}(\cos^{2}\varphi-\sin^{2}\varphi)].$
The last term of the second line vanishes, and the integration of the first
two terms over $d\varphi$ and $dx$ gives:
$I_{1}=\int_{0}^{1}x^{2}\left(\pi-2x\sqrt{1-x^{2}}-2\arcsin{x}\right)=\frac{8}{45}.$
Then Eq. (120) yields:
$\tilde{E}^{(1)}_{1}=\frac{8S}{45\pi^{2}}\frac{\hbar^{2}k_{F}^{4}}{m}k_{F}r_{*}=\frac{N^{2}}{S}\frac{128}{45}\frac{\hbar^{2}k_{F}^{2}}{m}k_{F}r_{*},$
(121)
which exactly coincides with the second term of the first line of Eq. (91).
Using Eq. (80) the contribution $\tilde{E}^{(1)}_{2}$ takes the form:
$\\!\\!\\!\\!\tilde{E}^{(\\!1\\!)}_{2}\\!\\!\\!=\\!\frac{S\hbar^{2}\\!k_{F}^{4}}{2\pi^{2}m}(k_{F}r_{*}\\!)^{2}\\!\left\\{\\!\left[\ln(\xi
k_{F}r_{*}\\!)\\!-\\!\frac{25}{12}\\!+\\!3\ln{2}\right]\\!I_{2}\\!+\\!I_{3}\\!\right\\}\\!\\!,\\!\\!\\!\\!\\!\\!\\!\\!$
(122)
where the integrals $I_{2}$ and $I_{3}$ are given by
$\displaystyle I_{2}$
$\displaystyle=\int_{0}^{2\pi}d\varphi\int_{0}^{1}x^{3}dx\int_{0}^{y_{0}(x,\varphi)}ydy=\frac{1}{2}\int_{0}^{2\pi}d\varphi\int_{0}^{1}x^{3}dx$
$\displaystyle\times\big{[}1-2|\cos\varphi|x\sqrt{1-x^{2}\sin^{2}\varphi}+x^{2}(\cos^{2}\varphi-\sin^{2}\varphi)\big{]}$
$\displaystyle=\frac{1}{2}\int_{0}^{1}x^{3}[2\pi-4x\sqrt{1-x^{2}}-4\arcsin{x}]dx=\frac{\pi}{32},$
and
$\displaystyle
I_{3}=\int_{0}^{2\pi}d\varphi\int_{0}^{1}x^{3}\ln{x}dx\int_{0}^{y_{0}(x,\varphi)}ydy=\frac{1}{2}\int_{0}^{2\pi}d\varphi\int_{0}^{1}dx$
$\displaystyle\times
x^{3}\ln{x}\big{[}1-2|\cos\varphi|x\sqrt{1-x^{2}\sin^{2}\varphi}+x^{2}(\cos^{2}\varphi-\sin^{2}\varphi)\big{]}$
$\displaystyle=\frac{1}{2}\int_{0}^{1}x^{3}\ln{x}[2\pi-4x\sqrt{1-x^{2}}-4\arcsin{x}]dx$
$\displaystyle=\frac{\pi}{32}\left(\frac{1}{6}-\ln{2}\right).$
Substituting the calculated $I_{2}$ and $I_{3}$ into Eq. (122) we obtain:
$\displaystyle\tilde{E}^{(1)}_{2}=\frac{S\hbar^{2}k_{F}^{4}}{64\pi
m}(k_{F}r_{*})^{2}\left[\ln(4\xi k_{F}r_{*})-\frac{23}{12}\right]$
$\displaystyle=\frac{N^{2}}{S}\frac{\pi\hbar^{2}}{4m}(k_{F}r_{*})^{2}\left[\ln(4\xi
k_{F}r_{*})-\frac{23}{12}\right].$ (123)
This exactly reproduces the third term of the first line of Eq. (91).
## Appendix B Calculation of the interaction function $\tilde{F}_{1}^{(2)}$
The interaction function $\tilde{F}_{1}^{(2)}$ is the second variational
derivative of the many-body contribution to the interaction energy,
$\tilde{E}^{(2)}_{1}$ (82), with respect to the momentum distribution
function. It can be expressed as
$\tilde{F}_{1}^{(2)}(\mathbf{k},\mathbf{k^{\prime}})=-\frac{2\hbar^{2}}{m}(k_{F}r_{*})^{2}(\tilde{I}_{1}+\tilde{I}_{2}+\tilde{I}_{3}),$
(124)
where
$\displaystyle\tilde{I}_{1}=2\\!\int_{|\mathbf{k_{1}}|<k_{F}}\\!\frac{d^{2}k_{1}}{k_{F}^{2}}\frac{|\mathbf{k}-\mathbf{k_{1}}|^{2}}{\mathbf{k^{2}}\\!+\\!\mathbf{k^{\prime
2}}\\!-\\!\mathbf{k^{2}_{1}}\\!-\\!\mathbf{k^{2}_{2}}}\delta_{\mathbf{k}\\!+\\!\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{1}}\\!-\\!\mathbf{k_{2}}},$
(125)
$\displaystyle\tilde{I}_{2}=2\\!\int_{|\mathbf{k_{1}}|<k_{F}}\\!\frac{d^{2}k_{1}}{k_{F}^{2}}\frac{|\mathbf{k}-\mathbf{k^{\prime}}|^{2}}{\mathbf{k^{2}}\\!+\\!\mathbf{k^{2}_{1}}\\!-\\!\mathbf{k^{\prime
2}}\\!-\\!\mathbf{k^{2}_{2}}}\delta_{\mathbf{k}\\!+\\!\mathbf{k_{1}}\\!-\\!\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{2}}},$
(126)
$\displaystyle\tilde{I}_{3}=2\\!\int_{|\mathbf{k_{1}}|<k_{F}}\\!\frac{d^{2}k_{1}}{k_{F}^{2}}\frac{|\mathbf{k_{1}}-\mathbf{k^{\prime}}|^{2}}{\\!\mathbf{k^{2}_{1}}\\!+\\!\mathbf{k^{2}}\\!-\\!\mathbf{k^{\prime
2}}\\!-\\!\mathbf{k^{2}_{2}}}\delta_{\mathbf{k_{1}}\\!+\\!\mathbf{k}\\!-\\!\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{2}}},$
(127)
and the presence of the Kronecker symbols $\delta_{\bf q}$ reflects the
momentum conservation law. On the Fermi surface we put $|{\bf k}|=|{\bf
k}^{\prime}|=k_{F}$ and denote the angle between ${\bf k}$ and ${\bf
k}^{\prime}$ as $\theta$. Due to the symmetry property:
$F(\mathbf{k},\mathbf{k^{\prime}})=F(\mathbf{k^{\prime}},\mathbf{k})$ we have
$F(\theta)=F(2\pi-\theta)$ and may consider $\theta$ in the interval from $0$
to $\pi$.
In order to calculate the integral $\tilde{I}_{1}$, we use the quantities
$\mathbf{s}=(\mathbf{k}+\mathbf{k^{\prime}})/2k_{F}$ and
$\mathbf{m}=(\mathbf{k}-\mathbf{k^{\prime}})/2k_{F}$ and turn to the variable
$\mathbf{x}=(\mathbf{k_{1}}-\mathbf{k_{2}})/2k_{F}=(2\mathbf{k_{1}}-\mathbf{s})/2k_{F}$.
For given vectors $\mathbf{k}$ and $\mathbf{k^{\prime}}$, the vectors
$\mathbf{s}$ and $\mathbf{m}$ are fixed and $|\mathbf{s}|=\cos(\theta/2)$,
$|\mathbf{m}|=\sin(\theta/2)$. The integral can then be rewritten as:
$\tilde{I}_{1}=\int\frac{m^{2}+x^{2}}{m^{2}-x^{2}}d^{2}x.$
The integration region is shown in Fig.1, where the distance between the
points $O_{1}$ and $O_{2}$ is ${\bf R}_{O_{1}O_{2}}=\mathbf{s}$. The distance
between the points $O_{1}$ and $N$ is ${\bf
R}_{O_{1}N}=\mathbf{k_{1}}/2k_{F}$, and ${\bf
R}_{NO_{2}}=\mathbf{k_{2}}/2k_{F}$, so that ${\bf R}_{ON}=\mathbf{x}/2$. The
quantity $|{\bf x}|$ changes from $0$ to $l_{1}(\varphi)$ where
$l_{1}^{2}(\varphi)+\cos^{2}\frac{\theta}{2}-2l_{1}(\varphi)\cos\frac{\theta}{2}\cos\varphi=1,$
and $l_{1}(\varphi)\cdot l_{1}(\varphi+\pi)=\sin^{2}(\theta/2)$, with
$\varphi$ being an angle between ${\bf m}$ and ${\bf x}$. In the polar
coordinates the integral $\tilde{I}_{1}$ takes the form:
$\tilde{I}_{1}=\int_{0}^{2\pi}d\varphi\int_{0}^{l_{1}(\varphi)}\left(-1+2\sin^{2}\frac{\theta}{2}\frac{1}{\sin^{2}(\theta/2)-x^{2}}\right)xdx,$
and after a straightforward integration we obtain:
$\tilde{I}_{1}=\pi\left(2\sin^{2}\frac{\theta}{2}\ln|\tan\frac{\theta}{2}|-1\right).$
(128)
Figure 1: (color online). Left: The integration area for $\tilde{I}_{1}$ (in
blue). The distance between the points $O_{1}$ and $P$ is ${\bf
R}_{O_{1}P}=\mathbf{k}/2k_{F}$, ${\bf R}_{PO_{2}}=\mathbf{k^{\prime}}/2k_{F},$
and ${\bf R}_{O_{1}N}=\mathbf{k_{1}}/2k_{F}$. Right: The integration area for
$\tilde{I}_{2}$ and $\tilde{I}_{3}$ (in red). The distance between the points
$O_{2}$ and $N$ is ${\bf R}_{O_{2}N}=\mathbf{k_{1}}/2k_{F}$, ${\bf
R}_{O_{1}P}={\bf k}/2k_{F}$, and ${\bf R}_{O_{2}P}={\bf k}^{\prime}/2k_{F}$.
In the integral $\tilde{I}_{2}$, using the variable
$\mathbf{y}=(\mathbf{k_{1}}+\mathbf{k_{2}})/2k_{F}$ we observe that it changes
from $0$ to $l_{2}(\tilde{\varphi})$ where
$l_{2}^{2}(\tilde{\varphi})+\sin^{2}\frac{\theta}{2}-2l_{2}(\tilde{\varphi})\sin\frac{\theta}{2}\cos(\tilde{\varphi})=1$
and
$l_{2}(\tilde{\varphi})-l_{2}(\tilde{\varphi}+\pi)=2\sin\frac{\theta}{2}\cos\tilde{\varphi}$,
with $\tilde{\varphi}$ being an angle between ${\bf y}$ and ${\bf m}$. We then
have:
$\displaystyle\tilde{I}_{2}$
$\displaystyle=-2\int\frac{m^{2}}{\mathbf{m}\cdot\mathbf{y}}d^{2}y=-2\sin\frac{\theta}{2}\int_{0}^{2\pi}d\tilde{\varphi}\int_{0}^{l_{2}(\tilde{\varphi})}\frac{dy}{\cos\tilde{\varphi}}$
$\displaystyle=-4\pi\sin^{2}\frac{\theta}{2}.$ (129)
For the integral $\tilde{I}_{3}$ we have:
$\displaystyle\tilde{I}_{3}=-\frac{1}{2}\int\frac{s^{2}+y^{2}-2\mathbf{s}\cdot\mathbf{y}}{\mathbf{m}\cdot\mathbf{y}}$
$\displaystyle=-\frac{1}{2\sin\frac{\theta}{2}}\int_{0}^{2\pi}\\!\\!\frac{d\tilde{\varphi}}{\cos\tilde{\varphi}}\\!\int_{0}^{l_{2}(\tilde{\varphi})}\\!\\!\\!dy\left[y^{2}+\cos^{2}\frac{\theta}{2}-2y\cos\frac{\theta}{2}\sin\tilde{\varphi}\right]$
$\displaystyle=-2\pi\left(\cos^{2}\frac{\theta}{2}+\frac{1}{3}\sin^{2}\frac{\theta}{2}\right),$
(130)
where we used the relation $l_{2}(\tilde{\varphi})=l_{2}(-\tilde{\varphi})$.
Using integrals $\tilde{I}_{1}$ (128), $\tilde{I}_{2}$ (129), and
$\tilde{I}_{3}$ (130) in Eq. (124), we obtain equation (84):
$\\!\\!\\!\tilde{F}_{1}^{(2)}(\theta)\\!\\!=\\!\\!\frac{2\hbar^{2}r^{2}_{*}k^{2}_{F}}{m}\\!\\!\left[\\!3\pi\\!+\\!2\pi\sin^{2}\frac{\theta}{2}\left(\frac{4}{3}\\!-\\!\ln|\tan\frac{\theta}{2}|\right)\\!\right].$
.
## Appendix C Calculation of the interaction function $\tilde{F}_{2}^{(2)}$
The interaction function $\tilde{F}_{2}^{(2)}$ is the second variational
derivative of the many-body contribution to the interaction energy,
$\tilde{E}^{(2)}_{2}$ (83), with respect to the momentum distribution. It
reads:
$\tilde{F}_{2}^{(2)}(\mathbf{k},\mathbf{k^{\prime}})=\frac{2\hbar^{2}}{m}(k_{F}r_{*})^{2}(I^{\prime}_{1}+I^{\prime}_{2}),$
(131)
where
$\displaystyle
I^{\prime}_{1}=2\\!\int_{|\mathbf{k_{1}}|<k_{F}}\\!\frac{d^{2}k_{1}}{k_{F}^{2}}\frac{|\mathbf{k}\\!-\\!\mathbf{k_{1}}|\cdot|\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{1}}|}{\mathbf{k^{2}}\\!+\\!\mathbf{k^{\prime
2}}\\!-\\!\mathbf{k^{2}_{1}}\\!-\\!\mathbf{k^{2}_{2}}}\delta_{\mathbf{k}\\!+\\!\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{1}}\\!-\\!\mathbf{k_{2}}},$
(132) $\displaystyle
I^{\prime}_{2}=4\\!\int_{|\mathbf{k_{1}}|<k_{F}}\\!\frac{d^{2}k_{1}}{k_{F}^{2}}\frac{|\mathbf{k}\\!-\\!\mathbf{k^{\prime}}|\cdot|\mathbf{k_{1}}\\!-\\!\mathbf{k^{\prime}}|}{\mathbf{k^{2}}\\!+\\!\mathbf{k^{2}_{1}}\\!-\\!\mathbf{k^{\prime
2}}\\!-\\!\mathbf{k^{2}_{2}}}\delta_{\mathbf{k}\\!+\\!\mathbf{k_{1}}\\!-\\!\mathbf{k^{\prime}}\\!-\\!\mathbf{k_{2}}}$
(133)
The integration area for $I^{\prime}_{1}$ is shown in Fig. 2, where the
distance between the points $O_{1}$ and $P$ is ${\bf R}_{O_{1}P}={\bf
k}/2k_{F}$, ${\bf R}_{PO_{2}}={\bf k}^{\prime}/2k_{F}$, ${\bf R}_{O_{1}N}={\bf
k}_{1}/2k_{F}$, and ${\bf R}_{ON}={\bf x}/2$. We thus have ${\bf R}_{NP}=({\bf
k}-{\bf k}_{1})/2k_{F}$ and ${\bf R}_{NP^{\prime}}=({\bf k}^{\prime}-{\bf
k}_{1})/2k_{F}$. In the region of integration we should have $|{\bf
R}_{O_{1}N}|=k_{1}/2k_{F}\leq 1/2$. This leads to
$\displaystyle I^{\prime}_{1}$ $\displaystyle=4\int\frac{|{\bf
R}_{NP}|\cdot|{\bf
R}_{NP^{\prime}}|}{m^{2}-x^{2}}d^{2}n=-\int_{0}^{2\pi}d\varphi\int_{0}^{l_{3}(\varphi)}xdx$
$\displaystyle\times\frac{\sqrt{[x^{2}+\sin^{2}(\theta/2)]^{2}-4x^{2}\sin^{2}(\theta/2)\cos^{2}\varphi}}{x^{2}-\sin^{2}(\theta/2)},$
(134)
where $\varphi$ is the angle between ${\bf x}$ and ${\bf m}$ (see Fig. 2), and
the quantity $l_{3}(\varphi)$ obeys the equation
$l_{3}^{2}(\varphi)-2\cos\frac{\theta}{2}\sin\varphi\cdot
l_{3}(\varphi)+\cos^{2}\frac{\theta}{2}=1.$
Turning to the variable $z=r^{2}-\sin^{2}(\theta/2)$ the integral
$I^{\prime}_{1}$ is reduced to
$\displaystyle
I^{\prime}_{1}=-\frac{1}{2}\int_{0}^{2\pi}d\varphi\int_{-\sin^{2}(\theta/2)}^{l_{3}^{2}(\varphi)-\sin^{2}(\theta/2)}\frac{\sqrt{R}}{z}dz,$
(135)
with
$R=z^{2}+4z\sin^{2}\frac{\theta}{2}\sin^{2}\varphi+4\sin^{4}\frac{\theta}{2}\sin^{2}\varphi.$
Figure 2: (color online). Left: The integration area for $I^{\prime}_{1}$ (in
blue): ${\bf R}_{O_{1}P}=\mathbf{k}/2k_{F}$, ${\bf
R}_{PO_{2}}=\mathbf{k^{\prime}}/2k_{F}$, ${\bf
R}_{O_{1}N}=\mathbf{k_{1}}/2k_{F}$, and $\varphi$ is the angle between the
vectors ${\bf R}_{OP}$ and ${\bf R}_{ON}$, which is the same as the angle
between ${\bf m}$ and ${\bf x}$. Right: The integration area for
$I^{\prime}_{2}$ (in red): ${\bf R}_{O_{1}P}=\mathbf{k}/2k_{F}$, ${\bf
R}_{O_{2}P}=\mathbf{k^{\prime}}/2k_{F}$, ${\bf
R}_{O_{2}N}=\mathbf{k_{1}}/2k_{F}$, $\alpha$ is the angle between ${\bf
R}_{PM}$ and ${\bf R}_{PN}$, and $\phi$ is the angle between ${\bf R}_{PN}$
and ${\bf R}_{OO_{2}}$.
It is easy to see that:
$\displaystyle
I_{r}=\int_{-\sin^{2}(\theta/2)}^{l_{3}^{2}(\varphi)-\sin^{2}(\theta/2)}\frac{\sqrt{R}}{z}dz$
$\displaystyle=\Big{\\{}\sqrt{R}-\sqrt{a}\ln\left(2a+bz+2\sqrt{aR}\right)$
$\displaystyle+\frac{b}{2}\ln\left(2\sqrt{R}+2z+b\right)\Big{\\}}\Big{|}^{l_{3}^{2}(\varphi)-\sin^{2}(\theta/2)}_{-\sin^{2}(\theta/2)}$
$\displaystyle+\sqrt{a}\cdot
P\int_{-\sin^{2}(\theta/2)}^{l_{3}^{2}(\varphi)-\sin^{2}(\theta/2)}\frac{dz}{z}=I_{r\uparrow}-I_{r\downarrow},$
where $a=4\sin^{4}(\theta/2)\sin^{2}\varphi$,
$b=4\sin^{2}(\theta/2)\sin^{2}\varphi$, and the symbol $P$ stands for the
principal value of the integral. The quantities $I_{r\uparrow}$ and
$I_{r\downarrow}$ denote the values of the integral at the upper and lower
bounds, respectively (in the last line we have to take the principal value of
the integral and, hence, if the upper bound of the integral is positive we
have to replace the lower bound with $\sin^{2}(\theta/2)$). Then
$I_{r\uparrow}$ and $I_{r\downarrow}$ are given by:
$\displaystyle I_{r\uparrow}=2|\sin\varphi|\cdot
l(\varphi)-2\sin^{2}\frac{\theta}{2}|\sin\varphi|\cdot\left[\ln\left(8\sin^{2}\frac{\theta}{2}\sin^{2}\varphi\right)+\ln\left(\sin^{2}\frac{\theta}{2}+\cos\frac{\theta}{2}\sin\varphi\cdot
l(\varphi)+l(\varphi)\right)\right]$
$\displaystyle+2\sin^{2}\frac{\theta}{2}|\sin\varphi|\cdot\ln|2\cos\frac{\theta}{2}\sin\varphi\cdot
l(\varphi)|+2\sin^{2}\frac{\theta}{2}\sin^{2}\varphi\left[\ln
4+\ln\left(|\sin\varphi|\cdot l(\varphi)+\cos\frac{\theta}{2}\sin\varphi\cdot
l(\varphi)+\sin^{2}\frac{\theta}{2}\sin^{2}\varphi\right)\right],$
$\displaystyle
I_{r\downarrow}=\sin^{2}\frac{\theta}{2}-2\sin^{2}\frac{\theta}{2}|\sin\varphi|\cdot\left[\ln\left(4\sin^{4}\frac{\theta}{2}\right)+\ln\left(\sin^{2}\varphi+|\sin\varphi|\right)\right]+2\sin^{2}\frac{\theta}{2}|\sin\varphi|\cdot\ln\left(\sin^{2}\frac{\theta}{2}\right)$
$\displaystyle+2\sin^{2}\frac{\theta}{2}\sin^{2}\varphi\cdot\ln\left(4\sin^{2}\frac{\theta}{2}\sin^{2}\varphi\right).$
The integral $I^{\prime}_{1}$ can be expressed as:
$I^{\prime}_{1}=-\frac{1}{2}\int_{0}^{2\pi}[I_{r\uparrow}-I_{r\downarrow}]d\varphi,$
and for performing the calculations we notice that $l_{3}(\varphi)\cdot
l_{3}(\varphi+\pi)=\sin^{2}\frac{\theta}{2}$,
$l_{3}(\varphi)-l_{3}(\varphi+\pi)=2\cos\frac{\theta}{2}\sin\varphi$, and
$l_{3}(\varphi)+l_{3}(\varphi+\pi)=2\sqrt{\cos^{2}\frac{\theta}{2}\sin^{2}\varphi+\sin^{2}\frac{\theta}{2}}$.
We then obtain:
$\displaystyle I^{\prime}_{1}=$
$\displaystyle-\sin^{2}\frac{\theta}{2}\left(\pi\ln
2+\pi/2-\pi\ln\sin\frac{\theta}{2}+4\ln|\cos\frac{\theta}{2}|-4\ln(1+\sin\frac{\theta}{2})+\mathcal{G}(\theta)-\frac{2\pi}{|\cos\frac{\theta}{2}|}-\frac{4\arcsin(\sin\frac{\theta}{2})}{|\cos\frac{\theta}{2}|}\right)$
$\displaystyle-\frac{k^{2}_{F}}{|\cos\frac{\theta}{2}|}\left(\pi-2\arcsin(\sin\frac{\theta}{2})+|\sin\theta|\right),$
(136)
with
$\mathcal{G}(\theta)=\int_{0}^{\pi}2\sin^{2}\varphi\ln\left(\sin\varphi+\sqrt{\sin^{2}\frac{\theta}{2}+\cos^{2}\frac{\theta}{2}\sin^{2}\varphi}\right)d\varphi.$
(137)
The integration area for $I^{\prime}_{2}$ is shown in Fig. 2, and we get:
$I^{\prime}_{2}=-4\int\frac{|\mathbf{m}|\cdot|{\bf
R}_{PN}|}{\mathbf{m}\cdot\mathbf{y}}d^{2}y=-8\int\frac{d^{2}\rho}{\cos\phi},$
(138)
where we denote ${\bf R}_{PN}=$ $\bm{\rho}$, and $\phi=\alpha-\theta/2$ is the
angle between the vectors $\mathbf{m}$ and $\bm{\rho}$, with $\alpha$ being
the angle between the vectors ${\bf R}_{PM}$ and ${\bf R}_{PN}$ (see Fig. 2).
We then have:
$\displaystyle I^{\prime}_{2}$
$\displaystyle=-8\int_{0}^{\pi}d\alpha\int_{0}^{\sin\alpha}\frac{\rho
d\rho}{\cos(\alpha-\theta/2)}$
$\displaystyle=-4\left[\cos^{2}\frac{\theta}{2}\ln\frac{1+\sin(\theta/2)}{1-\sin(\theta/2)}+2\sin\frac{\theta}{2}\right].$
(139)
Using $I^{\prime}_{1}$ (C) and $I^{\prime}_{2}$ (C) in Eq. (131) we get
equation (85) for the interaction function $\tilde{F}^{(2)}_{2}(\theta)$.
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* (37) K. Huang and C.N. Yang, Phys. Rev. 105, 767 (1957).
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* (41) L.D. Landau and E.M. Lifshitz, Quantum Mechanics, Non-Relativistic Theory (Butterworth-Heinemann, Oxford, 1999).
* (42) In fact, we arrived at this conclusion confining ourselves to the contributions up to $\beta^{2}$. In stead of considering higher order contributions we checked numerically that the solution with $s>1$ is absent if we omit the $(k_{F}r_{*})^{2}$-terms in the interaction function of quasiparticles.
* (43) M. J. H. Ku, A. T. Sommer, L. W. Cheuk, and M. W. Zwierlein, arXiv:1110.3309.
* (44) T. Yefsah, R. Desbuquois, L. Chomaz, K. J. Gunter, and J. Dalibard, Phys. Rev. Lett. 107, 130401 (2011).
* (45) D. Wang, B. Neyenhuis, M. H. G. de Miranda, K.-K. Ni, S. Ospelkaus, D. S. Jin, and J. Ye, Phys. Rev. A 81, 061404(R) (2010).
|
arxiv-papers
| 2011-11-30T10:34:54 |
2024-09-04T02:49:24.806793
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Zhen-Kai Lu and G. V. Shlyapnikov",
"submitter": "Zhenkai Lu Lu Zhenkai",
"url": "https://arxiv.org/abs/1111.7114"
}
|
1111.7188
|
¡html¿¡head¿ ¡meta http-equiv=”content-type” content=”text/html;
charset=ISO-8859-1”¿
¡title¿CERN-2010-002¡/title¿
¡/head¿
¡body¿
¡h1¿¡a href=”http://indico.cern.ch/conferenceDisplay.py?confId=97971”¿EuCARD-
AccNet-EuroLumi Workshop: The High-Energy Large Hadron Collider¡/a¿¡/h1¿
¡h2¿Villa Bighi, Malta, 14 - 16 Oct 2010¡/h2¿
¡h2¿Proceedings - CERN Yellow Report
¡a href=”http://cdsweb.cern.ch/record/1344820”¿CERN-2011-003¡/a¿¡/h2¿
¡h3¿editor: E. Todesco and F. Zimmermann ¡/h3¿ ¡p¿This report contains the
proceedings of the EuCARD-AccNet-EuroLumi Workshop on a High-Energy Large
Hadron Collider ‘HE-LHC10’ which was held on Malta from 14 to 16 October 2010.
This is the first workshop where the possibility of building a 33 TeV centre-
of-mass energy proton–proton accelerator in the LHC tunnel is discussed. The
key element of such a machine will be the 20 T magnets needed to bend the
particle beams: therefore much space was given to discussions about magnet
technologies for high fields. The workshop also discussed possible parameter
sets, issues related to beam dynamics and synchrotron radiation handling, and
the need for new injectors, possibly with 1 TeV energy. The workshop searched
for synergies with other projects and studies around the world facing similar
challenges or pushing related technologies, revisited past experience, and
explored a possible re-use of existing superconducting magnets. Last not
least, it reinforced the inter-laboratory collaborations within EuCARD,
especially between CERN and its European, US, and Japanese partners.¡/p¿
¡h2¿Lectures¡/h2¿
¡p¿ ¡!– Elements of a physics case for a High-Energy LHC –¿ Title: ¡a
href=”http://cdsweb.cern.ch/record/1403000”¿Elements of a physics case for a
High-Energy LHC¡/a¿ ¡br¿ Author: Wells, James D¡br¿ Journal-ref: CERN Yellow
Report CERN-2011-003, pp. 1-5¡br¿ ¡br¿
¡!–CERN Accelerator Strategy –¿ LIST:arXiv:1108.1115¡br¿
¡!– HE-LHC beam-parameters, optics and beam-dynamics issues –¿
LIST:arXiv:1108.1617¡br¿
¡!– Conceptual design of 20 T dipoles for high-energy LHC –¿
LIST:arXiv:1108.1619¡br¿ ¡br¿
¡!–What can the SSC and the VLHC studies tell us for the HE-LHC? –¿ Title: ¡a
href=”http://cdsweb.cern.ch/record/1403041”¿What can the SSC and the VLHC
studies tell us for the HE-LHC?¡/a¿ ¡br¿ Author: Wienands, U¡br¿ Journal-ref:
CERN Yellow Report CERN-2011-003, pp. 20-26¡br¿ ¡br¿
¡!–A high energy LHC machine: experiments ’first’ impressions –¿
LIST:arXiv:1108.1621¡br¿
¡!– Progress in high field accelerator magnet development by the US LHC
Accelerator Research Program–¿ LIST:arXiv:1108.1625¡br¿
¡!– LBNL high field core program –¿ LIST:arXiv:1108.1868¡br¿
¡!–KEK effort for high field magnets –¿ LIST:arXiv:1108.1626¡br¿
¡!– EUCARD magnet development –¿ LIST:arXiv:1108.1627¡br¿
¡!– Status of Nb3Sn accelerator magnet R&D at Fermilab –¿
LIST:arXiv:1108.1869¡br¿
¡!– Status of high temperature superconductor based magnets and the conductors
they depend upon –¿ LIST:arXiv:1108.1634¡br¿
¡!– 20 T dipoles and Bi-2212: the path to LHC energy upgrade –¿
LIST:arXiv:1108.1640¡br¿
¡!– Heat loads and cryogenics for HE-LHC –¿ LIST:arXiv:1108.1870¡br¿
¡!– HE-LHC: requirements from beam vacuum –¿ LIST:arXiv:1108.1642¡br¿
¡!– Beam screen issues –¿ LIST:arXiv:1108.1643¡br¿
¡!– Synchrotron radiation damping, intrabeam scattering and beam-beam
simulations for HE-LHC –¿ LIST:arXiv:1108.1644¡br¿
¡!– Beam-beam studies for the High-Energy LHC –¿ LIST:arXiv:1108.1871¡br¿
¡!– Preliminary considerations about the injectors of the HE-LHC –¿
LIST:arXiv:1108.1652¡br¿
¡!– Using tevatron magnets for HE-LHC or new ring in LHC tunnel –¿
LIST:arXiv:1108.1653¡br¿
¡!– Magnet design issues and concepts for the new injector –¿
LIST:arXiv:1108.1654¡br¿
¡!– Using LHC as injector and possible uses of HERA magnets/coils –¿
LIST:arXiv:1108.1655¡br¿
¡!– Intensity issues and machine protection of the HE-LHC –¿
LIST:arXiv:1108.1663¡br¿
¡!– Injection and dump considerations for a 16.5 TeV HE-LHC –¿
LIST:arXiv:1108.1664¡br¿
¡!– Radiation protection issues after 20 years of LHC operation –¿
LIST:arXiv:1108.1669¡br¿
¡!– 20 T dipoles and Bi-2212: the path to LHC energy upgrade –¿ ¡br¿
Title: ¡a href=”http://cdsweb.cern.ch/record/1373679”¿Summary of session 1:
introduction and overview¡/a¿ ¡br¿ Author: J.P. Koutchouk and R. Bailey ¡br¿
Journal-ref: CERN Yellow Report CERN-2011-003, pp. 137-139 ¡br¿ ¡br¿ Title: ¡a
href=”http://cdsweb.cern.ch/record/1373681”¿Summary of session 2: magnets for
the HE-LHC¡/a¿ ¡br¿ Author: L. Rossi and E. Todesco ¡br¿ Journal-ref: CERN
Yellow Report CERN-2011-003, pp. 140-142 ¡br¿ ¡br¿ LIST:arXiv:1202.3811¡br¿
¡br¿ Title: ¡a href=”http://cdsweb.cern.ch/record/1373683”¿Summary of session
4: HE-LHC injectors and infrastructure¡/a¿ ¡br¿ Author: E. Prebys and L.
Bottura ¡br¿ Journal-ref: CERN Yellow Report CERN-2011-003, pp. 145-148 ¡br¿
¡/p¿ ¡/body¿¡/html¿
|
arxiv-papers
| 2011-11-30T14:44:06 |
2024-09-04T02:49:24.822151
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "E. Todesco and F. Zimmermann",
"submitter": "Scientific Information Service CERN",
"url": "https://arxiv.org/abs/1111.7188"
}
|
1111.7211
|
# Disconnecting Open Solar Magnetic Flux
C.E. DeForest∗, T.A. Howard∗, and D.J. McComas∗† ∗Southwest Research
Institute, 1050 Walnut Street Suite 300, Boulder CO 80302 †University of
Texas at San Antonio, San Antonio TX 78249
Disconnection of open magnetic flux by reconnection is required to balance the
injection of open flux by CMEs and other eruptive events. Making use of recent
advances in heliospheric background subtraction, we have imaged many abrupt
disconnection events. These events produce dense plasma clouds whose
distinctive shape can now be traced from the corona across the inner solar
system via heliospheric imaging. The morphology of each initial event is
characteristic of magnetic reconnection across a current sheet, and the newly-
disconnected flux takes the form of a “U”-shaped loop that moves outward,
accreting coronal and solar wind material.
We analyzed one such event on 2008 December 18 as it formed and accelerated at
20 m s-2 to 320 km s-1, expanding self-similarly until it exited our field of
view 1.2 AU from the Sun. From acceleration and photometric mass estimates we
derive the coronal magnetic field strength to be $8\,\mu T$, 6$R_{\odot}$
above the photosphere, and the entrained flux to be $1.6\times 10^{11}Wb$
($1.6\times 10^{19}Mx$). We model the feature’s propagation by balancing
inferred magnetic tension force against accretion drag. This model is
consistent with the feature’s behavior and accepted solar wind parameters.
By counting events over a 36 day window, we estimate a global event rate of
$1\,d^{-1}$ and a global solar minimum unsigned flux disconnection rate of
$6\times 10^{13}Wb\,y^{-1}$ ($6\times 10^{21}Mx\,y^{-1}$) by this mechanism.
That rate corresponds to $\sim-0.2\,nT\,y^{-1}$ change in the radial
heliospheric field at 1 AU, indicating that the mechanism is important to the
heliospheric flux balance.
## 1 Introduction
The solar dynamo generates magnetic flux inside the Sun, whch is transported
outward and emerges through the Sun’s surface into the corona. Magnetic loops
build up “closed” magnetic flux (connected to the Sun at both ends) in the
corona. Some of these closed loops subsequently “open” into interplanetary
space – that is, they are connected to the Sun at only one end with the other
extending to great distances in the heliosphere or beyond. Owing to the very
high electrical conductivity, open magnetic flux is frozen into the solar wind
and carried out with it. The magnetized solar wind expands continuously
outward from the Sun in all directions, filling and inflating our heliosphere
and protecting the inner solar system from the vast majority of galactic
cosmic rays. The balance between the opening and closing of magnetic flux from
the Sun is thus critical and fundamental both to the solar wind and to the
radiation environment of our solar system.
Magnetic flux opens when coronal mass ejections (CMEs) erupt through the
corona, carrying previously closed magnetic loops beyond the critical point
where the solar wind exceeds the Alfvén speed (typically <$20\,R_{\odot}$) and
can no longer return to the Sun. CMEs were first studied in OSO-7 and Skylab
observations of the corona (e.g.Tousey 1973; Gosling et al. 1974; Hundhausen
1993), and since then continued work has provided an increasingly detailed
picture of these transient magnetic structures both during their formation and
ejection, and as they continue to evolve and interact with the solar wind.
Long lasting, radial "legs" are often observed along the flanks of a CME and
persisting behind it. These legs are generally interpreted as evidence for at
least some continued magnetic connection of CMEs back to the Sun and hence the
opening of new magnetic flux with CME ejections. That picture is further
supported by observatation, _in situ_ , of beamed suprathermal halo electrons
streaming in both directions along the local interplanetary magnetic field
(IMF) during the passage of an interplanetary CME (ICME) cloud (e.g. Gosling
1990, 1993 & references therein), which are commonly interpreted as signatures
of direct connection of the ICME magnetic field to the solar corona in both
directions, and hence of newly opening magnetic flux. However, less is known
about the equally necessary process of disconnection that must be present to
remove newly opened flux and prevent the IMF from growing without limit.
Because of the continual opening of magnetic flux through CMEs, McComas and
coworkers in the early 1990s pursued a series of studies to determine how
magnetic flux could be closed back off and avoid a so-called magnetic flux
“catastrophe” of ever increasing magnetic field strength in the interplanetary
magnetic field (McComas 1995& references therein). The amount of open magnetic
flux in interplanetary space can be approximated with the "total flux
integral" which removes the effects of variations in the solar wind speed in
determining the amount of magnetic flux crossing 1 AU (McComas et al. 1992a).
Using this integral, McComas et al. (1992a; 1995) showed that if all
counterstreaming electron events represent simply connected opening magnetic
loops, then for solar maximum CME rates, the amount of flux crossing 1 AU
would double over only ~9 months. For flux rope CMEs, significantly more
magnetic flux may be observed in the loops crossing 1 AU than what remains
attached to the Sun along the CMEs’ legs; however, it must be stressed that if
CMEs retain any solar attachment whatsoever, the flux catastrophe will
ultimately occur in the absence of some other process to close off previously
open fields.
Of course a magnetic flux catastrophe is not observed in the solar wind and,
in fact, the overall magnitude of the IMF and amount of open flux seems to
vary over the solar cycle. For cycle 21, the average magnitude varied by ~50%
(Slavin et al. 1986) while the total flux integral varied by ~60% (McComas et
al. 1992a; 1992b), with maxima shortly after solar maximum and minima shortly
after solar minimum McComas 1994. Since these studies, the solar wind has gone
through a prolonged (multi-cycle) reduction in both solar wind power (the
dynamic pressure of the solar wind that ultimately inflates the heliosphere)
(McComas et al. 2008) and magnetic field magnitude (Smith & Balogh 2008). The
lack of a flux catastrophe, solar cycle variation and now long-term reduction
in the open magnetic flux from the Sun all show that there must be some
process for closing off previously open field regions and returning magnetic
flux to the Sun.
Magnetic reconnection plays an important role in regulating the topology of
solar magnetic flux, however, once the top of a loop passes the critical
point, its magnetic flux remains open until some other process occurs to close
it off below the critical point. That is, reconnection above the critical
point can only rearrange the topology of open magnetic flux in the heliosphere
– only reconnection between two oppositely directed (inward and outward field)
regions of open flux close to the Sun can close off previously open magnetic
flux. The most obvious method of reducing the amount of magnetic flux open to
interplanetary space is via reconnection between oppositely directed,
previously open field lines (McComas et al. 1989), which creates closed field
loops that can return to the Sun and the release of disconnected U-shaped
field structures into interplanetary space. An example of such a coronal
disconnection event was shown by McComas et al. 1991 using SMM coronagraph
images from 1 June 1989\. An even older example of a likely coronal
disconnection event can be found as far back as the 16 April 1893 solar
eclipse (e.g., Cliver 1989), where sketches (data in 1893) made in time
ordered sequence from Chile, Brazil, and Senegal, indicate the outward motion
of a large U-shaped structure (McComas 1994).
For the opening and closing of the solar magnetic flux to maintain some sort
of equilibrium, there must be some type of feedback between these two
processes. McComas et al. (1989; 1991) suggested that this feedback occurs
through transverse magnetic pressure in the corona, where the expansion of
newly opened field regions must enhance transverse pressure and compress
already open flux elsewhere around the Sun. When enough pressure builds up,
reconnection between oppositely direct open flux would reduce the pressure and
amount of open flux. The sequence of images from the 27 June 1988 coronal
disconnection event, in fact showed just such a compression, indicated by the
deflection of the streamers in the corona, just prior to and appearing to
precipitate the coronal disconnection event. Another line of supporting
evidence was provided by numerical simulations (Linker et al. 1992), which
indicated that increased magnetic pressure could lead to reconnection across a
helmet streamer and the release of disconnected flux. Schwadron et al. (2010)
recently reexamined the flux balance issue in light of the anomalously long
solar minimum between cycles 23 and 24 and modeled the level of magnetic flux
in the inner heliosphere as a balance of that flux injected by CMEs, lost
through disconnection, and closed flux lost through interchange reconnection
near the Sun.
McComas et al. (1992c)conducted a statistical study of three months of SMM
coronagraph observations (Hundhausen (1993)) to assess the frequency of
coronal disconnection events. These authors found that while the initial
survey (St. Cyr & Burkepile 1990) found no obvious disconnections, six of the
53 transient events during this interval (11%) showed some evidence of
disconnection in more than one frame and 13 (23%) showed a single frame with
an outward "U" or "V" structure. Given the imaging and analysis technology of
the day, McComas et al. (1992c) concluded that magnetic disconnection events
on previously open field lines may be far more common than previously
appreciated. With today’s imaging and exceptional analysis capabilities, the
question of coronal disconnection events should finally be resolvable.
For this study we used image sequences, collected by the SECCHI (Howard et al.
2008) instrument suite on board NASA’s _STEREO-A_ spacecraft, of Thomson-
scattered sunlight from free electrons in the interplanetary plasma. The
observations span from the deep solar corona to beyond 1 AU at elongation
angles of up to $70^{\circ}$ from the solar disk; this continuous observation
is enabled by recently developed background subtraction techniques (DeForest
et al. 2011) operating on the _STEREO_ data. The signature U-shaped loops of
disconnected plasma are far clearer in the processed heliospheric data far
from the Sun than in the coronagraph data close to the Sun, and we detect 12
characteristic departing “V” or “U” events in 36 days – far more than the
expected number based on scaling the results of McComas et al. (1992c). For
this initial report, we focus on quantitative analysis of a single event. In
Sections 2.1-2.7, we describe the observations and calculate the geometry, the
mass evolution, and (by assuming the U-loop is accelerated by the tension
force) the coronal magnetic field sand entrained flux in the disconnecting
structure. As a plausibility check, we explore the tension force scenario and
its consequences for the long-term evolution of the feature, and find that the
scenario is consistent with accepted values for the solar wind density and
speed. In Section 3, we discuss broader consequences of the observation,
including estimating the disconnection rate based on the number of similar
events in our data set, and discuss implications for the global magnetic flux
balance.
## 2 Observations
The SECCHI suite on STEREO was intended to be used as a single integrated
imaging instrument (e.g. Howard et al. 2008). It consists of an EUV imager
(EUVI) observing the disk of the Sun, and four visible light imagers (COR-1,
COR-2, HI-1, and HI-2) with progressively wider overlapping fields of view, to
cover the entire range of angles between the solar disk and the Earth. The
visible light imagers view sunlight that has been Thomson scattered off of
free electrons in the corona and interplanetary space; the theory of Thomson
scattering observations has been recently reviewed by Howard & Tappin (2009a).
We set out to view coronal and heliospheric events in the weeks around 2008
December, using newly developed background subtraction techniques to observe
solar wind features in the HI-1 and HI-2 fields of view (DeForest et al.
(2011); Howard & DeForest (2011)). In the initial 36 day data set we prepared,
we observed 12 disconnection events identified by a clear “V” or “U” shaped
bright structure propagating outward in the heliosphere. We chose a
particuarly clearly presented one, which was easily traceable to its origin in
the low corona on 2008-Dec-12 at 04:00, for further detailed study.
### 2.1 Image Preparation
Data preparation followed standard and published techniques. For COR-1 and
COR-2, we downloaded Level 1 (photometrically calibrated) data, and further
processed them by fixed background subtraction: we acquired images for an
11-day period, and found the 10 percentile value of each pixel across the
entire 11-day dataset. This image was median filtered over a 5x5 pixel window
to generate a background image that included the F corona, any instrumental
stray light, and the smooth, steady portion of the K corona. Subtracting this
background from each image yielded familiar coronal images of excess feature
brightness compared to the smooth, steady background. We performed one
additional step: motion filtration to suppress stationary image components.
This step matches the motion filtration step used for HI-1 and HI-2 (below),
and suppresses the stationary streamer belt while not greatly affecting the
moving features under study.
The heliospheric imagers required further processing to remove the starfield,
which is quite bright compared to the faint Thomson scattering signal far from
the Sun in the image plane. We processed the STEREO-A HI-2 data as described
by DeForest et al. (2011). The HI-1 data used a similar process adapted to the
higher background gradients in that field of view and described by Howard &
DeForest (2011).
All the imagers yielded calibrated brightness data in physical units of the
mean solar surface brightness ( $B_{\odot}=2.3\times
10^{7}W\,m^{-2}\,SR^{-1}$). Because of the wide field of view, as a subsequent
processing step we distorted the images into azimuthal coordinates, in which
one coordinate is azimuth (solar position angle) in the image plane and the
other is either elongation angle $\epsilon$ (“radius” on the celestial sphere)
or its logarithm. The latter projection, if scaled properly, is conformal: it
preserves the shape of features that are small compared to their distance from
the Sun. To equalize brightness, we applied radial filters to the images for
presentation, with either a $\epsilon^{3.5}$ scaling (for coronal images) or a
$\epsilon^{3}$ scaling (for heliospheric images).
Figure 1 shows snapshots of the disconnected plasma and associated cusp, as
observed by four separate instruments over the course of four days as it
propagated outward. Shortly after 2008 Dec 18 04:00, the streamer belt at
$160^{\circ}$ ecliptic azimuth ($20^{\circ}$ CCW of the Sun-Earth line)
pinched and separated, forming a “U” loop that retracted outward, with a
trailing cusp, over the course of the following three days. The feature
remained visible in Thomson scattered light because of plasma scooped up
during the early acceleration period in the lower corona: this plasma remained
denser than the surrounding medium, yielding a bright feature throughout the
data set. The disconnected plasma completely missed the ecliptic plane and was
therefore not observed _in situ_ by any of the near-Earth or STEREO probes.
Figure 1: The disconnection event of 2008 Dec 18 in context \- 8 still images
showing formation and evolution of the U-loop: left to right, top to bottom.
### 2.2 Observing Geometry and 3-D structure
The observing geometry for the 2008 Dec 18 event is shown in Figure 2, from an
overhead (northward out-of-ecliptic) point of view co-rotating with the
STEREO-A orbit. The event departure angle was measured both using direct
triangulation between the coronagraphs in STEREO-A and STEREO-B. We used the
triangulation method described by Howard & Tappin (2008). Although the
disconnection event is small compared to most CMEs, subtending just a few
degrees in latitude, it is still large enough to cast doubt on the simple
triangulation results, so we also used “TH model” semi empirical transient
event reconstruction tool (Tappin & Howard 2009) to extract the departure
angle. TH was developed to reconstruct CME leading edge (“sheath”) overall
envelope and propagation speed, but is also applicable to smaller transient
events such as this one. Details, applications, and limitations of the TH
model are further described by Tappin & Howard 2009 and by Howard & Tappin
(2009b; 2010). Departure longitude was measured to be
$-10^{\circ}\text{\textpm}5^{\circ}$ in heliographic coordinates, with an
estimated event width of under $5^{\circ}.$
We took the disconnected feature’s trajectory to have constant radial motion
(the “Fixed-$\Phi$ aproximation”) in the solar inertial frame – this leads to
the slightly curved aspect to the trajectory in the co-rotating heliographic
ecliptic frame, which maintains the prime meridian at the Earth-Sun line.
Figure 2 shows an out-of-ecliptic projected view of the observing geometry,
including construction angles and distances used in Section 2.4 for trajectory
calculations.
Figure 2: Observing geometry in the ecliptic plane on 2008 Dec 18 - 2008 Dec
22
### 2.3 Feature Evolution
To analyze the feature’s evolution across a two-order-of-magnitude shift in
scale over its observed lifetime, we transformed the processed _STEREO-A_
source images into local heliographic radial coordinates – i.e. zero azimuth
is due solar West from the viewpoint of _STEREO-A_ , with azimuthal coordinate
increasing clockwise around the image plane; this follows early work by
DeForest et al. (2001) in imaging polar plumes. Distances from Sun center are
recorded as elongation angle $\epsilon$ from the center of the Sun, as a
reminder of the angular nature of the wide-field observations. To avoid
aliasing in the resampling process, we resampled the images using the
optimized resampling package described by DeForest (2004). Figures 3 and 4
show the liftoff and propagation of the feature across 65 degrees of
elongation from its origin in the solar streamer belt. Both figures have a
radial gain filter applied to equalize the feature’s brightness, which varies
by over seven orders of magnitude: from $1.5\times 10^{-9}B_{\odot}$ in the
low streamer belt at 2008 Dec 18 04:30 to $5.7\times 10^{-17}B_{\odot}$ six
days later, at $\epsilon=65^{\circ}$.
The bright feature takes the classic wishbone shape of reconnecting field
lines emerging from a current sheet (e.g. Chapter 4 of Priest et al. 2000).
The aspect ratio of the wishbone may be estimated by dividing the vertical
height from cusp to the top of the visible horns, by the width between the
horns. This aspect ratio varies from ~10:1 when the horns are first clearly
resolved near 2008 Dec 18 08:00, to approximately 2:1 some four hours later
and 1:1 by 2008 Dec 19 04:00 – one full day after the first pinch is observed
in the streamer belt. After 2008 Dec 19, the feature expands approximately
self-similarly as it propagates, subtending approximately $16^{\circ}$ of
azimuth and not changing its aspect ratio throughout the rest of its
trajectory.
Note that aspect ratio is _not_ preserved by the linear azimuthal mapping used
in Figure 3, which was selected to show the early acceleration clearly; aspect
ratio is preserved by the (conformal) logarithmic mapping used in Figure 4,
which shows nearly self-similar expansion in the image plane despite
perspective effects that come into play above about $\epsilon=30^{\circ}$
The scaling of brightness is reassuring because, in a uniformly propagating
wind with no acceleration, density must decrease as $r^{-2}$ and feature
column density must thus decrease as $r^{-1}$, while illumination decreases as
$r^{-2}$, so feature brightness is expected to decrease as $r^{-3}$. The fact
that brightness levels do not change much across Figure 4, which is scaled by
$\epsilon^{3}$, suggests that the disconnected flux and material entrained in
it are indeed propagating approximately uniformly. The fact that they _do_
change slightly, with brighter images to the right, indicates that the feature
is gaining intrinsic brightnsss by accumulating material as it propagates.
The horizontal positions and error bars in Figures 3 and 4 are the results of
manual feature location of the cusp, with a point-and-click interface. The
white error bars are based on the sharpness of the feature. In the excess
brightness plot, the feature is easy to see but blurs near the top of Figure 3
due to the higher levels of both photon noise and motion blur as the feature
accelerates to the top of the coronagraph field of view. The running
difference plot highlights fine scale feature and helps identify the cusp
location near the top of the COR-2 field of view.
Figure 3: Formation and early acceleration of the 2008 Dec 18 disconnection
event through the STEREO-A COR-1 and COR-2 fields of view The trailing edge of
the event is marked, with error bars based on feature identification. TOP:
direct excess-brightness images show feature formation and overall structure.
BOTTOM: running-difference images show detail. These stack plots include a
small image of the feature at each sampled time to show evolution. Intensities
are scaled with $\epsilon^{3.5}$ to equalize brightness vs. height. The
individual images have been resampled into linear azimuthal (radial)
coordinates, and the horizontal range is $160^{\circ}-174^{\circ}$ of azimuth.
Note that this projection does not preserve aspect ratio: despite appearances,
the event widens as it rises. Figure 4: Propagation and evolution of the 2008
Dec 18 disconnection event through the STEREO-A HI-1 and HI-2 fields of view.
The trailing edge of the event is marked, with error bars based on feature
identification. These stack plots include a small image of the feature at each
sampled time to show evolution. The individual images have been resampled into
logarithmic azimuthal (radial) coordinates, and the horizontal range is
$162^{\circ}-178\text{\textdegree}$of azimuth. This projection is conformal,
so the shape of the feature is preserved in each image. Intensities are scaled
with $\epsilon^{3}$ to equalize brightness vs. height. Note self-similar
expansion: the angular width and shape of the feature are preserved.
### 2.4 Acceleration profile
Converting angular observed coordinates to examine the inertial behavior of
the plasma requires triangulation using the Law of Sines. Using the “fixed
$\Phi"$ approximation (assuming the feature’s cusp is small and that it
propagates in a radial line from the Sun), the feature’s radius from the Sun
is easily calculated:
$r_{ev}=r_{A}\frac{{sin\left(\epsilon^{\prime}\right)}}{sin\left(\epsilon\right)},$
(1)
where the variables take the meanings in Figure 2: $\epsilon$ is the solar
elongation of the feature as seen from _STEREO-A_ , $\epsilon^{\prime}$ is the
solar elongation of _STEREO-A_ as seen from the feature), and the _STEREO-A_
solar distance $r_{A}$ is found by spacecraft tracking and is supplied by the
mission. Although no camera was present at the event itself,
$\epsilon^{\prime}$ is calculated by noting that
$\epsilon^{\prime}=180^{\circ}-\epsilon-\left(L-L_{ev}\right)$, as the
feature, _STEREO-A,_ and the Sun form a triangle. Figure 5 shows the results
of the tracking from Figures 3 and 4, propagated through Equation 1.
As expected, the event rapidly accelerates during the early phase, reaching a
peak acceleration of $20\,m\,s^{-2}$ as the aspect ratio changes in the
initial hours. The acceleration peaks 4-5 hours after the initial pinch in the
streamer belt, or 2-3 hours after the first observation of a well-formed cusp.
The feature reaches its final speed of $\sim 320\pm 15\,km\,s^{-1}$ within
just 8 hours of the initial pinch at 04:00 and within 6 hours of the first
observation of the well formed cusp at 06:00, and undergoes no further
significant acceleration nor deceleration during its obsered passage to beyond
1 AU over the next five days.
Figure 5: Inferred position, speed, and acceleration of the disconnected
plasma from the 2008-Dec-18 event, during onset (LEFT) and over the full
observation period (RIGHT). Error bars are derived by propagating _a priori_
location error and geometric error in the longitude of the event. The shaded
region indicates the full time range of the left-side plots.
### 2.5 Mass profile
We extracted photometric densities using the feature brightness in suitable
frames. The feature brightness is determined from the density via the Thomson
scattering equation (see, e.g., Howard & Tappin 2009a for a clear exposition).
Compact features can be treated as nearly point sources, and the line-of-sight
integral for the optically thin medium reduces to:
$B=B_{\odot}\Omega_{\odot}(r)\sigma_{e}\left(1+cos^{2}\chi\right)\rho\mu_{av}^{-1}d$
(2)
where $B$ is the measured feature brightness (in units of emissivity:
$Wm^{-2}SR^{-1}$), $B_{\odot}$ is (still) the solar surface brightness;
$\Omega_{\odot}(r)$ is the solid angle subtended by the Sun at the point of
scatter, well approximated by $\pi r_{\odot}^{2}r_{ev}^{-2}$ everywhere above
about 4 $r_{\odot}$; $\sigma_{e}$ is the differential Thomson scattering cross
section, given by half of the square of the classical electron radius
$r_{e}^{2}/2=4.0\times 10^{-30}m^{2}$; $\chi$ is the scattering angle (equal
to $\epsilon^{\prime}$ in Figure 2); $\rho$ is the mass density; $\mu_{av}$ is
the average mass per electron in the coronal plasma; and $d$ is the depth of
the feature.
$mu_{av}$ may be calculated from the spectroscopically measured 5% He/H number
ratio in the corona (Laming & Feldman 2000) and the assumption that the helium
is fully ionized (yielding two electrons per ion). This yields
$mu_{av}=1.1m_{p}=1.84\times 10^{-27}kg$.
Solving for the line-of-sight integrated mass surface density $\rho d$ gives:
$\rho
d=\mu_{av}\frac{B}{B_{\odot}}\Omega_{\odot}^{-1}(r)\sigma_{e}^{-1}\left(1+cos^{2}\epsilon^{\prime}\right)^{-1}$
(3)
and therefore
$m_{ev}=\left(\rho
d\right)wh=\mu_{av}\frac{B}{B_{\odot}}\Omega_{\odot}^{-1}(r)\sigma_{e}^{-1}\left(1+cos^{2}\epsilon^{\prime}\right)^{-1}\Omega_{ev}^{\left(1+cos^{2}\epsilon^{\prime}\right)}S^{2}$
(4)
where $w$ and $h$ are the dimensions shown in Figure 7; $\Omega_{ev}$ is the
solid angle subtended by the feature in the images; and $S$ is the calculated
spacecraft-feature distance, calculated by the law of sines as for $r_{ev}.$
To extract the mass profile from the data, we generated an image sequence
containing the feature, and marked the locus of the feature visually using a
pixel paint program. Using the generated masks, we summed masked pixels in the
feature for each photometric image, thereby integrating the feature brightness
over the solid angle represented by the corresponding pixels, to obtain an
intensity and an average brightness within the feature. To account for errors
in visual masking, we assigned error bars based on one-pixel dilation and one-
pixel contraction of the masked locus. We omitted frames with excessive noise,
encroachment of an image boundary, or a star or cosmic ray in or near the
feature. The results of the calculation are given in Figure 6, which shows
steady accretion of material through most of the journey through the
heliosphere.
Because our photometric analysis is based on subtraction of a calculated
background derived from the data set itself, we measure only excess feature
brightness (not absolute brightness) from Thomson scattering; thus our
brightness measurements and mass estimates are biased low, because we cannot
measure the absolute density of the background. The initial derived mass of
20-25 Tg translates to an electron number density of $2\times 10^{7}\,cm^{-3}$
in the lower corona, which is comparable to the density in bright coronal
features – so the total mass may be up to a factor of order two higher.
The final “feature excess” mass is $8\pm 2\times 10^{10}kg$, and the final
subtended solid angle is 0.028 SR, for a presented cross-section of $4.3\pm
0.1\times 10^{20}m^{2}$. Taking the depth to be the square root of the
observed cross section yields an estimated volume at 1 AU of $8.9\pm 0.3\times
10^{31}m^{3}$, for a total estimated excess electron density of $5\pm
1.4\,cm^{-3}$ at 1 AU, which is in good agreement with slow solar wind
densities ($3-10\,cm^{-3}$ when scaled to 1 AU) that were observed by
_Ulysses_ in situ in the same heliographic latitude range (e.g. McComas et al.
2000). Approximately 2/3 of this excess density appears to have been
accumulated enroute from the surroundign solar wind; this is further described
in Section 2.7.
Figure 6: Photometrically determined excess mass profile of the retracting
disconnected feature of 2008 Dec 18. Error bars are based on identification of
the feature boundary in the images. The trendline is extracted from regression
of the HI-1 and HI-2 data. The mass shown is excess mass in the feature
compared to the background solar wind (see text).
### 2.6 Entrained magnetic flux
Figure 7: Cartoon of the initial acceleration process of a disconnection
event. Tension force along newly released field lines is balanced by mass
entrained on the field lines. By measuring the acceleration and mass we infer
the amount of magnetic flux that was disconnected.
From the mass of the feature, and its acceleration, it is possible to extract
the entrained magnetic field by measuring the rate of change of momentum and
inferring a magnetic tension force via $f=ma$. The system is sketched in
Figure 7. The magnetic tension force is conserved along the open field lines,
so we can calculate it at any convenient cut plane including the one shown.
Tension force is frequently referred to as a “curvature force” and calculated
locally; here we integrate around the “U”, and notice that the integrated
force is just the unbalanced tension on the field lines contained in the “U”
shape. It is therefore given by
$f=m_{ev}a_{ev}=f_{B}=\frac{B^{2}A}{2\mu_{0}}=\frac{{\Phi^{2}}}{2\mu_{0}A}$
(5)
where, here, $B$ is the magnetic field strength (not brightness, as before).
Solving for $\Phi$,
$\Phi=\sqrt{2\mu_{0}dwm_{ev}a_{ev}}$ (6)
taking $m_{ev}$ to be $25Tg$ ($2.5\times 10^{10}kg$) during the peak of the
acceleration, and taking $w=d=0.2R_{\odot}$ (based on the measured width of
the feature’s fork during maximum acceleration, at $6R_{s}$ from the surface
($7R_{s}$ from Sun center) gives $\Phi=1.6\times 10^{11}Wb$ ($1.6\times
10^{19}Mx$ ), corresponding to an average field strength of 8$\mu T$ ( 0.08
Gauss) at that altitude, or an equivalent $r^{2}$-scaled field of 400$\mu T$
(4 Gauss) at the surface; this is comparable to accepted values of the open
flux density at the solar surface at solar minimum. Because of the way $m$ was
calculated (section 2.5 above) this figure is probably low by a factor of
order $\sqrt{2}$.
### 2.7 Accretion and force balance
As the disconnectioned structure travels outward, it accretes new material.
This effect is dramatic: as seen in Figure 6, the mass increases by a factor
of 3 from the corona to 1 AU. We conjecture that the material is accreted by
“snowplow” effects from the plasma ahead of the disconnected cusp as it
propagates. For the observed mass growth in the feature, new material must be
compressed to become visible in our Thomson scattering images, and the most
plausible way for it to be compressed is via ram effects. This scenario also
neatly explains the constant speed of the feature, by balancing the continued
tension force from the cusp with accretion momentum transfer. Here we explore
the concept of force balance between accretion and the tension force, to
identify whether some other model is required in addition to this simple one.
Extending Newton’s law to include momentum transfer by accretion, and
neglecting all but the tension force,
$\frac{\Phi^{2}}{2\mu_{0}A_{\Phi}}=m_{ev}a_{ev}+\frac{dm_{ev}}{dt}\Delta v,$
(7)
where the LHS is just the tension force from Equation 5, with the modification
that the cross section of the exiting field lines is written $A_{\Phi}$;
$a_{ev}=0$ after the initial acceleration; and the second term represents
momentum transfer into accreted material, with $\Delta v$ being the difference
between the feature speed and surrounding wind speed. The feature is thus in
equilibrium between accretion drag and continued acceleration by the tension
force. This accretion drag is important to the observed increase in feature
mass, because ram pressure against the surrounding wind material is what
compresses incoming material and renders it visible in the data against the
subtracted background.
Applying conservation of mass, we can relate $\Delta v$ and the average
density of the background solar wind through which the feature is propagating:
$\rho_{sw}=\frac{dm_{ev}/dt}{A_{ev}\Delta v}.$ (8)
where $A_{ev}$ is the geometrical area presented by the feature to the slow
wind ahead of it. Solving Equations 7 and 8 to eliminate $\Delta v$ gives
$\rho_{sw}=\left(\frac{dm_{ev}}{dt}\right)^{2}\left(\frac{2\mu_{0}}{\Phi^{2}}\right)\left(\frac{A_{\Phi}}{A_{ev}}\right),$
(9)
which gives the background solar wind density in terms of the accumulation
rate of mass in the observed feature, assuming constant outflow for both the
wind and the feature, and acceleration by the tension force. Given the
conservation of mass and the approximately constant speed of the solar wind,
$\rho_{sw}$ falls as approximately $r^{-2}$. Further, we observe nearly self-
similar expansion throughout most of the heliospheric range, so
$A_{\Phi}/A_{ev}$ is constant in that part of the trajectory – hence
$dm_{ev}/dt$ must also fall as $r^{-1}$ during the approximately constant
speed portion of the feature’s lifetime. Using this functional form, we can
extract an analytic expresson for the feature mass versus radius. We introduce
the $r^{-1}$ dependence by switching from the linear regression used in Figure
6, to a semi-log regression that assumes $dm_{ev}/d(log_{e}(r)$ to be
constant. Figure 8 shows such a regression, with the result that
$dm_{ev}/dr=34\pm 3\times 10^{9}\left(R_{\odot}/r\right)kgR_{\odot}^{-1}$.
Including the measured outflow speed of $315\pm 15\,km\,s^{-1}$, we find that
$dm_{ev}/dt=\left(1AU/r\right)\left(7.1\pm 1\times 10^{4}kg\,s^{-1}\right)$.
Including all of these values into Equation 9, together with the average
particle mass from Section 2.5, yields a background wind numeric density of
$n_{sw}\left(1\,AU\right)$ of $30\pm 6\,cm^{-3}\left(A_{\Phi}/A_{ev}\right)$.
From the morphology of the feature in Figure 4, we conservatively estimate
$A_{\Phi}/A_{ev}<0.25$, i.e. the forward cross section of the “horns” of the
vee appears to be well under 1/4 of the cross section of the vee itself. This
value yields a derived background solar wind density of $n_{sw}<8\,cm^{-3}$ at
1 AU to maintain the force balance in Equation 7. That figure is again in line
with the wind measurements from _Ulysses_ at $15^{\circ}$ heliographic
latitude (McComas et al. 2000), adding to the plausibility of the accretion
force balance picture. The corresponding mass density limit is
$\rho_{sw}<8\times 10^{-19}\,kg\,m^{-3}$ at 1 AU
As a sanity check, we can use this $\rho_{sw}$ limit and Equation 8 to find
that $\Delta v$ must then be a few tens of $km\,s^{-1},$ i.e. the background
wind speed must be close to $300\,km\,s^{-1}$.
We conclude that the picture of force balance between snowplow accretion and
the tension force is at least broadly consistent with the observed feature,
though further study of more events (preferably with corresponding _in situ_
measurements of the feature itself) is necessary.
Figure 8: Semilog regression fit of $m_{ev}(r)$
## 3 Discussion
Using data from STEREO/SECCHI, we have identified and measured the
characteristics of a single flux disconnection event and associated cusp
feature, similar to that discovered by McComas et al. (1992), from initial
detection in the lower corona to distances beyond 1 AU. The cusp feature is
formed in the classic X-point geometry and rapidly accelerates under the
tension force to approximately $320\,km\,s^{-1}$, which it reaches in under 4
hours at an altitude of approximately 10$R_{\odot}$. Thereafter the feature
continues to accumulate mass but maintains approximately constant speed until
it is lost to sight 1.2 AU from the Sun.
Based on photometry, we are able to estimate the onset mass of the event as
$25\,Tg$ and the entrained flux as $160\,GWb,$ corresponding to a coronal
field strength of $0.08\,G$ and an $r^{2}$-normalized surface open field of
$4\,G$ over the projected surface footprint of the feature. These estimates
are likely low by a factor of order $\sqrt{2}$, because they make use of
feature excess brightness rather than absolute Thomson-scattered brightness in
the coronagraph images; using polarized-brightness imagery could improve the
measurement by separating the non-transient component of the Thomson
scattering signal from the unwanted F coronal background.
Because our measurements are all based on morphology and photometry, we have
performed several consistency checks to build confidence in the calculated
parameters of the feature as it propagates. In particular, a model of simple
force balance between the tension force and mass accretion is consistent with
both the inferred magnetic field and accepted values for background slow solar
wind density and speed.
Simple accretion models such as we developed here demonstrate clearly why
ejected features such as U-loops or CMEs seem frequently to propagate at near
constant speed: under continuous weak driving, an equilibrium forms rapidly
between the driving force and momentum transfer by mass accretion. The
equilibrium outflow speed is the sum of a large, fixed (or at least driver-
independent) speed – that of the surrounding wind – with a smaller offset
speed that drives mass accretion. Thus the feature speed is quite insensitive
to the driver. In our case, doubling the tension force would only increase the
outflow rate by $\sim 10\%$.
The event under study is well presented, but is not unusual at all; such
events are easy to identify in heliospheric image sequences, because of their
distinctive “U” and cusp shape; they are readily traced back to the corona.
This technique represents a new, very effective way of finding these
disconnection events, which are small and hard to identify __ in the
coronagraph sequences alone, but are strongly and easily visible in the
processed heliospheric images.
In an initial reduced data set of 36 days near the deepest part of the recent
extended solar minimum (2008 Dec – 2009 Jan), we identified 12 such events;
all of them were identified by tracking “V” or “U” shapes back from the
heliospheric images to the corona. Assuming the present feature to be typical,
and considering that the single viewpoint affords clear coverage of about 1/4
of the circumference of the Sun, we estimate the global disconnection feature
rate at that time to be over $1\,event\,d^{-1}$, and the flux disconnection
rate to thus be at least of order $60\,TWb\,y^{-1}$. Expanded to a 1 AU
sphere, this amounts to a rate of change of the open field of order
$0.2\,nT\,y^{-1}$, which is a significant fraction of the observed cycle-
dependent rate of change of the open heliospheric field (e.g. Schwadron,
Connick, & Smith 2010). These figures are based on a single calculated flux
and an event rate obtained by initial visual inspection of a single 36-day
data set, and hence are merely rough estimates – but they indicate that flux
disconnections of this type are important to the global balance of open flux.
Further study, in the form of a systematic survey, is needed to determine
whether they are the primary mechanism of flux disconnection from the Sun.
The authors thank the STEREO instrument teams for making their data available.
Our image processing made heavy use of the freeware Perl Data Language
(http://pdl.perl.org). The work was enhanced by enlightening conversations
with J. Burkepile, C. Eyles, and N. Schwadron, to whom we are indebted. This
work was supported by NASA’s SHP-GI program, under grant NNG05GK14G.
## 4 References
## References
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* McComas et al. (1991) McComas, D. J., Phillips, J. L., Hundhausen, A. J., & Burkepile, J. T. 1991, Geophys. Res. Lett., 18, 73
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|
arxiv-papers
| 2011-11-30T15:37:46 |
2024-09-04T02:49:24.827543
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "C. E. DeForest, T. A. Howard, D. J. McComas",
"submitter": "Craig DeForest",
"url": "https://arxiv.org/abs/1111.7211"
}
|
1111.7243
|
Quantifying the dynamics of coupled networks of switches and oscillators
Matthew R. Francis1, Elana J. Fertig2,∗,
1 Physics Department, Randolph-Macon College, Ashland, VA, USA
2 Department of Oncology and Division of Oncology Biostatistics and
Bioinformatics, Sidney Kimmel Comprehensive Cancer Center, School of Medicine,
Johns Hopkins University, Baltimore, MD, USA
$\ast$ E-mail: Corresponding ejfertig@jhmi.edu
## Abstract
Complex network dynamics have been analyzed with models of systems of coupled
switches or systems of coupled oscillators. However, many complex systems are
composed of components with diverse dynamics whose interactions drive the
system’s evolution. We, therefore, introduce a new modeling framework that
describes the dynamics of networks composed of both oscillators and switches.
Both oscillator synchronization and switch stability are preserved in these
heterogeneous, coupled networks. Furthermore, this model recapitulates the
qualitative dynamics for the yeast cell cycle consistent with the hypothesized
dynamics resulting from decomposition of the regulatory network into dynamic
motifs. Introducing feedback into the cell-cycle network induces qualitative
dynamics analogous to limitless replicative potential that is a hallmark of
cancer. As a result, the proposed model of switch and oscillator coupling
provides the ability to incorporate mechanisms that underlie the synchronized
stimulus response ubiquitous in biochemical systems.
## Introduction
The dynamics in systems ranging from intercellular gene regulation to
organogenesis are driven by complex interactions (represented as edges) in
subcomponents (represented as nodes) in networks. If the structure of these
networks is known, network-wide models of coupled systems have been applied to
predict their qualitative dynamics. For example, models of coupled switches
based upon Glass networks [1] have been applied to model systems such as
neuronal synapses [2] and gene regulatory networks [3]. Similarly, models of
coupled oscillators along networks based upon the Kuramoto model [4] have been
used to model synchronization of oscillators in diverse systems reviewed in
[5]. In biochemical systems, in vivo oscillator synchronization has been
observed in synthetic oscillatory fluorescent bacteria [6, 7], yeast gene
regulatory networks [8, 9], and human cell fate decisions [10]. Such
spontaneous synchronization has also been attributed to the development of the
mammalian cardiac pacemaker cells (reviewed in [11]) and cortical systems
(reviewed in [12]) including notably the circadian pacemaker (e.g., [13]).
More recently, these network models have been found to be insufficient to
model more complex dynamics in neuronal information transfer [14, 15, 16, 12,
17] and cardiac arrhythmias [18, 19, 20, 21]. These limitations extend to
physical systems, such as the coupled lasers studied in [22]. Therefore,
numerous studies have modified these network models to account for evolving
networks [15, 23, 24, 25, 26, 27, 28], dynamic frequencies [15, 29, 30], or
phase delays [16, 31, 32, 33]. However, these mathematical modifications
typically do not encode the mechanism underlying the limitations in the
Kuramoto and Glass network models.
We hypothesize that the observed limitations in the standard Kuramoto and
Glass models arise from their exclusion of coupling components with
qualitatively different dynamics. Several studies have inferred that
biochemical systems contain “network motifs” with both oscillatory and switch-
like dynamics [34, 35]. The dynamics of these motifs are inferred from the
topology of subgraphs in the networks of these systems. Their structures are
statistically overrepresented in biochemical networks [36, 37] such as
intracellular regulatory networks [38], implicating evolutionary preservation
(and thus utility) of these network motifs [39]. The dynamics of these motifs
have been used to model yeast cell cycle regulation [40] and have been further
confirmed in synthetic, designed biochemical circuits (reviewed in [41]).
Because these heterogeneous network motifs are all identified as components
within a single biochemical network, their interactions must drive the global
dynamics of the network [42]. Previously, [43] have shown that coupling small
sets of heterogeneous network motifs ensures the robustness of motif dynamics
and [42] have shown that coupling networks changes their dynamics in
isolation. However, the network-level dynamics that result from coupling
oscillatory and switch-like components have not been studied comprehensively.
In this paper, we develop a theoretical framework to quantify the network-wide
dynamics resulting from coupling switches and oscillators. This model is based
upon introducing cross-coupling between the Kuramoto and Glass models, due to
their wide success in modeling the dynamics in networks of oscillators and
networks of switches, respectively. Simulations with the proposed model across
state-space in an all-to-all network yields four operational states: (1)
switches remain “on” and oscillators synchronize, (2) switches are “off” and
oscillators freeze, (3) switches fluctuate in sync with oscillators, and (4)
switches fluctuate transitionally until oscillators freeze. Further
application of our model to the network motifs identified in yeast in [44]
recapitulates the qualitative dynamics of the system observed in that study.
However, a simple rewiring of this cell-cycle network that introduces feedback
causes a cancer-like sustained re-activation of the cell cycle machinery
without regard for external signal growth signals. These dynamics suggest that
modeling cross-motif coupling may predict critical processes in the dynamics
of biochemical networks with minimal parameterization.
### The Kuramoto model of coupled oscillators
Quantitative studies of coupled oscillators often apply the Kuramoto model of
$M$ oscillators coupled in an all-to-all network. In this model, the change in
time $\dot{\theta}_{i}$ of the phase of the $i^{\mbox{th}}$ oscillator,
$\theta_{i}$, is governed by
$\dot{\theta_{i}}=\hat{\omega}_{i}+\frac{\kappa_{\theta,\theta}}{M}\sum_{j=1}^{M}\sin\left(\theta_{j}-\theta_{i}\right),$
(1)
where $\hat{\omega}_{i}$ is the natural frequency of the $i^{\mbox{th}}$
oscillator and $\kappa_{\theta,\theta}\geq 0$ is the coupling strength of the
oscillators [4]. Typically, the $\hat{\omega}_{i}$ values are drawn from a
normal distribution centered at 0 with variance $\sigma_{\omega}$.
In the Kuramoto model, the phases of the oscillators will synchronize if
$\kappa_{\theta,\theta}$ is above a threshold coupling strength
$\hat{\kappa}_{\theta,\theta}$. Such synchronization is quantified with the
mean field of the oscillators as
$r_{\theta}e^{i\psi}=\frac{1}{M}\sum_{j=1}^{M}e^{i\theta_{j}}.$ (2)
Here $\psi$ is the average phase of the oscillators and the coherence
$r_{\theta}$ represents the spread of the oscillators from that average phase.
Based upon eq. (2), $r_{\theta}=1$ if each $\theta_{i}=\psi$ and
$r_{\theta}=0$ if the values of $\theta_{i}$ are distributed uniformly between
$[0,2\pi)$ [45].
### Glass networks of coupled switches
Coupled sets of $N$ switches, which adopt one of a set of binary states, are
modeled with Glass networks [1]. These models describe the evolution of the
$i^{\mbox{th}}$ switch ($\tilde{x}_{i}$) as follows
$\dot{x}_{i}=-x_{i}+F_{i}\left(\tilde{x}_{1},\tilde{x}_{2},\ldots,\tilde{x}_{n}\right)\mbox{,
and}$ (3) $\tilde{x}_{i}=0\mbox{ if }x_{i}<0\mbox{; }1\mbox{ otherwise},$ (4)
where $\dot{x}_{i}$ represents the change in time of the value of each
$x_{i}$, which are unobservable continuous variables that control the time of
switching between observable, discrete states in $\tilde{x}_{i}$. In this
model, $F_{i}$ describes the change in state of the $i^{\mbox{th}}$ switch due
to the coupling with the other $N$ switches in the network [1]. In specified
network structures and functions $F_{i}$, such Glass networks can exhibit
complex dynamics, including periodic and aperiodic orbits (e.g., [46]).
One type of Glass network, called a Hopfield network [2], has dynamics
applicable to the smooth-decay of signal in biochemical switches [2]. The
Hopfield model lets
$F_{i}=\kappa_{x,x}\sum_{j=1}^{N}w_{ij}\tilde{x}_{j}-\tau_{i},$ (5)
where $w_{ij}$ takes values between $-1$ and $1$ representing the relative
strength of the connection between switches $i$ and $j$, $\kappa_{x,x}$ is the
magnitude of coupling strengths, and $\tau_{i}$ the threshold for switch
activation. Similar to the Kuramoto model, sets of the switches will
synchronize for $\kappa_{x,x}$ above a threshold $\hat{\kappa}_{x,x}$ in
appropriate network topologies.
## Results
### Network model of coupled oscillators and switches
By combining the established models for switches and oscillators, we model the
dynamics of the heterogeneous system of coupled switches and oscillators in
systems including biochemical networks with the following set of equations:
$\displaystyle\dot{x}_{i}$ $\displaystyle=$ $\displaystyle-
x_{i}+G_{i}\left(\tilde{x}_{1},\tilde{x}_{2},\ldots,\tilde{x}_{N},\theta_{1},\theta_{2},\ldots,\theta_{M}\right)$
(6) $\displaystyle\dot{\theta}_{l}$ $\displaystyle=$
$\displaystyle\omega_{l}\left(\tilde{x}_{1},\tilde{x}_{2},\ldots,\tilde{x}_{N}\right)$
$\displaystyle+$ $\displaystyle
H_{l}\left(\tilde{x}_{1},\tilde{x}_{2},\ldots,\tilde{x}_{N},\theta_{1},\theta_{2},\ldots,\theta_{M}\right).$
Here, eq. (6) is analogous to the Glass network in eq. (3) and $\tilde{x}_{i}$
is defined according to eq. (4).
In this study, we explore a case of the switch-oscillator model in eqs. (6)
and (Network model of coupled oscillators and switches) which contains an all-
to-all network that couples the Kuramoto model, eq. (1), and Hopfield network,
eqs. (3) - (5), as follows
$\displaystyle\dot{x}_{i}$ $\displaystyle=$ $\displaystyle-
x_{i}+\frac{\kappa_{x,x}}{N}\sum_{j\neq.~{}i}^{N}\tilde{x}_{j}+\frac{\kappa_{x,\theta}}{M}\sum_{k=1}^{M}\tilde{\theta}_{k}-\tau_{i},$
(8) $\displaystyle\dot{\theta}_{l}$ $\displaystyle=$
$\displaystyle\omega_{l}+\frac{\kappa_{\theta,\theta}}{M}\sum_{k=1}^{M}\sin\left(\theta_{k}-\theta_{l}\right),$
(9) $\displaystyle\dot{\omega}_{l}$ $\displaystyle=$
$\displaystyle\frac{\kappa_{\theta,x}}{N}\sum_{j=1}^{N}(\tilde{x}_{j}\hat{\omega}_{l}-\omega_{l}),$
(10)
where $\kappa_{x,\theta}$ and $\kappa_{\theta,x}$ are cross-component coupling
strengths. In eq. (10), $\omega_{l}$ is the time-varying frequency of the
$l^{\mbox{th}}$ oscillator resulting from switch coupling, with initial values
$\omega_{l}(t=0)=\hat{\omega}_{l}$, and
$\tilde{\theta}_{l}=\left\\{\begin{array}[]{ll}1&\mbox{ if
}0\leq\theta_{l}<\pi\\\ 0&\mbox{ otherwise}\\\ \end{array}\right..$ (11)
In this system, zero values of the cross-coupling parameters
$\kappa_{x,\theta}$ and $\kappa_{\theta,x}$ cause the model to reduce to the
standard uncoupled Kuramoto and Hopfield models. Similar decoupling of the
models occurs if the switch and oscillator systems are at vastly different
timescales, determined by the $\tau_{i}$ and $\hat{\omega}_{l}$ parameters,
respectively. The transformation in eq. (11) facilitates comparable switch-
like dynamics in the oscillators when they interact with switches in eq. (8).
Nonzero switch-oscillator ($\kappa_{x,\theta}$) interactions will cause an
oscillator in the “up” part ($\tilde{\theta}_{l}=1$) of its cycle to feed
energy into the switch in question, nudging it towards the “on” state if off
or delaying its decay if already on. Similarly, an “on” switch with a nonzero
oscillator-switch interaction ($\kappa_{\theta,x}$) will feed energy into the
oscillators causing them to cycle at their natural frequency if coupled to
that switch.
By thus incorporating coupling between switches and oscillators within the
framework established by the standard Kuramoto and Hopfield models, the
dynamics of our model in eqs. (8) - (10) can be analyzed within the framework
of these well established models. Similar to analysis of the Kuramoto model
and Glass network, we summarize the dynamics of our system using order
parameters. For oscillators, we utilize the order parameter defined in eq.
(2). We introduce a new order parameter
$r_{\omega}(t)=\frac{1}{M}\sum_{\mu}\frac{\omega_{\mu}(t)}{\hat{\omega}_{\mu}}$
(12)
that tracks how closely each individual oscillator’s frequency $\omega_{\mu}$
corresponds to the natural frequency $\hat{\omega}_{\mu}$. Analogously, we
measure the fraction of switches that are in the “on” position at a given time
using a switch-switch order parameter defined by
$r_{x}(t)=\frac{1}{N}\sum_{j}\tilde{x}_{j}(t).$ (13)
Both of these functions will have a maximum of 1 when all switches are on, and
minimum 0 if all switches are off.
### Simulation results in all-to-all networks
We first explore the qualitative dynamics of the heterogeneous system through
numerical simulations in all-to-all networks. We limited these simulations to
all-to-all networks, because of the ability of this network topology to
describe the qualitative dynamics from the Kuramoto model. These simulations
explore the majority of parameter space defined by $\kappa_{x,x}$,
$\kappa_{x,\theta}$, $\kappa_{\theta,\theta}$, $\kappa_{\theta,x}$, and
$\sigma_{\omega}$. Specifically, we select $\kappa_{x,x}=1<\hat{\kappa}_{x,x}$
to ensure that switches are able to turn off without appropriate stimulation
from the oscillators. We consider the effects of switches on oscillators for
values of $\kappa_{\theta,\theta}$ both above and below the Kuramoto threshold
$\hat{\kappa}_{\theta,\theta}$. Figure 1 plots the time-dependent order
parameters observed in the four qualitative states observed in simulations of
the coupled model eqs. (8) - (10) that are reflective of the qualitative
dynamics observed in simulations with these parameter values. Supplemental
videos S1-S4 further summarize the results of these simulations. We note that
these four states were the only qualitative states observed for our coupled
model in all-to-all networks simulated according to the description in the
Methods section. Because $\tau$, $\kappa_{x,\theta}$, and $\sigma_{\omega}$
all control the relative timing of switches and oscillators, their values were
selected in these simulations to optimize visualization in the supplemental
videos. When exploring the effect of timing on the system dynamics, we hold
$\tau$ and $\kappa_{x,\theta}$ fixed while varying $\sigma_{\omega}$. Figure 2
shows the probability of observing the states in Figures 1-1 in 100
simulations of all-to-all systems containing 100 switches and oscillators as a
function of $\kappa_{\theta,x}$ and $\sigma_{\omega}$. Because of their common
control of system timing, we would obtain comparable distributions when
varying either $\tau$ or $\kappa_{x,\theta}$ instead of $\sigma_{\omega}$.
#### The coupled system preserves synchronization in both oscillators and
switches.
Figure 1 shows a state of the model in which the switches are all in the “on”
state and oscillators are synchronized ($r_{x}$ near 1, $r_{\theta}$ near 1,
and $\psi$ oscillating between $[0,2\pi)$ periodically). While such
synchronization is observed in the uncoupled Hopfield and Kuramoto models, the
oscillator-switch cross coupling extends the region of parameter space over
which this synchronization occurs. Specifically, modest values of
$\kappa_{x,\theta}$ can induce sustained switch activity for parameter values
of $\kappa_{x,x}$ in which switches would decay in the uncoupled system.
Furthermore, this switch synchronization will occur for all values of
$\kappa_{x,x}$ in which synchronization occurs in the uncoupled Hopfield model
(i.e., all $\kappa_{x,x}$ larger than a threshold value $\hat{\kappa}_{x,x}$)
because the oscillators only contribute positively to the derivative in eq.
(8) in our model. On the other hand, no value of $\kappa_{\theta,x}$ will
cause oscillator phases to synchronize if $\kappa_{\theta,\theta}$ is below
the critical coupling parameter for the pure Kuramoto model
($\hat{\kappa}_{\theta,\theta}$). However, there are parameter regimes in
which this synchronization occurs stochastically, depending on the initial
values selected for $x_{i}$, $\theta_{j}$, and $\hat{\omega}_{j}$ (Figure 2).
In these cases, the average decrease in oscillator natural frequencies caused
by decreasing $\kappa_{\theta,x}$ or $\sigma_{\omega}$ will increase the
effective period of oscillators, thereby increasing the probability of
switches being locked in the “on” state and oscillator synchronization in the
heterogeneous system.
#### Coupling switches to unsynchronized oscillators can freeze network-wide
dynamics.
Figure 1 depicts a model state in which switches are all “off” ($r_{x}$ near
zero) and oscillators “freeze”: each
$\theta_{j}\left(t\right)=\psi\left(t\right)=\Psi$ for some constant values
$\Psi$ for all $t$ beyond the preliminary freezing time $t_{f}$. While the
decaying switches are observed in an uncoupled Hopfield model, the freezing
oscillators cannot be simulated in the uncoupled Kuramoto model. Such
oscillator freezing will occur whenever the oscillators decay to the “off”
state by virtue of the coupling of the oscillators to switches through the
$\omega_{j}$ in eq. (10). Specifically, this frozen state can occur whenever
$\kappa_{x,x}<\hat{\kappa}_{x,x}$ depending on the values of $\tilde{x}_{i}$,
$\theta_{j}$, and $\hat{\omega}_{j}$. However, the probability of selecting
these initial states is decreased when the heterogeneity of the oscillators
increases through incomplete synchronization ($r_{\theta}(t)<1$) or increased
$\sigma_{\omega}$ (Figure 2). In these cases, a single oscillator in the “up”
phase ($\tilde{\theta}=1$) can contribute positively to the switch states,
forcing the system out of this frozen state. The probability of obtaining this
frozen model state further depends on the relative timing of switch decay and
oscillator freezing. Specifically, the probability of obtaining the frozen
state decreases with the average oscillator frequency, determined
predominantly by the parameter $\kappa_{\theta,x}$ (Figure 2).
#### Coupling switches to synchronized oscillators can induce synchronized
oscillations in switches.
An additional consequence of coupling switches and oscillators in a state in
which switches vacillate between all “on” and all “off” along with the
synchronized oscillator frequency (Figure 1). This oscillatory synchronization
occurs when the pure Hopfield model would turn switches “off”
($\kappa_{x,x}<\hat{\kappa}_{x,x}$), the pure Kuramoto model would induce
oscillator synchronization
($\kappa_{\theta,\theta}>\hat{\kappa}_{\theta,\theta}$), and the timing
between the oscillators and switches are balanced such that the average period
of the coupled oscillators is slightly less than the average decay time of the
system of switches. Figure 2 shows that this balance in switch-oscillator
timescales increases with decreasing $\kappa_{\theta,x}$ and depends non-
monotonically on $\sigma_{\omega}$. As we see in the plot of
$r_{\omega}\left(t\right)$ in Figure 1, the average oscillator natural
frequencies will decrease towards the end of the “down” phase in response to
switches turning off, and then increase to their full natural values in the
“up” phase as switches turn back “on”. Therefore, if synchronized oscillator
period is too slow (i.e., $\sigma_{\omega}$ is too large), the system will
tend to be locked in the “on” state (Figure 2); if too fast (i.e.,
$\sigma_{\omega}$ too small) the system will tend to be locked in the “off”
state (Figure 2).
#### Synchronization of network-wide oscillations may be transitory.
Oscillatory behavior in the switches is also observed for unsynchronized
oscillators ($\kappa_{\theta,\theta}<\hat{\kappa}_{\theta,\theta}$) as
depicted in Figure 1. In this case, the value of $\kappa_{\theta,x}$ must be
large enough to enable switches to freeze the oscillators’ phases. However,
because the oscillators are uncoupled, a small subset of oscillators in the
“up” phase can drive the switches to turn on for large-enough values of
$\kappa_{x,\theta}$. These switch oscillations are transitory, ending when at
last the switch coupling dominates the system and induces all of the
oscillators to freeze. For unsynchronized oscillators in the parameter range
of Figure 1, the transitional oscillations in the switch state occurs
regularly in 21 of 100 simulations. In 8 of these 21 simulations, the switch
state turns “on” after decaying at least twice. More rarely, transitory
changes in switch state may be induced by a similar mechanism in simulations
for which $\kappa_{\theta,\theta}>\hat{\kappa}_{\theta,\theta}$ and switches
ultimately settle on the all “on” or all “off” states.
#### System size affects the distribution of qualitative dynamics
We also explored the dynamics of the coupled system for networks of sizes
ranging from $N=M=10$ to $N=M=500$ nodes, described in the methods. For
networks of all sizes, we observe that the dynamics of the system was limited
to the four qualitative behaviors observed for networks of size $N=M=100$
depicted in Figure 1. However, the system size does have a notable effect on
the frequency with which each of these behaviors occurs. Supplemental Figures
S5-S7 plot the observed frequencies for each of the network sizes as a
function of the $\kappa_{\theta,x}$ and $\sigma_{\omega}$ values considered in
Figure 2.
When $\kappa_{\theta,x}=0.01$, the observed frequencies of the system states
depend most strongly on network size in simulations using the smallest value
of $\sigma_{\omega}=1$ is also small (Supplemental Figure S5). In this case,
the probability of observing the system with synchronized oscillatory dynamics
in both switches and oscillators decays as the network grows. Both the state
in which the switches are on and oscillators are synchronized and the state in
which the switches are off and oscillators are frozen have with compensatory
increases in probability (Figure 3). The relative probability of obtaining the
frozen state increases, with notable decay in the probability of obtaining the
state in which switches are “on” and oscillators are synchronized in large
networks.
On the other hand, when $\kappa_{\theta,x}=1$, the system size has the
greatest influence on the resulting dynamics for large values of
$\sigma_{\omega}$ (Supplemental Figure S7). In this case, the system changes
from containing mostly switches in the on state and synchronized oscillators
to switches that are entirely in the “frozen” state for large network sizes
(Figure 4). We hypothesize that the system is forced into the frozen state in
larger networks because of increased oscillator synchronization in large
networks. Therefore, small networks would have a higher probability of having
few oscillators that are unsynchronized and in the “up” phase
($\tilde{\theta}=1$), causing the switches to turn “on” ($\tilde{x}=1$) due to
the structure of eq. (8) as was discussed previously. Furthermore, the rare
oscillations observed in both switches and oscillators when
$\kappa_{\theta,x}=1$ occur only when the network is small. Intermediate
values of $\kappa_{\theta,x}=0.1$ show similar changes to those described for
$\kappa_{\theta,x}=0.01$ when $\sigma_{\omega}=1$ and to those described for
$\kappa_{\theta,x}=1$ when $\sigma_{\omega}=10$ (Supplemental Figure S6).
### The heterogeneous network models qualitative dynamics of the yeast cell
cycle derived from network motifs.
Previous work by [47] make the cell cycle processes controlling mitotic
division of fission yeast Schizosaccharomyces pombe cells provides an optimal
system in which to apply our model. The biochemical reactions responsible for
driving the cell cycle are well understood and the resulting dynamics in each
of the stages of the cell cycle have been characterized extensively in [47,
44, 40]. The cell cycle machinery in mitosis is divided into four, sequential
stages: phase 1 is a gap or rest phase (G1); phase 2 is a DNA synthesis stage
(S); phase 3 is an additional gap stage (G2); and phase 4 is the mitotic
division stage (M). Previously, [47] observed that the dynamics of the yeast
cell cycle can be divided into three sequentially interacting modules,
triggered by a signal based upon cell size: (1) G1/S transitions with a
toggle-switch, (2) S/G2 transitions with a toggle-switch, and (3) G2/M
transitions with an oscillator. Although the specific timing differs from
[47], we observe similar qualitative dynamics to that observed in [47] when
applying our heterogeneous model to evolve the state of these cell cycle
stages (Figure 5) as described in the Methods section. We note that the
response in this system is consistent with the transitory oscillations
observed in Figure 1 in the case of all-to-all coupling. We also modeled this
cell-cycle system in a rewired-network, in which the G2/M transitions feedback
into G1/S (Figure 6). In this case, we observe sustained reactivation of the
cell cycle regardless of the external signal. These dynamics are analogous to
the synchronized dynamics in Figure 2 and consistent with cell growth arising
from re-wiring biochemical reactions in cancer cells [48].
## Discussion
Our model of coupled switches and oscillators in all-to-all networks
demonstrates that networks with components having heterogeneous dynamics can
exhibit synchronization similar to that observed in homogeneous systems. As is
the case in homogeneous models (e.g., [49, 50, 51, 52]), we expect analogous
synchronization to hold in small-world, biochemical network topologies (e.g.,
[53]). However, these alternative topologies would likely change the
probability of observing each of the qualitative model behaviors similar to
the observed dependence of probabilities in network size. In this alternative
network topologies, the qualitative states of the network model may have
greater variability in small network sizes in accordance with the findings of
[54]. Finally, in these topologies the heterogeneous model could yield
additional, complex qualitative dynamic states, resulting from the complex
dynamics that they cause in models of coupled switches alone [46].
While uncoupled network motifs may adopt switch-like or oscillatory dynamics,
coupling between these components can induce switch-like behavior in
oscillators and oscillatory behavior in switches. These qualitative changes in
component dynamics occur stochastically, depending on the distribution of
frequencies and switch states. They are more likely to occur in simulations
with an imbalance in relative timescales, in which the dynamics of the faster
network motif will dominate the system. Similarly, when $\kappa_{\theta,x}$
and $\sigma_{\omega}$ are both small, the coordinated oscillations in the
switches and oscillators that occur in frequently small networks are largely
eliminated in larger networks. We hypothesize that this larger network
effectively increases the range of natural frequencies and phases, making the
simulation less likely to have the constrained distribution required to obtain
such synchronized oscillations. We can expect that biological systems have
evolved components according to these distributions to ensure the robustness
of the dynamics in the system. For example, multiple proteins can often serve
similar functions in cell signaling pathways, which would increase the system
size and decrease the probability of transient behaviors in our model. This
robustness will be further ensured through the sheer size of most biochemical
systems. For example, in humans yeast two-hybrid maps and metabolic network
maps both contain on the order of thousands of interactions between thousands
of species [53].
Furthermore, we have also observed that the heterogeneous network model will
freeze the oscillator dynamics in the presence of inactive switches and then
subsequently activate in synchrony in the presence of active switches. As a
result, our model provides a natural mechanism for the coordination of complex
machinery such as the initiation of cell-cycle dynamics. For example, when we
apply our model to the yeast cell cycle motifs in [47], we recapitulate the
qualitative dynamics of delayed initiation of stages of the cell cycle
observed in simulations with differential equations of the regulatory dynamics
in [47]. Additional tuning of the model parameters or incorporation of
additional cell cycle checkpoints would facilitate a precise match of the
timing of [47]. Because parameters are defined for modules and their
interaction, our model requires far fewer rate parameters than any
differential equation model of sets of biochemical reactions of the yeast cell
cycle. Generally, the oscillator in the final G2/M step of the cell cycle is
active only when the series of switches in the previous steps of the cell
cycle are activated, consistent with the transient dynamics observed in our
network model. However, rewiring the network to introduce feedback from the
G2/M stage to the G1/S stage of the cell cycle will cause the modeled cell
cycle machinery to engage continually without regard to the external growth
signals, consistent with the malignant rewring in cancer cells [48]. Similar
to the oscillatory behavior induced in switches in simulations in all-to-all
networks, this small modification to the topology of cell cycle interactions
altered the resulting dynamics of the network motifs for the G1/S and S/G2
motifs. We, therefore, hypothesize that motif dynamics predicted by the
structure of subgraphs may not accurately describe their in vivo dynamics if
considered in isolation, consistent with the hypothesis in [55] and findings
of [42].
We observed that the switches in the cell cycle block activation of the yeast
cell cycle when no external signal is present. Similarly, when part of the
larger but sparse networks that compose biochemical systems [53], inactive
switches would effectively destroy links between nodes on the network. As a
result, the proposed heterogeneous model provides a potential mechanism for
Kuramoto-based models with evolving network topologies such as [15, 23, 24,
25, 26, 27, 28]. Similarly, we observed that the intermediate switches delay
the oscillations in the final G2/M motif in the simulated yeast cell cycle. As
a result, we hypothesize that coupling switches to oscillators through their
frequencies in this model also provides a natural mechanism for extensions of
the Kuramoto model with dynamic frequencies [15, 29, 30] or phase delays [16,
31, 32, 33].
The heterogeneous network model described in this paper facilitates
characterization of the dynamics of complex, biochemical systems by
abstracting the dynamics of their composite motifs such as the yeast cell
cycle based upon [47]. We note that the proposed heterogeneous network model
is deterministic once the initial values of all the switches and oscillator
frequencies have been specified. However, many intracellular reactions (e.g.,
[56]) and neuronal systems (reviewed in [57, 58]) evolve stochastically. In
these cases, the Hopfield networks used to model the switches could be
replaced with probabilistic Boolean networks [3] and the oscillators evolved
with stochastic solvers such as the stochastic simulation algorithm (reviewed
in [59]), integrated with the methodology developed in [60]. Similar
modifications could also extend the heterogeneous model to incorporate
coupling with components of additional dynamics pertinent to biochemical
systems, such as those of the network motifs enumerated in [44, 34, 35]. These
studies would also ideally consider the dynamics of the heterogeneous network
model in additional small-world and random network topologies, as well as the
topologies defined by neuronal systems and gene regulatory networks.
## Materials and Methods
### Numerical simulations in the all-to-all network
In this study, we analyze the range of possible dynamics of the coupled,
heterogeneous networks by applying this model to all-to-all networks. Analyses
were performed for networks with equal number of switches and oscillators
($N=M$) of sizes 10, 50, 100, 200, and 500. All simulations are run one
hundred times from random initial conditions for the state of switches
($x_{i},i=1,\ldots,N$) and oscillators ($\theta_{j},j=1,\ldots,M$), drawn from
a Gaussian distribution and a uniform distribution on $[0,2\pi)$,
respectively. Similarly, oscillator natural frequencies are drawn randomly
from a Gaussian distribution of mean zero and standard deviation parameter
$\sigma_{\omega}$. Simulations of 100 seconds (in the arbitrary units of the
model), with a time step of 0.01 seconds were found sufficient to reflect the
range of possible model behaviors and verify consistency across initial
conditions. The behavior of each simulation is summarized based on the time-
dependent order parameters $r_{\theta}\left(t\right)$ and
$\psi\left(t\right)$, $r_{\omega}\left(t\right)$, and $r_{x}\left(t\right)$.
### Numerical simulations of the yeast cell cycle
Based upon [47], we model the yeast cell cycle as an initiating external
signal (namely the cell size), coupled to a toggle switch representing the
transition between G1/S, a toggle switch representing the transition between
S/G2, and an oscillator representing the transition from G2/M. While the
external signal is incorporated into the model with coupling to the other
switches in eq. (6), its state is not updated by the model. The duration of
this external signal is set at 10 simulated minutes, based upon [47].
Similarly, the initial values of the hidden variable $x$ for the switches in
the G1/S and S/G2 modules are set at -0.5, $\tau$ to 1, and $\kappa_{x,x}$ to
2 to reproduce the approximate 10 minute duration of these switches in [47].
The natural frequency is for the G2/M module set to
$\frac{2\pi}{60}\mbox{min}^{-1}$ to likewise reflect the timescale reported in
[47], while the remainder of the coupling parameters are left untuned, set to
$\kappa_{\theta,x}=\kappa_{x,\theta}=\kappa_{\theta,\theta}=2$ because we
sought only to reproduce the qualitative dynamics of the [47] model. The
rewiring in the system with enduring cell cycle activation is achieved by
adding an edge from the module for G2/M to the switch in the G1/S module.
## Acknowledgments
This work was sponsored by NCI (CA141053). We thank LV Danilova, B Fertig, AV
Favorov, BR Hunt, MJ Stern, MF Ochs, E Ott, E Webster, and LM Weiner for
advice. Code available upon request.
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## Figure Legends
###### List of Figures
1. 1 Summary of the qualitative dynamics of the heterogeneous network model of eqs. (8) - (10). In all figures, top-panel shows temporal evolution of the mean field statistics ($r_{\theta}$ black, solid; $r_{x}$ green, dashed; and $r_{\omega}$ blue, dash-dotted) and the bottom-panel shows the evolution of the mean phase $\psi$ (red, solid). (a) Oscillators synchronize and switches stay “on” ($\kappa_{x,x}=11$, $\kappa_{x,\theta}=1.5$, $\kappa_{\theta,x}=1$, $\kappa_{\theta,\theta}=40$, and $\sigma_{\omega}=10$), (b) oscillators freeze (as evidenced by unchanging $\psi$) and switches stay “off” ($\kappa_{x,x}=1$, $\kappa_{x,\theta}=1.5$, $\kappa_{\theta,x}=1$, $\kappa_{\theta,\theta}=40$, and $\sigma_{\omega}=10$), (c) oscillators synchronize and switches oscillate ($\kappa_{x,x}=1$, $\kappa_{x,\theta}=160$, $\kappa_{\theta,x}=0.2$, $\kappa_{\theta,\theta}=42$, and $\sigma_{\omega}=3$), and (d) transitory oscillations in oscillators and switches ($\kappa_{x,x}=0.1$, $\kappa_{x,\theta}=1.4$, $\kappa_{\theta,x}=2$, $\kappa_{\theta,\theta}=1.8$, and $\sigma_{\omega}=10$).
2. 2 Percentage of simulations in which the qualitative dynamics in Figure 1 occur. In (a) oscillators synchronize and switches are “on”, in (b) oscillators freeze and switches are “off”, and in (c) switches vary with oscillators vs $\sigma_{\omega}$ for $\kappa_{\theta,x}=0.01$ (solid), $\kappa_{\theta,x}=0.1$ (dotted) and $\kappa_{\theta,x}=1$ (dashed). $\kappa_{x,x}=1<\hat{\kappa}_{x,x}$, $\kappa_{x,\theta}=1.5$, and $\kappa_{\theta,\theta}=40>\hat{\kappa}_{\theta,\theta}$.
3. 3 Dependence on network size for qualitative states for $\kappa_{\theta,x}=0.01$ and $\sigma_{\omega}=1$. Percentage of simulations in which the qualitative dynamics have switches off and oscillators frozen (blue, solid), switches on and oscillators synchronized (green, dashed), and oscillatory switches and synchronized oscillators (red, dotted).
4. 4 Dependence on network size for qualitative states for $\kappa_{\theta,x}=1$ and $\sigma_{\omega}=10$. Percentage of simulations with qualitative dynamics plotted as described in Figure 3.
5. 5 Simulated dynamics of the yeast cell cycle Evolution of the states of the cell cycle modules (G1/S top, green dashed; S/G2 top, red dotted; G2/M bottom, black) in response to an external stimulus to initiate the cell cycle (top, blue solid)
6. 6 Simulated dynamics of an aberrant cell cycle network. As for Figure 5 with a network topology linking the G2/M module to the G1/S module.
Figure 1: Summary of the qualitative dynamics of the heterogeneous network
model of eqs. (8) - (10). In all figures, top-panel shows temporal evolution
of the mean field statistics ($r_{\theta}$ black, solid; $r_{x}$ green,
dashed; and $r_{\omega}$ blue, dash-dotted) and the bottom-panel shows the
evolution of the mean phase $\psi$ (red, solid). (a) Oscillators synchronize
and switches stay “on” ($\kappa_{x,x}=11$, $\kappa_{x,\theta}=1.5$,
$\kappa_{\theta,x}=1$, $\kappa_{\theta,\theta}=40$, and $\sigma_{\omega}=10$),
(b) oscillators freeze (as evidenced by unchanging $\psi$) and switches stay
“off” ($\kappa_{x,x}=1$, $\kappa_{x,\theta}=1.5$, $\kappa_{\theta,x}=1$,
$\kappa_{\theta,\theta}=40$, and $\sigma_{\omega}=10$), (c) oscillators
synchronize and switches oscillate ($\kappa_{x,x}=1$, $\kappa_{x,\theta}=160$,
$\kappa_{\theta,x}=0.2$, $\kappa_{\theta,\theta}=42$, and
$\sigma_{\omega}=3$), and (d) transitory oscillations in oscillators and
switches ($\kappa_{x,x}=0.1$, $\kappa_{x,\theta}=1.4$, $\kappa_{\theta,x}=2$,
$\kappa_{\theta,\theta}=1.8$, and $\sigma_{\omega}=10$).
Figure 2: Percentage of simulations in which the qualitative dynamics in
Figure 1 occur. In (a) oscillators synchronize and switches are “on”, in (b)
oscillators freeze and switches are “off”, and in (c) switches vary with
oscillators vs $\sigma_{\omega}$ for $\kappa_{\theta,x}=0.01$ (solid),
$\kappa_{\theta,x}=0.1$ (dotted) and $\kappa_{\theta,x}=1$ (dashed).
$\kappa_{x,x}=1<\hat{\kappa}_{x,x}$, $\kappa_{x,\theta}=1.5$, and
$\kappa_{\theta,\theta}=40>\hat{\kappa}_{\theta,\theta}$. Figure 3: Dependence
on network size for qualitative states for $\kappa_{\theta,x}=0.01$ and
$\sigma_{\omega}=1$. Percentage of simulations in which the qualitative
dynamics have switches off and oscillators frozen (blue, solid), switches on
and oscillators synchronized (green, dashed), and oscillatory switches and
synchronized oscillators (red, dotted). Figure 4: Dependence on network size
for qualitative states for $\kappa_{\theta,x}=1$ and $\sigma_{\omega}=10$.
Percentage of simulations with qualitative dynamics plotted as described in
Figure 3. Figure 5: Simulated dynamics of the yeast cell cycle Evolution of
the states of the cell cycle modules (G1/S top, green dashed; S/G2 top, red
dotted; G2/M bottom, black) in response to an external stimulus to initiate
the cell cycle (top, blue solid) Figure 6: Simulated dynamics of an aberrant
cell cycle network. As for Figure 5 with a network topology linking the G2/M
module to the G1/S module.
### Supplemental Figure Captions
Figure S5. Dependence of dynamics on network size for
$\kappa_{\theta,x}=0.01$. Number of simulations (of 100) for which the
switches are off and oscillators are frozen (left panel), the switches are on
and the oscillators are synchronized (center panel), and both the oscillators
and switches have synchronized oscillations (right).
Figure S6. Dependence of dynamics on network size for $\kappa_{\theta,x}=0.1$.
As for Supplemental Figure S5.
Figure S7. Dependence of dynamics on network size for $\kappa_{\theta,x}=1$.
As for Supplemental Figure S5.
|
arxiv-papers
| 2011-11-30T17:33:02 |
2024-09-04T02:49:24.837198
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Matthew R. Francis and Elana J. Fertig",
"submitter": "Elana Fertig",
"url": "https://arxiv.org/abs/1111.7243"
}
|
1111.7302
|
# Self-Similar Solutions for Viscous and Resistive ADAF
Kazem Faghei
School of Physics e-mail: kfaghei@du.ac.ir Damghan University Damghan Iran
(Received 2011 January; accepted 2011 April)
###### Abstract
In this paper, the self-similar solution of resistive advection dominated
accretion flows (ADAF) in the presence of a pure azimuthal magnetic field is
investigated. The mechanism of energy dissipation is assumed to be the
viscosity and the magnetic diffusivity due to turbulence in the accretion
flow. It is assumed that the magnetic diffusivity and the kinematic viscosity
are not constant and vary by position and $\alpha$-prescription is used for
them. In order to solve the integrated equations that govern the behavior of
the accretion flow, a self-similar method is used. The solutions show that the
structure of accretion flow depends on the magnetic field and the magnetic
diffusivity. As, the radial infall velocity and the temperature of the flow
increase, and the rotational velocity decreases. Also, the rotational velocity
for all selected values of magnetic diffusivity and magnetic field is sub-
Keplerian. The solutions show that there is a certain amount of magnetic field
that the rotational velocity of the flow becomes zero. This amount of the
magnetic field depends on the gas properties of the disc, such as adiabatic
index and viscosity, magnetic diffusivity, and advection parameters. The
solutions show the mass accretion rate increases by adding the magnetic
diffusivity and in high magnetic pressure case, the ratio of the mass
accretion rate to the Bondi accretion rate decreases as magnetic field
increases. Also, the study of Lundquist and magnetic Reynolds numbers based on
resistivity indicates that the linear growth of magnetorotational instability
(MRI) of the flow decreases by resistivity. This property is qualitatively
consistent with resistive magnetohydrodynamics (MHD) simulations.
###### keywords:
accretion, accretion disks, magnetohydrodynamics: MHD
## 1 Introduction
The standard geometrically thin, optically thick accretion disc model can
successfully explain most of the observational features in active galactic
nuclei (AGN) and X-ray binaries (Shakura & Sunyaev 1973). In the standard thin
model, the motion of the matter in the accretion disc is nearly Keplerian, and
the gravitational energy released in the disc is radiated away locally. During
recent years, another type of accretion flow has been studied in which it is
assumed that energy released through dissipation processes in the disc may be
trapped within accreting gas, and only the small fraction of the energy
released in the accretion flow is radiated away due to inefficient cooling,
and most of the energy is stored in the accretion flow and advected to the
central object. This kind of accretion flow is known as advection-dominated
accretion flow (ADAF). The basis ideas of ADAF models have been developed by a
number of researchers (e.g. Ichimaru 1977; Rees et al. 1982; Narayan & Yi
1994, 1995; Abramowicz et al. 1995; Ogilvie 1999; Blandford & Begelman 1999).
It seems that accretion discs, whether in star-forming regions, in X-ray
binaries, in cataclysmic variables, or in the centers of active galactic
nuclei, are likely to be threaded by magnetic fields. Consequently, the role
of the magnetic fields on ADAFs has been analyzed in detail by a number of
investigators (Bisnovatyi-kogan & Lovelace 2001; Akizuki & Fukue 2006,
hereafter AF06; Shadmehri 2004, hereafter Sh04; Ghanbari et al., 2007; Bu,
Yung, & Xie 2009; Khesali & Faghei 2009, hereafter KF09). The existence of the
toroidal magnetic field has been proven in the outer regions of the discs of
young stellar objects (YSOs; Aitken et al. 1993; Wright et al. 1993) and in
the Galactic center (Novak et al. 2003; Chuss et al. 2003). Thus, considering
the accretion discs with a toroidal magnetic field have been studied by
several authors (AF06; Begelman & Pringle 2007; Khesali & Faghei 2008,
hereafter KF08; Bu, Yung, & Xie 2009; KF09).
The resistive diffusion of magnetic field is important in some systems, such
as the protostellar discs (Stone et al. 2000; Fleming & Stone 2003), discs in
dwarf nova systems (Gammie & Menou 1998), the discs around black holes (Kudoh
& Kaburaki 1996), and Galactic center (Melia & Kowalenkov 2001; Kaburaki et
al. 2010). Also, two and three dimensional simulations of local shearing box
have shown that resistive dissipation is one of the crucial processes that
determines the saturation amplitude of the magnetorotational instability
(MRI). As, linear growth rate of MRI can be reduced significantly because of
the suppression by ohmic dissipation (Fleming et al. 2000; Masada & Sano
2008).
AF06 proposed a self-similar advection-dominated accretion flow that the disc
plasma is highly ionized, so they assumed that resistivity of the plasma is
zero, and only viscosity is due to turbulence and dissipation in the disc.
However, recent works represent importance of magnetic diffusivity in
accretion discs (e.g. Kuwabara et al. 2000; Kaburaki 2000 ;Kaburaki et al.
2002; Sh04; Ghanbari et al. 2007; Krasnopolsky et al. 2010; Kaburaki et al.
2010). Sh04 studied a quasi-spherical accretion flow that dominant mechanism
of energy dissipation was assumed to be the magnetic diffusivity due to
turbulence in the accretion flow and the viscosity of the fluid was completely
neglected. The main focus of Sh04 was nonrotating accretion flow and ignored
from toroidal magnetic field. Also, Sh04 studied induction equation of
magnetic field in a steady state that is not according to anti-dynamo theorem
(e.g. Cowling 1981) and is usefull only in particular systems where the
magnetic dissipation time is very long, much longer than the age of the
system. Ghanbari, Salehi, & Abbassi (2007) considered an axisymmetric,
rotating, isothermal steady accretion flow, which contains a poloidal magnetic
field of the central star and from toroidal magnetic field of the flow
neglected. They assumed that the mechanisms of energy dissipation are the
turbulence viscosity and magnetic diffusivity due to the magnetic field of the
central star. They explored the effect of viscosity on a rotating disc in the
presence of constant magnetic diffusivity. Similar to Sh04 they considered the
flow in balance between escape and creation of the magnetic field, and ignored
from toroidal component of magnetic field.
As mentioned the observational evidences and the MHD simulations express that
the toroidal component of magnetic field and the magnetic diffusivity are
important in accretion discs. Thus in this paper by using AF06 technique we
will study the influence of presence of toroidal component of magnetic field
in a viscous and resistive accreting gas, and investigate the role of non-
constant magnetic diffusivity in systems that escaping and creating of
magnetic field are not balanced. We will show that the present model from some
aspects is qualitatively consistent with the observational evidences and the
resistive MHD simulation results. The paper is organized as follows. In
section 2, the basic equations of constructing a model for quasi-spherical
magnetized advection dominated accretion flow will be defined. In section 3,
self-similar method for solving equations which govern the behavior of the
accreting gas was utilized. The summary of the model will appear in section 4.
## 2 Basic Equations
We use spherical coordinate $(r,\theta,\varphi)$ centered on the accreting
object and make the following standard assumptions:
* (i)
The flow is assumed to be steady and axisymmetric
$\partial_{t}=\partial_{\varphi}=0$, so all flow variables are a function of
only $r$ ;
* (ii)
The magnetic field has only an azimuthal component;
* (iii)
The gravitational force on a fluid element is characterized by the Newtonian
potential of a point mass, $\Psi=-{GM_{*}}/{r}$, with $G$ representing the
gravitational constant and $M_{*}$ standing for the mass of the central star;
* (iv)
The equations written in spherical coordinates are considered in the
equatorial plane $\theta=\pi/2$ and terms with any $\theta$ and $\varphi$
dependence are neglected (Ogilvie 1999; KF09).
* (v)
For the sake of simplicity, the self-gravity and general relativistic effects
have been neglected;
The behavior of such system can be analyzed by magnetohydrodynamics (MHD)
equations. The general MHD equations are written as follows:
$\frac{\partial\rho}{\partial t}+\nabla\cdot(\rho{\bf v})=0,$ (1)
$\rho\left[\frac{\partial{\bf v}}{\partial t}+({\bf v}\cdot\nabla){\bf
v}\right]=-\nabla p_{gas}-\rho\nabla\Psi+\frac{1}{4\pi}{\bf
J}\times\bf{B}+{\bf{F}}_{vis},$ (2) ${\bf J}=\nabla\times{\bf B}$ (3)
$\frac{\partial{\bf B}}{\partial t}=\nabla\times\left({\bf v}\times{\bf
B}-\eta{\bf J}\right),$ (4)
$\rho\left[\frac{1}{\gamma-1}\frac{d}{dt}\left(\frac{p_{gas}}{\rho}\right)+\left(\frac{p_{gas}}{\rho}\right)(\nabla\cdot{\bf
v})\right]=Q_{\rm diss}-Q_{\rm cool}\equiv fQ_{\rm diss},$ (5)
$\nabla\cdot{\bf B}=0,$ (6)
where $\rho$, $\mathbf{v}$, and $p_{gas}$ are the density, the velocity field,
and the gas pressure, respectively; $\mathbf{F}_{vis}$ is the viscous force
per unit volume; $\mathbf{B}$ is the magnetic field; ${\bf J}$ is the current
density and $\eta$ represents the magnetic diffusivity; $\gamma$ is the
adiabatic index; The term on the right hand side of the energy equation,
$Q_{\rm diss}$, is the rate of heating of the gas by the dissipation and
$Q_{\rm cool}$ represents the energy loss through radiative cooling. The
advection factor, $f$ ($0\leq f\leq 1$), describes the fraction of the
dissipation energy which is stored in the accretion flow and advected into the
central object rather than radiated away. The advection factor of $f$ in
general depends on the details of the heating and cooling mechanism and will
vary with postion (e.g. Watari 2006, 2007; Sinha et al. 2009). However, we
assume a constant $f$ for simplicity. Clearly, the case $f=1$ corresponds to
the extreme limit of no radiative cooling and in the limit of efficient
cooling, we have $f=0$.
Under the assumptions (i)-(v), the approximation of quasi-spherical symmetry,
the equations (1)-(6) become
$\frac{1}{r^{2}}\frac{d}{dr}(r^{2}\rho v_{r})=\dot{\rho},$ (7)
$v_{r}\frac{dv_{r}}{dr}+\frac{1}{\rho}\frac{dp}{dr}+\frac{GM_{*}}{r^{2}}=r\Omega^{2}-\frac{B_{\varphi}}{4\pi
r\rho}\frac{d}{dr}(rB_{\varphi}),$ (8)
$\rho v_{r}\frac{d}{dr}(r^{2}\Omega)=\frac{1}{r^{2}}\frac{d}{dr}\left[\nu\rho
r^{4}\frac{d\Omega}{\partial r}\right],$ (9)
$\frac{1}{\gamma-1}\left[v_{r}\frac{dp}{dr}+\frac{\gamma
p}{r^{2}}\frac{d}{dr}\left(r^{2}v_{r}\right)\right]=fQ_{diss},$
$\frac{1}{r}\frac{d}{dr}\left[rv_{r}B_{\varphi}-\eta\frac{d}{dr}(rB_{\varphi})\right]=\dot{B}_{\varphi}.$
(10)
Here $v_{r}$ the radial velocity, $\dot{\rho}$ the mass-loss rate per unit
volume, $\Omega$ the angular velocity, $B_{\varphi}$ the toroidal component of
magnetic field, $\dot{B}_{\varphi}$ is the field escaping/creating rate due to
a magnetic instability or dynamo effect, $\nu$ the kinematic viscosity
coefficient. We assume both of the kinematic coefficient of viscosity and the
magnetic diffusivity due to turbulence in the accretion flow, it is reasonable
to use these parameters in analogy to the $\alpha$-prescription of Shakura &
Sunyaev (1973) for the turbulent,
$\nu=\alpha\frac{p_{gas}}{\rho\Omega_{K}}(1+\beta)^{1-\mu},$ (11)
$\eta=\eta_{0}\frac{p_{gas}}{\rho\Omega_{K}}(1+\beta)^{1-\mu}$ (12)
where $\Omega_{K}=({GM_{*}}/{r^{3}})^{1/2}$ is the Keplerian angular velocity.
Narayan & Yi (1995) used a similar form for kinematic coefficient of
viscosity, i.e. $\nu=\alpha(p_{gas}/\rho\Omega_{K})$, also Sh04 applied a
similar form for magnetic diffusivity, i.e.
$\eta=\eta_{0}(p_{gas}/\rho\Omega_{K})$. Thus in comparison to Narayan & Yi
(1995) and Sh04 prescriptions, we are using the above equations for the
kinematic coefficient of viscosity $\nu$ and the magnetic diffusivity $\eta$.
The parameters of $\alpha$ and $\eta_{0}$ are assumed to be positive
constants, and we assume that $\alpha$, $\eta_{0}\leq 1$ (Campbell 1999;
Kuwabara et al. 2000; Sh04; King et al. 2007). The parameter of
$\beta[=p_{mag}/p_{gas}]$ is the degree of magnetic pressure,
$p_{mag}=B^{2}_{\varphi}/8\pi$, to the gas pressure. Since, we will apply
steady self-similar method for solving system equations, this parameter will
be constant throughout the disc, but really this parameter varies by position
(see, KF08, KF09). The parameter of $\mu$ is a constant and states importance
degree of total pressure in the kinematic viscosity and the magnetic
diffusivity. Clearly, the case of $\mu=1$ corresponds to traditional
$\alpha$-model and in case of $\mu=0$ total pressure is used. Note the
kinematic coefficient of viscosity and the magnetic diffusivity are not
costant and depend on the physical quantities of the flow. we will show the
quantities of $\nu$ and $\eta$ in our self-similar solution vary with radius
$r^{1/2}$.
The ratio of the kinematic coefficient of viscosity to the magnetic
diffusivity is defined by the magnetic Prandtl number, $P_{m}=\nu/\eta$. By
using equations (12) and (13) the magnetic Prandtl number in present model is
$P_{m}=\alpha/\eta_{0}$. We will consider conditions that $P_{m}=\infty$ (in
case that magnetic diffusivity is zero), $P_{m}\geq 1$, and $P_{m}<1$.
For the heating term in equation (10), $Q_{diss}$, we use two sources of
dissipation: the viscous and resistive dissipations. Thus, for $Q_{diss}$ we
can write
$Q_{diss}=\nu\rho
r^{2}\left(\frac{d\Omega}{dr}\right)^{2}+\frac{\eta}{4\pi}J^{2}$ (13)
where first term is related to viscous dissipation and second term is related
to resistive dissipation. By converting the gas pressure and the magnetic
pressure in terms of sound speed ($c^{2}_{s}=p_{gas}/\rho$) and Alfvén speed
($c^{2}_{A}=B^{2}_{\varphi}/4\pi\rho$), and by using equations (12)-(14) for
the equations (7)-(10), we can write
$\frac{1}{r^{2}}\frac{d}{dr}(r^{2}\rho v_{r})=\dot{\rho},$ (14)
$v_{r}\frac{dv_{r}}{dr}+\frac{1}{\rho}\frac{d}{dr}(\rho
c^{2}_{s})+\frac{GM_{*}}{r^{2}}=r\Omega^{2}-\frac{c^{2}_{A}}{r}-\frac{1}{2\rho}\frac{d}{dr}(\rho
c^{2}_{A}),$ (15)
$\rho
v_{r}\frac{d}{dr}(r^{2}\Omega)=\frac{\alpha}{r^{2}}\frac{d}{dr}\left[\frac{\rho
c^{2}_{s}}{\Omega_{K}}(1+\beta)^{1-\mu}r^{4}\left(\frac{d\Omega}{dr}\right)\right],$
(16)
$\frac{1}{\gamma-1}\left[v_{r}\frac{d}{dr}(\rho c^{2}_{s})+\frac{\gamma\rho
c^{2}_{s}}{r^{2}}\frac{d}{dr}\left(r^{2}v_{r}\right)\right]=fQ_{diss},$
$\frac{1}{r}\frac{d}{dr}\left[\sqrt{4\pi\rho
c^{2}_{A}}\left(rv_{r}-\frac{\eta_{0}(1+\beta)^{1-\mu}}{4\rho\beta\Omega_{K}}\frac{d}{dr}(r^{2}\rho
c^{2}_{A})\right)\right]=\dot{B}_{\varphi},$ (17)
where
$Q_{diss}=\frac{(1+\beta)^{1-\mu}}{\Omega_{K}}\left[\alpha r^{2}\rho
c^{2}_{s}\left(\frac{d\Omega}{dr}\right)^{2}+\frac{\eta_{0}}{8r^{4}\rho\beta}\left(\frac{d}{dr}(r^{2}\rho
c^{2}_{A})\right)^{2}\right].$ (18)
## 3 Self-Similar Solutions
We seek self-similar solutions in the following forms:
$v_{r}=-c_{1}\alpha\sqrt{\frac{GM_{*}}{r}}$ (19)
$\Omega=c_{2}\sqrt{\frac{GM_{*}}{r^{3}}}$ (20)
$c^{2}_{s}=c_{3}\frac{GM_{*}}{r}$ (21)
$c^{2}_{A}=\frac{B^{2}_{\varphi}}{4\pi\rho}=2\beta c_{3}\frac{GM_{*}}{r}$ (22)
where coefficients $c_{1}$, $c_{2}$, and $c_{3}$ are determined later. we
assume a power-law relation for density
$\rho(r)=\rho_{0}r^{s}$ (23)
where $\rho_{0}$ and $s$ are constant. By using above self-similar quantities,
the mass-loss rate and the field escaping/creating rate must have the
following form:
$\dot{\rho}=\dot{\rho}_{0}r^{s-3/2}$ (24)
$\dot{B}_{\varphi}=\dot{B}_{0}r^{\frac{s-4}{2}}$ (25)
where $\dot{\rho}_{0}$ and $\dot{B}_{0}$ are constant.
Substituting the above solutions in the continuity, momentum, angular
momentum, energy, and induction equations [(15)-(20)], we can obtain the
following relations:
$\dot{\rho}_{0}=-\left(s+\frac{3}{2}\right)\alpha\rho_{0}c_{1}\sqrt{GM_{*}},$
(26)
$-\frac{1}{2}c^{2}_{1}\alpha^{2}+1-c^{2}_{2}+c_{3}\left(s-1+\beta(1+s)\right)=0,$
(27)
$c_{1}=3(s+2)(1+\beta)^{1-\mu}c_{3},$ (28)
$-\frac{1}{\gamma-1}\alpha
c_{1}(2s-2+3\gamma)=\frac{1}{2}f(1+\beta)^{1-\mu}\left(9\alpha
c^{2}_{2}+2\eta_{0}\beta c_{3}(1+s)^{2}\right),$ (29)
$\dot{B}_{0}=-\frac{s}{2}GM_{*}\sqrt{2\pi\rho_{0}\beta
c_{3}}\left(2c_{1}\alpha+\eta_{0}c_{3}(1+s)(1+\beta)^{1-\mu}\right).$ (30)
Above equations express for $s=-3/2$, there is no mass loss, while for
$s>-3/2$ mass loss (wind) exists. The escape and creation of magnetic fields
are balanced in $s=-2+(3\eta_{0})/(6\alpha+\eta_{0})$ that $\dot{B}_{0}$
becomes zero, solving it in $s=-3/2$ (no wind) implies that
$\eta_{0}=(6/5)\alpha$. This amount of $\eta_{0}$ corresponds to the magnetic
Prandtl number $5/6$. Thus, when the flow has $s=-3/2$ and
$\eta_{0}=(6/5)\alpha$, we expect the escape and creation of magnetic field
are balanced and there is no mass loss. In $\eta_{0}=0$ for balance of escape
and creation of magnetic field, we have $s=-2$. In this paper, only case of no
wind ($s=-3/2$) is considered that $\dot{\rho}_{0}=0$ and
$\dot{B}_{\varphi}\propto r^{-11/4}$ , we note $\dot{B}_{\varphi}$ in
$P_{m}=5/6$ will be zero, too. The equations (29)-(31) imply that for given
$\alpha$, $\eta_{0}$, $\beta$, $s$, and $f$ form a closed set of equations of
$c_{1}$, $c_{2}$, and $c_{3}$ which will determine behavior of the accretion
flow.
### 3.1 Analysis
By using the equations (29)-(31), the coefficients of $c_{i}$ have the
following forms
$c_{1}=\frac{1}{2\alpha^{2}}\left(D_{4}+\sqrt{D^{2}_{4}+8\alpha^{2}}\right),$
(31)
$c^{2}_{2}=-\frac{2}{9}c_{1}D_{1}D_{2},$ (32)
$c_{3}=\frac{1}{3(s+2)}c_{1}D_{1},$ (33)
where
$D_{1}=\frac{1}{(1+\beta)^{1-\mu}},$ (34)
$D_{2}=\frac{(2s-2+3\gamma)}{(\gamma-1)f}+\frac{\eta_{0}\beta}{3\alpha}\frac{(1+s)^{2}}{(s+2)},$
(35)
$D_{3}=\frac{s-1+\beta(1+s)}{s+2},$ (36)
$D_{4}=\frac{2}{3}D_{1}\left(\frac{2}{3}D_{2}+D_{3}\right).$ (37)
The obtained results imply that the model is parametrized by the ratio of
specific heat, $\gamma$, the standard viscous parameter, $\alpha$, the
magnetic diffusivity parameter, $\eta_{0}$, the degree of magnetic pressure to
gas pressure, $\beta$, the advection parameter, $f$, and the mass-loss rate
parameter, $s$.
The equation (37) implies that the value of $D_{2}$ for a value of $\beta$
will be zero. From equations (22) and (34) we can write $\Omega\propto
c_{2}^{2}\propto D_{2}$. Thus we conclude for a value of $\beta$ that we
define it $\beta_{b}$ (braking $\beta$), the angular velocity will be zero.
Over of $\beta_{b}$, $c_{2}^{2}$ becomes negative that it is not physical. By
solving $D_{2}=0$ for the $\beta$ parameter, we can write
$\beta_{b}=-\frac{3\alpha}{f\eta_{0}}\frac{(s+2)}{(s+1)^{2}}\frac{(2s-2+3\gamma)}{(\gamma-1)}.$
(38)
In the case of no mass-loss, $s=-3/2$, we can write
$\beta_{b}=\frac{18\alpha}{f\eta_{0}}\left(\frac{5/3-\gamma}{\gamma-1}\right).$
(39)
The above equation in terms of the magnetic Prandtl number becomes
$\beta_{b}=18(\frac{P_{m}}{f})\left(\frac{5/3-\gamma}{\gamma-1}\right).$ (40)
The equations (40)-(42) express that the amount of the $\beta_{b}$ parameter
depends on $f$, $P_{m}$, and $\gamma$. As $\beta_{b}$ decreases by adding the
advection degree, while increases by adding the magnetic Prandtl number. In a
flow with high conductivity, that $P_{m}$ is very large, $\beta_{b}$ becomes
very large. Also in the flow with high resistivity and low viscosity,
$\beta_{b}$ will be small. Since the magnetic pressure fraction always is
positive or equal to zero ($\beta\geq 0$), so in presence model $\gamma\leq
5/3$.
Examples of the coefficients $c_{i}$ in two cases of $\mu=0$ and $1.0$ are
shown in figures (1) and (2) as a function of the degree of magnetic pressure
to the gas pressure, ($\beta$), for different value of the magnetic
diffusivity, $\eta_{0}$. In figures (1) and (2), $c_{4}$ is the viscous torque
that is obtained from right hand side of equation (30), $c_{5}$ is total
energy dissipation by the viscosity and the resistivity, and is calculated by
right hand side of equation (31), and $c_{6}$ is the ratio of the resistive
dissipation to the viscous dissipation.
#### 3.1.1 First Case: $\mu=0$
The parameter of $\mu$ has appeared in the equations of the kinematic
coefficient of viscosity and the magnetic diffusivity to indicate important
degree of total pressure in them. In the case of $\mu=0$, the kinematic
coefficient of viscosity and the magnetic diffusivity become
$\displaystyle\nu=\alpha\frac{c^{2}_{s}}{\Omega_{K}}(1+\beta)$
$\displaystyle=\alpha c_{3}(1+\beta)\sqrt{GM_{*}}~{}r^{1/2}$ (41)
$\displaystyle\eta=\eta_{0}\frac{c^{2}_{s}}{\Omega_{K}}(1+\beta)$
$\displaystyle=\eta_{0}c_{3}(1+\beta)\sqrt{GM_{*}}~{}r^{1/2}.$ (42)
Figure 1: Physical quantities of the flow as a function of the degree of
magnetic pressure to the gas pressure, for several values of $\eta_{0}=0$,
$0.12$, $0.15$, and $0.2$ that corresponding to $P_{m}=\infty$, $5/6$, $2/3$,
and $5/10$. The disc density profile is set to be $s=-3/2$ (no wind), the
ratio of the specific heats is set to be $\gamma=1.3$, the viscous parameter
is $\alpha=0.1$, and the advection parameter is $f=1.0$.
In this case the function forms of the kinematic coefficient of viscosity and
the magnetic diffusivity deviate by factor $(1+\beta)$ from the function forms
of used by NY95 and Sh04. Also the profiles of non-resistive and non-magnetic
case is shown to compare the present model with canonical ADAF solutions (e.g.
NY95). The existence of $(1+\beta)$ in $\nu$ and $\eta$ causes the viscosity
and resistivity increase by adding the toroidal magnetic field. The $\nu$ and
$\eta$ have direct effects on the viscous torque ($c_{4}$) and the energy
dissipation ($c_{5}$). Thus, we expect increase of $c_{4}$ and $c_{5}$ by
adding the toroidal magnetic field that the profiles of $c_{4}$ and $c_{5}$
confirm it. Also, by adding the magnetic diffusivity ($\eta_{0}$) of the flow,
$c_{4}$ and $c_{5}$ increase. Because the ohmic dissipation increases due to
resistivity that it also increases the flow temperature (the profiles of
$c_{3}$ confirm it), and since the viscous torque is proportional with
temperature (sound speed), thus the viscous torque, $c_{4}$, increases by
adding the magnetic diffusivity. The profiles of $c_{4}$ and $c_{5}$ imply
that for all value of the $\beta$ and $\eta_{0}$, the total energy dissipation
and the viscous torque are larger than the canonical ADAF solutions. The
profiles of $c_{3}$ show that the temperature of the flow by adding the
toroidal component of magnetic field decreases. This property is qualitatively
consistent with the results of Bu et al. (2009) and KF09. The profiles of
$c_{6}$ show that the ratio of the resistive dissipation to the viscous
dissipation increases by adding the toroidal component of magnetic field
($\beta$) and the magnetic diffusivity ($\eta_{0}$). As in small amounts of
magnetic field and the magnetic diffusivity, the dominant heat generated is
the viscous dissipation, while in large values of magnetic field and the
magnetic diffusivity, the dominant heat generated will be the resistive
dissipation. Also, figure (1) shows by adding the $\beta$ and $\eta_{0}$
parameters, the radial infall velocity, $c_{1}$, increases, and the angular
velocity, $c_{2}$, decreases. It can be due to increase of the viscous torque
($c_{4}$) by parameters of $\beta$ and $\eta_{0}$. The raise of the viscous
torque by adding $\eta_{0}$ and $\beta$ parameters generates a larger negative
torque in angular momentum equation and causes the angular velocity of the
flow decreases, and the matter accretes with larger speed. In the present
model, the flow rotates slower than canonical ADAF and accretes speeder than
it. The increase of the radial infall velocity by adding the parameter of
$\beta$ is qualitatively consistent with the results by Bu et al. (2009) and
KF09.
#### 3.1.2 Second Case: $\mu=1$
In the case of $\mu=1$, for the kinematic coefficient of viscosity and the
magnetic diffusvity we can write
$\displaystyle\nu=\alpha\frac{c^{2}_{s}}{\Omega_{K}}$ $\displaystyle=\alpha
c_{3}\sqrt{GM_{*}}~{}r^{1/2}.$ (43)
$\displaystyle\eta=\eta_{0}\frac{c^{2}_{s}}{\Omega_{K}}$
$\displaystyle=\eta_{0}c_{3}\sqrt{GM_{*}}~{}r^{1/2}.$ (44)
In this case the function forms of the kinematic coefficient of viscosity and
the magnetic diffusivity are the same as Sh04 and NY95 used. The behavior of
the physical quantities of the flow in this case are shown in figure (2). The
absence of $(1+\beta)$ in this case for $\nu$ and $\eta$ in comparison to
previous case causes the value of them decrease by factor $(1+\beta)$. The
effects of absence of this factor increases by adding parameter of $\beta$.
The profiles of the angular momentum transport ($c_{4}$) and total energy
dissipation ($c_{5}$) imply that the amounts of them by factor $(1+\beta)$
compared with previous case decrease. As their behavior in terms of the
toroidal component of magnetic field have changed and they decrease by adding
the $\beta$ parameter. However, the magnetic diffusivity has the previous
effects and increase these two quantities. Also, the solutions show that the
viscous torque and the total dissipation in this case are smaller than
canonical ADAF. The weakening of $c_{4}$ and $c_{5}$ by adding the $\beta$
parameter reduces the radial infall velocity. Here, the radial infall velocity
increases by adding $\eta_{0}$ that is due to increase of $c_{4}$ and $c_{5}$.
The behavior of other physical quantities in terms of the $\beta$ and
$\eta_{0}$ parameters does not change, and represent small variations. To
compare the radial infall velocity profiles with canonical ADAF solutions
implies that the flow accretes slower than canonical ADAF that is different
with previous case. The decrease of the temperature and the viscous torque by
adding the parameter of $\beta$ is qualitatively consistent with the results
of Bu et al. (2009).
Figure 2: Physical quantities of the flow as a function of the degree of
magnetic pressure to the gas pressure, for several values of $\eta_{0}=0$,
$0.12$, $0.15$, and $0.2$ that corresponding to $P_{m}=\infty$, $5/6$, $2/3$,
and $5/10$. The disc density profile is set to be $s=-3/2$ (no wind), the
ratio of the specific heats is set to be $\gamma=1.3$, the viscous parameter
is $\alpha=0.1$, and the advection parameter is $f=1.0$.
### 3.2 Mass Accretion Rate
In according to assumptions of section (2) the mass accretion rate defines as
$\dot{M}=-4\pi r^{2}\rho v_{r}.$ (45)
The mass accretion rate under self-similar transformations of equations (21)
and (25) becomes
$\dot{M}=4\pi\alpha\rho_{0}c_{1}\sqrt{GM_{*}}~{}r^{s+3/2}.$ (46)
In our interesting case, $s=-3/2$ (no wind), for the mass accretion rate we
can write
$\dot{M}=4\pi\alpha\rho_{0}c_{1}\sqrt{GM_{*}}$ (47)
that implies the mass accretion rate does not vary by position. This result is
qualitatively consistent with the results by Sh04, Ghanbari et al (2007), and
AF06. Although the present model of accretion flow is different from Bondi
(1952) accretion in various aspect, we can define Bondi accretion rate as
$\dot{M}_{Bondi}=\pi
G^{2}M_{*}^{2}\left(\frac{\rho(\infty)}{c_{s}^{3}(\infty)}\right)\left[\frac{2}{5-3\gamma}\right]^{(5-3\gamma)/2(\gamma-1)}$
(48)
where $\rho(\infty)$ and $c_{s}(\infty)$ are the density and the sound speed
in the gas far away from the star (Frank et al. 2002). Bondi accretion rate in
terms of our self-similar transformations becomes
$\dot{M}_{Bondi}=\pi\sqrt{GM_{*}}\left(\frac{\rho_{0}}{c_{3}^{3/2}}\right)\left[\frac{2}{5-3\gamma}\right]^{(5-3\gamma)/2(\gamma-1)}.$
(49)
Thus, we can write the mass accretion rate in term of Bondi accretion rate as
follows
$\dot{M}/\dot{M}_{Bondi}=4\alpha
c_{1}c_{3}^{3/2}\left[\frac{2}{5-3\gamma}\right]^{(3\gamma-5)/2(\gamma-1)}r^{s+3/2}.$
(50)
In our interesting case, $s=-3/2$ (no wind), we can write
$c_{7}=\dot{M}/\dot{M}_{Bondi}=4\alpha
c_{1}c_{3}^{3/2}\left[\frac{2}{5-3\gamma}\right]^{(3\gamma-5)/2(\gamma-1)}.$
(51)
Here, we defined new parameter of $c_{7}$ that indicates the ratio of the mass
accretion rate to Bondi accretion rate. Examples of the coefficient of $c_{7}$
in two cases of $\mu=0$ and $1.0$ are shown in figures (3) as a function of
the degree of magnetic pressure to the gas pressure, ($\beta$), for different
value of the magnetic diffusivity, $\eta_{0}$. The profiles of $c_{7}$ show
that in the case $\mu=0$, the mass accretion rate increases by adding the
toroidal magnetic field and the magnetic diffusivity. While the solutions for
$\mu=1$ imply that the mass accretion rate decreases by adding the toroidal
magnetic field and increases by adding the magnetic diffusivity. On the other
hand, the magnetic diffusivity in two cases causes the mass accretion rate
increase. In high magnetic pressure, the $c_{7}$ profiles for both cases show
that the ratio of the mass accretion rate to the Bondi accretion rate is
decreased with an increase in magnetic pressure. This property is
qualitatively consistent with results of Kaburaki (2007). Comparision of the
present model with canonical ADAF solutions implies that in the case of
$\mu=0$ the flow accretes quicker than canonical ADAF, however in the case of
$\mu=1$ the flow accretes slower than canonical ADAF. Also, the profiles of
$c_{7}$ in both of cases show the mass accretion rate in our model is smaller
than Bondi accretion rate that is in adapting with observational evidences
from Sgr A∗, M87, and NGC 4261 (Kaburaki 2007).
Figure 3: The ratio of mass accretion rate to Bondi accretion rate
($c_{7}=\dot{M}/\dot{M}_{Bondi}$) as a function of the degree of magnetic
pressure to the gas pressure, for several values of $\eta_{0}=0$, $0.12$,
$0.15$, and $0.2$ that corresponding to $P_{m}=\infty$, $5/6$, $2/3$, and
$5/10$. The disc density profile is set to be $s=-3/2$ (no wind), the ratio of
the specific heats is set to be $\gamma=1.3$, the viscous parameter is
$\alpha=0.1$, and the advection parameter is $f=1.0$.
In section 3.1, the upper limit of the magnetic field obtained and mentioned
with $\beta_{b}$. By substituting $\beta_{b}$ and equations (33)-(39) in
equation (53) and assume of $s=-3/2$, the mass accretion rate to the Bondi
accretion rate ($c_{7}$) approximately is
$\displaystyle c_{7}=\dot{M}/\dot{M}_{Bondi}\approx
24\sqrt{2}~{}\alpha~{}g_{1}~{}\frac{\left(1+\beta_{b}\right)^{1-\mu}}{\left(5+\beta_{b}\right)^{5/2}}$
$\displaystyle=24\sqrt{2}~{}\alpha~{}g_{1}~{}\frac{\left(1+\frac{18~{}g_{2}P_{m}}{f}\right)^{1-\mu}}{\left(5+\frac{18~{}g_{2}P_{m}}{f}\right)^{5/2}}$
(52)
where
$\displaystyle
g_{1}=\left[\frac{2}{5-3\gamma}\right]^{\frac{(3\gamma-5)}{2(\gamma-1)}}$
$\displaystyle g_{2}=\left[\frac{5/3-\gamma}{\gamma-1}\right].$
To obtain the root of derivative of $c_{7}$ in terms of $f$, we can write
$\displaystyle\frac{dc_{7}}{df}=0~{}~{}\Rightarrow~{}~{}f_{max}=\left(\frac{18}{5}\right)\left(\frac{3+2\mu}{1-2\mu}\right)g_{2}P_{m},$
The above equation states the mass accretion rate in $f_{max}$ becomes maximum
and for $f>f_{max}$ the mass accretion rate decreases by advection degree. The
$f_{max}$ for $\mu>1/2$ becomes negative, while $0\leq f\leq 1$. Thus, the
relation of $f_{max}$ is valid only for the case of $\mu=0$. In the case of
$\mu=0$, $f_{max}=\left(54/5\right)g_{2}P_{m}$ and $\beta_{b}=5/3$. Thus, we
expect in dominant magnetic case ($\beta>1$), the accretion efficiency
decreases by advection degree parameter. In the case of $\mu=1$, due to the
lack of any extremum, the mass accretion rate only increases by advection
degree and does not show any decrease.
In the case of high magnetic pressure ($\beta\gg 1$), equation (54) becomes
$\displaystyle c_{7}=\dot{M}/\dot{M}_{Bondi}\approx
24\sqrt{2}~{}\alpha~{}g_{1}~{}\left(\frac{f}{18~{}g_{2}P_{m}}\right)^{3/2+\mu}.$
The above equation implies that the mass accretion rate to Bondi accretion is
strongly depends on viscosity parameter, Prandtl number and the advection
degree. The mass accretion rate increases by $f$ and decreases by $P_{m}$.
Thus, accretion efficiency increases by the advection degree parameter in high
magnetic pressure.
### 3.3 Timescales
To estimate the effect of viscosity and resistivity on the accretion discs, we
compare the viscous and resistive timescales with accretion timescale. The
accretion timescale, $t_{acc}$, and the viscous timescale, $t_{visc}$, are
given by
$t_{acc}=\frac{r}{-v_{r}},$ $t_{visc}=\frac{r^{2}}{\nu}.$
We are using a similar functional form of $t_{visc}$ for the resistive
timescale, $t_{resis}$, that is given by
$t_{resis}=\frac{r^{2}}{\eta}.$
By using self-similar forms of physical quantities, we can write
$\displaystyle\frac{t_{resis}}{t_{acc}}=\frac{\alpha}{\eta_{0}}\frac{c_{1}}{c_{3}}(1+\beta)^{\mu-1}$
$\displaystyle=3(\frac{\alpha}{\eta_{0}})(s+2).$ (53)
The equation (30) is used for fraction of $c_{1}/c_{3}$. As, we said in
previous section, in present model $\alpha/\eta_{0}$ is the magnetic Prandtl
number, $P_{m}$, so above equation becomes
$\frac{t_{resis}}{t_{acc}}=3P_{m}(s+2).$
For our interesting case, $s=-3/2$ (no wind), we can write
$\frac{t_{resis}}{t_{acc}}=\frac{3}{2}P_{m}.$
The above equation implies that for $P_{m}\leq 2/3$, the magnetic diffusivity
timescale is shorter than or equal to accretion timescale, while for
$P_{m}>2/3$ the accretion timescale is shorter.
Similar calculations for the viscous timescale express
$\frac{t_{visc}}{t_{acc}}=3(s+2),$
where in no wind case ($s=-3/2$) becomes
$\frac{t_{visc}}{t_{acc}}=(3/2).$
Thus, the viscosity timescale will be longer than the accretion timescale. To
compare the magnetic diffusivity with the viscous timescales, we can write
$\displaystyle\frac{t_{resis}}{t_{visc}}=\frac{r^{2}/\eta}{r^{2}/\nu}$
$\displaystyle=\frac{\nu}{\eta}~{}~{}$ $\displaystyle=P_{m}.$ (54)
Thus, the magnetic Prandtl number specifies which one is shorter. For example
in flow with high conductivity (e.g. AF06; KF09), $\eta\rightarrow 0$, the
magnetic Prandtl number limits to infinity, and so the magnetic diffusivity
timescale will be very longer than the viscous timescale. On the other hand,
for a flow with finite resistivity and tiny viscosity (e.g. Sh04), the
magnetic Prandtl number limits to zero, and so the magnetic diffusivity
timescale is very shorter than viscous timescale. When the resistivity and the
viscosity are approximately equal, $P_{m}\sim 1$, we expect $t_{resis}\sim
t_{visc}$. Also, in special case of $P_{m}=5/6$ and $s=-3/2$ that escape and
creation of magnetic field are balanced and there is no mass-loss,
$t_{resis}=(5/6)t_{visc}$.
## 4 Summary and Discussion
In this paper, the influences of the resistivity on the structure of the
advection-dominated accretion flow is investigated. It is used only azimuthal
component of magnetic field that is consistent with observational evidence of
Galactic center (Novak et al 2003; Chuss et al. 2003; Yuan 2006). The
$\alpha$-prescription is used for the kinematic coefficient of viscosity and
the magnetic diffusivity. The equations of the model are solved by a semi-
analytical self-similar method in comparison with the self-similar solution by
AF06.
The physical quantities of disc are sensitive to the amounts of the magnetic
pressure fraction ($\beta$) and the magnetic diffusivity ($\eta_{0}$)
parameters. As, the angular velocity of the flow by adding the $\beta$ and
$\eta_{0}$ parameters decreases. For a value of the magnetic pressure
fraction, the angular velocity of disc becomes zero. This amount of the
magnetic pressure fraction strongly depends on the properties of the accreting
gas, such as the viscosity, resistivity, adiabatic index, and advection
degree. The solutions represent the radial infall velocity increases by adding
the magnetic diffusivity. Also the solutions show that the temperature of the
flow decrease by adding the toroidal component of magnetic field. This result
qualitatively is consistent with the results of Bu et al. (2009) and KF09. The
profiles of the temperature of the flow show that it increases by adding the
magnetic diffusivity that is due to the raise of the resistive dissipation.
Comparison of the present model with Bondi accretion implies that for all
values of the $\beta$ and $\eta_{0}$ parameters, the Bondi accretion rate is
larger than the mass accretion rate that is in accord with observational
evidences of Sgr A∗, M87, and NGC 4261 (Kaburaki 2007). Also, the mass
accretion rate profiles at high magnetic field express that the magnetic field
reduces the mass accretion rate that is similar to results of Kaburaki (2007).
We found that in the small magnetic field, the more heat generated in the flow
is due to the viscous dissipation, while the ohmic dissipation will be
dominant in large amounts of magnetic field and resistivity.
As noted in the introduction, the MHD simulations show that linear growth of
MRI decreases significantly by ohmic dissipation. linear growth of the MRI in
the resistive fluid can be characterized by the Lundquist number
($S_{MRI}=c_{A}^{2}/\eta\Omega$) and magnetic Reynolds number
($Re_{M}=c_{s}^{2}/\eta\Omega$), where $c_{A}$, $c_{s}$, $\eta$, and $\Omega$
have usual meaning. In terms of our self-similar transformations, the
Lundquist number and magnetic Reynolds number become
$S_{MRI}=2\beta/\eta_{0}c_{2}(1+\beta)^{1-\mu}$ and
$Re_{M}=1/\eta_{0}c_{2}(1+\beta)^{1-\mu}$. The solutions of present model show
that $S_{MRI}$ and $Re_{M}$ decrease by resistivity. This property is
qualitatively consistent with MHD simulation results (Fleming et al. 2000;
Masada & Sano 2008).
Here, latitudinal dependence of physical quantities is ignored, while some
authors showed that latitudinal dependence is important in structure
consideration of a disc (Narayan & Yi 1995; Sh04; Ghanbari et al. 2007). One
can investigate latitudinal behavior of such discs. Furthermore, in a
realistic model the advection parameter $f$ is a function of position, one can
consider such discs.
## Acknowledgements
I would like to thank the referee for very useful comments that helped me to
improve the initial version of the paper. I would also like to thank Markus
Flaig for helpful discussion.
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2024-09-04T02:49:24.847327
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{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kazem Faghei",
"submitter": "Kazem Faghei",
"url": "https://arxiv.org/abs/1111.7302"
}
|
1112.0045
|
CytoITMprobe: a network information flow plugin for Cytoscape
Aleksandar Stojmirović , Alexander Bliskovsky and Yi-Kuo Yu***to whom
correspondence should be addressed
National Center for Biotechnology Information
National Library of Medicine
National Institutes of Health
Bethesda, MD 20894
United States
#### Background:
Cytoscape is a well-developed flexible platform for visualization, integration
and analysis of network data. Apart from the sophisticated graph layout and
visualization routines, it hosts numerous user-developed plugins that
significantly extend its core functionality. Earlier, we developed a network
information flow framework and implemented it as a web application, called ITM
Probe. Given a context consisting of one or more user-selected nodes, ITM
Probe retrieves other network nodes most related to that context. It requires
neither user restriction to subnetwork of interest nor additional and possibly
noisy information. However, plugins for Cytoscape with these features do not
yet exist. To provide the Cytoscape users the possibility of integrating ITM
Probe into their workflows, we developed CytoITMprobe, a new Cytoscape plugin.
#### Findings:
CytoITMprobe maintains all the desirable features of ITM Probe and adds
additional flexibility not achievable through its web service version. It
provides access to ITM Probe either through a web server or locally. The
input, consisting of a Cytoscape network, together with the desired origins
and/or destinations of information and a dissipation coefficient, is specified
through a query form. The results are shown as a subnetwork of significant
nodes and several summary tables. Users can control the composition and
appearance of the subnetwork and interchange their ITM Probe results with
other software tools through tab-delimited files.
#### Conclusions:
The main strength of CytoITMprobe is its flexibility. It allows the user to
specify as input any Cytoscape network, rather than being restricted to the
pre-compiled protein-protein interaction networks available through the ITM
Probe web service. Users may supply their own edge weights and
directionalities. Consequently, as opposed to ITM Probe web service,
CytoITMprobe can be applied to many other domains of network-based research
beyond protein-networks. It also enables seamless integration of ITM Probe
results with other Cytoscape plugins having complementary functionality for
data analysis.
## Background
Cytoscape [1, 2, 3] is a popular and flexible platform for visualization,
integration and analysis of network data. Apart from the sophisticated graph
layout and visualization routines, its main strength is in providing an API
that allows developers other than its core authors to produce extension
plugins. Over the last decade, a large number of plugins have been released,
supporting the features such as import and export of data, network analysis,
scripting and functional enrichment analysis. In this paper, we describe
CytoITMprobe a plugin that brings to Cytoscape new functionality founded on
information flow.
Numerous approaches for analyzing biological networks based on information
flow [4, 5, 6, 7, 8, 9, 10] have emerged in recent years. The main assumption
of all such methods is information transitivity: information can flow through
or can be exchanged via paths of biological interactions. Our contribution to
this area [11, 12] is a context-specific framework based on discrete-time
random walks (or equivalently, diffusion) over weighted directed graphs. In
contrast to most other approaches, our framework explicitly accommodates
directed networks as well as the information loss and leakage that generally
occurs in all networks. Apart from the network itself and a user-specified
context, it requires no prior restriction to the sub-network of interest nor
additional and possibly noisy information. We implemented our framework as an
application called ITM Probe [13] and made it available as a web service [14],
where users can query protein-protein interaction (PPI) networks from several
model organisms and visualize the results.
In addition to implementing network flow algorithms, the ITM Probe web service
possesses a number of useful features. Using the Graphviz [15] suite for
layout and visualization of graphs, it displays in the user’s web browser the
images of subnetworks consisting of nodes identified as significant by the
information flow models and offers a choice of multiple coloring schemes. The
entire query results can be retrieved in the CSV format or forwarded to a
functional enrichment analysis tool to facilitate their interpretation.
However, lacking a mechanism to decouple the algorithmic part from the
interaction graph, the ITM Probe web service restricts users to querying only
the few compiled PPI networks available on the website. Using a canned suite
for graph layout, ITM Probe limits the users’ ability to manipulate network
images. For example, the only way to change the layout of significant
subnetworks is to choose a different seed and re-compute the layout. Most
importantly, not having an adequate interface to a well-designed platform such
as Cytoscape, it is difficult to use the results of the ITM Probe service
within the workflows involving additional data and algorithms from other
sources. We thus developed CytoITMprobe to meet these challenges by (1)
providing an explicit decoupling between the algorithmic part and the
interaction graph, (2) utilizing the core graph manipulation functionality of
Cytoscape for a broader visualization choices, and (3) adding an appropriate
input/output interface for seamless integration with other resources available
in Cytoscape.
Figure 1: ITM Probe is based on discrete-time random walks with boundary nodes
and damping. As an example, consider the weighted directed network shown,
containing 19 nodes and 44 links. Single-directional links are assigned weight
2 and are indicated using arrows while bi-directional edges are assigned
weight 1 and are shown as lines. The first five graphs show the time progress
of a random walk in the presence of damping and two absorbing boundary nodes
(indicated by octagons). At $t=0$, 1000 random walkers start at a single point
in the network. At $t=1$, they have progressed one step from their origin to
the nodes adjacent to it, being distributed randomly in proportion to the
weights of the edges leading from the origin. Only 900 walkers remain in the
network at $t=1$ due to damping: the damping factor $\mu=0.9$ (dissipation
$0.1$) means that $10\%$ of walkers are dissipated at each step. At $t=60$,
most of the walks have terminated, either by dissipation, or by reaching one
of the two boundary nodes. The number of walkers terminating at each boundary
node depends on their starting location. The final graph shows the probability
$F_{ik}$ for a random walk starting at any transient node in the network
(indicated by circular shape) to terminate at the boundary node on the right-
hand side (scaled by 1000). Note that the value indicated in the final graph
for the starting node at $t=0$ (190) is the same as the final number of walks
shown at $t=60$ as terminating at the right boundary node.
## Information Flow Framework
ITM Probe extracts _context-specific_ information from networks. We elaborated
on the information flow framework underlying ITM Probe in our previous
publications [11, 12] and here we provide a non-technical explanation. Given a
context consisting of one or more user-selected network nodes, the aim is to
retrieve a set of other network nodes most related to that context. We model
networks as weighted directed graphs, where nodes are linked by directional
edges and each edge is assigned a positive weight. One can consider a random
walker that wanders among network nodes in discrete steps. The rule of the
walk is that the walker starts at a certain node and in each step moves
randomly to some adjacent node with probability proportional to the weight of
the edge linking these nodes (Fig. 1). If the graph is connected, that is, if
there is a directed path linking any two nodes, such a walk never terminates
and the walker will eventually visit every node in the graph.
Our main idea is to set termination or _boundary_ nodes for the walkers while
using random walks to explore the neighborhoods of the context nodes. Provided
there is a directed path linking any node to a boundary node, every random
walk here will eventually terminate. Furthermore, the nodes visited by a
walker before termination will vary depending on the origin of the walk. Since
a random walk is a stochastic process, and each walk is different, we are
interested in the cumulative behavior of infinitely many walkers following the
same rules. On average, we expect that the nodes more relevant to the context
will be more visited than those that are less relevant. Thus, the main
quantity of interest is the average number of visits to a node given the
selected origins and destinations of the walk.
A problem with the above approach is that random walkers may spend too much
time in the graph if the origins and destinations of the walk are far apart.
This could mean that the entire graph is visited so that the most significant
nodes are just those with the largest degree. To ensure that the significant
nodes are relatively close to the context nodes, our framework contains an
additional ingredient, _damping_ : at each step of a walk, we assign a certain
probability for the walker to dissipate, that is, to leave the network. We
still evaluate the average number of visits to each node, but now only count
the visits prior to the walker leaving the network. Evidently, the nodes that
are close to the walker’s origin will be significantly visited. In addition to
forcing locality, damping is also natural in physical or biological contexts.
If we treat random walkers as information propagating through the network, it
is natural to assume that some information is lost during transmission. For
protein-protein interaction networks, where nodes are proteins and links are
physical bindings between proteins, damping could be associated with protein
degradation by proteases, which would diminish the strength of information
propagation.
ITM Probe framework contains three models: _emitting_ , _absorbing_ and
_channel_. In the absorbing model (Fig. 1), the context nodes are interpreted
as destinations or _sinks_ of random walks, while every non-boundary or
_transient_ node is considered as a potential origin. For each transient node
$i$ and each sink $k$, the model computes $F_{ik}$, the average number of
visits to the terminating node $k$ by random walks originating at the node
$i$. Since a walk can either terminate at one sink or the other, $F_{ik}$ can
also be interpreted as the probability that a random walk from $i$ reaches
$k$. In the absence of damping, the sum of $F_{ik}$ over all sinks will be
exactly $1$ for any transient node $i$. However, in the presence of damping,
the sum of $F_{ik}$ over all sinks may be much less than $1$ (Fig. 1). The
emitting model (Fig. 2), offers a dual point of view. Here, the context nodes
are interpreted as origins or _sources_ of random walks. The walks terminate
by dissipating or by returning to the sources – the sources form an emitting
boundary. Since the origins of the walks are fixed, the quantity of interest
is the visits to the transient nodes. Specifically, for each source $s$ and
each transient node $i$, the emitting model returns $H_{si}$, the average
number of visits to $i$ by walkers originating at $s$.
Figure 2: The emitting model counts visits from sources. Using the example
network from Fig. 1 with the same damping factor, consider the case where 1000
random walkers start at the source node indicated by a hexagon. At each time
step, some random walkers leave the network due to damping or by moving back
to the source. In the first five graphs, the number in each node documents the
total number of visits to that node from all random walkers, dissipated or
not, up to the indicated time. The value of $H_{si}$ returned by the ITM Probe
emitting mode ($s$ here denotes the source node) yields the expected number of
visits to node $i$ per random walker that starts at $s$ over infinitely many
time steps. The final graph shows the values of $H_{si}$ for this context,
scaled by 1000. Note that the magnitude shown for one transient node is
greater than 1000 because a walker may visit the same node multiple times.
The values of $F_{ik}$ and $H_{si}$ can be efficiently computed by solving
(sparse) systems of linear equations. Let $W_{ij}$ denote the weight of the
directed link $i\to j$ and let $0<\mu<1$ denote the damping factor. For all
pairs of nodes $i,j$, construct the random walk evolution operator
$\mathbf{P}$, where $P_{ij}=\frac{\mu W_{ij}}{\sum_{j^{\prime}}W_{ij}}$. The
operator $\mathbf{P}$ includes damping and hence $\sum_{j}P_{ij}<1$. Let
$\mathbf{P}_{TT}$ denote the sub-operator of $\mathbf{P}$ with domain and
range restricted only to transient nodes and let
$\mathbf{G}=(\mathbb{I}-\mathbf{P}_{TT})^{-1}$, where $\mathbb{I}$ stands for
the identity matrix. Then, it can be shown [11], that
$\displaystyle F_{ik}$ $\displaystyle=\sum_{j}G_{ij}P_{jk},\qquad\text{and}$
$\displaystyle H_{si}$ $\displaystyle=\sum_{j}P_{sj}G_{ji}.$
More details, including the cases where $\mu=0$, $\mu=1$ or non-uniform
damping are covered in [11, 12].
Figure 3: The channel model highlights the directed flow from origins to
destinations. Consider once again the example network from Figs. 1 and 2, now
with a single source (hexagon) and two sinks (octagons). In common with the
case from Fig. 2, the walkers start at the source, but in this case can
terminate only by reaching the sinks. The damping factor is implicit: it
determines how far the walkers are allowed to deviate from the shortest path
towards one of the sinks. In the first five graphs, the number in each
transient node documents the total number of visits to that node from all
random walkers up to the indicated time. However, the value in each sink node
represents the likelihood to reach that sink from the source at the indicated
time. The value of $\hat{\varPhi}_{i,K}^{s}$ returned by the ITM Probe
normalized channel mode yields the expected number of visits to node $i$ per
random walker that starts at $s$ over infinitely many time steps. Note that
the sink nodes split the flow from the source depending on their location. In
this example, over infinitely many time steps, the node closer to the source
captures 970 walkers, while the further sink gets only the remaining 30.
The channel model combines the emitting and the absorbing model, with both
sources and sinks on the boundary. It illuminates the most likely paths from
sources to sinks. For each source node $s$, transient node $i$ and sink node
$k$, it computes $\varPhi_{i,k}^{s}=H_{si}F_{ik}$, the average number of
visits to $i$ by a random walker that originates at $s$ and terminates at $k$.
ITM Probe does not report $\varPhi_{i,k}^{s}$ directly, but instead shows a
simpler, _normalized_ quantity $\hat{\varPhi}_{i,K}^{s}$ (Fig. 3), which is
defined for each source $s$ and transient node $i$ by
$\hat{\varPhi}_{i,K}^{s}=\frac{\sum_{k}H_{si}F_{ik}}{\sum_{k^{\prime}}F_{sk^{\prime}}}.$
(1)
Here, the numerator $\sum_{k}H_{si}F_{ik}=\sum_{k}\varPhi_{i,k}^{s}$ gives the
average number of visits, in the presence of damping, to $i$ by a random
walker starting at $s$ and terminating at any sink. The denominator gives the
total probability of a walker starting at $s$ to terminate at any sink. Hence,
with the denominator off-setting the effect of damping, the value of
$\hat{\varPhi}_{i,K}^{s}$ counts the average number of visits to $i$ by
walkers that start at $s$ and terminate at any of the sinks as if no
dissipation is present. Generally, damping in the emitting or the absorbing
model determines how far the flow can reach away from its origins. In
contrast, the damping parameter for the normalized channel model plays a
different role (Fig. 4): it effectively determines the ‘width’ of the channel
from sources to sinks. When damping is very strong, only the nodes on the
shortest path from a source to its nearest sink will be visited.
Given the close relationship between random walks and diffusion, it is also
possible to interpret ITM Probe models through information diffusion (or
information flow). Within that paradigm, a fixed amount information is
constantly replenished at the source nodes while leaving the network at all
boundary nodes and through dissipation. At equilibrium, when the rate of flow
entering equals the rate of leaving, the amount of information occupying each
transient node is equivalent to the average number of visits to that node
(using the aforementioned non-replenishing random walk interpretation [11]).
We call the set of nodes most influenced by the flow an _Information
Transduction Module_ or ITM.
Figure 4: An example of the results of running different ITM Probe models.
Here we see the results of running the emitting (a,d,g), absorbing (b,e,h) and
channel (c,f,i) model of ITM Probe with the same sources and sinks but
different dissipation coefficients. The underlying undirected graph is derived
from a square lattice by removing random nodes and edges. Sources are shown as
hexagons, sinks as octagons, and transient nodes as squares. The top row
(a,b,c) shows the runs with damping factor $\mu=0.95$ (dissipation $0.05$),
the middle (d,e,f) with $\mu=0.75$ and the bottom with $\mu=0.25$. For the
emitting and channel model, each basic cyan, magenta or yellow color is
associated with a source. The coloring of each node arises by mixing the basic
color in proportion to the strength of information flow from their respective
sources. For the absorbing model, the nodes are shaded according to the total
probability of absorption at any sink on a logarithmic scale.
## Software architecture
CytoITMprobe architecture consists of two parts: the user interface front end
and computational back end. The user interface, written in Java [16] using
Cytoscape API, is accessed as a Cytoscape plugin. It consists of the query
form, the results viewer and the ITM subnetwork (Fig. 5). The back end is the
standalone ITM Probe program, written in Python, which can be installed
locally or accessed through a web service. In either configuration,
CytoITMprobe takes the user input through the graphical user interface,
validates it, and passes a query to the back end. Upon receiving from the back
end the entire query results, CytoITMprobe stores them as the node and network
attributes of the original network. Consequently, the query output can be
edited or manipulated within Cytoscape, as well as saved for later use.
Figure 5: CytoITMprobe interface. At startup from the Plugins menu,
CytoITMprobe embeds its query form into the Control Panel (left). After
performing a query or loading previously obtained search results, it creates
an ITM subnetwork showing significant nodes and a viewer embedded into Results
Panel (right). The viewer allows closer examination of the results and
manipulation of the contents and the look of the ITM subnetwork. The overall
visual styling of CytoITMprobe components closely resembles that of the ITM
Probe web version.
Standalone ITM Probe is a part of the qmbpmn-tools Python package, which also
contains the code supporting the ITM Probe and SaddleSum web services, as well
as the scripts for constructing the underlying datasets. The ITM Probe part
depends on Numpy and Scipy [17] packages for numerical computations. The
performance of ITM Probe critically depends on the routines for computing
direct solutions of large, sparse, nonsymmetric systems of linear equations.
Scipy supports two sparse direct solver libraries (both written in C): SuperLU
[18] as default and UMFPACK [19] as an optional add on through SciKits
collection [20]. In our experience, UMFPACK runs faster than SuperLU and Scipy
always uses it if available. However, for optimal performance, UMFPACK
requires well-tuned Basic Linear Algebra Subroutines (BLAS) libraries and may
not be easy to install. To support users who are inclined not to install
UMFPACK or Scipy, CytoITMprobe supports remote queries by default.
## Input
CytoITMprobe requires as input a weighted directed graph and the ITM Probe
model parameters that include a selection of boundary nodes and a dissipation
probability.
### Step one: defining a query graph
In CytoITMprobe graph connectivity is specified by selecting a Cytoscape
network. In addition, each link must be assigned a weight and a direction
through the query form. Edge weights are set using the _Weight attribute_
dropdown box, which lists all available floating-point edge attributes of the
selected network and the default option (_NONE_). If the default option is
selected, CytoITMprobe assumes a weight $2$ for any self-pointing edge and $1$
for all other edges. If an attribute is selected, the weight of an edge is set
to the value of the selected attribute for that edge. Null attribute values
are treated as zero weights.
Since Cytoscape edges are always internally treated as directed, the user must
also indicate the directedness of each edge type through the query form.
Whenever a new Cytoscape network is selected, CytoITMprobe updates the query
form and places all of the network’s edge types into the _undirected_
category. The user can use arrow buttons to move some edge types to the
_directed_ or _ignored_ category. Undirected edges are treated as
bidirectional, with the same weight in both directions. Directed edges have a
specified weight assigned only in the forward direction, with the backward
direction receiving the zero weight. Ignored edges have zero weight in both
directions. Since Cytoscape allows multiple edges of different types between
the same nodes, CytoITMprobe collapses multiple edges in each direction into a
single edge by appropriately summing their weights (Fig. 6).
Figure 6: Edge weights example Consider the following example: Suppose A and B
are nodes in a Cytoscape network linked by three edges of two types with shown
edge weights. Assume two type I edges (lighter gray), $A\to B$ and $B\to A$
are directed, while a single type II edge (darker gray) $A\to B$ is
undirected. At query time, CytoITMprobe creates two directed edges, $A\to B$
and $B\to A$, with weights $3$ and $6$, respectively.
### Step two: selecting a model and boundary nodes
In addition to a weighted directed graph, ITM Probe requires an information
flow model (emitting, absorbing or normalized channel), a selection of sources
and/or sinks, and dissipation probability. The choice of the model determines
the types of boundary nodes that need to be specified, as well as the ways in
which the damping factor can be set (see ‘Step three: specifying dissipation
probability’ below). The query form also allows users to specify _excluded
nodes_. Any flow reaching excluded nodes is fully dissipated. This is a way to
remove those nodes that do not participate in information propagation in the
desired context or that introduce undesirable shortcuts.
### Step three: specifying dissipation probability
The values of $H$, $F$, and $\hat{\varPhi}_{,}$ all implicitly depend on the
dissipation probability. In ITM Probe the user can set the dissipation
probability directly or specify a related quantity that can, using Newton’s
method, determine the dissipation probability. The choice of the alternative
quantity depends on the selected model. For the emitting model, this quantity
is the average path length before termination, which we denote by $\bar{t}$.
For example, the user can require a random walker to make on average three
steps before terminating. The formula for $\bar{t}$ is
$\bar{t}=1+\frac{1}{n_{S}}\sum_{s}\sum_{j}H_{sj},$ (2)
where $n_{S}$ denotes the number of sources. For the normalized channel model,
the path length before termination is given by
$\bar{t}=1+\frac{1}{n_{S}}\sum_{s}\sum_{j}\hat{\varPhi}_{j,K}^{s}.$ (3)
Since the normalized channel model counts only the random walkers actually
terminating at sinks, $\bar{t}$ is in this case bounded below by the length of
the shortest path from any source to any sink. Hence, ITM Probe accepts the
desired value of $\bar{t}$ in terms of length deviation from the shortest
path. There are two ways to set the average path-length deviation: in absolute
units (steps) or as a proportion of the length of the shortest path. The
absorbing model allows users to obtain the dissipation probability by setting
the average absorption probability, denoted $\bar{r}$. The formula for
$\bar{r}$ is
$\bar{r}=\frac{1}{n_{T}}\sum_{i}\sum_{k}F_{ik},$ (4)
where $k$ ranges over all sinks, $i$ ranges over all transient nodes _that are
connected to at least one sink_ , and $n_{T}$ is the total number of such
nodes. The value of $\bar{r}$ represents the likelihood of a random walk
starting at a randomly selected point in the network to reach a sink. The
dissipation probability obtained in this way is larger if the sinks are well-
connected hubs near the center of the network, in contrast to the case when
the chosen sinks are not as well connected.
### Step four: submitting a query
After specifying all necessary input, the user submits a query by pressing the
_QUERY_ button on the query form. The time required for a run depends on
whether the query is local or remote, as well as on the size of the submitted
graph and the number of selected sources and/or sinks.
## Output
For every completed query, CytoITMprobe displays its results in a viewer
embedded in Cytoscape Results Panel and a new Cytoscape network consisting of
significant nodes (ITM subnetwork). The results viewer has five tabs: _Top
Scoring Nodes_ , _Summary_ , _Input Parameters_ , _Excluded Nodes_ , and
_Display Options_. The first four tabs contain information about the query and
the results, while the last one contains a form that allows users to
manipulate the ITM subnetwork. The form controls two aspects of the
subnetwork: composition (what nodes are selected and how many) and node
coloring.
### Displaying significant nodes
Subnetwork nodes are selected through a _ranking attribute_ , which assigns a
numerical value from ITM Probe results to each node. The nodes are listed in
descending order of the ranking attribute and top nodes are displayed as the
ITM subnetwork. The number of top nodes is determined by specifying a
_selection criterion_ , which can be simply a number of nodes to show, a
cutoff value or the ‘participation ratio’. Specifying a cutoff value $x$
selects the nodes with their ranking attribute greater than $x$. Participation
ratio estimates the number of ‘significant’ nodes by considering all values of
the ranking attribute in a scale-independent manner [11]. The available
choices for the ranking attribute depend on the ITM Probe model and the number
of boundary points. For the emitting and normalized channel model, the user
can select visits to a node from each source or the sum of visits from all
sources. It is also possible to use _interference_ [11], which denotes the
minimum number of node visits, taken over all sources. For the absorbing
model, the available attributes are absorbing probabilities to each sink and
the total probability of termination at a sink. The values of all attributes
for the subnetwork nodes are displayed in the _Top Scoring Nodes_ tab.
The colors of the subnetwork nodes are determined by selecting _coloring
attributes_ , a _scaling function_ and a _color map_. The list of coloring
attributes is the same as the list of ranking attributes but the user can
select up to three coloring attributes. If a single attribute is selected,
node colors are determined by the selected eight-category ColorBrewer [21]
color map. Otherwise, they are resolved by color mixing: each coloring
attribute is assigned a single basic color (cyan, magenta or yellow), and the
final node color is obtained by mixing the three basic colors in proportion to
the values of their associated attributes at that node. The scaling function
serves to scale and discretize the coloring attributes to the ranges
appropriate for color maps. Figure 4 shows examples of mixed color scheme with
three boundary points (left and right columns) and of a coloring using a
single attribute (center column).
### Manipulating node attributes
Since the ITM Probe query results are saved as Cytoscape attributes of the
original network, they can be arbitrarily modified through Cytoscape. Any
changes made are reflected in the results viewer and the corresponding ITM
subnetwork after pressing the _RESET_ button on the Display Options form.
Using the CytoITMprobe attribute nomenclature, users can create additional
attributes to be used for ranking or coloring. Consider the following usage
example. A user has run an emitting model query with three sources, S1, S2,
and S3, and obtained the results in a viewer labeled ITME243. At the end of
the run, CytoITMprobe created the attributes ITME243[S1], ITME243[S2] and
ITME243[S3] for the nodes of the input network and saved the results as their
values. The user creates a new floating-point node attribute with a label
ITME243[avgS1S2] and fills it with an average of ITME243[S1] and ITME243[S2].
After resetting the Display Options form, an item ‘Custom [avgS1S2]’ is
available for selection as a ranking or coloring attribute. This gives the
user the flexibility to reinterpret S1 and S2 as if they were a single source
of equal weight as S3. Another possibility is to combine the results of
queries with different boundaries and display them together on the same
subnetwork.
### Saving and restoring results
The query network together with its attributes containing ITM Probe results
can be saved as a Cytoscape session and later retrieved. After reloading the
session, the user can regenerate the results viewer and the corresponding
subnetwork for a stored ITM by pressing the _LOAD_ button on the CytoITMprobe
query form and selecting the desired ITM from a list.
Alternatively, the ITM Probe results can be exported to tab-delimited text
files through the Cytoscape _Export_ menu. Each exported tab-delimited file
contains all the information necessary to restore the results except the query
network and can be easily manipulated both by humans and by external programs
or scripts. The results from tab-delimited files can be imported into any
selected Cytoscape network through the _Import_ menu. Since the selected
network may be different from the original query network, only the results for
the nodes in the selected network whose IDs match the IDs from the imported
file will be loaded. After importing the results, CytoITMprobe generates a new
results viewer and a subnetwork, as if the results originated from a direct
ITM Probe query.
## Discussion
The main function of ITM Probe, also applicable to domains other than PPI
networks, is to retrieve information from large and complex networks by
discovering the possible interface between network nodes that are hypothesized
to be related. This paradigm works best with large networks, where such
information cannot be easily accessed by other means. For examples of
applications of the ITM Probe frameworks to protein-protein interaction
networks, consult our earlier papers [11, 13, 12].
With a network as an _encyclopedia_ of domain-specific knowledge, ITM Probe
enables a direct access to its specific portions related to a specified
context. The user can learn about the objects representing individual nodes by
setting them as sources and/or sinks and retrieving information about the most
significant objects in the resulting ITM. This approach not only extracts a
relevant subnetwork but also produces context-specific weights for each node.
With their interpretation as average numbers of node visits, or equivalently,
as average numbers of paths passing through a node, the ITM weights signify
the relative importance of network nodes in the context of the query and thus
can be used to refine its interpretation as a whole. For example, a single
node with a large weight in an ITM resulting from a normalized channel model
query represents a choke point _in the particular context of the query_. The
same node need not have a high global centrality.
Containing both sources and sinks, the normalized channel model offers the
users the ability to formulate and evaluate network based hypotheses in
silico. Since information flow that reaches one sink cannot subsequently
terminate at any other, sink nodes can be associated with alternative
hypotheses, such as different biological functions if the network is PPI. The
information flow from each source will then, depending on the dissipation
coefficient used, mainly trace the path towards the sink most likely to be
reached first from that source (see Fig. 4, right column). The ITM Probe
framework considers all weighted paths from sources to sinks and hence
produces more robust results than approaches involving only the shortest
paths. The path weights are tunable using the dissipation probability.
Compared to the previously described web interface to ITM Probe [13],
CytoITMprobe significantly benefits from being a part of the Cytoscape
platform. Although the _Display Options_ form is very similar to the web
version, the sophisticated network visualization functionality provided by
Cytoscape allows significantly more versatility in displaying ITMs. For
example, Cytoscape GUI allows users to manually alter node placements, rotate
network views, or arbitrarily change the look of a network. In addition,
Cytoscape interface enables users to directly manipulate node attributes
representing ITM Probe results and possibly create new node summary variables
appropriate to their problem. The newly created variables can be immediately
reflected in the graphical representation of an ITM, which is not possible in
the web setting. Most importantly, the results of ITM Probe can be integrated
into workflows involving other Cytoscape plugins that provide complementary
functionality. For instance, output ITMs can be related to terms from
controlled vocabularies such as Gene Ontology [22] using functional enrichment
analysis plugins such as PinGO [23] or our own recently released CytoSaddleSum
[24]. The graph-theoretic structure of ITM subnetworks can be analyzed using a
variety of algorithms such as MCODE [25] or GraphletCounter [26, 27].
The architecture of CytoITMprobe with a Cytoscape front end and an ITM Probe
back end offers flexibility for a variety of usage scenarios. In contrast to
the web version, it allows users to use ITM Probe with arbitrary networks and
edge weights, rather than being limited to compiled PPIs from few model
organisms. Most users will be content with accessing ITM Probe through the web
server. However, the option to download and install the qmbpmn-tools package
provides not only faster running times for queries but also the ability to use
the command line interface for ITM Probe to perform batch queries and to
locally reproduce its web service. The separation of the presentation layers
(web or Cytoscape) from the ‘business’ layer (standalone ITM Probe)
facilitates easy future updates to any components.
## Conclusion
CytoITMprobe is a plugin that brings the previously unavailable network flow
algorithms of ITM Probe to the Cytoscape platform. It enables users to extract
context-specific subnetworks from large networks by specifying the origins
and/or destinations of information flow. CytoITMprobe significantly extends
the features of the previously released web version of ITM Probe.
The main novelty of CytoITMprobe is that it allows the user to specify as
input any Cytoscape network, rather than being restricted to the PPI networks
available through the ITM Probe web service. Using Cytoscape attributes to
hold their desired values, users may easily supply their own edge weights and
denote edge directionalities. Additionally, the ability to manipulate and add
new node attributes through Cytoscape reduces the workload required for
visualizing various combinations of ITM components. In the context of
biological cellular networks, this additional flexibility may lead to
constructions of new node attributes that can better reflect biological
significance, hence facilitating more educated hypothesis forming.
By bringing ITM Probe to Cytoscape, CytoITMprobe enables seamless integration
of ITM Probe results with other Cytoscape plugins having complementary
functionality for data analysis. By decoupling the query network from the
information flow algorithm, the newly developed CytoITMprobe can be applied to
many other domains of network-based research beyond protein-networks.
## Availability and requirements
### CytoITMprobe plugin
Project name: CytoITMprobe
Project home page:
http://www.ncbi.nlm.nih.gov/CBBresearch/Yu/downloads/itmprobe.html
Documentation:
http://www.ncbi.nlm.nih.gov/CBBresearch/Yu/mn/itm_probe/doc/cytoitmprobe.html
Video tutorial: http://www.youtube.com/watch?v=4Cdf-mSKtWo
Operating system(s): Platform independent
Programming language: Java
Other requirements: Java SE 6 or higher and Cytoscape 2.7 or higher
License: All components written by the authors at the NCBI are released into
Public Domain. Components included from elsewhere are available under their
own open source licenses and attributed in the source code.
### Standalone ITM Probe (optional for CytoITMprobe)
Project name: qmbpmn-tools
Project home page:
http://www.ncbi.nlm.nih.gov/CBBresearch/Yu/downloads/itmprobe.html
Documentation: http://www.ncbi.nlm.nih.gov/CBBresearch/Yu/mn/itm_probe/doc/
Operating system(s): Platform independent
Programming language: Python
Other requirements: Python 2.6 or 2.7, Numpy 1.3 or higher and Scipy 0.7 or
higher. UMFPACK Scikit is recommended for good performance.
License: All components written by the authors at the NCBI are released into
Public Domain. Components included from elsewhere are available under their
own open source licenses and attributed in the source code.
## Acknowledgments
This work was supported by the Intramural Research Program of the National
Library of Medicine at the National Institutes of Health.
## References
* [1] Cline MS, Smoot M, Cerami E, Kuchinsky A, Landys N, Workman C, Christmas R, Avila-Campilo I, Creech M, Gross B, Hanspers K, Isserlin R, Kelley R, Killcoyne S, Lotia S, Maere S, Morris J, Ono K, Pavlovic V, Pico AR, Vailaya A, Wang PL, Adler A, Conklin BR, Hood L, Kuiper M, Sander C, Schmulevich I, Schwikowski B, Warner GJ, Ideker T, Bader GD: Integration of biological networks and gene expression data using Cytoscape. _Nat Protoc_ 2007, 2(10):2366–82.
* [2] Shannon P, Markiel A, Ozier O, Baliga NS, Wang JT, Ramage D, Amin N, Schwikowski B, Ideker T: Cytoscape: a software environment for integrated models of biomolecular interaction networks. _Genome Res_ 2003, 13(11):2498–504.
* [3] Smoot ME, Ono K, Ruscheinski J, Wang PL, Ideker T: Cytoscape 2.8: new features for data integration and network visualization. _Bioinformatics_ 2011, 27(3):431–2.
* [4] Nabieva E, Jim K, Agarwal A, Chazelle B, Singh M: Whole-proteome prediction of protein function via graph-theoretic analysis of interaction maps. _Bioinformatics_ 2005, 21 Suppl 1:302–310.
* [5] Tu Z, Wang L, Arbeitman M, Chen T, Sun F: An integrative approach for causal gene identification and gene regulatory pathway inference. _Bioinformatics_ 2006, 22:e489–496.
* [6] Suthram S, Beyer A, Karp R, Eldar Y, Ideker T: eQED: an efficient method for interpreting eQTL associations using protein networks. _Mol. Syst. Biol._ 2008, 4:162.
* [7] Zotenko E, Mestre J, O’Leary DP, Przytycka TM: Why do hubs in the yeast protein interaction network tend to be essential: reexamining the connection between the network topology and essentiality. _PLoS Comput Biol_ 2008, 4(8):e1000140.
* [8] Missiuro P, Liu K, Zou L, Ross B, Zhao G, Liu J, Ge H: Information flow analysis of interactome networks. _PLoS Comput Biol_ 2009, 5(4):e1000350.
* [9] Voevodski K, Teng S, Xia Y: Spectral affinity in protein networks. _BMC Syst Biol_ 2009, 3:112.
* [10] Kim YA, Przytycki JH, Wuchty S, Przytycka TM: Modeling information flow in biological networks. _Phys Biol_ 2011, 8(3):035012.
* [11] Stojmirović A, Yu YK: Information flow in interaction networks. _J Comput Biol_ 2007, 14(8):1115–43.
* [12] Stojmirović A, Yu YK: Information flow in interaction networks II: channels, path lengths and potentials. _J Comput Biol_ 2012\. in press.
* [13] Stojmirović A, Yu YK: ITM Probe: analyzing information flow in protein networks. _Bioinformatics_ 2009, 25(18):2447–9.
* [14] ITM Probe Web Service [http://www.ncbi.nlm.nih.gov/CBBresearch/Yu/mn/itm_probe].
* [15] Gansner ER, North SC: An open graph visualization system and its applications to software engineering. _Software — Practice and Experience_ 2000, 30(11):1203–1233.
* [16] Java [http://www.java.com].
* [17] Jones E, Oliphant T, Peterson P, et al.: SciPy: Open source scientific tools for Python 2001–. [http://www.scipy.org/].
* [18] Demmel JW, Eisenstat SC, Gilbert JR, Li XS, Liu JWH: A supernodal approach to sparse partial pivoting. _SIAM J. Matrix Analysis and Applications_ 1999, 20(3):720–755.
* [19] Davis TA: Algorithm 832: UMFPACK V4.3—an unsymmetric-pattern multifrontal method. _ACM Trans. Math. Softw._ 2004, 30(2).
* [20] SciKits [http://scikits.appspot.com/].
* [21] Harrower M, Brewer C: ColorBrewer.org: An Online Tool for Selecting Colour Schemes for Maps. _Cartogr J_ 2003, 40:27–37.
* [22] Ashburner M, Ball CA, Blake JA, Botstein D, Butler H, Cherry JM, Davis AP, Dolinski K, Dwight SS, Eppig JT, Harris MA, Hill DP, Issel-Tarver L, Kasarskis A, Lewis S, Matese JC, Richardson JE, Ringwald M, Rubin GM, Sherlock G: Gene ontology: tool for the unification of biology. The Gene Ontology Consortium. _Nat Genet_ 2000, 25:25–29.
* [23] Smoot M, Ono K, Ideker T, Maere S: PiNGO: a Cytoscape plugin to find candidate genes in biological networks. _Bioinformatics_ 2011, 27(7):1030–1.
* [24] Stojmirovic A, Bliskovsky A, Yu YK: CytoSaddleSum: a functional enrichment analysis plugin for Cytoscape based on sum-of-weights scores. _Bioinformatics_ 2012\. [Doi://10.1093/bioinformatics/bts041].
* [25] Bader GD, Hogue CWV: An automated method for finding molecular complexes in large protein interaction networks. _BMC Bioinformatics_ 2003, 4:2.
* [26] Whelan C, Sönmez K: Computing graphlet signatures of network nodes and motifs in Cytoscape with GraphletCounter. _Bioinformatics_ 2012, 28(2):290–1.
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|
arxiv-papers
| 2011-11-30T22:10:50 |
2024-09-04T02:49:24.861303
|
{
"license": "Public Domain",
"authors": "Aleksandar Stojmirovi\\'c, Alexander Bliskovsky and Yi-Kuo Yu",
"submitter": "Aleksandar Stojmirovi\\'c",
"url": "https://arxiv.org/abs/1112.0045"
}
|
1112.0051
|
# QCD measurements at the Tevatron
Dmitry Bandurin
(for the D0 and CDF Collaborations)
Florida State University Department of Physics
Tallahassee FL 32306 USA
###### Abstract
Selected quantum chromodynamics (QCD) measurements performed at the Fermilab
Run II Tevatron $p\bar{p}$ collider running at $\sqrt{s}=1.96$ TeV by CDF and
D0 Collaborations are presented. The inclusive jet, dijet production and
three-jet cross section measurements are used to test perturbative QCD
calculations, constrain parton distribution function (PDF) determinations, and
extract a precise value of the strong coupling constant,
$\alpha_{s}(m_{Z})=0.1161^{+0.0041}_{-0.0048}$. Inclusive photon production
cross-section measurements reveal an inability of next-to-leading-order (NLO)
perturbative QCD (pQCD) calculations to describe low-energy photons arising
directly in the hard scatter. The diphoton production cross-sections check the
validity of the NLO pQCD predictions, soft-gluon resummation methods
implemented in theoretical calculations, and contributions from the parton-to-
photon fragmentation diagrams. Events with $W/Z$+jets productions are used to
measure many kinematic distributions allowing extensive tests and tunes of
predictions from pQCD NLO and Monte-Carlo (MC) event generators. The charged-
particle transverse momenta ($p_{T}$) and multiplicity distributions in the
inclusive minimum bias events are used to tune non-perturbative QCD models,
including those describing the multiple parton interactions (MPI). Events with
inclusive production of $\gamma$ and 2 or 3 jets are used to study
increasingly important MPI phenomenon at high $p_{T}$, measure an effective
interaction cross section, $\sigma_{\rm eff}=16.4\pm 2.3$ mb, and limit
existing MPI models.
## 1 Introduction
QCD, the theory of the strong interaction between quarks and gluons, is
heavily tested in experimental studies at hadron colliders. QCD results from
the CDF and D0 collaborations obtained with integrated luminosity up to 8 fb-1
are reviewed in this paper. These results provide a crucial tests for pQCD,
PDFs, the strong coupling constant, non-perturbative models describing parton
fragmentation, and MPI phenomena. At the same time, the results are used to
search for new phenomena and impose limits on the corresponding models. The
performed extensive studies prepare a solid base for the LHC era of $pp$
collisions.
## 2 Jet production
Thorough testing of pQCD at short distances is provided through measurements
of differential inclusive jet, dijet and three-jet cross sections. The
measurements of the inclusive jet cross sections done by the D0 [1] and CDF
[2] collaborations are in agreement with pQCD predictions in a few jet
rapidity regions. However, data with uncertainties lower than theoretical
(mostly PDF) ones, favor a smaller gluon content at high Feynman $x$ ($>$0.2).
The D0 inclusive jet data has also been used to extract values of the strong
coupling constant $\alpha_{s}$ in the interval of $50<p_{T}^{\rm jet}<145$ GeV
[3]. The best fit over 22 data points leads to
$\alpha_{s}(m_{Z})=0.1161^{+0.0041}_{-0.0048}$ with improved accuracy from the
Run I CDF result [4] and also in agreement with result from HERA jet data [5].
The inclusive dijet cross sections have been measured in the D0 [6] and CDF
[7] experiments. Both measurements cover the mass range up to about 1.2 TeV
with good agreement with pQCD, and no indication of any new physics. CDF
imposed limits on some models with exotic particles decaying into two jets
[7]. D0 results are compared to mstw2008 [9] and cteq6.6m [8] PDF sets. They
are in a better agreement with mstw2008 and are systematically lower than the
central pQCD prediction at high rapidities ($1.2<|y|<2.4$).
The D0 collaboration has also measured the three-jet mass cross sections using
jets with leading (in $p_{T}$) jet $p_{T}>150$ GeV, and considering three
regions with different lower cut on the 3rd jet $p_{T}$ ($40,70,100$ GeV) and
three different jet rapidity regions ($|y|<2.4,1.6,0.8$) [10]. Results are
shown in Fig. 2 and in agreement with NLO pQCD predictions, which use
mstw2008, ct10 [11], nnpdf2.1 [12], hera1.0 [13] and abkm09 [14] PDF sets.
Results favor more mstw2008 and nnpdf2.1 PDFs. D0 has also presented a
measurement of the ratio of the inclusive 3-jet to 2-jet production cross
sections [15]. The shape of the ratio is well described by NLO QCD and is
practically independent of the PDF set. The results can potentially be used to
test the running of $\alpha_{s}$ up to a $p_{T}$ scale of 500 GeV.
The CDF collaboration studied structure of high $p_{T}$ jets by selecting only
events with at least one jet having $p_{T}>400$ GeV, $0.1<|y|<0.7$ and
considering jets with cone sizes $R=0.4,0.7$ and $1.0$ [16]. Such studies can
be used to tune parton showering and search for heavy resonances decaying
hadronically. The jet mass is calculated using 4-vectors of calorimeter towers
in a jet. Fig. 3 shows the jet mass distribution for $R=0.7$ at high masses.
The data are in agreement pythia predictions and interpolate between the QCD
LLA predictions [17] for quark and gluon jets, and confirm that the high mass
jets are mostly caused by quark fragmentation.
Figure 1: Ratios of data/theory for the dijet mass cross sections measured in
D0 and CDF are shown on the left and right plots. Figure 2: The 3-jet mass
differential cross section in the three jet rapidity intervals are compared to
NLO pQCD with different PDF sets. Figure 3: Distribution of jet mass for jets
with $R=0.7$; black crosses are data, red dashed is QCD MC, theoretical “all
quarks” and “all gluons” curves are presented as well; the inset plot compares
the results with Midpoint/SC and Anti-kT jet algorithms.
## 3 $W/Z$ \+ jets production
Both collaborations have extensively studied the $W/Z$ \+ jet productions
since these events are the main background to top-quark, Higgs boson, SUSY and
many other new physics production channels. In this section we review some of
the latest results.
Fig. 4 shows the inclusive cross section for $Z/\gamma^{\ast}$+jets production
measured by CDF [18] as a function of dijet mass and jet multiplicity. The
measurements are compared to LO and NLO pQCD predictions obtained with MCFM
[19] and are in good agreement with the NLO theory predictions. D0 measured
jet $p_{T}$ inclusive cross sections of $W+n$-jet production for jet
multiplicities $n=1-4$ [20]. The measurements are compared to the NLO
predictions for $n=1-3$ and to LO predictions for $n=4$. The measured cross
sections are generally found to agree with the NLO calculation although
certain regions of phase space are identified where the calculations could be
improved.
D0 recently published the measured cross section ratio
$\sigma(Z+b)/\sigma(Z+{\rm jet})=0.0193\pm 0.0022({\rm stat})\pm 0.0015({\rm
syst})$ for events with jet $p_{T}>20$ GeV and $|\eta|<2.5$ [21]. This most
precise measurement of the $Z+b$ fraction is consistent with the NLO theory
prediction, $0.0192\pm 0.0022$, (done with MCFM, renormalization and
factorization scales set at $m_{Z}$) and the CDF result [22]. the CDF
collaboration measured the cross section of $W+b$-jet production
$\sigma(W+b)\cdot Br(W\to l\nu)=2.74\pm 0.27({\rm stat})\pm 0.42({\rm syst})$
pb with jet $p_{T}>20$ GeV, $|\eta|<2.0$ and $l=e,\mu$. The measurement
significantly exceeds the NLO prediction $1.2\pm 0.14$ pb.
Figure 4: Left: measured inclusive cross section for $Z/\gamma^{*}$+jets
production as a function of dijet mass compared to NLO pQCD predictions as
determined using MCFM. Right: measured cross section as a function of
inclusive jet multiplicity compared to LO and NLO pQCD predictions as
determined using MCFM.
Figure 5: Left: measured $W+n$ jet differential cross section as a function of
jet $p_{T}$ for $n=1-4$, normalized to the inclusive $W\to e\nu$ cross
section. The $W+1$ jet inclusive spectra are shown by the top curve, the $W+4$
jet inclusive spectra by the bottom curve. The measurements are compared to
the fixed-order NLO predictions for $n=1-3$ and to LO predictions for $n=4$.
Right: (a) total inclusive $n$-jet cross sections $\sigma_{n}$ as a function
of $n$, (b) the ratio of the theory predictions to the measurements, and (c)
$\sigma_{n}/\sigma_{n-1}$ ratios for data, Blackhat+Sherpa and Rocket+MCFM.
The hashed areas represent the theoretical uncertainty arising from the choice
of renormalization and factorization scale.
Figure 6: Left (D0): the distributions of the $b$, $c$, light jets and data
over the $b-$jet discriminant; MC templates are weighted by the fractions
found from the fit to data. Right (CDF): the secondary vertex mass fit for the
tagged jets in the selected sample.
## 4 Photon production
Since high $p_{T}$ photons emerge directly from $p\bar{p}$ collisions and
provide a direct probe of the parton hard scattering dynamics, they are of
permanent interest in high energy physics. The inclusive photon production
cross sections measured by D0 and CDF in the central rapidity region [24, 25]
are in agreement within experimental uncertainties and both indicate a
difficulty for NLO pQCD to describe the low $p_{T}$ behavior.
In light of the Higgs boson search and other possible resonances decaying to a
photon pair, both the collaborations performed a thorough study of the
diphoton production. D0 measured the diphoton cross sections (see Fig. 7) as a
function of the diphoton mass $M_{\gamma\gamma}$, the transverse momentum of
the diphoton system $p_{T}^{\gamma\gamma}$, the azimuthal angle between the
photons $\Delta\phi_{\gamma\gamma}$, and the polar scattering angle of the
photons. The latter three cross sections are also measured in the three
$M_{\gamma\gamma}$ bins, $30-50,50-80$ and $80-350$ GeV. The measurements are
compared to NLO QCD and pythia [32] predictions. The results show that the
largest discrepancies between data and NLO predictions for each of the
kinematic variables originate from the lowest $M_{\gamma\gamma}$ region
($M_{\gamma\gamma}$ $<50$ GeV), where the contribution from
$gg\to\gamma\gamma$ is expected to be largest [26]. The discrepancies between
data and the theory predictions are reduced in the intermediate
$M_{\gamma\gamma}$ region, and a quite satisfactory description of all
kinematic variables is achieved for the $M_{\gamma\gamma}$$>80$ GeV region,
the relevant region for the Higgs boson and new phenomena searches. The CDF
collaboration measured the diphoton production cross sections functions of
$M_{\gamma\gamma}$, $p_{T}^{\gamma\gamma}$and $\Delta\phi_{\gamma\gamma}$.
They are shown in Fig. 8. None of the models describe the data well in all
kinematic regions, in particular at low diphoton mass ($M_{\gamma\gamma}$$<60$
GeV), low $\Delta\phi_{\gamma\gamma}$($<1.7$ rad) and moderate
$p_{T}^{\gamma\gamma}$($20-50$ GeV).
Figure 7: The measured differential diphoton production cross sections as
functions of (a) $M_{\gamma\gamma}$, (b) $p_{T}^{\gamma\gamma}$and (c)
$\Delta\phi_{\gamma\gamma}$ in D0 experiment. The data are compared to the
theoretical predictions from resbos, diphox, and pythia. The ratio of
differential cross sections between data and resbos are displayed as black
points with uncertainties in the bottom plots. The solid (dashed) line shows
the ratio of the predictions from diphox (pythia) to those from resbos. In the
bottom plots, the scale uncertainties are shown by dash-dotted lines and the
PDF uncertainties by shaded regions.
Figure 8: The measured differential diphoton production cross sections as
functions of (from left to right) $M_{\gamma\gamma}$,
$p_{T}^{\gamma\gamma}$and $\Delta\phi_{\gamma\gamma}$ in CDF experiment.
## 5 Multiple parton interactions
The CDF and D0 collaborations comprehensively studied the phenomenon of MPI
events in a few Run II measurements. In this section we mention some of the
recently published results. CDF studied charged-particle $p_{T}$ sum densities
and multiplicities in Drell-Yan and jet events [28]. Specifically, both
distributions have been analyzed in the three regions: towards the total
lepton pair (Z-boson) $\vec{p}_{T}$, opposite to this direction (“away”
region) and in the region transverse to the lepton pair/jet $\vec{p}_{T}$. The
charged-particle $p_{T}$ sum density in the three regions in the Drell-Yan
events is shown on the left plot of Fig. 9 as a function of the lepton pair
$p_{T}$. The same quantity is also plotted for the transverse region on the
right plot. One can see a similar trend in both Drell-Yan and jet events which
can be considered as MPI universality. The tuned pythia describes the data
very well.
Figure 9: Left: the comparison of the total charged-particle $p_{T}$ sum
density $dp_{T}/d\eta d\phi$ in the three regions in the Drell-Yan events:
“transverse”, “away” and “toward”. Right: the $p_{T}$ sum density in the
transverse region in Drell-Yan and jet production as a function of the lepton
pair or leading jet $p_{T}$.
D0 has studied events with double parton (DP) scattering in $\gamma+3$ jet
events [29], in which two pairs of partons undergo two hard interactions in a
single $p\bar{p}$ collision. The DP events can be a background to many rare
processes but they also provide insight into the spatial distribution of
partons in the colliding hadrons. D0 measured the so-called effective cross
section, that characterizes rates of the DP events, $\sigma^{\gamma j,jj}_{\rm
DP}=\sigma^{\gamma j}\sigma^{jj}/\sigma_{\rm eff}$. The measurement was done
in the three bins of the 2nd (ordered in $p_{T}$) jet $p_{T}$. The results are
shown in Fig. 11. Using these three (almost uncorrelated0 points, the obtained
average effective cross section is $\sigma_{\rm eff}=16.4\pm 0.3({\rm
stat})\pm 2.3({\rm syst})$. It is in agreement with the previous CDF result
[30]. To tune MPI models, D0 also measured cross sections for the azimuthal
angle defined between the $p_{T}$ vectors of the $\gamma+$jet and dijet
systems in the three $p_{T}$ bins of the 2nd jet $p_{T}$ [31]. Comparison of
data with a few MPI and two “no MPI” models are shown in Fig. 11. One can see
that data clearly contain DP events and favor more Perugia MPI tunes.
Figure 10: The measured effective cross section vs 2nd jet $p_{T}$.
Figure 11: Left: Normalized differential cross section in the $\gamma+3$-jet +
X events, $(1/\sigma_{\gamma 3j})\sigma_{\gamma 3j}/d\Delta S$, in data
compared to MC models and the ratio of data over theory, only for models
including MPI, in the range $15<p_{T}^{jet2}<30$ GeV. Right: Normalized
differential cross section in $\gamma+2$-jet + X events, $(1/\sigma_{\gamma
2j})\sigma_{\gamma 2j}/d\Delta\phi$, in data compared to MC models and the
ratio of data over theory, only for models including MPI, in the range
$15<p_{T}^{jet2}<20$ GeV.
## 6 Summary
The Tevatron experiments provide precision QCD measurements of many
fundamental observables. In most cases, the results are mutually consistent
and/or complementary to each other. Jet measurements show good agreement with
pQCD, sensitivity to PDF sets, the strongest constraint on high-$x$ gluon PDF,
provide detailed studies of different jet algorithms, are used to extract
$\alpha_{s}$, study jet substructure, and provide limits on many new phenomena
models. The $W/Z+$jets results provide extensive tests of pQCD and tune
existing MC models. The photon results test fixed order NLO pQCD predictions
accounting for resummation and fragmentation effects and show that the theory
should be better understood. Measurements of underlying/MPI events impose
strong constraints and improve phenomenological MPI models at low and high
$p_{T}$ regimes.
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* [24] V.M. Abazov et al. (D0), Phys.Lett. B 639, 151 (2006); ibid. 658, 285 (2008).
* [25] T. Aaltonen et al. (CDF), Phys. Rev. D 80, 111106 (2009).
* [26] V. M. Abazov et al. (D0), Phys. Lett. B 690, 108 (2010) and references therein.
* [27] T. Aaltonen et al. (CDF), CDF Note CDF-10160 (2011).
* [28] T. Aaltonen et al. (CDF), Phys.Rev. D 82, 034001 (2010).
* [29] V. M. Abazov et al. (D0), Phys.Rev. D 81, 052012 (2010).
* [30] F. Abe et al. (CDF), Phys.Rev. D 56, 3811 (1997).
* [31] V. M. Abazov et al. (D0), Phys.Rev. D 83, 052008 (2011).
* [32] T. Sjostrand, S. Mrenna, P. Z. Skands, JHEP 0605, 026 (2006).
|
arxiv-papers
| 2011-11-30T23:12:38 |
2024-09-04T02:49:24.872809
|
{
"license": "Public Domain",
"authors": "Dmitry Bandurin (for the D0 and CDF Collaborations)",
"submitter": "Dmitry Bandurin V",
"url": "https://arxiv.org/abs/1112.0051"
}
|
1112.0343
|
# Ontological Queries: Rewriting and Optimization (Extended Version)111This is
an extended and revised version of the paper [1].
Georg Gottlob1,2, Giorgio Orsi1,3, Andreas Pieris1
$~{}^{1}$Department of Computer Science, University of Oxford, UK
$~{}^{2}$Oxford-Man Institute of Quantitative Finance, University of Oxford,
UK
$~{}^{3}$Institute for the Future of Computing, University of Oxford, UK
{georg.gottlob,giorgio.orsi,andreas.pieris}@cs.ox.ac.uk
###### Abstract
Ontological queries are evaluated against an ontology rather than directly on
a database. The evaluation and optimization of such queries is an intriguing
new problem for database research. In this paper we discuss two important
aspects of this problem: query rewriting and query optimization. Query
rewriting consists of the compilation of an ontological query into an
equivalent query against the underlying relational database. The focus here is
on soundness and completeness. We review previous results and present a new
rewriting algorithm for rather general types of ontological constraints. In
particular, we show how a conjunctive query against an ontology can be
compiled into a union of conjunctive queries against the underlying database.
Ontological query optimization, in this context, attempts to improve this
process so to produce possibly small and cost-effective UCQ rewritings for an
input query. We review existing optimization methods, and propose an effective
new method that works for linear Datalog±, a class of Datalog-based rules that
encompasses well-known description logics of the _DL-Lite_ family.
## 1 Introduction
This paper is about ontological query processing, an important new challenge
to database research. We will review existing methods and propose new
algorithms for compiling an ontological query, that is, a query against an
ontology on top of a relational database, into a direct query against this
database, and we will deal with optimization issues related to this process so
as to obtain possibly small and efficient compiled queries. In this section,
we first discuss a number of relevant concepts, and then illustrate query
rewriting and optimization processes in the context of a small but non-trivial
example.
#### Ontologies.
The use of ontologies and ontological reasoning in companies, governmental
organizations, and other enterprises has become widespread in recent years. An
ontology is an explicit specification of a conceptualization of an area of
interest [2], and consists of a formal representation of knowledge as a set of
concepts within a domain, and the relationships between those concepts [3]. To
distinguish an enterprise ontology from a data dictionary, Dave McComb
explicitly refers to the formal semantics of ontologies that enables automated
processing and inferencing, while the interpretation of a data dictionary is
strictly done by humans [4]. Moreover, ontologies have been adopted as high-
level conceptual descriptions of the data contained in data repositories that
are sometimes distributed and heterogeneous in the data models. Due to their
high expressive power, ontologies are also substituting more traditional
conceptual models such as UML class-diagrams and E/R schemata.
#### Description Logics.
Description logics (DLs) are logical languages for expressing and modelling
ontologies. The best known DLs are those underlying the _OWL_
language222http://www.w3.org/TR/owl2-overview/. The main ontological reasoning
and query answering tasks in the complete OWL language, called _OWL Full_ ,
are undecidable. For the most well-known decidable fragments of OWL,
ontological reasoning and query answering is still computationally very hard,
typically 2exptime-complete.
In description logics, the ontological axioms are usually divided into two
sets: The ABox (assertional box), which essentially contains factual knowledge
such as “IBM is a company”, denoted by $\mathit{company}(\mathit{ibm})$, or
“IBM is listed on the NASDAQ”, which could be represented as a fact of the
form $\mathit{list\\_comp}(\mathit{ibm},\mathit{nasdaq})$, and a TBox
(terminological box) which contains axioms and constraints that allow us, on
the one hand, to infer new facts from those given in the ABox, and, on the
other hand, to express restrictions such as keys. For example, a TBox may
contain an axiom stating that for each fact $\mathit{list\\_comp}(X,Y)$, $Y$
must be a financial index, which in DL is expressed as
$\exists\mathit{list\\_comp}^{-}\sqsubseteq\mathit{fin\\_idx}$. If the fact
$\mathit{fin\\_idx}(\mathit{nasdaq})$ is not already present in the ABox, it
can be derived via the above axiom from
$\mathit{list\\_comp}(\mathit{ibm},\mathit{nasdaq})$. Thus, the atomic query
“$q(X)\leftarrow\mathit{fin\\_idx}(X)$” would return $\mathit{nasdaq}$ as one
of the answers. Note that the axiom
$\exists\mathit{list\\_comp}^{-}\sqsubseteq\mathit{fin\\_idx}$, which
corresponds to an inclusion dependency, is enforced by adding new tuples,
rather than just being checked. This is one main difference between
ontological constraints and classical database dependencies. In database
terms, the above axiom is to be interpreted more like a trigger than a
classical constraint.
#### Ontology Based Data Access (OBDA).
We are currently witnessing the marriage of ontological reasoning and database
technology. In fact, this amalgamation consists in the realization of the
obvious idea that ABoxes shall be implemented in form of a relational
database, or even stored in classical RDBMSs. Moreover, very large existing
databases are semantically enriched with ontological constraints. There are a
number of recent commercial systems and experimental prototypes that extend
RDBMSs with the possibility of querying an ontology that is rooted in a
database (for examples, see Section 2). The main problem here is how to couple
these two different types of technology smoothly and efficiently, and this is
also the main theme of the present paper.
One severe obstacle to efficient OBDA is the already mentioned high
computational complexity of query answering with description logics. The
situation clearly worsens when the ABoxes of enterprise ontologies are very
large databases. To tackle this problem, new, lightweight DLs have been
designed, that guarantee polynomial-time data complexity for conjunctive query
answering. This means that based on a fixed TBox, a fixed query can be
answered in polynomial time over variable databases. The best-known and best-
studied examples of such lightweight DLs are the _DL-Lite_ [5] and
$\mathcal{E}\mathcal{L}$ (see, e.g., [6]) families. These languages can be
considered tractable subclasses of OWL. It was convincingly argued that simple
DLs such as DL-Lite or $\mathcal{E}\mathcal{L}$ are sufficient for modelling
an overwhelming number of applications.
More recently, the Datalog± family of description logics was introduced [7, 8,
9, 10]. Its syntax is based on classical first-order logic, more specifically,
on variants of the well-known Datalog language [11]. The basic Datalog± rules
are known as _tuple-generating dependencies_ (TGDs) in the database literature
[12]. Tractable DLs in this framework are guarded Datalog±, which is
noticeably more general than both DL-Lite and $\mathcal{E}\mathcal{L}$, and
the DLs linear Datalog± and sticky-join Datalog±, which both encompass DL-
Lite.
Besides being more expressive than DL-Lite, suitable Datalog± languages offer
a more compact representation of the attributes of concepts and roles, since
description logics are usually restricted to unary and binary predicates only.
Consider, as an example, a relation $\mathit{stock}(\underline{{\sf id}},{\sf
name},{\sf unit}$-${\sf price})$. Representing this relation in DL would
require the introduction of a concept symbol $\mathit{stock}$, and of three
attribute symbols $\mathit{id}$, $\mathit{name}$ and
$\mathit{unit}$-$\mathit{price}$. These entities must be then weaved together
by the TBox formula $\mathit{stock}\sqsubseteq\exists
id\sqcap\exists\mathit{name}\sqcap\exists\mathit{unit}$-$\mathit{price}$.
Datalog± represents the relation in a natural way by means of a ternary
predicate $\mathit{stock}$. In the same way, Datalog± provides a more natural
syntax for many other DL formulae; for example, an inverse role assertion
$r\sqsubseteq s^{-}$ is represented as a (full) TGD $r(X,Y)\rightarrow
s(Y,X)$, while an existential restriction $p\sqsubseteq\exists r.q$ is
represented as a (partial) TGD $p(X)\rightarrow\exists Y\,r(X,Y),q(Y)$.
#### First-Order Rewritability.
Polynomial-time tractability is often considered not to be good enough for
efficient query processing. Ideally, one would like to achieve the same
complexity as for processing SQL queries, or, equivalently, first-order (FO)
queries. An ontology language $\mathcal{L}$ is _first-order rewritable_ if,
for every TBox $\Sigma$ expressed in $\mathcal{L}$ and a query $q$, a first-
order query $q_{\Sigma}$ (called the perfect rewriting) can be constructed
such that, given a database $D$, $q_{\Sigma}$ evaluated over $D$ yields
exactly the same result as $q$ evaluated against $D$ and $\Sigma$. Since
answering first-order queries is in the class ac0 in data complexity [13], it
immediately follows that under FO-rewritable TGDs, query answering is also in
ac0 in data complexity
This notion was first introduced by Calvanese et al. [5] in the concept of
description logics. If a DL guarantees the FO-rewritability of each query
under every TBox, we simply say that the logic is FO-rewritable. FO-
rewritability is a most desirable property since it ensures that the reasoning
process can be largely decoupled from data access. In fact, to answer query
$q$, a separate software can compile $q$ into $q_{\Sigma}$, and then just
submit $q_{\Sigma}$ as a standard SQL query to the DBMS holding $D$, where it
is evaluated and optimized in the usual way.
Excitingly, it was shown that the members of the DL-Lite family, as well as
the slightly more expressive language linear Datalog± are FO-rewritable.
Moreover, even the much more expressive language of sticky-join Datalog± is
FO-rewritable. For these languages, a pair $\langle\Sigma,q\rangle$, where $q$
is a CQ, is rewritten as an SQL expression equivalent to a UCQ $q_{\Sigma}$.
The research challenge we address in this paper is precisely the question of
how to rewrite $\langle\Sigma,q\rangle$ to $q_{\Sigma}$ correctly and
efficiently. Let us illustrate this process by a small, but comprehensive
example.
Consider the following relational schema $\mathcal{R}$ representing financial
information about companies and their stocks:
$\begin{array}[]{rcl}&&\mathit{stock}({\sf id},{\sf name},{\sf unit}$-${\sf
price})\\\ &&\mathit{company}({\sf name},{\sf country},{\sf segment})\\\
&&\mathit{list\\_comp}({\sf stock},{\sf list})\\\ &&\mathit{fin\\_idx}({\sf
name},{\sf type},{\sf ref}$-${\sf mkt})\\\ &&\mathit{stock\\_portf}({\sf
company},{\sf stock},{\sf qty}).\end{array}$
The $\mathit{stock}$ relation contains information about stocks such as the
name, and the price per unit. The relation $\mathit{company}$ contains
information about companies; in particular, the name, the country, and the
market segment of a company. The relation $\mathit{list\\_comp}$ relates a
stock to a financial index (i.e., NASDAQ, FTSE, NIKKEI) represented by the
relation $fin\\_idx$ which, in turn, contains information about the types of
stocks in the index, and the reference market (e.g., London Stock Exchange).
Finally, $\mathit{stock\\_portf}$ relates companies to their stocks along with
an indication of the amount of the investment.
Datalog± provides the necessary expressive power to extend $\mathcal{R}$ with
ontological constraints in an easy and intuitive way. Examples of such
constraints follow:
$\begin{array}[]{rcl}\sigma_{1}&:&\mathit{stock\\_portf(X,Y,Z)}\rightarrow\exists
V\exists W\ \mathit{company(X,V,W)}\\\
\sigma_{2}&:&\mathit{stock\\_portf(X,Y,Z)}\rightarrow\exists V\exists W\
\mathit{stock(Y,V,W)}\\\
\sigma_{3}&:&\mathit{list\\_comp(X,Y)}\rightarrow\exists Z\exists W\
\mathit{fin\\_idx(Y,Z,W)}\\\
\sigma_{4}&:&\mathit{list\\_comp(X,Y)}\rightarrow\exists Z\exists W\
\mathit{stock(X,Z,W)}\\\
\sigma_{5}&:&\mathit{stock\\_portf(X,Y,Z)}\rightarrow\mathit{has\\_stock(Y,X)}\\\
\sigma_{6}&:&\mathit{has\\_stock(X,Y)}\rightarrow\exists Z\
\mathit{stock\\_portf(Y,X,Z)}\\\
\sigma_{7}&:&\mathit{stock(X,Y,Z)}\rightarrow\exists V\exists W\
\mathit{stock\\_portf(V,X,W)}\\\
\sigma_{8}&:&\mathit{stock(X,Y,Z)}\rightarrow\mathit{fin\\_ins(X)}\\\
\sigma_{9}&:&\mathit{company(X,Y,Z)}\rightarrow\mathit{legal\\_person(X)}\\\
\delta_{1}&:&\mathit{legal\\_person(X,Y,Z),fin\\_ins(X,V,W)\rightarrow\bot}.\end{array}$
Figure 1: A (partial) rewriting for the Stock Exchange example.
$q^{[0]}(A,B,C)\leftarrow\mathit{fin\\_ins}(A),\mathit{stock\\_portf}(B,A,D),\mathit{company}(B,E,F),\mathit{list\\_comp}(A,C),\mathit{fin\\_idx}(C,G,H)$
---
$q^{[1]}(A,B,C)\leftarrow\mathit{fin\\_ins}(A),\underline{\mathit{has\\_stock}(A,B)},\mathit{company}(B,E,F),\mathit{list\\_comp}(A,C),\mathit{fin\\_idx}(C,G,H)$
$q^{[2]}(A,B,C)\leftarrow\mathit{fin\\_ins}(A),\mathit{has\\_stock}(A,B),\underline{\mathit{stock\\_portf}(B,E,F)},\mathit{list\\_comp}(A,C),\mathit{fin\\_idx}(C,G,H)$
$q^{[3]}(A,B,C)\leftarrow\underline{\mathit{stock}(A,J,K)},\mathit{has\\_stock}(A,B),\mathit{stock\\_portf}(B,E,F),\mathit{list\\_comp}(A,C),\mathit{fin\\_idx}(C,G,H)$
$\ldots$
The first four TGDs set the “domain” and the “range” of the
$\mathit{stock\\_portf}$ and $\mathit{list\\_comp}$ relations, respectively.
TGDs $\sigma_{5}$ and $\sigma_{6}$ assert that $\mathit{stock\\_portf}$ and
$\mathit{has\\_stock}$ are “inverse relations”, while $\sigma_{7}$ expresses
that each stock must belong to some stock portfolio. TGDs $\sigma_{8}$ and
$\sigma_{9}$ model taxonomic relationships such as the facts that each stock
is a financial instrument, and each company is a legal person. Finally, the
negative constraint $\delta_{1}$ (where $\bot$ denotes the truth constant
$\mathit{false}$) states that legal persons and financial instruments are
disjoint sets.
Consider now the following conjunctive query $q$ asking for all the triples
$\langle a,b,c\rangle$, where $a$ is a financial instrument owned by the
company $b$ and listed on $c$:
$\begin{array}[]{rcl}\mathit{q(A,B,C)}&\leftarrow&\mathit{fin\\_ins}(A),\mathit{stock\\_portf}(B,A,D),\mathit{company}(B,E,F),\\\
&&\mathit{list\\_comp}(A,C),\mathit{fin\\_idx}(C,G,H).\end{array}$
Since $\Sigma=\\{\sigma_{1},\ldots,\sigma_{9}\\}$ is a set of linear TGDs,
i.e., TGDs with single body-atom, query answering under $\Sigma$ is FO-
rewritable. Thus, it is possible to reformulate $\langle\Sigma,q\rangle$ to a
first-order query $q_{\Sigma}$ such that, for every database $D$,
$D\cup\Sigma\models q$ iff $D\models q_{\Sigma}$. A naive rewriting procedure
would use the TGDs of $\Sigma$ as rewriting rules for the atoms in $q$ to
generate all the CQs of the perfect rewriting. Figure 1 shows a (partial)
rewriting for $q$, where the query obtained at the $i$-th step is denoted as
$q^{[i]}$, and the newly introduced atoms are underlined. In particular,
$q^{[0]}$ is the given query $q$, $q^{[1]}$ is obtained from $q^{[0]}$ by
using $\sigma_{6}$, $q^{[2]}$ is obtained from $q^{[1]}$ by applying
$\sigma_{1}$, and $q^{[3]}$ is obtained from $q^{[2]}$ by using $\sigma_{8}$.
The complete perfect rewriting contains more than 200 queries executing more
than 1000 joins. However, by exploiting the set of constraints, it is possible
to eliminate redundant atoms in the generated queries, and thus reduce the
number of the CQs in the rewritten query. For example, in the given query $q$
above it is possible to eliminate the atom $\mathit{fin\\_ins(A)}$ since, due
to the existence of the TGDs $\sigma_{2}$ and $\sigma_{8}$ in $\Sigma$, if the
atom $\mathit{stock\\_portf}(B,A,D)$ is satisfied, then immediately the atom
$\mathit{fin\\_ins(A)}$ is also satisfied. Notice that by eliminating a
redundant atom from a query, we also eliminate all the atoms that are
generated starting from it during the rewriting process. Moreover, due to the
TGD $\sigma_{3}$, if the atom $\mathit{list\\_comp(A,C)}$ in $q$ is satisfied,
then the atom $\mathit{fin\\_idx(C,G,H)}$ is also satisfied, and therefore can
be eliminated. Finally, due to the TGD $\sigma_{1}$, if the atom
$\mathit{stock\\_portf}(B,A,D)$ is satisfied, then the atom
$\mathit{company(B,E,F)}$ is also satisfied, and hence is redundant. The query
that has to be considered as input of the rewriting process is therefore
$q(A,B,C)\leftarrow\mathit{stock\\_portf}(B,A,D),\mathit{list\\_comp}(A,C)$
that produces a perfect rewriting containing the following two queries
executing only two joins:
$\begin{array}[]{rcl}q(A,B,C)&\leftarrow&\mathit{list\\_comp}(A,C),\mathit{stock\\_portf}(B,A,D)\\\
q(A,B,C)&\leftarrow&\mathit{list\\_comp}(A,C),\mathit{has\\_stock}(A,B).\end{array}$
#### Contributions and Roadmap.
After a review of previous work on ontology based data access in the next
section, and some formal definitions and preliminaries in Section 3, we
present a short overview of the Datalog± family in Section 4. We then proceed
with new research results. In Section 5, we propose a new rewriting algorithm
that improves the one stated in [14] by substantially reducing the number of
redundant queries in the perfect rewriting. In Section 6, we present a
polynomial-time optimization strategy based on the early-pruning of redundant
atoms produced during the rewriting process. An implementation and
experimental evaluation of the new method is discussed in Section 7. We also
discuss the relationship between our optimization technique and optimal query
minimization algorithms such as the _chase & back-chase_ algorithm [15]. We
conclude with a brief outlook on further research.
## 2 Ontology Based Data Access
Answering queries under constraints and the related optimization techniques
are important topics in data management beyond the obvious research interest.
In particular, they are profitable opportunities for companies that need to
deliver efficient and effective data management solutions to their customers.
This trend is becoming even more evident as a plethora of robust systems and
APIs for Semantic Web data management proposed in the recent years. These
systems span from open-source solutions such as
Virtuoso333http://virtuoso.openlinksw.com/, Sesame444http://www.openrdf.org/,
RDFSuite [16], KAON555http://kaon.semanticweb.org/ and
Jena666http://jena.sourceforge.net/, to commercial implementations such as the
semantic extensions implemented in Oracle Database 11g R2 [17] and
BigOWLLim777http://www.ontotext.com/owlim/. In this Section we briefly analyze
the systems providing rewriting-based access to databases under ontological
constraints, and we highlight some crucial points that we want to address in
this work.
We first present the class of constraints identified by the members of the DL-
Lite family [5], namely, DL-LiteA, DL-LiteF, and DL-LiteR, underlying the W3C
OWL-QL profile of the OWL language. These constraints correspond to unary and
binary _inclusion dependencies_ combined with a restricted form of _key
constraints_. In order to perform query answering under this class of
constraints, a rewriting algorithm, introduced in [5] and implemented in the
QuOnto system, reformulates the given query into unions of conjunctive
queries. The size of the reformulated query is unnecessarily large due to a
number of reasons. In the first place, _(i)_ basic optimization techniques
such as the identification of the connected components in the body of the
input query, or the computation of any form of query decomposition [18], are
not applied. Moreover, _(ii)_ the fact that the given set of constraints can
be used to identify existential joins in the reformulated query which can be
eliminated is not exploited. Finally, _(iii)_ the factorization step (which is
needed to guarantee completeness) is applied exhaustively, and as a result
many superfluous queries are generated.
Peréz-Urbina et al. [19] proposed an alternative resolution-based rewriting
algorithm, implemented in the Requiem system, that addressed the issue of the
useless factorizations (and therefore of the redundant queries generated due
to this weakness) by directly handling existential quantification through
proper functional terms. The algorithm has then been extended to more
expressive DL languages [19]. In this case the output of the rewriting is a
Datalog program.
Rosati et al. [20] recently proposed a very sophisticated rewriting technique,
implemented in the Presto system, that addresses some of the issues described
above. In particular, _(i)_ the unnecessary existential joins are eliminated
by resorting to the concept of _most-general subsumees_ , which also avoids
the unnecessary factorizations, and _(ii)_ the connectivity of the given query
is checked before executing the algorithm; in case the query is not connected,
Presto splits the query in connected components and rewrites them separately.
Notice that Presto produces a non-recursive Datalog program, and not a union
of conjunctive queries. This allows the “hiding” of the exponential blow-up
inside the rules instead of generating explicitly the disjunctive normal form.
The final rewriting is exponential only in the number of non-eliminable
existential joins, but not in the size of the input query.
The approaches presented above have been proven very effective when applied to
very particular classes of description logic constraints. Following a more
general approach for ontological query answering, Calì et al. [14] presented a
backward-chaining rewriting algorithm which is able to deal with arbitrary
sets of TGDs, providing that the class of TGDs under consideration satisfies
suitable syntactic restrictions that guarantee the termination of the
algorithm. However, this algorithm is inspired by the original QuOnto
algorithm and inherits all its drawbacks.
Despite the possibly exponential number of queries to be constructed, we know
that all such queries are independent from each other, and thus they can be
easily executed in parallel threads and distributed on multiple processors.
Notice that a non-recursive Datalog program is not equally easy to distribute.
Moreover, the optimizations implemented in current DBMS systems for (unions
of) conjunctive queries are much more advanced than those implemented for the
positive existential first-order queries resulting from the translation of a
non-recursive Datalog program into a concrete query language such as SQL. It
is clear that a trade-off between these two approaches must be found in order
to exploit as much as possible the current optimization techniques, while
keeping the size of the rewriting reasonably small in order to make the
execution of it feasible in practice.
A related research field is that of query minimization [21], in particular, in
presence of views and constraints [22, 15]. Given a conjunctive query $q$, and
a set of constraints $\Sigma$, the goal is to find all the minimal equivalent
reformulations of $q$ under the constraints of $\Sigma$. The most interesting
approach in this respect is the chase & back-chase algorithm (C&B) [15],
implemented in the MARS system [23]. The algorithm freezes the atoms of
$\mathit{body}(q)$ and, by considering them as a database $D_{q}$, applies the
following two steps. During the chase-step, the chase of $D_{q}$ w.r.t.
$\Sigma$ is constructed, and then the atoms of $\mathit{chase}(D_{q},\Sigma)$
are considered as the body-atoms of a query $q_{u}$, called the universal
plan. The back-chase step considers all the possible subsets of the atoms of
$\mathit{body}(q_{u})$, starting from those with a single-atom, which are then
considered as the body of a query $q^{\prime}$. Whenever there exists a
containment mapping from $\mathit{body}(q_{u})$ to
$\mathit{chase}(D_{q^{\prime}},\Sigma)$, where $D_{q^{\prime}}$ is the
database obtained by freezing $\mathit{body}(q^{\prime})$, then $q^{\prime}$
is an equivalent reformulation of $q$. Moreover, every time an equivalent
reformulation $q^{\prime}$ is found, the back-chase does not consider any of
the supersets of the atoms of $\mathit{body}(q^{\prime})$ because they will be
automatically implied by the atoms of $q^{\prime}$, and therefore the produced
query would be redundant. This particular exploration strategy guarantees the
minimality of the reformulations. A non-negligible drawback of this approach
is the fact that we need to compute the chase of $D_{q}$ w.r.t. $\Sigma$, and
also the chase for the (exponentially many) databases $D_{q^{\prime}}$ w.r.t
$\Sigma$. Clearly, this makes the procedure computationally expensive.
## 3 Preliminaries
In this section we recall some basics on relational databases, conjunctive
queries, tuple-generating dependencies, and the chase procedure.
### 3.1 Relational Databases and Conjunctive Queries
Consider two pairwise disjoint (infinite) sets of symbols $\Delta_{c}$ and
$\Delta_{z}$ such that: $\Delta_{c}$ is a set of _constants_ (which
constitutes the domain of a database), and $\Delta_{z}$ is a set of _labeled
nulls_ (used as placeholders for unknown values). Different constants
represent different values (_unique name assumption_), while different nulls
may represent the same value. Throughout the paper, we denote by $\mathbf{X}$
sequences of variables $X_{1},\ldots,X_{k}$, where $k\geq 0$, and by $[n]$ the
set $\\{1,\ldots,n\\}$, for any $n\geq 1$.
A _relational schema_ $\mathcal{R}$ (or simply _schema_) is a set of
_relational symbols_ (or _predicate symbols_), each with its associated arity.
A _position_ $r[i]$ (or $\langle r,i\rangle$) is identified by a predicate
$r\in\mathcal{R}$ and its $i$-th argument. A _term_ $t$ is a constant, labeled
null, or variable. An _atomic formula_ (or simply _atom_) has the form
$r(t_{1},\ldots,t_{n})$, where $r\in\mathcal{R}$ has arity $n$, and
$t_{1},\ldots,t_{n}$ are terms. Conjunctions of atoms are often identified
with the sets of their atoms.
A _substitution_ from one set of symbols $S_{1}$ to another set of symbols
$S_{2}$ is a function $h:S_{1}\rightarrow S_{2}$. A _homomorphism_ from a set
of atoms $A_{1}$ to a set of atoms $A_{2}$, both over the same schema
$\mathcal{R}$, is a substitution $h$ from the set of terms of $A_{1}$ to the
set of terms of $A_{2}$ such that: (i) if $t\in\Delta_{c}$, then $h(t)=t$, and
(ii) if $r(t_{1},\ldots,t_{n})$ is in $A_{1}$, then
$h(r(t_{1},\ldots,t_{n}))=r(h(t_{1}),\ldots,h(t_{n}))$ is in $A_{2}$. The
notion of homomorphism naturally extends to conjunctions of atoms.
A _relational instance_ (or simply _instance_) $I$ for a schema $\mathcal{R}$
is a (possibly infinite) set of atoms of the form $r({\bf t})$, where
$r\in\mathcal{R}$ has arity $n$ and ${\bf
t}\in(\Delta_{c}\cup\Delta_{z})^{n}$. A _database_ is a finite relational
instance. A _conjunctive query_ (CQ) $q$ of arity $n$ over a schema
$\mathcal{R}$ is a formula of the form
$q(\mathbf{X})\leftarrow\phi(\mathbf{X},\mathbf{Y})$, where
$\phi(\mathbf{X},\mathbf{Y})$ is a conjunction of atoms over $\mathcal{R}$,
and $q$ is an $n$-ary predicate. $\phi(\mathbf{X},\mathbf{Y})$ is called the
body of $q$, denoted as $body(q)$, and $q(\mathbf{X})$ is the head of $q$,
denoted as $head(q)$. A Boolean conjunctive query (BCQ) is a CQ of arity zero.
The answer to a CQ $q$ of arity $n$ over an instance $I$, denoted as $q(I)$,
is the set of all $n$-tuples $\mathbf{t}\in(\Delta_{c})^{n}$ for which there
exists a homomorphism
$h:\mathbf{X}\cup\mathbf{Y}\rightarrow\Delta_{c}\cup\Delta_{z}$ such that
$h(\phi(\mathbf{X},\mathbf{Y}))\subseteq I$ and $h(\mathbf{X})=\mathbf{t}$. A
BCQ has only the empty tuple $\langle\rangle$ as possible answer, in which
case it is said that has positive answer. Formally, a BCQ has _positive_
answer over $I$, denoted as $I\models q$, iff $\langle\rangle\in q(I)$. A
_union of CQs_ (UCQ) $Q$ of arity $n$ is a set of CQs, where each $q\in Q$ has
the same arity $n$ and uses the same predicate symbol in the head. The answer
to $Q$ over an instance $I$, denoted as $Q(I)$, is defined as the set of
tuples $\\{\mathbf{t}~{}|~{}\textrm{there~{}exists~{}}q\in
Q\textrm{~{}such~{}that~{}}\mathbf{t}\in q(I)\\}$.
### 3.2 Tuple-Generating Dependencies
A tuple-generating dependency (TGD) $\sigma$ over a schema $\mathcal{R}$ is a
first-order formula $\forall{\bf X}\forall{\bf Y}\phi({\bf X},{\bf
Y})\rightarrow\exists\mathbf{Z}\,\psi({\bf X},{\bf Z})$, where
$\phi(\mathbf{X},\mathbf{Y})$ and $\psi(\mathbf{X},\mathbf{Z})$ are
conjunctions of atoms over $\mathcal{R}$, called the body and the head of
$\sigma$, denoted as $\mathit{body}(\sigma)$ and $\mathit{head}(\sigma)$,
respectively. Henceforth, to avoid notational clutter, we will omit the
universal quantifiers in TGDs. Such $\sigma$ is satisfied by an instance $I$
for $\mathcal{R}$ iff, whenever there exists a homomorphism $h$ such that
$h(\phi(\mathbf{X},\mathbf{Y}))\subseteq I$, there exists an extension
$h^{\prime}$ of $h$ (i.e., $h^{\prime}\supseteq h$) such that
$h^{\prime}(\psi(\mathbf{X},\mathbf{Z}))\subseteq I$.
We now define the notion of _query answering_ under TGDs. Given a database $D$
for $\mathcal{R}$, and a set $\Sigma$ of TGDs over $\mathcal{R}$, the _models_
of $D$ w.r.t. $\Sigma$, denoted as $\mathit{mods}(D,\Sigma)$, is the set of
all instances $I$ such that $I\models D\cup\Sigma$, which means that
$I\supseteq D$ and $I$ satisfies $\Sigma$. The _answer_ to a CQ $q$ w.r.t. $D$
and $\Sigma$, denoted as $\mathit{ans}(q,D,\Sigma)$, is the set
$\\{\mathbf{t}~{}|~{}\mathbf{t}\in
q(I){\rm~{}for~{}each~{}}I\in\mathit{mods}(D,\Sigma)\\}$. The _answer_ to a
BCQ $q$ w.r.t. $D$ and $\Sigma$ is _positive_ , denoted as $D\cup\Sigma\models
q$, iff $\mathit{ans}(q,D,\Sigma)\neq\varnothing$. Note that query answering
under general TGDs is undecidable [24], even when the schema and the set of
TGDs are fixed [25]. We recall that the two problems of answering CQs and BCQs
under TGDs are equivalent [21, 26]. Roughly speaking, we can enumerate the
polynomially many tuples of constants which are possible answers to $q$, and
then, instead of answering the given query $q$, we answer the polynomially
many BCQs that we obtain by replacing the variables in the body of $q$ with
the appropriate constants. A certain tuple $\mathbf{t}$ of constants is in the
answer of $q$ iff the answer to the BCQ that we obtain from $\mathbf{t}$ is
positive. Henceforth, we thus focus only on the BCQ answering problem.
### 3.3 The TGD Chase
The _chase procedure_ (or simply _chase_) is a fundamental algorithmic tool
introduced for checking implication of dependencies [27], and later for
checking query containment [28]. Informally, the chase is a process of
repairing a database w.r.t. a set of dependencies so that the resulted
database satisfies the dependencies. We shall use the term chase
interchangeably for both the procedure and its result. The chase works on an
instance through the so-called TGD _chase rule_.
TGD Chase Rule: Consider a database $D$ for a schema $\mathcal{R}$, and a TGD
$\sigma:\phi(\mathbf{X},\mathbf{Y})\rightarrow\exists\mathbf{Z}\,\psi(\mathbf{X},\mathbf{Z})$
over $\mathcal{R}$. If $\sigma$ is applicable to $D$, i.e., there exists a
homomorphism $h$ such that $h(\phi(\mathbf{X},\mathbf{Y}))\subseteq D$ then:
_(i)_ define $h^{\prime}\supseteq h$ such that $h^{\prime}(Z_{i})=z_{i}$, for
each $Z_{i}\in\mathbf{Z}$, where $z_{i}\in\Delta_{z}$ is a “fresh” labeled
null not introduced before, and _(ii)_ add to $D$ the set of atoms in
$h^{\prime}(\psi(\mathbf{X},\mathbf{Z}))$, if not already in $D$.
Given a database $D$ and a set of TGDs $\Sigma$, the chase algorithm for $D$
and $\Sigma$ consists of an exhaustive application of the TGD chase rule in a
breadth-first fashion, which leads as result to a (possibly infinite) chase
for $D$ and $\Sigma$, denoted as $\mathit{chase}(D,\Sigma)$. For the formal
definition of the chase algorithm we refer the reader to [8].
The (possibly infinite) chase for $D$ and $\Sigma$ is a _universal model_ of
$D$ w.r.t. $\Sigma$, i.e., for each instance $I\in\mathit{mods}(D,\Sigma)$,
there exists a homomorphism from $\mathit{chase}(D,\Sigma)$ to $I$ [26, 29].
Using this fact it can be shown that $D\cup\Sigma\models q$ iff
$\mathit{chase}(D,\Sigma)\models q$, for every BCQ $q$.
## 4 The Datalog± Family
In this section we present the main Datalog± languages under which query
answering is decidable, and (almost in all cases) also tractable in data
complexity.
### 4.1 Decidability Paradigms
We first discuss the three main paradigms for ensuring decidability of query
answering, namely, chase termination, guardedness and stickiness.
#### Chase Termination.
In this case the chase always terminates and produces a finite universal model
$U$. Thus, given a query we just need to evaluate it over the finite database
$U$. The most notable syntactic restriction of TGDs guaranteeing chase
termination is _weak-acyclicity_ , which is defined by means of a graph-based
condition, for which we refer the reader to the landmark paper [29]. Roughly
speaking, in the chase constructed under a weakly-acyclic set of TGDs over a
schema $\mathcal{R}$, only a finite number of distinct values can appear at
any position of $\mathcal{R}$, and thus after finitely many steps the chase
procedure terminates. It is known that query answering under a weakly-acyclic
set of TGDs is ptime-complete [29] and 2exptime-complete [10] in data and
combined complexity, respectively. More general syntactic restrictions that
guarantee chase termination were proposed in [26] and [30].
#### Guardedness.
_Guarded_ TGDs, introduced and studied in [25], have an atom in their body,
called the _guard_ , that contains all the universally quantified variables.
For example, the TGD $r(X,Y),s(X,Y,Z)\rightarrow\exists Ws(Z,X,W)$ is guarded
via the guard atom $s(X,Y,Z)$, while the TGD $r(X,Y),r(Y,Z)\rightarrow r(X,Z)$
is not. Decidability of query answering follows from the fact that the chase
constructed under a set of guarded TGDs has the bounded treewidth property,
i.e., is a “tree-like” structure. The data and combined complexity of query
answering under a set of guarded TGDs is ptime-complete [7] and 2exptime-
complete [25], respectively.
_Linear_ TGDs, proposed in [7], is a FO-rewritable variant of guarded TGDs. A
TGD is linear iff it contains only one atom in its body. Obviously a linear
TGD is trivially guarded since the singleton body-atom is automatically a
guard. Linear TGDs are more expressive than the well-known class of inclusion
dependencies. Query answering under linear TGDs is in the highly tractable
class ac0 in data complexity [7]. The same problem is pspace-complete in
combined complexity; this result is immediately implied by results in [28].
An expressive class, which forms a generalization of guarded TGDs, is the
class of _weakly-guarded_ sets of TGDs introduced in [25]. Intuitively
speaking, a set $\Sigma$ of TGDs is weakly-guarded iff in the body of each TGD
of $\Sigma$ there exists an atom, called the _weak-guard_ , that contains all
the universally quantified variables that appear only at positions where a
“fresh” null of $\Delta_{z}$ can appear during the construction of the chase.
Query answering under a weakly-guarded set of TGDs is exptime-complete [25]
and 2exptime-complete [25] in data and combined complexity, respectively.
#### Stickiness.
In this paragraph we present a Datalog± language (and its extensions), which
hinges on a paradigm that is very different from guardedness. _Sticky_ sets of
TGDs are defined formally by an efficiently testable condition involving
variable-marking [9]. In what follows we just give an intuitive definition of
this class. For every database $D$, assume that during the construction of
chase of $D$ under a sticky set of TGDs, we apply a TGD $\sigma\in\Sigma$ that
has a variable $V$ appearing more than once in its body; assume also that $V$
maps (via homomorphism) on the symbol $z$, and that by virtue of this
application the atom $\underline{a}$ is introduced. In this case, for each
atom $\underline{b}$ in $\mathit{body}(\sigma)$, we say that $\underline{a}$
is _derived_ from $\underline{b}$. Then, we have that $z$ appears in
$\underline{a}$ and in all atoms resulting from some chase derivation sequence
starting from $\underline{a}$, “sticking” to them (hence the name “sticky”
sets of TGDs). Interestingly, sticky sets of TGDs are FO-rewritable, and thus
query answering is feasible in ac0 in data complexity [9]. Combined complexity
of query answering is known to be exptime-complete [9].
In [10] the FO-rewritable class of _sticky-join_ sets of TGDs, that captures
both linear TGDs and sticky sets of TGDs, is introduced. Similarly to sticky
sets of TGDs, sticky-join sets are defined formally by a testable condition
based on variable-marking. However, this variable-marking procedure is more
sophisticated than the one used for sticky sets, and due to this fact the
problem of identifying sticky-join sets of TGDs is harder than the one of
identifying sticky sets. In particular, given a set $\Sigma$ of TGDs, we can
decide in ptime whether $\Sigma$ is sticky, while the problem whether $\Sigma$
is sticky-join is pspace-complete. Note that the data and combined complexity
of query answering under sticky and sticky-join sets of TGDs coincide.
### 4.2 Additional Features
In this subsection we briefly discuss how the languages presented above can be
combined with negative constraints and key dependencies, without altering the
complexity of query answering.
#### Negative Constraints.
A _negative constraint_ (NC) $\nu$ over a schema $\mathcal{R}$ is a first-
order formula $\forall{\bf X}\,\phi({\bf X})\rightarrow\bot$, where $\bot$
denotes the truth constant false. NCs are vital when representing ontologies
(see, e.g., [7, 9]), as well as conceptual schemas such as Entity-Relationship
diagrams (see, e.g., [31, 32]). With NCs we can assert, for example, that
students and professors are disjoint sets: $\forall
X\mathit{student}(X),\mathit{professor}(X)\rightarrow\bot$. Also, we can state
that a student cannot be the leader of a research group: $\forall X\forall
Y\mathit{student}(X),\mathit{leads}(X,Y)\rightarrow\bot$.
It is known that checking NCs is tantamount to query answering [7]. In
particular, given an instance $I$, a set $\Sigma_{\bot}$ of NCs, and a set
$\Sigma$ of TGDs, for each NC $\nu$ of the form
$\forall\mathbf{X}\,\phi(\mathbf{X})\rightarrow\bot$, we answer the BCQ
$q_{\nu}()\leftarrow\phi(\mathbf{X})$. If at least one of such queries answers
positively, then $I\cup\Sigma\cup\Sigma_{\bot}\models\bot$ (i.e., the theory
is inconsistent), and therefore $I\cup\Sigma\cup\Sigma_{\bot}\models q$, for
every BCQ $q$; otherwise, given a BCQ $q$, we have
$I\cup\Sigma\cup\Sigma_{\bot}\models q$ iff $I\cup\Sigma\models q$, i.e., we
can answer $q$ by ignoring the set of NCs.
#### Key Dependencies.
It is well-known that the interaction of general TGDs and key dependencies
(KDs) leads to undecidability of query answering [33]; we assume that the
reader is familiar with the notion of KD (see, e.g., [34]). Thus, the classes
of TGDs presented above cannot be combined arbitrarily with KDs. Suitable
syntactic restrictions are needed in order to ensure decidability of query
answering.
A crucial concept towards this direction is separability [35], which
formulates a controlled interaction of TGDs and KDs. Formally speaking, a set
$\Sigma=\Sigma_{T}\cup\Sigma_{K}$ over a schema $\mathcal{R}$, where
$\Sigma_{T}$ and $\Sigma_{K}$ are sets of TGDs and KDs, respectively, is
_separable_ iff for every instance $I$ for $\mathcal{R}$, either $I$ violates
$\Sigma_{K}$, or for every BCQ $q$ over $\mathcal{R}$, $I\cup\Sigma\models q$
iff $I\cup\Sigma_{T}\models q$. Notice that separability is a semantic notion.
A sufficient syntactic criterion for separability of TGDs and KDs is given in
[7]; TGDs and KDs satisfying the criterion are called _non-conflicting_.
Obviously, in case of non-conflicting sets of TGDs and KDs, we just need to
perform a preliminary check whether the given instance satisfies the KDs, and
if this is the case, then we eliminate them, and proceed by considering only
the set of TGDs. This preliminary check can be done using negative
constraints. For example, to check whether the KD $\mathit{key}(r)=\\{1\\}$,
stating that the first attribute of the binary relation $r$ is a key
attribute, is satisfied by the database $D$, we just need to check whether the
database $D_{\neq}$ obtained by adding to $D$ the set of atoms
$\\{\mathit{neq}(a,b)~{}|~{}a\neq
b,\textrm{~{}and~{}}a,b\textrm{~{}are~{}constants~{}occurring~{}in~{}}D\\}$,
where $\mathit{neq}$ is an auxiliary predicate, satisfies the negative
constraint $r(X,Y),r(X,Z),\mathit{neq}(Y,Z)\rightarrow\bot$. The atom
$\mathit{neq}(a,b)$ implies that $a$ and $b$ are different constants. Since,
as already mentioned, checking NCs is tantamount to query answering, we
immediately get that the complexity of query answering under non-conflicting
sets of TGDs and KDs is the same as in the case of TGDs only.
Interestingly, by combining non-conflicting linear (or sticky) sets of TGDs
and KDs with NCs, we get strictly more expressive formalisms than the most
widely-adopted tractable ontology languages, in particular DL-LiteA, DL-LiteF
and DL-LiteR, without loosing FO-rewritability, and consequently high
tractability of query answering in data complexity. For more details, we refer
the interested reader to [7, 9].
## 5 Datalog± for OBDA
In this section we consider the problem of BCQ answering under the FO-
rewritable members of the Datalog± family, namely, linear, sticky and sticky-
join sets of TGDs. Given a BCQ $q$ and a set $\Sigma$ of TGDs, the actual
computation of the rewriting is done by applying a backward-chaining
resolution procedure using the rules of $\Sigma$ as rewriting rules. Our
algorithm optimizes the algorithm presented in [14] by greatly reducing the
number of BCQs in the rewriting, and therefore improves the overall
performance of query answering. Before going into the details of the rewriting
algorithm, we first give some useful notions.
A set of atoms $A=\\{\underline{a}_{1},\ldots,\underline{a}_{n}\\}$, where
$n\geqslant 2$, _unifies_ if there exists a substitution $\gamma$, called
_unifier_ for $A$, such that
$\gamma(\underline{a}_{1})=\ldots=\gamma(\underline{a}_{n})$. A _most general
unifier (MGU)_ for $A$ is a unifier for $A$, denoted as $\gamma_{A}$, such
that for each other unifier $\gamma$ for $A$, there exists a substitution
$\gamma^{\prime}$ such that $\gamma=\gamma^{\prime}\circ\gamma_{A}$. Notice
that if a set of atoms unify, then there exists a MGU. Furthermore, the MGU
for a set of atoms is unique (modulo variable renaming). The MGU for a
singleton set $\\{\underline{a}\\}$ is defined as the identity substitution on
the set of terms that occur in $\underline{a}$.
Let us now give some auxiliary results which will allow us to simplify our
later technical definitions and proofs. The first such lemma states that we
can restrict our attention on TGDs that have only one head-atom.
###### Lemma 1
BCQ answering under (general) TGDs and BCQ answering under TGDs with just one
head-atom are logspace-equivalent problems.
Proof. It suffices to show that BCQ answering under (general) TGDs can be
reduced in logspace to BCQ answering under TGDs with just one head-atom.
Consider a BCQ $q$ over a schema $\mathcal{R}$, a database $D$ for
$\mathcal{R}$, and a set $\Sigma$ of TGDs over $\mathcal{R}$. We construct
$\Sigma^{\prime}$ from $\Sigma$ by applying the following procedure. For each
TGD $\sigma\in\Sigma$, where
$\mathit{head}(\sigma)=\\{\underline{a}_{1},\ldots,\underline{a}_{k}\\}$ and
$\mathbf{X}$ is the set of variables that occur in $\mathit{head}(\sigma)$,
replace $\sigma$ with the following set of TGDs:
$\begin{array}[]{rcl}\mathit{body}(\sigma)&\rightarrow&r_{\sigma}(\mathbf{X}),\\\
r_{\sigma}(\mathbf{X})&\rightarrow&\underline{a}_{1},\\\
r_{\sigma}(\mathbf{X})&\rightarrow&\underline{a}_{2},\\\ &\vdots&\\\
r_{\sigma}(\mathbf{X})&\rightarrow&\underline{a}_{k},\end{array}$
where $r_{\sigma}$ is an auxiliary predicate not occurring in $\mathcal{R}$
having the same arity as the number of variables in $\mathbf{X}$. It is not
difficult to see that the above construction is feasible in logspace. By
construction, except for the atoms with an auxiliary predicate,
$\mathit{chase}(D,\Sigma)$ and $\mathit{chase}(D,\Sigma^{\prime})$ coincide.
The auxiliary predicates, being introduced only during the above
transformation, do not match any predicate symbol in $q$, and hence
$\mathit{chase}(D,\Sigma)\models q$ iff
$\mathit{chase}(D,\Sigma^{\prime})\models q$, or, equivalently,
$D\cup\Sigma\models q$ iff $D\cup\Sigma^{\prime}\models q^{\prime}$.
The next lemma implies that we can restrict our attention on TGDs that have
only one existentially quantified variable which occurs only once.
###### Lemma 2
BCQ answering under (general) TGDs and BCQ answering under TGDs with at most
one existentially quantified variable that occurs only once are logspace-
equivalent problems.
Proof. It suffices to show that BCQ answering under (general) TGDs can be
reduced in logspace to BCQ answering under TGDs that have at most one
existentially quantified variable which occurs only once. Consider a BCQ $q$
over a schema $\mathcal{R}$, a database $D$ for $\mathcal{R}$, and a set
$\Sigma$ of TGDs over $\mathcal{R}$. We construct $\Sigma^{\prime}$ from
$\Sigma$ by applying the following procedure. For each TGD $\sigma\in\Sigma$,
where $\\{X_{1},\ldots,X_{n}\\}$, for $n\geqslant 1$, is the set of variables
that occur both in $\mathit{body}(\sigma)$ and $\mathit{head}(\sigma)$, and
$\\{Z_{1},\ldots,Z_{m}\\}$, for $m>1$, is the set of the existentially
quantified variables of $\sigma$, replace $\sigma$ with the following set of
TGDs:
$\begin{array}[]{rcl}\mathit{body}(\sigma)&\rightarrow&\exists
Z_{1}\,r_{\sigma}^{1}(X_{1},\ldots,X_{n},Z_{1}),\\\
r_{\sigma}^{1}(X_{1},\ldots,X_{n},Z_{1})&\rightarrow&\exists
Z_{2}\,r_{\sigma}^{2}(X_{1},\ldots,X_{n},Z_{1},Z_{2}),\\\ &\vdots&\\\
r_{\sigma}^{m-1}(X_{1},\ldots,X_{n},Z_{1},\ldots,Z_{m-1})&\rightarrow&\exists
Z_{m}\,r_{\sigma}^{m}(X_{1},\ldots,X_{n},Z_{1},\ldots,Z_{m}),\\\
r_{\sigma}^{m}(X_{1},\ldots,X_{n},Z_{1},\ldots,Z_{m})&\rightarrow&\mathit{head}(\sigma),\end{array}$
where $r_{\sigma}^{i}$ is an auxiliary predicate of arity $n+i$, for each
$i\in[m]$. It is easy to see that the above procedure can be carried out in
logspace. By construction, except for the atoms with an auxiliary predicate,
$\mathit{chase}(D,\Sigma)$ and $\mathit{chase}(D,\Sigma^{\prime})$ are the
same (modulo bijective variable renaming). The auxiliary predicates, being
introduced only during the above construction, do not match any predicate
symbol in $q$, and hence $\mathit{chase}(D,\Sigma)\models q$ iff
$\mathit{chase}(D,\Sigma^{\prime})\models q$, or, equivalently,
$D\cup\Sigma\models q$ iff $D\cup\Sigma^{\prime}\models q$.
Since the transformations given above preserve the syntactic condition of
linear, sticky and sticky-join sets of TGDs, henceforth we assume w.l.o.g.
that every TGD has just one atom in its head which contains only one
existentially quantified variable that occurs only once. In the rest of the
paper, for notational convenience, given a TGD $\sigma$, we denote by
$\pi_{\sigma}$ the position in $\mathit{head}(\sigma)$ at which the
existentially quantified variable occurs.
We now give the notion of _applicability_ of a TGD to a set of body-atoms of a
query. Let us assume w.l.o.g that the variables that appear in the query, and
those that appear in the TGD, constitute two disjoint sets. Given a BCQ $q$, a
variable is called _shared_ in $q$ if it occurs more than once in
$\mathit{body}(q)$. Notice that in the case of (non-Boolean) CQs, a variable
is shared in $q$ if it occurs more than once in $q$ (considering also the head
of $q$ and not just its body).
###### Definition 1 (Applicability)
Consider a BCQ $q$ over a schema $\mathcal{R}$, and a TGD $\sigma$ over
$\mathcal{R}$. Given a set of atoms $A\subseteq\mathit{body}(q)$ that unifies,
we say that $\sigma$ is _applicable_ to $A$ if the following conditions are
satisfied: (i) the set $A\cup\\{\mathit{head}(\sigma)\\}$ unifies, and (ii)
for each $\underline{a}\in A$, if the term at position $\pi$ in
$\underline{a}$ is either a constant or a shared variable in $q$, then
$\pi\neq\pi_{\sigma}$.
Let us now introduce the notion of _factorizability_ which, as we explain
below, makes one of the main differences between our algorithm and the one
presented in [14], due to which a perfect rewriting with less BCQs is
obtained.
###### Definition 2 (Factorizability)
Consider a BCQ $q$ over a schema $\mathcal{R}$, and a TGD $\sigma$ over
$\mathcal{R}$ which contains an existentially quantified variable. A set of
atoms $A\subseteq\mathit{body}(q)$, where $|A|\geqslant 2$, that unifies is
_factorizable_ w.r.t. $\sigma$ if there exists a variable $V$ that occurs in
every atom of $S$ only at position $\pi_{\sigma}$, and also $V$ does not occur
in $\mathit{body}(q)\setminus S$.
It is important to clarify that in the case of (non-Boolean) CQs, the notion
of factorizability is defined as above, except that the variable $V$ does not
occur in $(\\{\mathit{head}(\sigma)\\}\cup\mathit{body}(\sigma))\setminus S$.
###### Example 1 (Factorization)
Consider the BCQs
$\begin{array}[]{rcl}q_{1}&:&q()\,\leftarrow\,\underbrace{t(A,B,C),t(A,E,C)}_{S_{1}}\\\
q_{2}&:&q()\,\leftarrow\,s(C),\underbrace{t(A,B,C),t(A,E,C)}_{S_{2}}\\\
q_{3}&:&q()\,\leftarrow\,\underbrace{t(A,B,C),t(A,C,C)}_{S_{3}}\end{array}$
and the TGD $\sigma:s(X),r(X,Y)\,\rightarrow\,\exists Z\,t(X,Y,Z).$ Clearly,
$S_{1}$ is factorizable w.r.t. $\sigma$ since the substitution
$\\{E\rightarrow B\\}$ is a unifier for $S_{1}$, and also $C$ appears in both
atoms of $S_{1}$ only at position $\pi_{\sigma}$. The factorization results in
the query $q()\leftarrow t(A,B,C)$; notice that $\sigma$ is not applicable to
$S_{1}$, but it is applicable to $\\{t(A,B,C)\\}$. On the contrary, despite
the fact that $S_{2}$ unifies, it is not factorizable w.r.t. $\sigma$ since
$C$ occurs also in $\mathit{body}(q_{2})\setminus S_{2}$. Finally, even if
$S_{3}$ unifies, it is not factorizable w.r.t. $\sigma$ since $C$ appears in
$S_{3}$, not only at position $\pi_{\sigma}$, but also at position $t[2]$.
We are now ready to describe the algorithm TGD-rewrite, depicted in Algorithm
1, which is based on the rewriting algorithm presented in [14]. The perfect
rewriting of a BCQ $q$ w.r.t. a set of TGDs $\Sigma$ is computed by
exhaustively applying (i.e., until a fixpoint is reached) two steps:
_factorization_ and _rewriting_.
Input: a BCQ $q$ over a schema $\mathcal{R}$, a set $\Sigma$ of TGDs over
$\mathcal{R}$
Output: the FO-rewriting $Q_{\textsc{fin}}$ of $q$ w.r.t. $\Sigma$
$Q_{\textsc{rew}}:=\\{\langle q,1\rangle\\}$;
repeat
$Q_{\textsc{temp}}:=Q_{\textsc{rew}}$;
foreach _$\\{\langle q,x\rangle\\}\in Q_{\textsc{temp}}$ , where
$x\in\\{0,1\\}$,_ do
/* factorization step */
foreach _$\sigma\in\Sigma$_ do
$q^{\prime}:=\mathit{factorize}(q,\sigma)$;
if _$\mathit{notExists}(\langle q^{\prime},y\rangle,Q_{\textsc{rew}})$ ,
where $y\in\\{0,1\\}$,_ then
$Q_{\textsc{rew}}:=Q_{\textsc{rew}}\cup\\{\langle q^{\prime},0\rangle\\}$;
/* rewriting step */
foreach _$A\subseteq\mathit{body}(q)$_ do
foreach _$\sigma\in\Sigma$_ do
if _$\mathit{isApplicable}(\sigma,A,q)$_ then
$q^{\prime}:=\gamma_{A\cup\\{\mathit{head}(\sigma)\\}}(q[A/\mathit{body}(\sigma)])$;
if _$\mathit{notExists}(\langle q^{\prime},1\rangle,Q_{\textsc{rew}})$_ then
$Q_{\textsc{rew}}:=Q_{\textsc{rew}}\cup\\{\langle q^{\prime},1\rangle\\}$;
until _$Q_{\textsc{temp}}=Q_{\textsc{rew}}$_ ;
$Q_{\textsc{fin}}:=\\{q~{}|~{}\langle q,x\rangle\in
Q_{\textsc{rew}}\textrm{~{}and~{}}x=1\\}$;
return _$Q_{\textsc{fin}}$_
Algorithm 1 The algorithm TGD-rewrite
Factorization Step. The function $\mathit{factorize}(q,\sigma)$, providing
that there exists a subset of $\mathit{body}(q)$ which is factorizable w.r.t.
$\sigma$ (otherwise, the query $q$ is returned), first selects such a set
$S\subseteq\mathit{body}(q)$. Then, the query $q^{\prime}$ is constructed by
applying the MGU $\gamma_{S}$ for $S$ on $q$. Providing that there is no pair
$\langle q^{\prime\prime},y\rangle$, where $y\in\\{0,1\\}$, in
$Q_{\textsc{rew}}$ such that $q^{\prime}$ and $q^{\prime\prime}$ are the same
(modulo bijective variable renaming), the pair $\langle q^{\prime},0\rangle$
is added to $Q_{\textsc{rew}}$; the label $0$ keeps track of the queries
generated by the factorization step that must be excluded from the final
rewriting. This is carried out by the $\mathit{notExists}$ function.
Rewriting Step. If there exists a pair $\langle q,y\rangle$ and a TGD
$\sigma\in\Sigma$ which is applicable to a set of atoms
$A\subseteq\mathit{body}(q)$, then the algorithm constructs a new query
$q^{\prime}=\gamma_{A\cup\\{\mathit{head}(\sigma)\\}}(q[A/\mathit{body}(\sigma)])$,
that is, the BCQ obtained from $q$ by replacing $A$ with
$\mathit{body}(\sigma)$ and then applying the MGU for the set
$A\cup\\{\mathit{head}(\sigma)\\}$. Providing that there is no pair $\langle
q^{\prime\prime},1\rangle$ in $Q_{\textsc{rew}}$ such that $q^{\prime}$ and
$q^{\prime\prime}$ are the same (modulo bijective variable renaming), the pair
$\langle q^{\prime},1\rangle$ is added to $Q_{\textsc{rew}}$; the label $1$
keeps track of the queries generated by the rewriting step which will be the
final rewriting.
###### Example 2 (Rewriting)
Consider the set $\Sigma$ of TGDs
$\begin{array}[]{rcl}\sigma_{1}&:&s(X)\,\rightarrow\exists Z\ \,t(X,X,Z)\\\
\sigma_{2}&:&t(X,Y,Z)\,\rightarrow\,r(Y,Z)\end{array}$
and the query $q()\leftarrow t(A,B,C),r(B,C).$ TGD-rewrite first applies
$\sigma_{2}$ to $\\{r(B,C)\\}$ since $\sigma_{1}$ is not applicable. The query
$q_{1}:q()\leftarrow t(A,B,C),t(V^{1},B,C)$ is produced. Clearly,
$\mathit{body}(q_{1})$ is factorizable w.r.t. $\sigma_{1}$ and the query
$q_{2}:q()\leftarrow t(A,B,C)$ is obtained. Now, $\sigma_{1}$ is applicable to
$\\{t(A,B,C)\\}$ and the query $q_{3}:q()\leftarrow s(A)$ is obtained. The
perfect rewriting constructed by the algorithm is the set
$\\{q,q_{1},q_{3}\\}$.
The next example shows that dropping the applicability condition, then TGD-
rewrite may produce unsound rewritings.
###### Example 3 (Loss of soundness)
Suppose that we ignore the applicability condition during the rewriting
process. Consider the set $\Sigma$ of TGDs given in Example 2, and also the
BCQ $q_{1}:q()\leftarrow t(A,B,c)$, where $c$ is a constant of $\Delta_{c}$. A
BCQ $q^{\prime}$ of the form $q()\leftarrow s(V)$ is obtained, where the
information about the constant $c$ is lost. Consider now the database
$D=\\{s(b),t(a,b,d)\\}$ for $\mathcal{R}$. The query $q^{\prime}$ maps to the
atom $s(b)$ which implies that $D\models q^{\prime}$. However, the original
query $q$ does not map to $\mathit{chase}(D,\Sigma)$, and thus
$D\cup\Sigma\not\models q$. Therefore, any rewriting containing $q^{\prime}$
is not a sound rewriting of $q$ given $\Sigma$. Consider now the query
$q^{\prime\prime}:q()\leftarrow t(A,B,B)$. The same query $q^{\prime}$ mapping
to the atom $s(b)$ of $D$ is obtained. However, during the construction of
$\mathit{chase}(D,\Sigma)$ it is not possible to get an atom of the form
$t(X,Y,Y)$, where at positions $t[2]$ and $t[3]$ the same value occurs. This
implies that there is no homomorphism that maps $q$ to
$\mathit{chase}(D,\Sigma)$, and hence $D\cup\Sigma\not\models q$. Therefore,
any rewriting containing $q^{\prime}$ is again unsound.
The applicability condition may prevent the generation of queries that are
vital to guarantee completeness of the rewritten query, as shown by the
following example. This is exactly the reason why the factorization step is
also needed.
###### Example 4 (Loss of completeness)
Consider the set $\Sigma$ of TGDs
$\begin{array}[]{rcl}\sigma_{1}&:&p(X)\,\rightarrow\,\exists Y\,t(X,Y)\\\
\sigma_{2}&:&t(X,Y)\,\rightarrow\,s(Y)\end{array}$
and the query $q:q()\leftarrow t(A,B),s(B).$ The only viable strategy in this
case is to apply $\sigma_{2}$ to $\\{s(B)\\}$, since $\sigma_{1}$ is not
applicable to $\\{t(A,B)\\}$ due to the shared variable $B$. The query that we
obtain is $q^{\prime}:q()\leftarrow t(A,B),t(V^{1},B)$, where $V^{1}$ is a
fresh variable. Notice that in $q^{\prime}$ the variable $B$ remains shared
thus it is not possible to apply $\sigma_{1}$. It is obvious that without the
factorization step there is no way to obtain the query
$q^{\prime\prime}:q()\leftarrow p(A)$ during the rewriting process. Now,
consider the database $D=\\{p(a)\\}$. Clearly,
$\mathit{chase}(D,\Sigma)=\\{p(a),t(a,z_{1}),s(z_{1})\\}$, and therefore
$\mathit{chase}(D,\Sigma)\models q$, or, equivalently, $D\cup\Sigma\models q$.
However, the rewritten query is not entailed by the given database $D$, since
$q^{\prime\prime}$ does not belong to it, which implies that it is not
complete.
We proceed now to establish soundness and completeness of the proposed
algorithm. Towards this aim we need two auxiliary technical lemmas. The first
one, which is needed for soundness, states that once the chase entails the
rewritten query constructed by the rewriting algorithm, then the chase entails
also the given query. In the sequel, for brevity, given a BCQ $q$ over a
schema $\mathcal{R}$ and a set $\Sigma$ of TGDs over $\mathcal{R}$, we denote
by $q_{\Sigma}$ the rewritten query $\textsf{TGD-rewrite}(q,\Sigma)$.
###### Lemma 3
Consider a BCQ $q$ over a schema $\mathcal{R}$, a database $D$ for
$\mathcal{R}$, and a set $\Sigma$ of TGDs over $\mathcal{R}$. If
$\mathit{chase}(D,\Sigma)\models q_{\Sigma}$, then
$\mathit{chase}(D,\Sigma)\models q$.
Proof. The proof is by induction on the number of applications of the
rewriting step. We denote by $q_{\Sigma}^{[i]}$ the part of the rewritten
query $q_{\Sigma}$ obtained by applying $i$ times the rewriting step.
Base Step. Clearly, $q_{\Sigma}^{0}=q_{\Sigma}$, and the claim holds
trivially.
Inductive Step. Suppose now that $\mathit{chase}(D,\Sigma)\models
q_{\Sigma}^{[i]}$, for $i\geq 0$. This implies that there exists $p\in
q_{\Sigma}^{[i]}$ such that $\mathit{chase}(D,\Sigma)\models p$, and thus
there exists a homomorphism $h$ such that
$h(\mathit{body}(p))\subseteq\mathit{chase}(D,\Sigma)$. If $p\in
q_{\Sigma}^{[i-1]}$, then the claim follows by induction hypothesis. The
interesting case is when $p$ was obtained during the $i$-th application of the
rewriting step from a BCQ $p^{\prime}\in q_{\Sigma}^{[i-1]}$, i.e.,
$q_{\Sigma}^{[i]}=q_{\Sigma}^{[i-1]}\cup\\{p\\}$. By induction hypothesis, it
suffices to show that $\mathit{chase}(D,\Sigma)\models q_{\Sigma}^{[i-1]}$.
Clearly, there exists a TGD $\sigma\in\Sigma$ of the form
$\phi(\mathbf{X},\mathbf{Y})\rightarrow\exists Z\,r(\mathbf{X},Z)$ which is
applicable to a set $A\subseteq\mathit{body}(p^{\prime})$, and $p$ is such
that $\mathit{body}(p)=\gamma(p^{\prime}[A/\phi(\mathbf{X},\mathbf{Y})])$,
where $\gamma$ is the MGU for the set $A\cup\\{\mathit{head}(\sigma)\\}$.
Observe that
$h(\gamma(\phi(\mathbf{X},\mathbf{Y})))\subseteq\mathit{chase}(D,\Sigma)$, and
hence $\sigma$ is applicable to $\mathit{chase}(D,\Sigma)$; let
$\mu=h\circ\gamma$. Thus,
$\mu^{\prime}(r(\mathbf{X},Z))\in\mathit{chase}(D,\Sigma)$, where
$\mu^{\prime}\supset\mu$. We define the substitution
$h^{\prime}=h\cup\\{\gamma(Z)\rightarrow\mu^{\prime}(Z)\\}$.
Let us first show that $h^{\prime}$ is a well-defined substitution. It
suffices to show that $\gamma(Z)$ is not a constant, and also that $\gamma(Z)$
does not appear in the left-hand side of an assertion of $h$. Towards a
contradiction, suppose that $\gamma(Z)$ is either a constant or appears in the
left-hand side of an assertion of $h$. It is easy to verify that in this case
there exists an atom $\underline{a}\in A$ such that at position $\pi_{\sigma}$
in $\underline{a}$ occurs either a constant or a variable which is shared in
$p^{\prime}$. But this contradicts the fact that $\sigma$ is applicable to
$A$. Consequently, $h^{\prime}$ is well-defined. It remains to show that the
substitution $h^{\prime}\circ\gamma$ maps $\mathit{body}(p^{\prime})$ to
$\mathit{chase}(D,\Sigma)$, and thus $\mathit{chase}(D,\Sigma)\models
q_{\Sigma}^{[i-1]}$. Clearly, $\gamma(\mathit{body}(p^{\prime})\setminus
A)\subseteq\mathit{body}(p)$. Since
$h(\mathit{body}(p))\subseteq\mathit{chase}(D,\Sigma)$, we get that
$h^{\prime}(\gamma(\mathit{body}(p^{\prime})\setminus
A))\subseteq\mathit{chase}(D,\Sigma)$. Moreover,
$\begin{array}[]{rcl}h^{\prime}(\gamma(A))&=&h^{\prime}(\gamma(r(\mathbf{X},Z)))\\\
&=&r(h^{\prime}(\gamma(\mathbf{X})),h^{\prime}(\gamma(Z)))\\\
&=&r(\mu(\mathbf{X}),\mu^{\prime}(Z))\\\ &=&\mu^{\prime}(r(\mathbf{X},Z))\\\
&\in&\mathit{chase}(D,\Sigma).\end{array}$
The proof is now complete.
The second auxiliary lemma, which is needed for completeness, asserts that
once the chase entails the rewritten query constructed by the rewriting
algorithm, then the given database also entails the rewritten query.
###### Lemma 4
Consider a BCQ $q$ over a schema $\mathcal{R}$, a database $D$ for
$\mathcal{R}$, and a set $\Sigma$ of TGDs over $\mathcal{R}$. If
$\mathit{chase}(D,\Sigma)\models q_{\Sigma}$, then $D\models q_{\Sigma}$.
Proof. We proceed by induction on the number of applications of the chase
step.
Base Step. Clearly, $\mathit{chase}^{[0]}(D,\Sigma)=D$, and the claim holds
trivially.
Inductive Step. Suppose now that $\mathit{chase}^{[i]}(D,\Sigma)\models
q_{\Sigma}$, for $i\geq 0$. This implies that there exists $p\in q_{\Sigma}$
such that $\mathit{chase}^{[i]}(D,\Sigma)\models p$, and thus there exists a
homomorphism $h$ such that
$h(\mathit{body}(p))\subseteq\mathit{chase}^{[i]}(D,\Sigma)$. If
$h(\mathit{body}(p))\subseteq\mathit{chase}^{[i-1]}(D,\Sigma)$, then the claim
follows by induction hypothesis. The non-trivial case is when the atom
$\underline{a}$, obtained during the $i$-th application of the chase step due
to a TGD $\sigma\in\Sigma$ of the form
$\phi(\mathbf{X},\mathbf{Y})\rightarrow\exists Z\,r(\mathbf{X},Z)$, belongs to
$h(\mathit{body}(p))$. Clearly, there exists a homomorphism $\mu$ such that
$\mu(\phi(\mathbf{X},\mathbf{Y}))\subseteq\mathit{chase}^{[i-1]}(D,\Sigma)$
and $\underline{a}=\mu^{\prime}(r(\mathbf{X},\mathbf{Y}))$, where
$\mu^{\prime}\supseteq\mu$. By induction hypothesis, it suffices to show that
$\mathit{chase}^{[i-1]}(D,\Sigma)\models q_{\Sigma}$. Before we proceed
further, we need to establish an auxiliary technical claim.
###### Claim 5
There exists a BCQ $p^{\prime}\in q_{\Sigma}$ and a set of atoms
$A\subseteq\mathit{body}(p^{\prime})$ such that $\sigma$ is applicable to $A$,
and also there exists a homomorphism $\lambda$ such that
$\lambda(\mathit{body}(p^{\prime})\setminus
A)\subseteq\mathit{chase}^{[i-1]}(D,\Sigma)$ and $\lambda(A)=\underline{a}$.
Proof. Clearly, there exists a set of atoms $B$ such that
$h(\mathit{body}(p)\setminus B)\subseteq\mathit{chase}^{[i-1]}(D,\Sigma)$ and
$h(B)=\underline{a}$. Observe that the null value that occurs in
$\underline{a}$ at position $\pi_{\sigma}$ does not occur in
$\mathit{chase}^{[i-1]}(D,\Sigma)$ or in $\underline{a}$ at some position
other than $\pi_{\sigma}$. Therefore, the variables that occur in the atoms of
$B$ at $\pi_{\sigma}$ do not appear at some other position. Consequently, $B$
can be partitioned into the sets $B_{1},\ldots,B_{m}$, where $m\geq 1$, and
the following holds: for each $i\in[m]$, in the atoms of $B_{i}$ at position
$\pi_{\sigma}$ the same variable $V_{i}$ occurs, and also $V_{i}$ does not
occur in some other set $B\in\\{B_{1},\ldots,B_{m}\\}\setminus\\{B_{i}\\}$ or
in $B_{i}$ at some position other than $\pi_{\sigma}$. It is easy to verify
that each set $B_{i}$ is factorizable w.r.t. $\sigma$.
Suppose that we factorize $B_{1}$. Then, the query $p_{1}=\gamma_{1}(p)$,
where $\gamma_{1}$ is the MGU for $B_{1}$, is obtained. Observe that $h$ is a
unifier for $B_{1}$. By definition of the MGU, there exists a substitution
$\theta_{1}$ such that $h=\theta_{1}\circ\gamma_{1}$. Clearly,
$\begin{array}[]{rcl}\theta_{1}(\mathit{body}(p_{1})\setminus\gamma_{1}(B))&=&\theta_{1}(\gamma_{1}(\mathit{body}(p))\setminus\gamma_{1}(B))\\\
&=&h(\mathit{body}(p)\setminus B)\\\
&\subseteq&\mathit{chase}^{[i-1]}(D,\Sigma),\end{array}$
and $\theta_{1}(\gamma_{1}(B))=h(B)=\underline{a}$.
Now, observe that the set $\gamma_{1}(B_{2})\subseteq\mathit{body}(p_{1})$ is
factorizable w.r.t. $\sigma$. By applying factorization we get the query
$p_{2}=\gamma_{2}(p_{1})$, where $\gamma_{2}$ is the MGU for
$\gamma_{1}(B_{2})$. Since $\theta_{1}$ is a unifier for $\gamma_{1}(B_{2})$,
there exists a substitution $\theta_{2}$ such that
$\theta_{1}=\theta_{2}\circ\gamma_{2}$. Clearly,
$\begin{array}[]{rcl}\theta_{2}(\mathit{body}(p_{2})\setminus\gamma_{2}(\gamma_{1}(B)))&=&\theta_{2}(\gamma_{2}(\mathit{body}(p_{1}))\setminus\gamma_{2}(\gamma_{1}(B)))\\\
&=&\theta_{1}(\gamma_{1}(\mathit{body}(p))\setminus\gamma_{1}(B))\\\
&=&h(\mathit{body}(p)\setminus B)\\\
&\subseteq&\mathit{chase}^{[i-1]}(D,\Sigma),\end{array}$
and
$\theta_{2}(\gamma_{2}(\gamma_{1}(B)))=\theta_{1}(\gamma_{1}(B))=h(B)=\underline{a}$.
Eventually, by applying the factorization step as above, we will get the BCQ
$p_{m}\ =\ \gamma_{m}\circ\ldots\circ\gamma_{1}(p),$
where $\gamma_{j}$ is the MGU for the set
$\gamma_{j-1}\circ\ldots\circ\gamma_{1}(B_{j})$, for $j\in\\{2,\ldots,m\\}$
(recall that $\gamma_{1}$ is the MGU for $B_{1}$), such that
$\theta_{m}(\mathit{body}(p_{m})\setminus\gamma_{m}\circ\ldots\circ\gamma_{1}(B))\subseteq\mathit{chase}^{[i-1]}(D,\Sigma)$
and $\theta_{m}(\gamma_{m}\circ\ldots\circ\gamma_{1}(B))=\underline{a}$.
It is easy to verify that $\sigma$ is applicable to $A$. The claim follows
with $p^{\prime}=p_{m}$, $A=\gamma_{m}\circ\ldots\circ\gamma_{1}(B)$ and
$\lambda=\theta_{m}$.
The above claim implies that during the rewriting process eventually we will
get a BCQ $p^{\prime\prime}$ such that
$\mathit{body}(p^{\prime\prime})=\gamma(\mathit{body}(p^{\prime})\setminus
A)\cup\gamma(\phi(\mathbf{X},\mathbf{Y}))$, where $\gamma$ is the MGU for the
set $A\cup\\{\mathit{head}(\sigma)\\}$. It remains to show that there exists a
homomorphism that maps $\mathit{body}(p^{\prime\prime})$ to
$\mathit{chase}^{[i-1]}(D,\Sigma)$. Since $\lambda\cup\mu^{\prime}$ is a well-
defined substitution, we get that $\lambda\cup\mu^{\prime}$ is a unifier for
$A\cup\\{\mathit{head}(\sigma)\\}$. By definition of the MGU, there exists a
substitution $\theta$ such that $\lambda\cup\mu^{\prime}=\theta\circ\gamma$.
Observe that
$\begin{array}[]{rcl}\theta(\mathit{body}(p^{\prime\prime}))&=&\theta(\gamma(\mathit{body}(p^{\prime})\setminus
A)\cup\gamma(\phi(\mathbf{X},\mathbf{Y})))\\\
&=&(\lambda\cup\mu^{\prime})(\mathit{body}(p^{\prime})\setminus
A)\cup(\lambda\cup\mu^{\prime})(\phi(\mathbf{X},\mathbf{Y}))\\\
&=&\lambda(\mathit{body}(p^{\prime})\setminus
A)\cup\mu^{\prime}(\phi(\mathbf{X},\mathbf{Y}))\\\
&\subseteq&\mathit{chase}^{[i-1]}(D,\Sigma).\end{array}$
Consequently, the desired homomorphism is $\theta$, and the claim follows.
We are now ready to establish soundness and completeness of the algorithm TGD-
rewrite.
###### Theorem 6
Consider a BCQ $q$ over a schema $\mathcal{R}$, a database $D$ for
$\mathcal{R}$, and a set $\Sigma$ of TGDs over $\mathcal{R}$. It holds that,
$D\models q_{\Sigma}$ iff $D\cup\Sigma\models q$.
Proof. Suppose first that $D\models q_{\Sigma}$. Since
$D\subseteq\mathit{chase}(D,\Sigma)$, we get that
$\mathit{chase}(D,\Sigma)\models q_{\Sigma}$, and the claim follows by Lemma
3. Suppose now that $D\cup\Sigma\models q_{\Sigma}$. Since $q\in q_{\Sigma}$,
we get that $\mathit{chase}(D,\Sigma)\models q_{\Sigma}$, and the claim
follows by Lemma 4.
Notice that the above result holds for arbitrary TGDs. However, termination of
TGD-rewrite is guaranteed if we consider linear, sticky or sticky-join sets of
TGDs since, during the rewriting process, only finitely many queries (modulo
bijective variable renaming) are generated.
###### Theorem 7
The algorithm TGD-rewrite terminates under linear, sticky or sticky-join sets
of TGDs.
Approaches such as those of [5] and [14] resort to exhaustive factorizations
of the atoms in the queries generated by the rewriting algorithm. By
factorizing a query $q$ we obtain a subquery $q^{\prime}$, that is, $q$
implies $q^{\prime}$ (w.r.t. the given set of TGDs). Observe that by computing
the factorized query $q^{\prime}$ we eliminate unnecessary shared variables,
in the body of $q$, due to which the applicability condition is violated.
Consider for example the query $q^{\prime}$ of Example 4. By factorizing the
body of $q^{\prime}$ we obtain the query $q()\leftarrow t(A,B)$ which is a
subquery (w.r.t. to the given set $\Sigma$ of TGDs) of $q^{\prime}$ (in this
case equivalent to $q^{\prime}$), where the variable $B$ is no longer shared.
Thus, the rewriting step can now apply $\sigma_{1}$ to $\\{t(A,B)\\}$ and
produce the query $q()\leftarrow p(A)$ which is needed to ensure completeness.
The exhaustive factorization produces a non-negligible number of redundant
queries as demonstrated by the simple example above. It is thus necessary to
apply a restricted form of factorization that generates a possibly small
number of BCQs that are necessary to guarantee completeness of the rewritten
query. This corresponds to the identification of all the atoms in the query
whose shared existential variables come from the same atom in the chase, and
they can be thus unified with no loss of information. The key principle behind
our factorization process is that, in order to be applied, there must exist a
TGD that can be applied to the output of the factorization.
### 5.1 Exploiting Negative Constraints
It is well-known that negative constraints (NCs) of the form
$\forall\mathbf{X}\,\phi(\mathbf{X})\rightarrow\bot$ are vital for
representing ontologies. As already explained in Subsection 4.2, given a
database $D$ for a schema $\mathcal{R}$, a set $\Sigma$ of TGDs over
$\mathcal{R}$, and a set $\Sigma_{\bot}$ of NCs over $\mathcal{R}$, once the
theory $D\cup\Sigma\cup\Sigma_{\bot}$ is consistent, then we are allowed to
ignore the NCs since, for every BCQ $q$, $D\cup\Sigma\cup\Sigma_{\bot}\models
q$ iff $D\cup\Sigma\models q$. However, as shown in the following example, by
exploiting the given set of NCs it is possible to further reduce the size of
the final rewriting.
###### Example 5
Consider the TGD $\sigma:t(X),s(Y)\rightarrow\exists Z\,p(Y,Z)$, the NC
$\nu:r(X,Y),s(Y)\rightarrow\bot$, and the BCQ $q()\leftarrow r(A,B),p(B,C)$.
Clearly, due to the rewriting step, the query $p:q()\leftarrow
r(A,B),t(V^{1}),s(B)$ is obtained during the rewriting process. However, this
query is not really needed since, for any database $D$ for $\mathcal{R}$,
$D\not\models p$; otherwise, $D$ violates the NC $\nu$ which is a
contradiction since we always assume that the theory $D\cup\\{\sigma,\nu\\}$
is consistent.
It is not difficult to show that, given a BCQ $q$, and a set $\Sigma$ of TGDs,
if a query $p\in q_{\Sigma}$ is not entailed by $\mathit{chase}(D,\Sigma)$,
for an arbitrary database $D$, then any query $p^{\prime}\in q_{\Sigma}$
obtained during the rewriting process starting from $p$, also it is not
entailed by $\mathit{chase}(D,\Sigma)$. Assume now that the set
$\Sigma_{\bot}$ of NCs is part of the input. If we obtain a query $p\in
q_{\Sigma}$ such that there exists a homomorphism that maps
$\mathit{body}(\nu)$, for some NC $\nu\in\Sigma_{\bot}$, to
$\mathit{body}(p)$, then we can safely ignore $p$ since
$\mathit{chase}(D,\Sigma)$ does not entail $p$.
From the above informal discussion, we conclude that we can further reduce the
size of the final rewriting by modifying our algorithm as follows. During the
execution of the rewriting algorithm TGD-rewrite (see Algorithm 1), after the
factorization step (resp., rewriting step) we check whether there exists a
homomorphism that maps $\mathit{body}(\nu)$, for some NC $\nu$ of the given
set of NCs, to the body of the generated query $q^{\prime}$. If there exists
such a homomorphism, then the pair $\langle q^{\prime},0\rangle$ (resp.,
$\langle q^{\prime},1\rangle$) is not added to the set $Q_{\textsc{rew}}$.
Furthermore, the pair $\langle q,1\rangle$ is added to $Q_{\textsc{rew}}$ (see
the first line of the algorithm) only if there is no homomorphism that maps
$\mathit{body}(\nu)$, for some NC $\nu$ of the given set of NCs, to
$\mathit{body}(q)$. If there exists such a homomorphism, then the algorithm
terminates and returns the emptyset, which means that
$\mathit{chase}(D,\Sigma)\not\models q$, for every database $D$ for
$\mathcal{R}$.
## 6 Rewriting Optimization
It is common knowledge that the perfect rewriting obtained by applying a
backward-chaining rewriting algorithm (like TGD-rewrite) is, in general, not
very well-suited for execution by a DB engine due to the large number of
queries to be evaluated. In this section we propose a technique, called _query
elimination_ , aiming at optimizing the obtained rewritten query under the
class of linear TGDs. As we shall see, query elimination (which is an
additional step during the execution of the algorithm TGD-rewrite) reduces
_(i)_ the number of BCQs of the perfect rewriting, _(ii)_ the number of atoms
in each query of the rewriting as well as _(iii)_ the number of joins. Note
that in the rest of the paper we restrict our attention on linear TGDs. Recall
that linear TGDs are TGDs with just one atom in their body. Since we also
assume, as explained in the previous section, TGDs with just one atom in their
head, henceforth, when using the term TGD, we shall refer to TGDs with just
one body-atom and one head-atom.
By exploiting the given set of TGDs, it is possible to identify atoms in the
body of a certain query that are logically implied (w.r.t. the given set of
TGDs) by other atoms in the same query. In particular, for each BCQ $q$
obtained by applying the rewriting step of TGD-rewrite, the atoms of
$\mathit{body}(q)$ that are logically implied (w.r.t. the given set of TGDs)
by some other atom of $\mathit{body}(q)$ are eliminated. Roughly speaking, the
elimination of an atom from the body of a query implies the avoidance of the
construction of redundant queries during the rewriting process. Thus, this
step greatly reduces the number of BCQs in the perfect rewriting. Before going
into the details, let us first introduce some necessary technical notions.
###### Definition 3 (Dependency Graph)
Consider a set $\Sigma$ of TGDs over a schema $\mathcal{R}$. The _dependency
graph_ of $\Sigma$ is a labeled directed multigraph $\langle
N,E,\lambda\rangle$, where $N$ is the node set, $E$ is the edge set, and
$\lambda$ is a labeling function $E\rightarrow\Sigma$. The node set $N$ is the
set of positions of $\mathcal{R}$. If there is a TGD $\sigma\in\Sigma$ such
that the same variable appears at position $\pi_{b}$ in
$\mathit{body}(\sigma)$ and at position $\pi_{h}$ in $\mathit{head}(\sigma)$,
then in $E$ there is an edge $e=(\pi_{b},\pi_{h})$ with $\lambda(e)=\sigma$.
Intuitively speaking, the dependency graph of a set $\Sigma$ of TGDs describes
all the possible ways of propagating a term from a position to some other
position during the construction of the chase under $\Sigma$. More precisely,
the existence of a path $P$ from $\pi_{1}$ to $\pi_{2}$ implies that it is
possible (but not always) to propagate a term from $\pi_{1}$ to $\pi_{2}$. The
existence of $P$ guarantees the propagation of a term from $\pi_{1}$ to
$\pi_{2}$ if, for each pair of consecutive edges $e=(\pi,\pi^{\prime})$ and
$e^{\prime}=(\pi^{\prime},\pi^{\prime\prime})$ of $P$, where $e$ and
$e^{\prime}$ are labeled by the TGDs $\sigma$ and $\sigma^{\prime}$,
respectively, the atom obtained during the chase by applying $\sigma$ triggers
$\sigma^{\prime}$. To verify whether this holds we need an additional piece of
information, the so-called _equality type_ , about the body-atom and the head-
atom of each TGD that occurs in $P$.
###### Definition 4 (Equality Type)
Consider an atom $\underline{a}$ of the form $r(t_{1},\ldots,t_{n})$, where
$n\geq 1$. The _equality type_ of $\underline{a}$ is the set of equalities
$\displaystyle\left\\{r[i]=r[j]~{}|~{}t_{i},t_{j}\not\in\Delta_{c}\textrm{~{}and~{}}t_{i}=t_{j}\right\\}$
$\displaystyle\bigcup$
$\displaystyle\left\\{r[i]=c~{}|~{}c\in\Delta_{c}\textrm{~{}and~{}}t_{i}=c\right\\}.$
We denote the above set as $\mathit{eq}(\underline{a})$.
It is straightforward to see that, given a pair of TGDs $\sigma$ and
$\sigma^{\prime}$, if
$\mathit{eq}(\mathit{body}(\sigma^{\prime}))\subseteq\mathit{eq}(\mathit{head}(\sigma))$,
then there exists a substitution $\mu$ such that
$\mu(\mathit{body}(\sigma^{\prime}))=\mathit{head}(\sigma)$. This allows us to
show that the atom obtained by applying $\sigma$ during the construction of
the chase triggers $\sigma^{\prime}$. Consequently, the existence of a path
$P$ (as above) guarantees the propagation of a term from $\pi_{1}$ to
$\pi_{2}$ if, for each pair of consecutive edges $e$ and $e^{\prime}$ of $P$
which are labeled by $\sigma$ and $\sigma^{\prime}$, respectively,
$\mathit{eq}(\mathit{body}(\sigma^{\prime}))\subseteq\mathit{eq}(\mathit{head}(\sigma))$.
###### Example 6 (Dependency Graph)
Consider the set $\Sigma$ of TGDs
$\begin{array}[]{rcl}\sigma_{1}&:&p(X,Y)\rightarrow\exists Zr(X,Y,Z)\\\
\sigma_{2}&:&r(X,Y,c)\rightarrow s(X,Y,Y)\\\ \sigma_{3}&:&s(X,X,Y)\rightarrow
p(X,Y).\end{array}$
The equality type of the body-atoms and head-atoms of the TGDs of $\Sigma$ are
as follows:
$\begin{array}[]{rcl}\mathit{eq}(\mathit{body}(\sigma_{1}))&=&\varnothing\\\
\mathit{eq}(\mathit{head}(\sigma_{1}))&=&\varnothing\\\
\mathit{eq}(\mathit{body}(\sigma_{2}))&=&\\{r[3]=c\\}\\\
\mathit{eq}(\mathit{head}(\sigma_{2}))&=&\\{s[2]=s[3]\\}\\\
\mathit{eq}(\mathit{body}(\sigma_{3}))&=&\\{s[1]=s[2]\\}\\\
\mathit{eq}(\mathit{head}(\sigma_{3}))&=&\varnothing.\end{array}$
The dependency graph of $\Sigma$ is shown in Figure 2.
We are now ready, by exploiting the dependency graph of a set of TGDs, and the
equality type of an atom, to introduce _atom coverage_.
Figure 2: Dependency graph for Example 6.
###### Definition 5 (Atom Coverage)
Consider a BCQ $q$ over a schema $\mathcal{R}$, and a set $\Sigma$ of TGDs
over $\mathcal{R}$. Let $\underline{a}$ and $\underline{b}$ be atoms of
$\mathit{body}(q)$, where $\\{t_{1},\ldots,t_{n}\\}$, for $n\geq 0$, is the
set of shared variables and constants that occur in $\underline{b}$. Also, let
$G_{\Sigma}$ be the dependency graph of $\Sigma$. We say that $\underline{a}$
_covers_ $\underline{b}$ w.r.t. $q$ and $\Sigma$, written as
$\underline{a}\prec_{\Sigma}^{q}\underline{b}$, if for each $i\in[n]$: (i)
the term $t_{i}$ occurs also in $\underline{a}$, and (ii) if $t_{i}$ occurs
in $\underline{a}$ and $\underline{b}$ at positions $\Pi_{\underline{a},i}$
and $\Pi_{\underline{b},i}$, respectively, then, there exists an integer
$k\geq 2$ and a set of TGDs
$\\{\sigma_{1},\ldots,\sigma_{k-1}\\}\subseteq\Sigma$, where
$\mathit{eq}(\mathit{body}(\sigma_{1}))\subseteq\mathit{eq}(\underline{a})$
and, for each $j\in[k-2]$,
$\mathit{eq}(\mathit{body}(\sigma_{j+1}))\subseteq\mathit{eq}(\mathit{head}(\sigma_{j}))$,
such that, for each $\pi\in\Pi_{\underline{b},i}$, in $G_{\Sigma}$ there
exists a path $\pi_{i_{1}}\pi_{i_{2}}\ldots\pi_{i_{k}}$, where
$\pi_{i_{1}}\in\Pi_{\underline{a},i}$, $\pi_{i_{k}}=\pi$, and
$\lambda((\pi_{i_{j}},\pi_{i_{j+1}}))=\sigma_{j}$, for each $j\in[k-1]$.
Condition _(i)_ ensures that by removing $\underline{b}$ from $q$ we do not
loose any constant, and also all the joins between $\underline{b}$ and the
other atoms of $\mathit{body}(q)$, except $\underline{a}$, are preserved.
Condition _(ii)_ guarantees that the atom $\underline{b}$ is logically implied
(w.r.t. $\Sigma$) by the atom $\underline{a}$, and therefore can be
eliminated.
###### Lemma 8
Consider a BCQ $q$ over a schema $\mathcal{R}$, and a set $\Sigma$ of linear
TGDs over $\mathcal{R}$. Suppose that
$\underline{a}\prec_{\Sigma}^{q}\underline{b}$, where
$\underline{a},\underline{b}\in\mathit{body}(q)$, and $q^{\prime}$ is the BCQ
obtained from $q$ by eliminating the atom $\underline{b}$. Then, $I\models q$
iff $I\models q^{\prime}$, for each instance $I$ that satisfies $\Sigma$.
Proof (Sketch). ($\Rightarrow$) By hypothesis, there exists a homomorphism $h$
such that $h(\mathit{body}(q))\subseteq I$. Since, by definition of
$q^{\prime}$, $\mathit{body}(q^{\prime})\subset\mathit{body}(q)$, we
immediately get that $h(\mathit{body}(q^{\prime}))\subseteq I$, which implies
that $I\models q^{\prime}$.
($\Leftarrow$) Conversely, there exists a homomorphism $h$ such that
$h(\mathit{body}(q^{\prime}))\subseteq I$, and thus
$h(\mathit{body}(q)\setminus\\{\underline{b}\\})\subseteq I$. It suffices to
show that there exists an extension of $h$ which maps $\underline{b}$ to $I$.
Since $\underline{a}\prec_{\Sigma}^{q}\underline{b}$, it is not difficult to
verify that there exists an atom $\underline{c}\in I$ such that
$\mathit{eq}(\underline{b})=\mathit{eq}(\underline{c})$, which implies that
there exists a substitution $\mu$ such that
$\mu(\underline{b})=\underline{c}$, and also $\mu$ is compatible with $h$.
Consequently, $(h\cup\mu)(\mathit{body}(q))\subseteq I$, and thus $I\models
q$.
An _atom elimination strategy_ for a BCQ is a permutation of its body-atoms.
Given a BCQ $q$ and a set $\Sigma$ of linear TGDs, the set of atoms of
$\mathit{body}(q)$ that cover $\underline{a}\in\mathit{body}(q)$ w.r.t.
$\Sigma$, denoted as $\mathit{cover}(\underline{a},q,\Sigma)$, is the set
$\\{\underline{b}~{}|~{}\underline{b}\in\mathit{body}(q)\textrm{~{}and~{}}\underline{b}\prec_{\Sigma}^{q}\underline{a}\\}$;
when $q$ and $\Sigma$ are obvious from the context, we shall denote the above
set as $\mathit{cover}(\underline{a})$. By exploiting the cover set of the
atoms of $\mathit{body}(q)$, we associate to each atom elimination strategy
$S$ for $q$ a subset of $\mathit{body}(q)$, denoted
$\mathit{eliminate}(q,S,\Sigma)$, which is the set of atoms of
$\mathit{body}(q)$ that can be safely eliminated (according to $S$) in order
to obtain a logically equivalent query (w.r.t. $\Sigma$) with less atoms in
its body. Formally, $\mathit{eliminate}(q,S,\Sigma)$ is computed by applying
the following procedure; in the sequel, let
$S=[\underline{a}_{1},\ldots,\underline{a}_{n}]$, where
$\\{\underline{a}_{1},\ldots,\underline{a}_{n}\\}=\mathit{body}(q)$:
* $A:=\varnothing$;
* foreach $i:=1$ to $n$ do
* $\underline{a}:=S[i]$;
* if $\mathit{cover}(\underline{a})\neq\varnothing$ then
* $A:=A\cup\\{\underline{a}\\}$;
* foreach $\underline{b}\in\mathit{body}(q)\setminus A$ do
* $\mathit{cover}(\underline{b}):=\mathit{cover}(\underline{b})\setminus\\{\underline{a}\\}$;
* return $A$.
By exploiting the fact that the binary relation $\prec_{\Sigma}^{q}$ is
transitive, it is possible to establish the uniqueness (w.r.t. the number of
the eliminated atoms) of the atom elimination strategy for a BCQ. In
particular, the following lemma can be shown.
###### Lemma 9
Consider a BCQ $q$ over a schema $\mathcal{R}$, and a set $\Sigma$ of linear
TGDs over $\mathcal{R}$. Let $S_{1}$ and $S_{2}$ be arbitrary elimination
strategies for $q$. It holds that,
$|\mathit{eliminate}(q,S_{1},\Sigma)|=|\mathit{eliminate}(q,S_{2},\Sigma)|$.
Since the elimination strategy for a query is unique (w.r.t. the number of the
eliminated atoms), in the rest of this section we refer to the set of atoms
that can be safely eliminated from a query $q$ (w.r.t. a set $\Sigma$ of
linear TGDs) by $\mathit{eliminate}(q,\Sigma)$.
We are now ready to describe how query elimination works. During the execution
of the rewriting algorithm TGD-rewrite (see Algorithm 1), after the
factorization step and the rewriting step the so-called _elimination_ step is
applied. In particular, the factorized query $q^{\prime}$ obtained during the
factorization step is the query
$\mathit{eliminate}(\mathit{factorize}(q,\sigma),\Sigma)$, while the rewritten
query obtained during the rewriting step is the query
$\mathit{eliminate}(\gamma_{A\cup\\{\mathit{head}(\sigma)\\}}(q[A/\mathit{body}(\sigma)]),\Sigma)$.
Moreover, instead of adding the given query $q$ in $Q_{\textsc{rew}}$, we add
the eliminated query. In particular, the first line of the algorithm is
replaced by $Q_{\textsc{rew}}:=\langle\mathit{eliminate}(q),1\rangle$. An
example of query elimination follows.
###### Example 7 (Query Elimination)
Consider the set $\Sigma$ of TGDs of Example 6, and the BCQ
$\begin{array}[]{rcl}q()&\leftarrow&\underbrace{p(A,B)}_{\underline{a}},\underbrace{r(A,B,C)}_{\underline{b}},\underbrace{s(A,A,D)}_{\underline{c}}.\end{array}$
Based on the Definition 5, it is an easy task to verify that
$\mathit{cover}(\underline{a})=\varnothing$,
$\mathit{cover}(\underline{b})=\\{\underline{a}\\}$ and
$\mathit{cover}(\underline{c})=\varnothing$. Therefore, the output of the
function $\mathit{eliminate}(q,\Sigma)$ is the singleton set
$\\{\underline{b}\\}$. Consequently, by applying the elimination step we get
the BCQ $q()\leftarrow p(A,B),s(A,A,D)$.
As already mentioned, the fact that an atom $\underline{a}$ covers some atom
$\underline{b}$, means that $\underline{b}$ is logically implied (w.r.t. the
given set of TGDs) by $\underline{a}$. However, as shown by the following
example, this fact is not also necessary for the implication of
$\underline{b}$ by $\underline{a}$.
###### Example 8 (Atom Implication)
Consider the set $\Sigma$ of TGDs of Example 6, and the BCQ $q$
$\begin{array}[]{rcl}q()&\leftarrow&\underbrace{r(A,A,c)}_{\underline{a}},\underbrace{p(A,A)}_{\underline{b}},\end{array}$
where $c$ is a constant of $\Delta_{c}$. Observe that $\underline{a}$ does not
cover $\underline{b}$ since, despite the existence of the paths $r[1]s[1]p[1]$
and $r[2]s[3]p[2]$ in the dependency graph of $\Sigma$,
$\mathit{eq}(\mathit{body}(\sigma_{3}))\not\subseteq\mathit{eq}(\mathit{head}(\sigma_{2}))$.
However, $\underline{b}$ is logically implied (w.r.t. $\Sigma$) by
$\underline{a}$. In particular, for every instance $I$ that satisfies
$\Sigma$, if $I\models\underline{a}$, which implies that an atom of the from
$r(V,V,c)$ exists in $I$, then due to the TGDs $\sigma_{2}$ and $\sigma_{3}$
there exists also an atom $p(V,V)$, and thus $I\models\underline{b}$. Note
that such cases are identified by the C&B algorithm [15]. Nevertheless, as
already criticized in Section 2, this requires to pay a price in the number of
queries in the rewritten query.
It is not difficult to see that the function eliminate runs in quadratic time
in the number of atoms of $\mathit{body}(q)$ (by considering the given set of
TGDs as fixed). In particular, to compute the cover set of each body-atom of
$q$ we need to consider all the pairs of atoms of $\mathit{body}(q)$. Note
that the problem whether a certain atom $\underline{a}$ covers some other atom
$\underline{b}$ is feasible in constant time since the given set of TGDs (and
thus its dependency graph) is fixed.
The following result implies that the rewriting algorithm $\textsf{TGD-
rewrite}^{\star}$, obtained from TGD-rewrite by applying the additional step
of elimination, is still sound and complete.
###### Theorem 10
Consider a BCQ $q$ over a schema $\mathcal{R}$, a database $D$ for
$\mathcal{R}$, and a set $\Sigma$ of linear TGDs over $\mathcal{R}$. Then,
$D\models\textsf{TGD-rewrite}^{\star}(\mathcal{R},\Sigma,q)$ iff
$D\cup\Sigma\models q$.
Proof (Sketch). This result follows from the fact that the algorithm TGD-
rewrite is sound and complete under linear TGDs (see Theorem 6) and Lemma 8.
It is important to clarify that the above result does not hold if we consider
arbitrary TGDs (as in Theorem 6). This is because the proof of Lemma 8, which
states that atom coverage implies logical implication (w.r.t. the given set of
TGDs), is based heavily on the linearity of TGDs. Termination of $\textsf{TGD-
rewrite}^{\star}$ follows immediately from the fact that TGD-rewrite
terminates under linear TGDs (see Theorem 7).
## 7 Implementation and Experimental Setting
TGD-rewrite (without the additional check described in Subsection 5.1) and the
query elimination technique presented in Section 6 have been implemented in
the prototype system Nyaya [36] available at
http://mais.dia.uniroma3.it/Nyaya. The reasoning and query answering engine is
based on the IRIS Datalog engine888http://www.iris-reasoner.org/. extended to
support the FO-rewritable fragments of the Datalog± family. In particular, we
extended IRIS to natively support existential variables in the head without
introducing function symbols and to support the constant $\mathit{false}$ as
head of a rule (used to represent negative constraints). Both IRIS and our
extension are implemented in Java.
Since TGD-rewrite is designed for reasoning over ontologies with large ABoxes,
we put ourselves in a similar experimental setting such that of [19]. Thus, we
use DL-LiteR ontologies with a varying number of axioms. The queries under
consideration are based on canonical examples used in the research projects
where these ontologies have been developed. VICODI (V) is an ontology of
European history, and developed in the EU-funded VICODI
project999http://www.vicodi.org.. STOCKEXCHANGE (S) is an ontology for
representing the domain of financial institutions of the European Union.
UNIVERSITY (U) is a DL-LiteR version of the LUBM
Benchmark101010http://swat.cse.lehigh.edu/projects/lubm/., developed at Lehigh
University, and describes the organizational structure of universities.
ADOLENA (A) (Abilities and Disabilities OntoLogy for ENhancing Accessibility)
is an ontology developed for the South African National Accessibility Portal,
and describes abilities, disabilities and devices. The Path5 (P5) ontology is
a synthetic ontology encoding graph structures and used to generate an
exponential-blowup of the size of the rewritten queries. Recall that the
transformation of a set of TGDs into an equivalent set of single-head TGDs
with a single existential variable can introduce auxiliary predicates and
rules (see Lemmas 1 and 2). The ontologies UX, AX and P5X are equivalent
ontologies to U, A and P5 where the auxiliary predicates are considered part
of the schema. These ontologies allow to study the impact of such
transformations on the size of the rewriting.
We compared our implementation with two other rewriting-based query answering
systems for FO-rewritable ontologies:
QuOnto111111http://www.dis.uniroma1.it/quonto/., based on [5] and developed by
the University of Rome La Sapienza, and
Requiem121212http://www.comlab.ox.ac.uk/projects/requiem/home.html., based on
[19] and developed by the Knowledge Representation group of the University of
Oxford.
Table 1: Evaluation of Nyaya System. | | Size | Length | Width
---|---|---|---|---
| | QO | RQ | NY | NY⋆ | QO | RQ | NY | NY⋆ | QO | RQ | NY | NY⋆
V | $q_{1}$ | 15 | 15 | 15 | 15 | 15 | 15 | 15 | 15 | 0 | 0 | 0 | 0
$q_{2}$ | 11 | 10 | 10 | 10 | 32 | 30 | 30 | 30 | 31 | 30 | 30 | 30
$q_{3}$ | 72 | 72 | 72 | 72 | 216 | 216 | 216 | 216 | 144 | 144 | 144 | 144
$q_{4}$ | 185 | 185 | 185 | 185 | 555 | 555 | 555 | 555 | 370 | 370 | 370 | 370
$q_{5}$ | 150 | 30 | 30 | 30 | 900 | 210 | 210 | 210 | 1,110 | 270 | 270 | 270
S | $q_{1}$ | 6 | 6 | 6 | 6 | 6 | 6 | 6 | 6 | 0 | 0 | 0 | 0
$q_{2}$ | 204 | 160 | 160 | 2 | 566 | 480 | 480 | 2 | 362 | 320 | 320 | 0
$q_{3}$ | 1,194 | 480 | 480 | 4 | 5,026 | 2,400 | 2,400 | 8 | 4,778 | 2,400 | 2,400 | 4
$q_{4}$ | 1,632 | 960 | 960 | 4 | 7,384 | 4,800 | 4,800 | 8 | 7,112 | 4,800 | 4,800 | 4
$q_{5}$ | 11,487 | 2,880 | 2,880 | 8 | 67,664 | 20,160 | 20,160 | 24 | 84,064 | 25,920 | 25,920 | 24
U | $q_{1}$ | 5 | 2 | 2 | 2 | 10 | 4 | 4 | 4 | 5 | 2 | 2 | 2
$q_{2}$ | 287 | 148 | 148 | 1 | 813 | 444 | 444 | 1 | 526 | 296 | 296 | 0
$q_{3}$ | 1,260 | 224 | 224 | 4 | 7,296 | 1,344 | 1,344 | 16 | 10,812 | 2,016 | 2,016 | 20
$q_{4}$ | 5,364 | 1,628 | 1,628 | 2 | 15,723 | 4,884 | 4,884 | 2 | 10,393 | 3,256 | 3,256 | 0
$q_{5}$ | 9,245 | 2,960 | 2,960 | 10 | 35,710 | 11,840 | 11,840 | 20 | 52,970 | 17,760 | 17,760 | 20
A | $q_{1}$ | 783 | 402 | 402 | 247 | 1,540 | 779 | 779 | 197 | 757 | 377 | 377 | 86
$q_{2}$ | 1,812 | 103 | 103 | 92 | 5,350 | 256 | 256 | 234 | 3,538 | 153 | 153 | 142
$q_{3}$ | 4,763 | 104 | 104 | 104 | 23,804 | 520 | 520 | 520 | 23,804 | 520 | 520 | 520
$q_{4}$ | 7,251 | 492 | 492 | 454 | 21,406 | 1,288 | 1,288 | 1,212 | 14,155 | 796 | 796 | 758
$q_{5}$ | 66,068 | 624 | 624 | 624 | 195,042 | 3,120 | 3,120 | 3,120 | 128,974 | 3,120 | 3,120 | 3,120
P5 | $q_{1}$ | 14 | 6 | 6 | 6 | 14 | 6 | 6 | 6 | 0 | 0 | 0 | 0
$q_{2}$ | 86 | 10 | 10 | 10 | 156 | 16 | 16 | 16 | 70 | 6 | 6 | 6
$q_{3}$ | 538 | 13 | 13 | 13 | 1,413 | 29 | 29 | 29 | 900 | 16 | 16 | 16
$q_{4}$ | 3,620 | 15 | 15 | 15 | 14,430 | 44 | 44 | 44 | 10,260 | 29 | 29 | 29
$q_{5}$ | 25,256 | 16 | 16 | 16 | 107,484 | 60 | 60 | 60 | 103,361 | 44 | 44 | 44
UX | $q_{1}$ | 5 | 5 | 5 | 5 | 10 | 10 | 10 | 10 | 5 | 5 | 5 | 5
$q_{2}$ | 286 | 240 | 240 | 1 | 156 | 147 | 147 | 1 | 70 | 70 | 70 | 0
$q_{3}$ | 1,248 | 1,008 | 1,008 | 12 | 1,397 | 1,125 | 1,125 | 48 | 892 | 735 | 735 | 60
$q_{4}$ | 5,358 | 5,000 | 5,000 | 5 | 12,006 | 7,578 | 7,578 | 5 | 9,828 | 5,625 | 5,625 | 0
$q_{5}$ | 9,220 | 8,000 | 8,000 | 25 | 101,652 | 47,656 | 47,656 | 50 | 96,677 | 37,890 | 37,890 | 50
AX | $q_{1}$ | 783 | 782 | 782 | 555 | 1,543 | 1,541 | 1,541 | 1,084 | 763 | 761 | 761 | 529
$q_{2}$ | 1,812 | 1,781 | 1,781 | 1,737 | 3,589 | 3,528 | 3,528 | 3,514 | 3,576 | 3,516 | 3,516 | 3,401
$q_{3}$ | 4,763 | 4,752 | 4,752 | 4,741 | 27,705 | 23,760 | 23,760 | 23,760 | 23,824 | 23,815 | 23,815 | 23,694
$q_{4}$ | 7,251 | 7,100 | 7,100 | 6,564 | 7,739 | 7,578 | 7,578 | 6,178 | 5,744 | 5,625 | 5,625 | 5,201
$q_{5}$ | - | - | 76,032 | 76,032 | - | - | 81,173 | 81,173 | - | - | 95,942 | 95,942
P5X | $q_{1}$ | 14 | 14 | 14 | 14 | 14 | 14 | 14 | 14 | 0 | 0 | 0 | 0
$q_{2}$ | 86 | 77 | 77 | 66 | 156 | 147 | 147 | 121 | 70 | 70 | 70 | 55
$q_{3}$ | 530 | 390 | 390 | 329 | 1,397 | 1,125 | 1,125 | 925 | 892 | 735 | 735 | 596
$q_{4}$ | 3,476 | 1,953 | 1,953 | 1,644 | 12,006 | 7,578 | 7,578 | 6,263 | 9,828 | 5,625 | 5,625 | 4,619
$q_{5}$ | 23,744 | 9,766 | 9,766 | 8,219 | 101,652 | 47,656 | 47,656 | 39,531 | 96,677 | 37,890 | 37,890 | 31,312
Since TGD-rewrite, as well as the algorithms presented in [5] and [19], are
proven to be sound and complete, the most relevant way of judging the quality
of the rewriting is the _size_ of the perfect rewriting, i.e., the number of
CQs in the perfect UCQ rewriting. In addition, we use two additional metrics,
namely, the _length_ of the rewriting, i.e., the number of atoms in the
perfect rewriting, and the _width_ , i.e., the number of joins to be performed
when the rewritten query is executed. We believe these metrics to be more
appropriate than the number of symbols in the rewritten query used, for
example, in [19], since they allow to establish in a more precise way the cost
of executing the rewriting on a database system. Table 1 reports the results
of our experiments131313Additional data can be found on the Nyaya’s Web site.
while Table 2 shows the queries used in the experiments. We use the symbol “-”
to denote those cases where the tool did not complete the rewriting within 15
minutes. By QO and RQ we refer to the QuOnto and Requiem systems,
respectively, while NY and NY⋆ refer to Nyaya with factorisation and Nyaya
with both factorisation and query elimination, respectively. All the tests
have been performed on an Intel Core 2 Duo Processor at 2.50 GHz and 4GB of
RAM. The OS is Ubuntu Linux 9.10 carrying a Sun JVM Standard Edition with
maximum heap size set at 2GB of RAM.
Table 2: Test Queries TBox | Queries
---|---
V | $q_{1}(A)\leftarrow\mathit{Location(A).}$
$q_{2}(A,B)\leftarrow\mathit{Military\\_Person(A),hasRole(B,A),related(A,C).}$
$q_{3}(A,B)\leftarrow\mathit{Time\\_Dependant\\_Relation(A),hasRelationMember(A,B),Event(B).}$
$q_{4}(A,B)\leftarrow\mathit{Object(A),hasRole(A,B),Symbol(B).}$
$q_{5}(A)\leftarrow\mathit{Individual(A),hasRole(A,B),Scientist(B),hasRole(A,C),Discoverer(C),hasRole(A,D),Inventor(D).}$
S | $q_{1}(A)\leftarrow\mathit{StockExchangeMember(A).}$
$q_{2}(A,B)\leftarrow\mathit{Person(A),hasStock(A,B),Stock(B).}$
$q_{3}(A,B,C)\leftarrow\mathit{FinantialInstrument(A),belongsToCompany(A,B),Company(B),hasStock(B,C),Stock(C).}$
$q_{4}(A,B,C)\leftarrow\mathit{Person(A),hasStock(A,B),Stock(B),isListedIn(B,C),StockExchangeList(C).}$
$q_{5}(A,B,C,D)\leftarrow\mathit{FinantialInstrument(A),belongsToCompany(A,B),Company(B),hasStock(B,C),Stock(C),}$
| $\mathit{isListedIn(B,D),StockExchangeList(D).}$
U(X) | $q_{1}(A)\leftarrow\mathit{worksFor(A,B),affiliatedOrganizationOf(B,C).}$
$q_{2}(A,B)\leftarrow\mathit{Person(A),teacherOf(A,B),Course(B).}$
$q_{3}(A,B,C)\leftarrow\mathit{Student(A),advisor(A,B),FacultyStaff(B),takesCourse(A,C),teacherOf(B,C),Course(C).}$
$q_{4}(A,B)\leftarrow\mathit{Person(A),worksFor(A,B),Organization(B).}$
$q_{5}(A)\leftarrow\mathit{Person(A),worksFor(A,B),University(B),hasAlumnus(B,A).}$
A(X) | $q_{1}(A)\leftarrow\mathit{Device(A),assistsWith(A,B).}$
$q_{2}(A)\leftarrow\mathit{Device(A),assistsWith(A,B),UpperLimbMobility(B).}$
$q_{3}(A)\leftarrow\mathit{Device(A),assistsWith(A,B),Hear(B),affects(C,B),Autism(C).}$
$q_{4}(A)\leftarrow\mathit{Device(A),assistsWith(A,B),PhysicalAbility(B).}$
$q_{5}(A)\leftarrow\mathit{Device(A),assistsWith(A,B),PhysicalAbility(B),affects(C,B),Quadriplegia(C).}$
P5(X) | $q_{1}(A)\leftarrow\mathit{edge(A,B).}$
$q_{2}(A)\leftarrow\mathit{edge(A,B),edge(B,C).}$
$q_{3}(A)\leftarrow\mathit{edge(A,B),edge(B,C),edge(C,D).}$
$q_{4}(A)\leftarrow\mathit{edge(A,B),edge(B,C),edge(C,D),edge(D,E).}$
$q_{4}(A)\leftarrow\mathit{edge(A,B),edge(B,C),edge(C,D),edge(D,E),edge(E,F).}$
As it can be seen, query elimination provides a substantial advantage in terms
of the size of the perfect rewriting for the real-world ontologies A, U and S.
In particular, for the queries denoted as Q2 in U and S, our procedure
eliminates all the redundant atoms in the input query, and drastically reduces
the number of queries in the final rewriting. On the other side, query
elimination is not particularly effective in the synthetic test case P5 and
P5X, since these cases have been intentionally created in order to generate
perfect rewritings of exponential size.
## 8 Future Work
We plan to investigate rewriting and optimization techniques for sticky-join
sets of TGDs, and alternative forms of rewriting such as positive-existential
queries. We also plan to develop improved techniques for rewriting an
ontological query into a non-recursive Datalog program, rather than into a
union of conjunctive queries (recall the discussion in Section 2). While the
current approaches yield exponentially large non-recursive Datalog programs,
it is possible to rewrite queries and TBoxes into non-recursive Datalog
programs whose size is simultaneously polynomial in the query and the TBox.
This will be dealt in a forthcoming paper.
## Acknowledgments
G. Gottlob’s work was funded by the European Research Council under the
European Community’s Seventh Framework Programme (FP7/2007-2013)/ERC grant no.
246858 – DIADEM. Gottlob gratefully acknowledges a Royal Society Wolfson
Research Merit Award. G. Orsi and G. Gottlob also acknowledge the Oxford
Martin School - Institute for the Future of Computing. A. Pieris’ work was
funded by the EPSRC project “Schema Mappings and Automated Services for Data
Integration and Exchange” (EP/E010865/1). We thank Michaël Thomazo for his
useful and constructive comments on the conference version of this paper.
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|
arxiv-papers
| 2011-12-01T22:19:24 |
2024-09-04T02:49:24.885695
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Georg Gottlob and Giorgio Orsi and Andreas Pieris",
"submitter": "Giorgio Orsi PhD",
"url": "https://arxiv.org/abs/1112.0343"
}
|
1112.0408
|
# Late-time symmetry near black hole horizons
Kentaro Tanabe Yukawa Institute for Theoretical Physics, Kyoto University,
Kyoto 606-8502, Japan Tetsuya Shiromizu Shunichiro Kinoshita Department of
Physics, Kyoto University, Kyoto 606-8502, Japan
###### Abstract
It is expected that black holes are formed dynamically under gravitational
collapses and approach stationary states. In this paper, we show that the
asymptotic Killing vector at late time should exist on the horizon and then
that it can be extended outside black holes under the assumption of the
analyticity of spacetimes. This fact implies that if there is another
asymptotic Killing vector which becomes a stationary Killing at a far region
and spacelike in the “ergoregion,” the rotating black holes may have the
asymptotically axisymmetric Killing vector at late time. Thus, we may expect
that the asymptotic rigidity of the black holes holds.
###### pacs:
04.20.-q, 04.20.Ha
††preprint: YITP-11-99
## I introduction
Black holes in our Universe are expected to be formed under gravitational
collapses, and to finally approach stationary and vacuum states by radiating
and absorbing energy, momentum, and angular momentum. Then the uniqueness
theorem Israel guarantees that the black hole candidates in our Universe are
the Kerr black holes. The key ingredient for the proof of the uniqueness
theorem is the rigidity theorem Hawking:1973uf ; Hawking:1971vc ;
Friedrich:1998wq ; Racz:1999ne . The rigidity theorem shows that the
stationary rotating black holes have axisymmetric Killing vectors. The outline
of the proof is as follows: The stationarity of spacetime implies that there
are no gravitational waves around a black hole. Then the expansion and shear
of the event horizon must vanish by virtue of the Raychaudhuri equation and
stationarity. Using the Einstein equations, then, we can find that the null
geodesic generator of the event horizon is a Killing vector. If the black hole
is rotating, this new Killing vector may deviate from the stationary Killing
vector which becomes spacelike in the ergoregion. This means that the
stationary black holes might have two Killing vectors, that is, the stationary
and axisymmetric Killing vectors. Hence the stationary black holes should
rotate rigidly.
However, the late-time phase of black holes produced by gravitational
collapses would not be exactly stationary but nearly stationary. “Nearly
stationary” means that the black holes are surrounded by the gravitational
waves at late time. Then we cannot apply the rigidity theorem to such black
holes because the expansion and shear of the null geodesic generator on the
event horizon do not vanish due to the presence of the gravitational waves on
the event horizon. Note that the late-time behaviors of the perturbations
around the Schwarzschild and Kerr black holes were examined and it was shown
that the perturbations decay at late time both on the horizon and null
infinity at the same rate (for examples, see Refs. Gundlach:1993tp ;
Barack:1999ma ; Barack:1999ya ). In the dynamical processes in gravitational
collapses, however, it is quite nontrivial whether the formed black hole
approaches the Kerr black hole. In this paper, we show that the asymptotic
Killing vector, which will asymptotically approach a Killing vector at late
time, should exist on the horizon without assuming any symmetries and then
that it can be extended outside of the event horizon using the Einstein
equations. This indicates that if there is another asymptotic Killing vector
which will be an asymptotically stationary Killing vector at a far region, the
rotating black hole may have the asymptotically axisymmetric Killing vector at
late time.
This paper is organized as follows. In the next section we provide the
Einstein equations near the future event horizons. In Sec. III, we investigate
the late-time behaviors of the metric on the event horizon and find the late-
time symmetry. Then, under the assumption of the analyticity, we show that
this late-time symmetry can be extended outside of the event horizon using the
Einstein equation. In Sec. IV, we provide a summary and discussion. In
Appendix A we perform the $(n+1)$ decompositions of the Einstein equation.
Using them, we provide the explicit form of the Einstein equation in the
current form of the metric in Appendix B. Almost a similar formulation is
developed in the study on null infinity Tanabe2011 . In Appendix C, we present
the details of the discussion associated with gauge freedom.
## II Bondi-like coordinate and Einstein equations
In this section, introducing the Bondi-like coordinate near the event horizon,
we investigate the initial value problem. For the details of these
derivations, see Appendixes A and B.
### II.1 Bondi-like coordinate near event horizon
We consider a dynamical black hole and investigate the late-time behavior of
the near horizon geometry at late time. The black hole is defined as the
region which is not contained in $J^{-}(\mathscr{I}^{+})$ where
$\mathscr{I}^{+}$ is the future null infinity. $J^{-}(S)$ is the chronological
past connecting a set $S$ by causal curves from $S$. The event horizon is the
boundary of the black hole and it is a null hypersurface. See Ref.
Hawking:1973uf for the precise definitions. Then, we introduce the Bondi-like
(Gaussian null) coordinates $x^{A}=(u,r,x^{a})$ near the event horizon as
$ds^{2}=g_{AB}dx^{A}dx^{B}=-Adu^{2}+2dudr+h_{ab}(dx^{a}+U^{a}du)(dx^{b}+U^{b}du),$
(1)
where $u$ is a time coordinate. In this coordinate, the horizon position is
taken to be $r=0$. We assume that a cross section of the event horizon is
compact in the $u=\text{constant}$ hypersurface and its topology is $S^{2}$.
$x^{a}$ are coordinates on $S^{2}$. Since the event horizon is a null
hypersurface, $g_{uu}$ must vanish on the event horizon and
$l=\partial/\partial u$ is the null geodesic generator on the event horizon.
Furthermore we can choose the coordinate $u$ so that $l_{a}=0$ on the event
horizon. Then, we have
$A\,\hat{=}\,0,$ (2)
and
$U^{a}\,\hat{=}\,0,$ (3)
where $\hat{=}$ means the evaluation on the event horizon $r=0$.
To solve the vacuum Einstein equations, we formulate the initial value problem
in the Bondi coordinates. For the convenience of our study in the following
sections, we solve the Einstein equations as the evolution equations in the
direction of $r$. Thus, the initial value of the metric should be set on the
$r=\text{constant}$ surface. In this paper we take the event horizon $r=0$ as
the initial surface.
The evolution equations are given by
$\displaystyle R_{rr}$ $\displaystyle=$ $\displaystyle-\frac{1}{2}(\log
h)^{\prime\prime}-\frac{1}{4}h^{ac}h^{bd}(h_{ab})^{{}^{\prime}}(h_{cd})^{{}^{\prime}}\,=\,0,$
(4) $\displaystyle R_{rB}h^{aB}$ $\displaystyle=$
$\displaystyle\frac{1}{2}U^{a^{\prime\prime}}+\frac{1}{2}h^{ac}h_{bc}^{{}^{\prime}}U^{b^{\prime}}+\frac{1}{4}(\log
h)^{\prime}U^{a^{\prime}}+\frac{1}{2}h^{ab}\bar{D}^{c}[h_{bc}^{{}^{\prime}}-h_{bc}(\log
h)^{\prime}]\,=\,0$ (5)
and
$\displaystyle R_{AB}h^{A}_{~{}a}h^{B}_{~{}b}$ $\displaystyle=$
$\displaystyle-\frac{1}{2}Ah^{{}^{\prime\prime}}_{ab}-\frac{1}{2}A^{{}^{\prime}}h^{{}^{\prime}}_{ab}-(\dot{h}_{ab})^{\prime}+{}^{(h)}R_{ab}+\frac{A}{2}h^{cd}h^{{}^{\prime}}_{ac}h^{{}^{\prime}}_{bd}+\frac{1}{2}h^{cd}(h^{{}^{\prime}}_{ac}\dot{h}_{bd}+h^{{}^{\prime}}_{bd}\dot{h}_{ac})+\mathcal{L}_{U}h^{\prime}_{ab}$
(6)
$\displaystyle~{}~{}-\frac{1}{2}h^{{}^{\prime}}_{ac}(\bar{D}_{b}U^{c}+\bar{D}^{c}U_{b})-\frac{1}{2}h^{{}^{\prime}}_{bc}(\bar{D}_{a}U^{c}+\bar{D}^{c}U_{a})-\frac{1}{2}h_{ac}h_{bd}U^{c}{}^{{}^{\prime}}U^{d}{}^{{}^{\prime}}-\frac{1}{4}[\dot{(\log{h})}-2\bar{D}_{a}U^{a}]h^{{}^{\prime}}_{ab}$
$\displaystyle~{}~{}-\frac{1}{4}(\log{h})^{{}^{\prime}}(Ah^{{}^{\prime}}_{ab}+\dot{h}_{ab}-\bar{D}_{a}U_{b}-\bar{D}_{b}U_{a})+\frac{1}{2}(h_{bc}\bar{D}_{a}U^{c}{}^{{}^{\prime}}+h_{ac}\bar{D}_{b}U^{c}{}^{{}^{\prime}})\,=\,0,$
where the prime and dot denote the $r$ and $u$ derivative, respectively,
$\bar{D}_{a}$ is a covariant derivative with $h_{ab}$ and $h=\det{h_{ab}}$.
Also, ${}^{(h)}R_{ab}$ is the Ricci tensor with respect to $h_{ab}$. The
evolutions of the metric functions $A$, $U^{a}$ and $h_{ab}$ are determined by
Eqs. (4), (5) and (6) completely. Note that the trace part of Eq. (6) gives us
$\displaystyle A^{{}^{\prime}}(\log{h})^{{}^{\prime}}$ $\displaystyle=$
$\displaystyle
2{}^{(h)}R-\frac{A}{2}[(\log{h})^{{}^{\prime}}]^{2}-[\dot{(\log{h})}-2\bar{D}_{a}U^{a}](\log{h})^{{}^{\prime}}-h_{ab}U^{a^{\prime}}U^{b^{\prime}}$
(7)
$\displaystyle~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}-A(\log{h})^{{}^{\prime\prime}}-2\dot{(\log{h})}{}^{{}^{\prime}}+2(\bar{D}_{a}U^{a})^{{}^{\prime}}+U^{a}\bar{D}_{a}(\log
h)^{\prime},$
which is used as the evolution equation for $A$.
The other components of the Einstein equations, $R_{uu}=0$, $R_{ur}=0$ and
$R_{ua}=0$, are related to the above evolution equations by the Bianchi
identities. Therefore, once they are satisfied at the initial surface $r=0$,
we do not need to solve them any more. In fact, if the evolution equations are
satisfied, the Bianchi identities lead to
$\left\\{\begin{aligned} (\sqrt{h}R_{ua})^{\prime}&=\,\sqrt{h}R_{ur},\\\
(\sqrt{h}R^{r}{}_{u})^{\prime}&=\,-\sqrt{h}\bar{D}_{a}R^{a}{}_{u}-\dot{(\sqrt{h})}R_{ur},\\\
(\log h)^{\prime}R_{ur}&=\,0,\end{aligned}\right.$ (8)
where $R^{r}{}_{u}=R_{uu}+AR_{ur}-U^{a}R_{ua}$ and
$R^{a}{}_{u}=-U^{a}R_{ur}+h^{ab}R_{ub}$. Thus, the evolution equations
guarantee that $R_{uu}=0$, $R_{ur}=0$ and $R_{ua}=0$ will always be satisfied
at $r\neq 0$ if $R_{uu}\,\hat{=}\,0$, $R_{ur}\,\hat{=}\,0$ and
$R_{ua}\,\hat{=}\,0$ at the initial surface $r=0$. For convenience, at $r=0$
we will use the following constraint equations:
$R_{uu}\hat{=}-\frac{1}{2}\ddot{(\log h)}+\frac{A^{\prime}}{4}\dot{(\log
h)}-\frac{1}{4}h^{ac}h^{bd}\dot{h}_{ab}\dot{h}_{cd}=0,$ (9)
and
$R^{ra}\hat{=}-\frac{1}{2}(\dot{U}^{a})^{\prime}-\frac{1}{4}\dot{(\log
h)}{U^{a}}^{\prime}-\frac{1}{2}h^{ab}{U^{c}}^{\prime}\dot{h}_{bc}-\frac{1}{2}\bar{D}^{a}A^{\prime}+\frac{1}{2}h^{ab}\bar{D}^{c}\dot{h}_{bc}-\frac{1}{2}\bar{D}^{a}\dot{(\log
h)}=0$ (10)
[see Eqs. (79) and (80)]. Moreover, since $A$ should vanish on the initial
surface $r=0$ [see Eq. (2)], the evolution equations Eq. (6) become singular
on $r=0$. Analyticity of $h_{ab}$ on $r=0$ gives us the regularity condition
$\displaystyle-\frac{1}{2}A^{{}^{\prime}}h^{{}^{\prime}}_{ab}$
$\displaystyle-$
$\displaystyle\dot{h}^{{}^{\prime}}_{ab}+{}^{(h)}R_{ab}+\frac{1}{2}h^{cd}(h^{{}^{\prime}}_{ac}\dot{h}_{bd}+h^{{}^{\prime}}_{bd}\dot{h}_{ac})$
(11) $\displaystyle-$
$\displaystyle\frac{1}{2}h_{ac}h_{bd}U^{c}{}^{{}^{\prime}}U^{d}{}^{{}^{\prime}}-\frac{1}{4}\dot{(\log{h})}h^{{}^{\prime}}_{ab}$
$\displaystyle-$
$\displaystyle\frac{1}{4}(\log{h})^{{}^{\prime}}\dot{h}_{ab}+\frac{1}{2}(h_{bc}\bar{D}_{a}U^{c}{}^{{}^{\prime}}+h_{ac}\bar{D}_{b}U^{c}{}^{{}^{\prime}})\,\hat{=}\,0.$
If we give $h_{ab}|_{r=0}$ on $r=0$, we can determine
$h^{\prime}_{ab}|_{r=0}$, ${U^{a}}^{\prime}|_{r=0}$, and $A^{\prime}|_{r=0}$
by solving Eqs. (9), (10) and (11). As a result, we will solve the evolution
equations (4), (5) and (6) using the above initial values on $r=0$.
### II.2 Some explicit solutions near event horizon
In this subsection, it is shown that we explicitly solve the constraint
equations and evolution equations near the event horizon using power series
expansion around $r=0$. To do this, we expand the metric functions with $r$
near the event horizon as
$A\,=\,rA^{(1)}+\sum_{i\geq 2}r^{i}A^{(i)},$ (12)
$U^{a}\,=\,rU^{(1)a}+\sum_{i\geq 2}r^{i}U^{(i)a},$ (13)
and
$h_{ab}\,=\,h^{(0)}_{ab}+rh^{(1)}_{ab}+\sum_{i\geq 2}r^{i}h^{(i)}_{ab},$ (14)
where from the gauge conditions Eqs. (2) and (3), the expansions of $A$ and
$U^{a}$ start from the first order of $r$. In the following, the tensor index
of $h^{(i)}_{ab}$ and $U^{(i)a}$ is raised and lowered by $h^{(0)}_{ab}$. The
trace and traceless parts are also taken by $h^{(0)}_{ab}$.
First let us solve the constraint equations and regularity conditions in order
to determine initial values. On the initial surface $r=0$, $h^{(0)}_{ab}$,
$h^{(1)}_{ab}$, $A^{(1)}$ and $U^{(1)a}$ should be set on $r=0$ as initial
values. The constraint equations for these initial values are
$R_{uu}\,\hat{=}\,0$ and $R^{ra}\,\hat{=}\,0$. Now $R_{uu}\,\hat{=}\,0$ [Eq.
(9)] is rewritten as
$\ddot{h}^{(0)}-\frac{1}{2}A^{(1)}\dot{h}^{(0)}+\frac{1}{2}h^{(0)ac}h^{(0)bd}\dot{h}_{ab}^{(0)}\dot{h}_{cd}^{(0)}\,=\,0,$
(15)
where $h^{(0)}=\det{h^{(0)}_{ab}}$. Thus, $A^{(1)}$ should be given to satisfy
this equation for given $h^{(0)}_{ab}$. Also, $R^{ra}\,\hat{=}\,0$ [Eq. (10)]
becomes
$\dot{U}^{(1)a}\,=\,-D^{a}A^{(1)}+h^{(0)ac}D^{b}\dot{h}_{bc}^{(0)}-D^{a}\dot{(\log
h^{(0)})}-h^{(0)ac}U^{(1)b}\dot{h}^{(0)}_{bc}-\frac{1}{2}U^{(1)a}\dot{(\log
h^{(0)})},$ (16)
where $D_{a}$ is the covariant derivative with respect to $h^{(0)}_{ab}$. Then
the initial value $U^{(1)a}$ is given to satisfy the above. The regularity
condition Eq. (11) becomes
$\displaystyle-\dot{h}^{(1)}_{ab}$ $\displaystyle-$
$\displaystyle\frac{1}{2}A^{(1)}h^{(1)}_{ab}+{}^{(h^{(0)})}R_{ab}+\frac{1}{2}h^{(0)cd}(h^{(1)}_{ac}\dot{h}_{bd}^{(0)}+h^{(1)}_{bd}\dot{h}_{ac}^{(0)})-\frac{1}{2}h^{(0)}_{ac}h^{(0)}_{bd}U^{(1)c}U^{(1)d}$
(17)
$\displaystyle~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}-\frac{1}{4}\dot{(\log{h^{(0)}})}h^{(1)}_{ab}-\frac{1}{4}h^{(0)cd}h^{(1)}_{cd}\dot{h}^{(0)}_{ab}+\frac{1}{2}(D_{a}U^{(1)}_{b}+D_{b}U^{(1)}_{a})=0,$
where ${}^{(h^{(0)})}R_{ab}$ is the Ricci tensor of $h^{(0)}_{ab}$. We obtain
$h^{(1)}_{ab}$ satisfying this equation. Hence, if we give $h^{(0)}_{ab}$ on
the initial surface, the constraint equations and the regularity condition
yield $h^{(1)}_{ab}$, $A^{(1)}$ and $U^{(1)a}$.
Next we solve the evolution equations. Equation (4) becomes near the event
horizon as
$R_{rr}\,=\,-h^{(0)ab}h^{(2)}_{ab}+\frac{1}{4}(h^{(1)}_{ab})^{2}+O(r)\,=\,0.$
(18)
This equation gives us the trace part of $h^{(2)}_{ab}$ as
$h^{(0)ab}h^{(2)}_{ab}\,=\,\frac{1}{4}(h^{(1)}_{ab})^{2}.$ (19)
The evolution equation of $U^{a}$ as $R_{rB}h^{aB}=0$ [Eq. (5)] can be
expanded as
$R_{rB}h^{aB}\,=\,U^{(2)a}+\frac{1}{2}h^{(0)ac}h^{(1)}_{bc}U^{(1)b}+\frac{1}{4}U^{(1)a}h^{(0)bc}h^{(1)}_{bc}+\frac{1}{2}h^{(1)ac}D^{b}\left(h^{(1)}_{bc}-h_{bc}^{(0)}h^{(0)de}h^{(1)}_{de}\right)+O(r).$
(20)
Then $U^{(2)a}$ is given by
$U^{(2)a}\,=\,-\frac{1}{2}h^{(0)ac}h^{(1)}_{bc}U^{(1)b}-\frac{1}{4}U^{(1)a}h^{(0)bc}h^{(1)}_{bc}-\frac{1}{2}h^{(1)ac}D^{b}\left(h^{(1)}_{bc}-h_{bc}^{(0)}h^{(0)de}h^{(1)}_{de}\right).$
(21)
In the same way, expanding the evolution equations Eq. (6) near the event
horizon, we can obtain the traceless part $h^{(2)}_{\langle ab\rangle}$ and
$A^{(2)}$ in terms of $h^{(0)}_{ab},h^{(1)}_{ab},U^{(1)a}$ and $A^{(1)}$. Note
that $A^{(2)}$ is given by the trace part of Eq. (6), namely Eq. (7). However,
we do not provide its explicit form because its form is very cumbersome.
Hence, we can determine $h^{(2)}_{ab}$, $U^{(2)a}$, and $A^{(2)}$ using the
evolution equations. To determine the higher order quantities, $h^{(i)}_{ab}$,
$U^{(i)a}$, and $A^{(i)}$ ($i>2$), we have to repeat the same procedure.
## III Late-time symmetry on and near event horizon
In this section, we show that there is late-time symmetry on the event
horizon. Then we will extend it outside of black hole regions.
### III.1 Late-time behaviors on event horizon
To investigate late-time behaviors of the event horizon, we introduce the
expansion and shear of the null geodesic generator $l=\partial/\partial u$ of
the event horizon. The expansion $\theta$ and shear $\sigma_{ab}$ are defined
as
$\displaystyle\sigma_{ab}+\frac{1}{2}\theta h^{(0)}_{ab}$
$\displaystyle\hat{=}$ $\displaystyle h_{a}^{~{}A}h_{b}^{~{}B}\nabla_{A}l_{B}$
(22) $\displaystyle\hat{=}$ $\displaystyle\frac{1}{2}\dot{h}^{(0)}_{ab},$
where $\sigma_{ab}$ is the traceless part of $\dot{h}_{ab}^{(0)}$ with respect
to $h^{(0)}_{ab}$. Then we can rewrite Eq. (9), one of the constraint
equations, using $\theta$ and $\sigma_{ab}$ as
$\dot{\theta}-\frac{1}{2}A^{(1)}\theta\,=\,-\frac{1}{2}\theta^{2}-\sigma_{ab}\sigma^{ab}.$
(23)
We can regard $A^{(1)}/2$ as a surface gravity of the black hole with respect
to the time coordinate $u$ because
$l^{A}\nabla_{A}l^{B}\,\hat{=}\,\frac{A^{(1)}}{2}l^{B}$ (24)
holds. Using the affine parameter $w$ defined by
$\frac{dw}{du}\,=\,\exp\Bigl{(}\int^{u}\\!\frac{A^{(1)}}{2}du^{\prime}\Bigr{)},$
(25)
we can obtain the Raychaudhuri equation
$\partial_{w}\theta_{(w)}\,=\,-\frac{1}{2}\theta_{(w)}^{2}-\sigma_{(w)ab}\sigma_{(w)}^{ab},$
(26)
where $\theta_{(w)}$ and $\sigma_{(w)ab}$ are expansion and shear with respect
to $w$. We can see the relation between $\theta$, $\sigma_{ab}$ and
$\theta_{(w)}$, $\sigma_{(w)ab}$ as
$\theta\,=\,\theta_{(w)}\exp{\Bigl{(}\int^{u}\\!\frac{A^{(1)}}{2}du^{\prime}\Bigr{)}}\,,\quad\sigma_{ab}\,=\,\sigma_{(w)ab}\exp{\Bigl{(}\int^{u}\\!\frac{A^{(1)}}{2}du^{\prime}\Bigr{)}}.$
(27)
Here we remember that the area law of the event horizon holds for spacetimes
satisfying the null energy condition, that is, $\theta\geq 0$ and
$\theta_{(w)}\geq 0$. Since $\sigma_{(w)ab}\sigma_{(w)}^{ab}\geq 0$, Eq. (26)
implies the inequality
$\partial_{w}\theta_{(w)}+\frac{1}{2}\theta_{(w)}^{2}\leq 0.$ (28)
Then the integration over $w$ gives us
$\theta_{(w)}\leq\frac{1}{1/\theta_{(0)}+(w-w_{0})/2}\to
0~{}~{}({\text{as}}~{}~{}w\to\infty),$ (29)
where we used the fact of $\theta_{(0)}=\theta_{(w)}(w=w_{0})\geq 0$. In
addition, Eq. (26) shows that the shear $\sigma_{(w)ab}$ should also vanish as
$w\to\infty$. This is shown as a part of the proof of another theorem HSN .
From now on, we assume that $w\to\infty$ corresponds to $u\to\infty$. Then, we
see that $\theta$ and $\sigma_{ab}$ should also vanish as $u\to\infty$. It is
natural to assume that the cross section of the event horizon is compact. Then
the vanishing of the expansion implies that the horizon area approaches a
constant and finite value.
Altogether we see the behavior of the metric at late time as
$\mathcal{L}_{l}g_{AB}|_{\rm horizon}\to 0~{}~{}(u\to\infty).$ (30)
Here we impose the following decaying condition on the event horizon for the
metric:
$\mathcal{L}_{l}g_{AB}\,\hat{=}\,O\left(\frac{1}{u^{n}}\right),$ (31)
which explicitly means $\dot{h}^{(0)}_{ab}=O(u^{-n})$. This equation means
that the null geodesic generator of the event horizon $l$ should be an
asymptotic Killing vector at late time ($u\rightarrow\infty$). Thus there is a
late-time symmetry on the event horizon.
Since we consider the dynamics only near the event horizon, we cannot
determine the decaying rate of the metric. However, the details of the
decaying properties are not important for our argument, that is, our result
does not depend on $n$. Note that it will be determined by the boundary
conditions at a far region from the black holes as in Refs. Gundlach:1993tp ;
Barack:1999ma ; Barack:1999ya .
It should be remembered that the horizon which satisfies this condition is
called a slowly evolving horizon in Refs. Booth:2003ji ; Booth:2006bn . If the
right-hand side of Eq. (31) vanishes, the event horizon will be identical to
the isolated horizon Ashtekar:2004cn .
### III.2 Extension of late-time symmetry
The purpose of this subsection is to show that the decaying condition of Eq.
(31) can be extended outside of the event horizon as
$\mathcal{L}_{l}g_{AB}\,=\,O\left(\frac{1}{u^{n}}\right).$ (32)
In the following we assume the analyticity of $g_{AB}$.
Under the presence of the analyticity of spacetimes, the above is equivalent
with
$(\mathcal{L}_{n})^{m}\mathcal{L}_{l}g_{AB}\,\hat{=}\,O\left(\frac{1}{u^{n}}\right),$
(33)
where $n=\partial/\partial r$ and $m=0,1,2,\cdots$. Note that “$\hat{=}$”
means the evaluation on the event horizon $(r=0)$ again.
First we will show the $m=1$ case:
$\mathcal{L}_{n}\mathcal{L}_{l}g_{AB}\,\hat{=}\,O\left(\frac{1}{u^{n}}\right).$
(34)
Substituting the explicit form of $g_{AB}$ to the above, we rewrite Eq. (34)
as
$-\mathcal{L}_{n}\mathcal{L}_{l}(A-U_{a}U^{a})(du)_{A}(du)_{B}+2\mathcal{L}_{n}\mathcal{L}_{l}U_{a}(du)_{(A}(dx^{a})_{B)}+\mathcal{L}_{n}\mathcal{L}_{l}h_{ab}(dx^{a})_{A}(dx^{b})_{B}\hat{=}\,O\left(\frac{1}{u^{n}}\right).$
(35)
Using Eqs. (12), (13) and (14), the above will be equivalent with
$\mathcal{L}_{l}A^{(1)}\,=\,O\left(\frac{1}{u^{n}}\right),$ (36)
$\mathcal{L}_{l}U^{(1)a}\,=\,O\left(\frac{1}{u^{n}}\right),$ (37)
and
$\mathcal{L}_{l}h^{(1)}_{ab}\,=\,O\left(\frac{1}{u^{n}}\right).$ (38)
Let us examine these conditions. First, we focus on Eq. (36). As a result,
using the gauge freedom, we can show that the slightly strong condition,
$\mathcal{L}_{l}A^{(1)}\,=\,O\ (1/u^{n+1})$, holds. To see this, we decompose
$A^{(1)}$ into the $u$-independent term and others as
$A^{(1)}\,=\,A^{(1)}_{0}(x^{a})+\tilde{A}^{(1)}(u,x^{a}).$ (39)
As shown in Ref. Hollands:2006rj , we can choose the gauge so that
$A^{(1)}_{0}$ is a constant. Furthermore, using the residual gauge, we can
take $\tilde{A}^{(1)}=O(1/u^{n})$. Then, using the gauge freedom in our
coordinate, we can make $A^{(1)}$ satisfy stronger decaying property as
$\mathcal{L}_{l}A^{(1)}\,=\,O\left(\frac{1}{u^{n+1}}\right).$ (40)
For the details, see Appendix C. Of course, Eq. (36) is satisfied.
Now $U^{(1)a}$ should satisfy the constraint equation of Eq. (16) as
$\mathcal{L}_{l}U^{(1)a}\,=\,-D^{a}\tilde{A}^{(1)}+h^{(0)ac}D^{b}\dot{h}_{bc}^{(0)}-D^{a}\dot{(\log
h^{(0)})}-h^{(0)ac}U^{(1)b}\dot{h}^{(0)}_{bc}-\frac{1}{2}U^{(1)a}\dot{(\log
h^{(0)})}.$ (41)
Together with Eqs. (31) and (40), we can see that Eq. (37) holds from the
above.
Furthermore, $h^{(1)}_{ab}$ satisfy the following equation [see Eq. (17)] as a
regularity condition
$\displaystyle\dot{h}^{(1)}_{ab}+\frac{1}{2}A^{(1)}h^{(1)}_{ab}$
$\displaystyle=$
$\displaystyle{}^{(h^{(0)})}R_{ab}+\frac{1}{2}h^{(0)cd}(h^{(1)}_{ac}\dot{h}_{bd}^{(0)}+h^{(1)}_{bd}\dot{h}_{ac}^{(0)})-\frac{1}{2}h^{(0)}_{ac}h^{(0)}_{bd}U^{(1)c}U^{(1)d}$
(42)
$\displaystyle~{}~{}-\frac{1}{4}\dot{(\log{h^{(0)}})}h^{(1)}_{ab}-\frac{1}{4}h^{(0)cd}h^{(1)}_{cd}\dot{h}^{(0)}_{ab}+\frac{1}{2}(D_{a}U^{(1)}_{b}+D_{b}U^{(1)}_{a}),$
Then multiplying $\mathcal{L}_{l}$ to the above and using Eq. (16), we can see
$\frac{1}{2}A^{(1)}\mathcal{L}_{l}h^{(1)}_{ab}=\mathcal{L}_{l}{}^{(h^{(0)})}R_{ab}+O(u^{-n})$
(43)
holds. In the above, the higher-order terms like
$\mathcal{L}_{l}^{2}{h}^{(1)}_{ab},(\mathcal{L}_{l}h^{(0)}_{ab})^{2},(\mathcal{L}_{l}h^{(1)}_{ab})^{2}$,
$\mathcal{L}_{l}h^{(0)}_{ab}\mathcal{L}_{l}h^{(1)}_{ab}$ and so on are
contained in $O(u^{-n})$. Thus Eq. (31) tells us that
$\mathcal{L}_{l}h^{(1)}_{ab}=O(u^{-n})$ holds. As a consequence, we can show
the $m=1$ case of Eq. (33).
For the $m>1$ cases of Eq. (33), we perform the same procedure inductively.
Let $m>1$ be an integer and assume that the metric satisfies
$(\mathcal{L}_{n})^{k}\mathcal{L}_{l}g_{AB}\,\hat{=}O\left(\frac{1}{u^{n}}\right)$
(44)
for all $k(<m)$. Then
$(\mathcal{L}_{n})^{m}\mathcal{L}_{l}g_{AB}\,\hat{=}O(u^{-n})$ is equivalent
with
$\displaystyle\mathcal{L}_{l}A^{(m)}\,=\,O\left(\frac{1}{u^{n}}\right),$ (45)
$\displaystyle\mathcal{L}_{l}U^{(m)a}\,=\,O\left(\frac{1}{u^{n}}\right),$ (46)
and
$\mathcal{L}_{l}h^{(m)}_{ab}\,=\,O\left(\frac{1}{u^{n}}\right).$ (47)
Since $U^{(m)a}$ is written by $U^{(1)a}$,
$A^{(1)},h^{(0)}_{ab},\cdots,h^{(m-1)}_{ab}$ from Eq. (5) like Eq. (21), the
induction assumption
$\mathcal{L}_{l}h^{(i)}_{ab}=O\left(\frac{1}{u^{n}}\right)~{}~{}(i\leq m-1)$
(48)
immediately shows us that Eq. (46) holds.
From Eq. (7), $A^{(m)}$ is written by $U^{(1)},A^{(1)},h^{(i)}_{ab}(i<m)$ and
the trace part of $h^{(m)}_{ab}$. Since the trace part of $h^{(m)}_{ab}$ is
written by $U^{(1)},A^{(1)}$ and $h^{(i)}_{ab}$ ($i<m$) through Eq. (4),
$A^{(m)}$ is written only by $U^{(1)},A^{(1)},h^{(i)}_{ab}~{}(i<m)$ in the
end. Then, by the assumption of the induction, Eq. (45) holds. Next, the
evolution equation for $h^{(m)}_{ab}$ is described by Eq. (6). Expanding Eq.
(6) near the event horizon and multiplying, $(\mathcal{L}_{n})^{m-1}$, Eq. (6)
becomes the following form on the event horizon
$\dot{h}^{(m)}_{ab}+\frac{1}{2}A^{(1)}h^{(m)}_{ab}=F_{ab},$ (49)
where $F_{ab}$ is a function of $h^{(i)}_{ab}\,(i\leq m)$,
$\dot{h}^{(i)}_{ab}\,(i<m)$ and so on. Acting $\mathcal{L}_{l}$ to Eq. (49),
then, we can see that
$(1/2)A^{(1)}\dot{h}^{(m)}_{ab}=\dot{F}_{ab}\sim\dot{h}^{(j)}_{ab}=O(1/u^{n})$
for $j<m$. Thus Eq. (47) holds.
Now we can confirm that the induction loop is closed. Then we can show that
Eq. (33), equivalently Eq. (32) holds if the spacetime is real analytic.
Therefore we can show that the null geodesic generator $l$ of the event
horizon is the asymptotic Killing vector at late time in the sense of Eq.
(32).
## IV Summary and discussion
We have confirmed that the expansion and shear of the event horizon should
decay at late time in the vacuum spacetimes. Then, assuming the compactness of
the cross sections of the event horizon, the null geodesic generators on the
horizon give us an asymptotic Killing vector $l$ at late time. This means that
the horizon has late-time symmetry. By solving the Einstein equations, then,
we have found that this late-time symmetry can be extended outside of the
black holes. Therefore, at late time, there is the asymptotic symmetry outside
of black holes.
If the black hole rotates and there is another asymptotic Killing vector at
late time, $k$, which will be a stationary Killing vector at a far distance
and spacelike near the horizon, $k-l$ is also an asymptotic Killing vector
expected to correspond to axisymmetry. In this sense, we would expect that the
rigidity holds in gravitational collapse at late time. In these discussions,
we assume the compactness of the event horizon. Thus this result cannot be
applied to other null hypersurfaces which do not have a compact cross section.
There is a remaining issue. In the proof of the rigidity theorem, the exact
stationarity does show that the null geodesic generators of the horizon are
Killing orbits. On the other hand, our argument could show us that the null
geodesic generators of the horizon is a Killing orbit without assuming the
presence of asymptotically stationary Killing vectors. It is likely that this
difference suggests the existence of important and new points in black hole
physics.
## Acknowledgment
KT is supported by JSPS Grant-in-Aid for Scientific Research (No. 21-2105). TS
is supported by Grant-in-Aid for Scientific Research from the Ministry of
Education, Science, Sports and Culture of Japan (No. 21244033, No. 21111006,
No. 20540258, and No. 19GS0219). This work is also supported by the Grant-in-
Aid for the Global COE Program “The Next Generation of Physics, Spun from
Universality and Emergence” from the Ministry of Education, Culture, Sports,
Science and Technology (MEXT) of Japan.
## Appendix A Einstein equations near event horizon
In Appendix A, using the suitable variables, we will show the derivations of
the Einstein equations near the event horizon.
### A.1 $(n+1)$ decomposition
First we describe the formula of the $(n+1)$ decomposition. The metric can be
written as
$g_{AB}\,=\,\epsilon n_{A}n_{B}+\gamma_{AB},$ (50)
where $\gamma_{AB}$ is an $n$-dimensional induced metric. $n_{A}$ is the unit
normal vector with $n_{A}n^{A}=\epsilon=+1$ ($n^{A}$: spacelike) or $-1$
($n^{A}$: timelike).
We define the extrinsic curvature as
$K_{AB}\,=\,\frac{1}{2}\mathcal{L}_{n}\gamma_{AB}.$ (51)
Now $n_{A}$ can be written as $n_{A}=\epsilon N(d\Omega)_{A}$, where $\Omega$
is a function which describes the hypersurface as $\Omega=\text{const.}$ and
$N$ is the “lapse” function. Then the Riemann tensor is decomposed into
$R_{EFGH}\gamma_{A}{}^{E}\gamma_{B}{}^{F}\gamma_{C}{}^{G}\gamma_{D}{}^{H}={}^{(\gamma)}R_{ABCD}-\epsilon
K_{AC}K_{BD}+\epsilon K_{AD}K_{BC},$ (52)
$R_{EFGD}\gamma_{A}{}^{E}\gamma_{B}{}^{F}\gamma_{C}{}^{G}n^{D}=\nabla_{A}K_{BC}-\nabla_{B}K_{AC},$
(53)
$R_{ACBD}n^{C}n^{D}=-\mathcal{L}_{n}K_{AB}+K_{AC}K_{B}{}^{C}-\epsilon\frac{1}{N}\nabla_{A}\nabla_{B}N,$
(54)
where $\nabla_{A}$ denotes the covariant derivative with respect to
$\gamma_{AB}$.
The Ricci tensor is decomposed into
$R_{AB}n^{A}n^{B}=-\mathcal{L}_{n}K-K_{AB}K^{AB}-\epsilon\frac{1}{N}\nabla^{2}N,$
(55) $R_{AC}n^{A}\gamma_{B}{}^{C}=\nabla^{A}K_{AB}-\nabla_{B}K,$ (56)
$R_{CD}\gamma_{A}{}^{C}\gamma_{B}{}^{D}={}^{(\gamma)}R_{AB}-\epsilon\mathcal{L}_{n}K_{AB}-\epsilon
KK_{AB}+2\epsilon K_{AC}K_{B}{}^{C}-\frac{1}{N}\nabla_{A}\nabla_{B}N.$ (57)
The Ricci scalar is written as
$\displaystyle R=$
$\displaystyle{}^{(\gamma)}R-2\epsilon\mathcal{L}_{n}K-\epsilon K^{2}-\epsilon
K_{AB}K^{AB}-\frac{2}{N}\nabla^{2}N$ (58) $\displaystyle=$
$\displaystyle{}^{(\gamma)}R+\epsilon K^{2}-\epsilon
K_{AB}K^{AB}-\frac{2}{N}\nabla^{2}N-2\epsilon\nabla_{A}(Kn^{A}).$
The components of the Einstein tensor are
$G_{AB}n^{A}n^{B}=\frac{1}{2}(-\epsilon{}^{(\gamma)}R+K^{2}-K_{AB}K^{AB}),$
(59) $G_{AC}n^{A}\gamma_{B}{}^{C}=\nabla^{A}K_{AB}-\nabla_{B}K,$ (60)
$\displaystyle G_{CD}\gamma_{A}{}^{C}\gamma_{B}{}^{D}=$
$\displaystyle{}^{(\gamma)}G_{AB}-\epsilon KK_{AB}+2\epsilon
K_{AC}K_{B}{}^{C}+\frac{\epsilon}{2}\gamma_{AB}(K_{CD}K^{CD}+K^{2})$ (61)
$\displaystyle-\epsilon\mathcal{L}_{n}K_{AB}+\epsilon\gamma_{AB}\mathcal{L}_{n}K-\frac{1}{N}\nabla_{A}\nabla_{B}N+\frac{1}{N}\gamma_{AB}\nabla^{2}N.$
### A.2 Einstein equations
We apply the $(n+1)$ decomposition presented in the previous section to the
$r$-constant surface in our current four dimensional metric form:
$ds^{2}\,=\,-Adu^{2}+2dudr+h_{ab}(dx^{a}+U^{a}du)(dx^{b}+U^{b}du).$ (62)
We express the above in the following form
$ds^{2}\,=\,N^{2}dr^{2}+q_{\mu\nu}(dx^{\mu}+N^{\mu}dr)(dx^{\nu}+N^{\nu}dr),$
(63)
where
$N^{2}\,=\,\frac{1}{A},$ (64) $N^{u}\,=\,-\frac{1}{A},$ (65)
$N^{a}\,=\,\frac{1}{A}U^{a}.$ (66)
$q_{\mu\nu}$ is the induced metric on the $r$-const. hypersurface as
$\displaystyle q_{\mu\nu}=\begin{pmatrix}-A+U^{a}U_{a}&U_{b}\\\
U_{a}&h_{ab}\end{pmatrix}.$ (67)
Note that the Latin indices $a,b,...$ and the Greek indices $\mu,\nu,..$ are
raised and lowered by $h_{ab}$ and $q_{\mu\nu}$ respectively. The unit normal
vector to the $r$-const. hypersurface is given by $m_{A}=N(dr)_{A}$ and
$m^{A}=N^{-1}(\partial_{r}-N^{\mu}\partial_{\mu})^{A}$.
The extrinsic curvature on the $r$-const. hypersurface is defined as
$K_{\mu\nu}\,=\,\frac{1}{2}\mathcal{L}_{m}q_{\mu\nu}\,=\,\frac{1}{2N}(\partial_{r}q_{\mu\nu}-\mathcal{D}_{\mu}N_{\nu}-\mathcal{D}_{\nu}N_{\mu}),$
(68)
where $\mathcal{D}_{\mu}$ is the covariant derivative with respect to
$q_{\mu\nu}$.
We rewrite the induced metric as
$q_{\mu\nu}dx^{\mu}dx^{\nu}=-\alpha^{2}du^{2}+h_{ab}(dx^{a}+\beta^{a}du)(dx^{b}+\beta^{b}du),$
(69)
where
$\alpha^{2}\,=\,A\,,\,\beta^{a}\,=\,U^{a}.$ (70)
The timelike unit vector to the $u$-const. surface is given by $u=-\alpha du$
and $u^{\mu}=\alpha^{-1}(\partial_{u}-\beta^{a}\partial_{a})^{\mu}$.
The extrinsic curvature on the $u$-const. surface is given by
$k_{ab}\,=\,\frac{1}{2}\mathcal{L}_{u}h_{ab}\,=\,\frac{1}{2\alpha}(\partial_{u}h_{ab}-\bar{D}_{a}\beta_{b}-\bar{D}_{b}\beta_{a}),$
(71)
where $\bar{D}_{a}$ is the covariant derivative with respect to $h_{ab}$.
We introduce the following quantities for later convenience
$\Theta\equiv
K_{\,u\nu}u^{\mu}u^{\nu}\,=\,-\frac{1}{N}\partial_{r}(\log{\alpha})-\mathcal{L}_{u}\log{N},$
(72) $\rho^{a}\equiv
K^{a}_{\mu}u^{\mu}\,=\,\frac{1}{2}\partial_{r}\beta^{a}+\bar{D}^{a}\log{\alpha},$
(73) $\kappa_{ab}\equiv
K_{cd}h_{a}{}^{c}h_{b}{}^{d}\,=\,\frac{\alpha}{2}\partial_{r}h_{ab}+k_{ab},$
(74)
and
$\kappa\,=\,\kappa_{ab}h^{ab}\,=\,\frac{\alpha}{2}(\log{h})^{{}^{\prime}}+k,$
(75)
where $\rho_{\mu}u^{\mu}=0=\kappa_{\mu\nu}u^{\mu}$, $h=\det{h}_{ab}$ and
$k=k_{ab}h^{ab}$. The prime denotes the $r$ differentiation. Then $K_{\mu\nu}$
can be written as
$K_{\mu\nu}\,=\,\Theta u_{\mu}u_{\nu}-2\rho_{(\mu}u_{\nu)}+\kappa_{\mu\nu}.$
(76)
Using these quantities, we can decompose the four dimensional Ricci tensor
$R_{AB}$ into the quantities on two dimensional space:
$\displaystyle R_{AB}m^{A}m^{B}$ $\displaystyle=$
$\displaystyle\frac{1}{N}(\Theta-\kappa)^{{}^{\prime}}+\mathcal{L}_{u}(\Theta-\kappa)-\Theta^{2}+2\rho_{a}\rho^{a}-\kappa_{ab}\kappa^{ab}$
(77)
$\displaystyle~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}+\frac{1}{N}(\mathcal{L}_{u}\mathcal{L}_{u}N+k\mathcal{L}_{u}N-\bar{D}^{2}N-\bar{D}^{a}N\bar{D}_{a}\log{\alpha}),$
$R_{AB}m^{A}q^{B}{}_{C}u^{C}\,=\,-\mathcal{L}_{u}\kappa+\bar{D}^{a}\rho_{a}-k_{ab}\kappa^{ab}+2\rho^{a}\bar{D}_{a}\log{\alpha}-\Theta
k,$ (78) $\displaystyle R_{AB}q^{A}{}_{C}q^{B}{}_{D}u^{C}u^{D}$
$\displaystyle=$
$\displaystyle-\frac{1}{N}\Theta^{{}^{\prime}}-\mathcal{L}_{u}\Theta+\Theta^{2}-\Theta\kappa-2\rho^{a}\rho_{a}-2\rho^{a}\bar{D}_{a}\log{\frac{N}{\alpha}}$
(79)
$\displaystyle~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}+\bar{D}^{a}\log{\alpha}\bar{D}_{a}\log{N}+\frac{1}{\alpha}\bar{D}^{2}\alpha-\mathcal{L}_{u}k-k_{ab}k^{ab}-\frac{1}{N}\mathcal{L}_{u}\mathcal{L}_{u}N,$
$R_{AB}m^{A}h^{Ba}\,=\,\Theta\bar{D}^{a}\log{\alpha}-2\rho_{b}k^{ab}-k\rho^{a}+\bar{D}_{b}\kappa^{ab}-\bar{D}^{a}\kappa+\bar{D}^{a}\Theta+\kappa^{ab}\bar{D}_{b}\log{\alpha}-\frac{1}{\alpha}(\partial_{u}\rho^{a}-\mathcal{L}_{\beta}\rho^{a}),$
(80) $\displaystyle R_{AB}q^{A}{}_{C}u^{C}h^{Bb}$ $\displaystyle=$
$\displaystyle\bar{D}_{a}k^{ab}-\bar{D}^{b}k-\rho^{b}\kappa-2\rho_{a}\kappa^{ab}-\frac{1}{N}\bar{D}^{b}\mathcal{L}_{u}N+k^{ab}\bar{D}_{a}\log{N}-\frac{1}{N}(\rho^{b})^{{}^{\prime}}$
(81)
$\displaystyle~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}~{}-\Theta\bar{D}^{b}\log{\frac{N}{\alpha}}-\kappa^{ab}\bar{D}_{a}\log{\frac{N}{\alpha}}-\frac{1}{\alpha}(\partial_{u}\rho^{b}-\mathcal{L}_{\beta}\rho^{b}),$
and
$\displaystyle
R_{AB}h^{A}{}_{a}h^{B}{}_{b}\,=\,{}^{(h)}R_{ab}+\mathcal{L}_{u}k_{ab}+kk_{ab}$
$\displaystyle-$ $\displaystyle
2k_{ac}k_{b}{}^{c}-\frac{1}{\alpha}\bar{D}_{a}\bar{D}_{b}\alpha-\frac{1}{N}\kappa_{ab}^{{}^{\prime}}-\frac{1}{\alpha}\dot{\kappa}_{ab}+\frac{1}{\alpha}\mathcal{L}_{\beta}\kappa_{ab}-2\rho_{(a}\bar{D}_{b)}\log{\frac{N}{\alpha}}$
(82) $\displaystyle+$
$\displaystyle(\Theta-\kappa)\kappa_{ab}-2\rho_{a}\rho_{b}+2\kappa_{ac}\kappa_{b}{}^{c}-\frac{1}{N}\bar{D}_{a}\bar{D}_{b}N+k_{ab}\mathcal{L}_{u}\log{N},$
where the dot denotes $\partial_{u}$. The vacuum Einstein equation is given by
$R_{AB}=0$.
## Appendix B Explicit form of Einstein equations
In Appendix B, we describe the components of the Einstein equations in terms
of the metric functions explicitly.
Using
$u=\alpha^{-1}(l-U^{a}\partial_{a})=\alpha^{-1}(\partial_{u}-U^{a}\partial_{a}),$
(83) $\Theta\,=\,-(A^{1/2})^{{}^{\prime}}+\frac{1}{2}\mathcal{L}_{u}\log{A},$
(84)
$\rho^{a}\,=\,\frac{1}{2}U^{a}{}^{{}^{\prime}}+\frac{1}{2}\bar{D}^{a}\log{A},$
(85)
and
$\kappa_{ab}\,=\,\frac{A^{1/2}}{2}h^{{}^{\prime}}_{ab}+\frac{1}{2A^{1/2}}(\dot{h}_{ab}-\bar{D}_{a}U_{b}-\bar{D}_{b}U_{a}),$
(86)
we can rewrite the Einstein equation $R_{AB}=0$ in terms of our metric form.
We will not provide all components of the Einstein equation explicitly.
For example, $R_{AB}h^{A}{}_{a}h^{B}{}_{b}=0$ becomes
$\displaystyle\mathcal{L}_{l}h^{{}^{\prime}}_{ab}+\frac{1}{2}A^{{}^{\prime}}h^{{}^{\prime}}_{ab}$
$\displaystyle=$
$\displaystyle{}^{(h)}R_{ab}+\frac{A}{2}h^{cd}h^{{}^{\prime}}_{ac}h^{{}^{\prime}}_{bd}+\frac{1}{2}h^{cd}(h^{{}^{\prime}}_{ac}\dot{h}_{bd}+h^{{}^{\prime}}_{bd}\dot{h}_{ac})-\frac{1}{2}h^{{}^{\prime}}_{ac}(\bar{D}_{b}U^{c}+\bar{D}^{c}U_{b})+\mathcal{L}_{U}h^{\prime}_{ab}$
(87)
$\displaystyle~{}~{}~{}~{}~{}-\frac{1}{2}h^{{}^{\prime}}_{bc}(\bar{D}_{a}U^{c}+\bar{D}^{c}U_{a})-\frac{1}{2}h_{ac}h_{bd}U^{c}{}^{{}^{\prime}}U^{d}{}^{{}^{\prime}}-\frac{1}{4}[\dot{(\log{h})}-2\bar{D}_{a}U^{a}]h^{{}^{\prime}}_{ab}$
$\displaystyle~{}~{}~{}~{}~{}-\frac{1}{4}(\log{h})^{{}^{\prime}}(Ah^{{}^{\prime}}_{ab}+\dot{h}_{ab}-\bar{D}_{a}U_{b}-\bar{D}_{b}U_{a})-\frac{1}{2}Ah^{{}^{\prime\prime}}_{ab}+\frac{1}{2}(h_{bc}\bar{D}_{a}U^{c}{}^{{}^{\prime}}+h_{ac}\bar{D}_{b}U^{c}{}^{{}^{\prime}}).$
This determines the evolutions of $h_{ab}$. For $R_{AB}m^{A}h^{Ba}=0$, we have
$\displaystyle\mathcal{L}_{l}U^{a}{}^{{}^{\prime}}$ $\displaystyle=$
$\displaystyle-\bar{D}^{a}A^{{}^{\prime}}-h^{ac}U^{b}{}^{{}^{\prime}}(\dot{h}_{bc}-\bar{D}_{a}U_{b}-\bar{D}_{b}U_{a})+h^{ab}\bar{D}^{c}(Ah^{{}^{\prime}}_{bc}+\dot{h}_{bc}-\bar{D}_{b}U_{c}-\bar{D}_{c}U_{b})-A\bar{D}^{a}(\log{h})^{{}^{\prime}}$
(88)
$\displaystyle~{}~{}~{}~{}-\frac{1}{2}U^{a}{}^{{}^{\prime}}[\dot{(\log{h})}-2\bar{D_{a}}U^{a}]-\frac{1}{2}(\log{h})^{{}^{\prime}}\bar{D}^{a}A-\bar{D}^{a}[\dot{(\log{h})}-2\bar{D}_{b}U^{b}]+\mathcal{L}_{U}U^{a}{}^{{}^{\prime}}.$
## Appendix C The gauge issue for $A^{(1)}$ [Eq. (40)]
In Appendix C we will show the presence of the gauge where Eq. (40) is
satisfied. In our coordinate $x^{A}=(u,r,x^{a})$, the metric can be written as
$ds^{2}\,=\,-Adu^{2}+2dudr+h_{ab}(dx^{a}+U^{a}du)(dx^{b}+U^{b}du).$ (89)
The components of the metric are expanded near the event horizon $(r=0)$ as
$A\,=\,rA^{(1)}+O(r^{2}),$ (90) $U^{a}\,=\,rU^{(1)a}+O(r^{2})$ (91)
and
$h_{ab}\,=\,h^{(0)}_{ab}+O(r).$ (92)
Here $A^{(1)}$ can be decomposed as
$A^{(1)}\,=\,A_{0}^{(1)}+\tilde{A}^{(1)}(u,x^{a}),$ (93)
where $A_{0}^{(1)}$ is set to be a constant as shown in Ref. Hollands:2006rj .
When we consider the gauge transformation $x^{A}\rightarrow x^{A}+\xi^{A}$,
the metric is transformed as
$g_{AB}\rightarrow g_{AB}+\mathcal{L}_{\xi}g_{AB}\equiv g_{AB}+\delta g_{AB}.$
(94)
To keep our gauge, the following conditions will be imposed:
$\displaystyle\delta g_{ur}\,=\,0\,,\,\delta g_{rr}\,=\,0\,,\,\delta
g_{ra}\,=\,0\,,$ $\displaystyle\delta g_{uu}\,=\,O(r)\,,\,\delta
g_{ua}\,=\,O(r)\,,\,\delta g_{ab}\,=\,O(1).$ (95)
From $\delta g_{rr}=2\partial_{r}\xi^{u}=0$, we have $\xi^{u}=f(u,x^{a})$.
Since $\delta g_{ra}$ and $\delta g_{ur}$ are given by
$\delta
g_{ra}\,=\,\partial_{a}\xi^{u}+U_{a}\partial_{r}\xi^{u}+h_{ab}\partial_{r}\xi^{b},$
(96) $\delta
g_{ur}\,=\,(-A+U^{a}U_{a})\partial_{r}\xi^{u}+\partial_{r}\xi^{r}+\partial_{u}\xi^{u}+U_{a}\partial_{r}\xi^{a},$
(97)
$\delta g_{ra}=0$ and $\delta g_{ur}=0$ lead to
$\xi^{r}=-r\partial_{u}f+\partial_{a}f\int^{r}U^{a}dr^{\prime},\quad\xi^{a}=-\partial_{b}f\int^{r}h^{ab}dr^{\prime}.$
(98)
Then $\delta g_{uu}$ becomes
$\delta g_{uu}\,=\,r[-\partial_{u}(fA^{(1)})-\partial^{2}_{u}f]+O(r^{2}),$
(99)
where we used Eqs. (90) and (91). This means that $A^{(1)}$ is transformed
under the gauge transformation as
$A^{(1)}\rightarrow A^{(1)}+\partial_{u}(fA^{(1)})+\partial_{u}^{2}f.$ (100)
Thus if we choose $f$ as
$f\,=\,-\frac{1}{A^{(1)}_{0}}\int^{u}du^{{}^{\prime}}\tilde{A}^{(1)},$ (101)
$A^{(1)}$ is transformed as
$A^{(1)}\rightarrow\bar{A}^{(1)}\,=\,A^{(1)}_{0}-\frac{1}{A^{(1)}_{0}}[\partial_{u}\tilde{A}^{(1)}+\partial_{u}(f\tilde{A}^{(1)})].$
(102)
Let assume that $\tilde{A}^{(1)}$ decays as $u\to\infty$. For the moment, we
write $\tilde{A}^{(1)}=O(1/u^{m})$, where $m$ is an integer. If $m\geq n$, Eq.
(40) is already satisfied. Therefore, we suppose that $m$ is smaller than $n$.
In the current gauge transformation, the transformed $\bar{A}^{(1)}$ satisfies
$\bar{A}^{(1)}\,=\,A^{(1)}_{0}+O(1/u^{m+1}).$ (103)
Repeating this procedure, we can always choose the gauge satisfying
$\mathcal{L}_{l}A^{(1)}\,=\,O\left(\frac{1}{u^{n+1}}\right).$ (104)
If one wishes, one can continue the same procedure and then achieve an
arbitrary faster decaying rate. But, the above is enough for our current
purpose.
## References
* (1) W. Israel, Phys. Rev. 164, 1776 (1967); B. Carter, Phys. Rev. Lett. 26, 331 (1971); S. W. Hawking, Commun. Math. Phys. 25, 152 (1972); D. C. Robinson, Phys. Rev. Lett. 34, 905 (1975); P. O. Mazur, J. Phys. A15, 3173 (1982); For a review, see M. Heusler, Black Hole Uniqueness Theorems, (Cambridge University Press, London, 1996); P. O. Mazur, hep-th/0101012; G. L. Bunting, PhD thesis, Univ. of New England, Armidale (1983).
* (2) S. W. Hawking and G. F. R. Ellis, The Large scale structure of space-time, (Cambridge Univ. Press, Cambridge, 1973).
* (3) S. W. Hawking, Commun. Math. Phys. 25, 152-166 (1972).
* (4) H. Friedrich, I. Racz, R. M. Wald, Commun. Math. Phys. 204, 691-707 (1999).
* (5) I. Racz, Class. Quant. Grav. 17, 153-178 (2000).
* (6) C. Gundlach, R. H. Price, J. Pullin, Phys. Rev. D49, 883-889 (1994).
* (7) L. Barack, A. Ori, Phys. Rev. Lett. 82, 4388 (1999).
* (8) L. Barack, A. Ori, Phys. Rev. D60, 124005 (1999).
* (9) K. Tanabe, S. Kinoshita, T. Shiromizu, Phys. Rev. D84, 044055 (2011).
* (10) S. A. Hayward, T. Shiromizu, K. -i. Nakao, Phys. Rev. D49, 5080-5085 (1994).
* (11) I. Booth, S. Fairhurst, Phys. Rev. Lett. 92, 011102 (2004).
* (12) I. Booth, S. Fairhurst, Phys. Rev. D75, 084019 (2007).
* (13) A. Ashtekar, B. Krishnan, Living Rev. Rel. 7, 10 (2004).
* (14) S. Hollands, A. Ishibashi, R. M. Wald, Commun. Math. Phys. 271, 699-722 (2007).
|
arxiv-papers
| 2011-12-02T09:30:17 |
2024-09-04T02:49:24.901658
|
{
"license": "Public Domain",
"authors": "Kentaro Tanabe, Tetsuya Shiromizu, Shunichiro Kinoshita",
"submitter": "Kentaro Tanabe",
"url": "https://arxiv.org/abs/1112.0408"
}
|
1112.0580
|
# Perfect Absorption in Ultrathin Epsilon-Near-Zero Metamaterials Induced by
Fast-Wave Non-Radiative Modes
Simin Feng simin.feng@navy.mil Klaus Halterman Michelson Lab, Physics
Division, Naval Air Warfare Center, China Lake, California 93555
###### Abstract
Above-light-line surface plasmon polaritons can arise at the interface between
a metal and $\epsilon$-near-zero metamaterial. This unique feature induces
unusual fast-wave non-radiative modes in a $\epsilon$-near-zero material/metal
bilayer. Excitation of this peculiar mode leads to wide-angle perfect
absorption in low-loss ultrathin metamaterials. The ratio of the perfect
absorption wavelength to the thickness of the $\epsilon$-near-zero
metamaterial can be as high as $10^{4}$; the electromagnetic energy can be
confined in a layer as thin as $\lambda/10000$. Unlike conventional fast-wave
leaky modes, these fast-wave non-radiative modes have quasi-static capacitive
features that naturally match with the space-wave field, and thus are easily
accessible from free space. The perfect absorption wavelength can be tuned
from mid- to far-infrared by tuning the $\epsilon\approx 0$ wavelength while
keeping the thickness of the structure unchanged.
###### pacs:
42.25.Bs, 78.67.Pt, 42.82.Et
A metamaterial is a composite structure with an electromagnetic (EM) response
not readily observed in naturally occurring materials. Many remarkable
phenomenon have been predicted Vesalago ; Feng ; Lai ; nguyen by tuning the
permittivity $\epsilon$ and permeability $\mu$ in extraordinary ways. If the
dielectric response is made vanishingly small, creating an epsilon-near-zero
(ENZ) material, interesting radiative effects are expected to occur Enoch ;
Silveirinha ; Halterman . On the other hand, by manipulating the EM response
to achieve a small transmittance ($T$) and reflectance ($R$), enhanced
absorption can ensue. Strong absorption in a thin layer typically requires
high loss. One of the earliest absorbers, the Salisbury screen Salisbury , is
based on the phenomenon of destructive wave interference, and is thus limited
to a minimum thickness of one quarter wavelength. To overcome the thickness
constraint, absorbing screens using metamaterials Engheta and high impedance
ground planes Sievenpiper have recently been proposed. By exploring resonant
enhancement, thin metamaterial and nanoplasmonic absorbers were demonstrated
in structures having localized resonances Landy ; XLiu ; Brown ; Li ; Avitzour
; Ye ; Kazemzadeh ; Costa ; Kravets ; NLiu ; Kang ; Chern ; atwater . In those
structures, the geometrical quality factor (GQF), i.e. the ratio of the
perfect absorption wavelength to the thickness of the medium, is significantly
improved compared to the standard Salisbury screen. The best GQF of those
structures is about 40 Li .
Based on a fundamentally different mechanism, in this paper we demonstrate
wide-angle perfect absorption in a low-loss ultrathin ENZ-metamaterial/metal
bilayer as shown in Fig. 1. In our structure, the GQF can be as high as
$10^{4}$. This structure possesses fast-wave non-radiative (FWNR) modes due to
unconventional above-light-line surface plasmon polaritons (ALL-SPPs) at the
ENZ-metal interface, which is easily accessible from free space. Fast waves
have phase velocities exceeding the speed of light in vacuum. Conventional
fast waves in planar structures are radiative leaky modes that cannot be
excited by plane wave incidence due to wave-vector mismatch at the boundary.
For our ENZ-metal structure, the exotic FWNR mode is naturally matched to free
space at the ENZ-air interface, and thus is easily accessible from free space.
Unlike the Salisbury screen where the thickness of the absorbers increases
with the absorption wavelength, in our structure the wavelength can be tuned
from mid- to far-infrared (IR) by tuning the $\epsilon\approx 0$ wavelength
while maintaining the thickness.
Figure 1: (Color online) A plane wave is incident (at an angle $\theta$) on an
anisotropic ($\epsilon_{1x}=1$ and $\epsilon_{1z}\approx 0$) ENZ/metal bilayer
with thicknesses $d_{1}$ and $d_{2}$ for the ENZ medium and metal (Ag),
respectively. Above-light-line surface plasmon polaritons (ALL-SPPs) are
excited at the ENZ-metal interface.
To investigate this absorption phenomenon in a general fashion, an anisotropic
ENZ metamaterial ($\epsilon_{1x}=1$, $\epsilon_{1z}\approx 0$) is assumed (see
Fig. 1). In the following, the subscripts 1 and 2 refer, respectively, to the
ENZ medium and metal. The areas above the ENZ material and below the metal are
free space and refer to regions 0 and 3, respectively. Since the
$\epsilon_{1x}$ has no significant impact on the result, $\epsilon_{1x}=1$ is
assumed throughout the paper. The permittivity of silver in the infrared
region is obtained by curve fitting experimental data Palik with a Drude
model, $\epsilon_{m}=1-\omega_{p}^{2}/\omega^{2}$, where the “plasma
frequency” $\omega_{p}=5.38\,\mu$m-1. The propagation constant of SPPs at the
interface of two semi-infinite anisotropic media is given by
$\frac{k_{spp}}{k_{0}}=\sqrt{\dfrac{\epsilon_{1z}\epsilon_{2z}\bigl{(}\epsilon_{2x}\mu_{1y}-\epsilon_{1x}\mu_{2y}\bigr{)}}{\epsilon_{2x}\epsilon_{2z}-\epsilon_{1x}\epsilon_{1z}}}\approx\sqrt{\epsilon_{1z}\mu_{1y}\left(1-\delta\right)}\,,$
(1)
where $k_{0}=\omega/c$ and
$\delta=\epsilon_{1x}\mu_{2y}/(\epsilon_{2x}\mu_{1y})$. The approximation is
taken when $\epsilon_{1z}\rightarrow 0$. In our case, $\mu_{1y}=\mu_{2y}=1$
and $\epsilon_{2x}=\epsilon_{2z}=\epsilon_{m}\sim-10^{5}$ (in the IR region).
Thus, $k_{spp}/k_{0}\approx\sqrt{\epsilon_{1z}\mu_{1y}}<1$, which
characterizes above light-line SPP dispersion, and is clearly different from
conventional SPP dispersion.
The absorption can be calculated from Maxwell’s equations. Assuming a harmonic
time dependence $\exp(-i\omega t)$ for the EM field, we have
$\displaystyle\begin{split}\nabla\times\bigl{(}\bar{\bar{\mu}}_{n}^{-1}\cdot\nabla\times{\bm{E}}\bigr{)}&=\,k_{0}^{2}\bigl{(}\bar{\bar{\epsilon}}_{n}\cdot{\bm{E}}\bigr{)}\,,\\\
\nabla\times\bigl{(}\bar{\bar{\epsilon}}_{n}^{-1}\cdot\nabla\times{\bm{H}}\bigr{)}&=\,k_{0}^{2}\bigl{(}\bar{\bar{\mu}}_{n}\cdot{\bm{H}}\bigr{)}\,,\end{split}$
(2)
where $\bar{\bar{\epsilon}}_{n}$ and $\bar{\bar{\mu}}_{n}$ are, respectively,
the permittivity and permeability tensors for a given (uniform) region
($n=0,1,2,\cdots$), which in the principal coordinates are described by,
$\displaystyle\bar{\bar{\epsilon}}_{n}$ $\displaystyle=$
$\displaystyle\epsilon_{nx}\hat{\bm{x}}\hat{\bm{x}}+\epsilon_{ny}\hat{\bm{y}}\hat{\bm{y}}+\epsilon_{nz}\hat{\bm{z}}\hat{\bm{z}}\,,$
(3) $\displaystyle\bar{\bar{\mu}}_{n}$ $\displaystyle=$
$\displaystyle\mu_{nx}\hat{\bm{x}}\hat{\bm{x}}+\mu_{ny}\hat{\bm{y}}\hat{\bm{y}}+\mu_{nz}\hat{\bm{z}}\hat{\bm{z}}\,.$
(4)
We consider TM modes, corresponding to non-zero field components $H_{y}$,
$E_{x}$, and $E_{z}$. The magnetic field $H_{y}$ satisfies the following wave
equation:
$\frac{1}{\epsilon_{z}}\frac{\partial^{2}H_{y}}{\partial
x^{2}}+\frac{1}{\epsilon_{x}}\frac{\partial^{2}H_{y}}{\partial
z^{2}}+k_{0}^{2}\mu_{y}H_{y}=0\,,$ (5)
which admits solutions of the form $\psi(z)\exp(i\beta x)$. The parallel wave-
vector $\beta$ is determined by the incident wave, and is conserved across the
interface,
$\beta^{2}=k_{0}^{2}\epsilon_{nz}\mu_{ny}-\alpha_{n}^{2}\frac{\epsilon_{nz}}{\epsilon_{nx}}\,,\hskip
14.45377pt(n=0,1,2,\cdots)\,,$ (6)
where $\alpha_{n}$ is the wave number in the $z$ direction. The functional
form of $\psi(z)$ is either a simple exponential $\exp(i\alpha_{n}z)$ for the
semi-infinite region or a superposition of $\cos(\alpha_{n}z)$ and
$\sin(\alpha_{n}z)$ terms for the bounded region (along the $z$ direction).
The other two components $E_{x}$ and $E_{z}$ are found using Maxwell’s
equations. By matching boundary conditions at the interface, i.e., continuity
of $H_{y}$ and $E_{x}$, the transmittance ($T$) and reflectance ($R$) can be
calculated via the Poynting vector $\bm{S}$, given by
$\bm{S}=c/8\pi\Re(\bm{E}\times\bm{H}^{*})$. The fraction of energy that is
absorbed by the system is determined by the absorptance ($A$), where
$A=1-T-R$, consistent with energy conservation.
Figure 2: (Color online) Absorptance vs. angle of incidence (see Fig. 1) at
the wavelength $\lambda=10.94\,\mu$m (left panels) and $\lambda=201.66\,\mu$m
(right panels) for two different thicknesses of the ENZ medium:
$d_{1}=0.02\,\mu$m (top panels) and $d_{1}=0.2\,\mu$m (bottom panels). The Ag-
layer thickness ($d_{2}$) is $20\,$nm.
Figure 2 shows the absorptance of the ENZ-metal structure versus the angle of
incidence (AOI) in the mid- to far-IR regime. Unless stated otherwise, the
results that follow correspond to $\epsilon_{1z}=0.001+i0.01$. The AOI at
which the perfect absorption occurs depends on the wavelength and the
thickness of the ENZ layer. The full-width-half-maximum (FWHM) angular width
can be as high as $\sim 60^{\circ}$ (top-left panel), while the GQF can reach
$\sim 10^{4}$ (top-right panel). When the thickness increases, the angular
bandwidth varies differently in the mid- and far-IR regions.
Figure 3: (Color online) Absorptance (right panels) and direction of perfect
absorption (left panels) when the ENZ-thickness $d_{1}=0.02\,\mu$m (top
panels) and $d_{1}=0.2\,\mu$m (bottom panels). Color-bars represent the
magnitude of the absorptance. The curves in the left panels are computed using
three different methods (see text). Solid (blue): numerically extracted from
the corresponding 2D plot in the right panels. Dashed (green): obtained from
Eq. (7). Circles (red): computed from Eq. (8). The Ag thickness is $20\,$nm.
In Fig. 3, we show how the absorptance varies with the AOI and incident
wavelength (right panels). The left set of panels corresponds to the incident
angle $\theta_{p}$ that results in perfect absorption, as a function of
wavelength. Perfect absorption occurs when the $\theta$ dependent effective
impedance (${\cal Z}_{e}$) of the ENZ-metal structure matches that of free
space (${\cal Z}_{0}$), i.e. ${\cal Z}_{e}={\cal Z}_{0}$, in which case the
reflection coefficient is zero. The effective impedance of the ENZ-metal
structure can be expressed as
${\cal Z}_{e}={\cal
Z}_{1}\frac{1-r_{12}\exp\bigl{(}i2\phi\bigr{)}}{1+r_{12}\exp\bigl{(}i2\phi\bigr{)}}={\cal
Z}_{0}\,,$ (7)
where $r_{12}=({\cal Z}_{1}-{\cal Z}_{2})/({\cal Z}_{1}+{\cal Z}_{2})$, ${\cal
Z}_{j}=\alpha_{j}/\epsilon_{jx}\,(j=0,1,2)$, and $\phi=\alpha_{1}d_{1}$.
Perfect absorption occurs at the particular AOI that satisfy Eq. (7). It is
also informative to analyze the corresponding waveguide modes, which are
solutions to the transcendental equation,
$\tan\phi=-i\frac{{\cal Z}_{1}({\cal Z}_{0}+{\cal Z}_{2})}{{\cal
Z}_{1}^{2}+{\cal Z}_{0}{\cal Z}_{2}},$ (8)
which implicitly depends on the parallel complex propagation constant $\beta$
(Eq. (6)). The mode solutions of Eq. (8) have four branches due to the $\pm$
signs inherent to the square root of $\alpha_{0}$ and $\alpha_{2}$ in Eq. (6).
The branch leading to perfect absorption corresponds to both $\alpha_{0}<0$
and $\alpha_{2}<0$. In the long wavelength regime and for ultrathin slabs,
this branch provides fast-wave ($\beta/k_{0}<1$), non-radiative ($\beta$ real)
modes that are guided along the ENZ layer. In the left panels of Fig. 3 we
compute the direction of perfect absorption as a function of wavelength in
three complimentary ways: The solid (blue) curves are numerically extracted
from the absorptance corresponding to the right panels, the dashed (green)
curves are obtained from Eq. (7), where the structure is impedance matched to
free space, and the circles (red) are calculated using Eq. (8) and a root
searching algorithm. The overlap of each set of data further demonstrates that
the underlying mechanism of perfect absorption is the excitation of FWNR
modes. We see also that perfect absorption occurs only at oblique incidence
where the electric field has a nonzero component in the normal direction.
The energy density, $U$, in lossy and dispersive anisotropic media is given
by,
$U=\frac{1}{16\pi}\Re\left[{\bm{E}}^{\dagger}\cdot\frac{\partial(\omega\bar{\bar{\epsilon}})}{\partial\omega}{\bm{E}}+{\bm{H}}^{\dagger}\cdot\frac{\partial(\omega\bar{\bar{\mu}})}{\partial\omega}{\bm{H}}\right].$
(9)
Figure 4: (Color online) Spatial distribution of the energy density computed
from Eq. (9) (normalized so that $H_{y}$ is unity at the ENZ-air interface)
for $\lambda=10\,\mu$m and AOI = $16.1^{\circ}$ (left) and $\lambda=155\,\mu$m
and AOI = $56.3^{\circ}$ (right). The ENZ-layer with $d_{1}=0.2\,\mu$m is
located at $z\in[-0.1\ 0.1]$. The regions $z\in[0.1\ 0.15]$ and $z\in[-0.2\
-0.1]$ are, respectively, silver and air. Color-bars represent the magnitude
($\times 0.01$) of the energy density, indicating the energy is more confined
to the ENZ medium at far-IR wavelengths. Figure 5: (Color online) Direction
(top left) and FWHM bandwidth (bottom left) of perfect absorption vs. ENZ
thickness. The curves are numerically extracted from the absorptance when
$\lambda=10\,\mu$m (dots (blue)), $\lambda=100\,\mu$m (solid (green)), and
$\lambda=300\,\mu$m (dashed (red)). The circles in the left-top panel are
obtained from Eq. (8). The normalized propagation constant $\beta/k_{0}$
corresponding to these modes is shown in the right panels. Top-right: real
part of $\beta/k_{0}$. Bottom-right: imaginary part of $\beta/k_{0}$. The Ag
thickness ($d_{2}$) is $10\,$nm.
We show in Fig. 4 the spatial distribution of the energy density in the ENZ-
metal structure under conditions of perfect absorption at mid- and far-IR
wavelengths. In the mid-IR (left panel), the field inside the ENZ-layer is
concentrated at the two interfaces, and a very thin layer of ALL-SPPs can be
seen at the ENZ-metal interface. The strong field at the ENZ-air interface is
distinguishable from conventional surface waves, since in this case the
exponential decay of the field occurs strictly inside the ENZ medium. Outside
of the ENZ medium, the field neither decays like a surface wave nor does it
increase like a leaky wave. The observed uniformity indicates that the field
outside the structure is plane wave matching with the free-space wave, and
thus can be easily excited by plane wave incidence – a striking difference
from surface waves and leaky waves. In the far-IR (right panel), the field is
even more uniform and confined to the ENZ-medium and also matched to the
space-wave field at the ENZ-air boundary. Around the point where the real part
of $\epsilon$ approaches zero, even with a small loss
$\Im(\epsilon_{1z})=0.01$, the presence of the FWNR mode reinforces the
increasing absorption due to the strong field inside the ENZ medium (by virtue
of the electric displacement continuity) leading to perfect absorption in the
ultrathin ENZ medium.
The minimum thickness of ENZ materials is generally limited by fabrication
methods except in some cases of naturally occurring resonances. Figure 5 shows
the direction, $\theta_{p}$ (top left), and FWHM angular bandwidth (bottom
left) of perfect absorption versus ENZ thickness for three wavelengths, as
well as the corresponding normalized complex propagation constants (right
panels). The circles in the top left panel represent solutions to Eq. (8),
which are consistent with the absorption calculated using the Poynting vector.
These curves also correlate with the solutions from the matched impedance
expression (Eq. (7)) (not shown). When the thickness increases, the direction
of perfect absorption is shifted towards the surface normal for both mid- and
far-IR regions. There is an optimal thickness that yields the largest
bandwidth. The real part of $\beta/k_{0}$ is above the vacuum light-line,
indicating a fast-wave character for all three wavelengths shown. The
imaginary component of $\beta$ is zero for the far-IR wavelengths. For
$\lambda=10\,\mu$m, $\Im(\beta)=0$ when $d_{1}<0.2\,\mu$m. When
$d_{1}>0.2\,\mu$m, $\beta$ starts to pick up a small imaginary part which
increases with thickness: The fast-wave non-radiative mode transforms into a
fast-wave leaky mode Halterman as the thickness increases.
Figure 6: (Color online) Average energy density (left top), transport energy
($S_{x}/c$) per unit area (left bottom), voltage ($V$) across the ENZ-layer
(middle top), the ratio of power to dissipation rate (middle bottom), the
effective surface charge density $Q$ (right top) and the normalized
capacitance $C/C_{0}$ (right bottom) vs. $d_{1}$ for $\lambda=10\,\mu$m (dots
(blue)), $\lambda=100\,\mu$m (solid (green)), and $\lambda=300\,\mu$m (dashed
(red)). Figure 7: (Color online) Color map of the absorptance versus AOI and
loss ($\Im(\epsilon_{1z})$) for $\lambda=10\,\mu$m (top panels) and
$\lambda=200\,\mu$m (bottom panels). The left and right set of panels
correspond to $\Re(\epsilon_{1z})=0.001$ and $\Re(\epsilon_{1z})=0.01$
respectively. Here, $d_{1}=0.15\,\mu$m and $d_{2}=10\,$nm. The color bars
represent the magnitude of the absorptance.
To further understand the FWNR modes, we integrated the transport power
$S_{x}$, dissipation rate $D$, and $\Re(E_{z})$ over the ENZ-thickness
(normalized by the input $E_{z}$). The dissipation rate is given by,
$D=\frac{\omega}{8\pi}\Im\Bigl{[}{\bm{E}}^{\dagger}\cdot\bar{\bar{\epsilon}}{\bm{E}}+{\bm{H}}^{\dagger}\cdot\bar{\bar{\mu}}{\bm{H}}\Bigr{]}.$
(10)
The average energy density in the ENZ medium is also calculated (Eq. (9)), as
well as the effective surface charge density $Q=\epsilon_{1z}E_{z}/4\pi$
(normalized by the input $E_{z}$) and capacitance $C=Q/V$ per unit area, which
is much larger than the capacitance ($C_{0}=\epsilon_{1z}/d_{1}$) of a
parallel-plate capacitor (of unit-area), due in part to the field enhancement
inside the ENZ material. These results are presented in Fig. 6. The
confinement is stronger for far-IR and decreases with the thickness (left-
top). The linear responses of the voltage (middle-top) and constant behavior
of the charge (right-top) as a function of thickness are hallmarks of
capacitors. The deviation from the linear behavior for $\lambda=10\,\mu$m can
be understood as a balance between radiative loss, which acquires more
“charge” stored in the “capacitor” (right- and middle-top) and energy
transport (left-bottom) (see right-bottom panel of Fig. 5 where
$\Im(\beta)\neq 0$ when $d_{1}>0.2\,\mu$m for $\lambda=10\,\mu$m). The
radiative loss reduces the enhancement of the capacitance at
$\lambda=10\,\mu$m as shown in the right-bottom panel. In essence, the higher
the confinement, the larger the capacitance. The effective capacitance is
found to be equal to $cQ^{2}/(2S_{x})$ or $2S_{x}/(cV^{2})$, which implies the
transport energy is driven by quasi-static “surface charges”. Thus, FWNR modes
can be considered quasi-static, capacitively driven fast waves. The middle-
bottom panel of Fig. 6 shows the ratio of transport power to dissipation
power, suggesting longer wavelengths have better power handling capability if
the structure is used as an ultrathin channel to transport light. The
influence of loss on the absorption is illustrated in Fig. 7 for different
$\epsilon_{1z}$ (real part) in the mid- and far-IR. In the far-IR, the smaller
$\Im(\epsilon_{1z})$ has larger angular bandwidth, which is opposite to what
occurs in the mid-IR.
In summary, wide-angle perfect absorbers have been demonstrated in ultrathin
ENZ-metal structures, and which have extraordinary high geometric quality
factors. Such structures may have applications involving ultra light-weight
infrared absorbers, cloaking, EM shielding, photovoltaic systems, and highly
efficient infrared sensors and detectors. The authors gratefully acknowledge
the sponsorship of ONR N-STAR and NAVAIR’s ILIR programs.
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* (8) W. W. Salisbury, U.S. Patent No. 2,599,944 (1952).
* (9) N. Engheta, IEEE Antennas and Propagation Society (AP-S) Int. Symp. v.2, p.392 (2002).
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|
arxiv-papers
| 2011-12-01T17:29:16 |
2024-09-04T02:49:24.915826
|
{
"license": "Public Domain",
"authors": "Simin Feng and Klaus Halterman",
"submitter": "Simin Feng",
"url": "https://arxiv.org/abs/1112.0580"
}
|
1112.0715
|
# A dynamical characterization of $C$ sets
John H. Johnson
Department of Mathematics, James Madison University, Harrisonburg, VA 22807,
USA
john.j.jr@gmail.com
###### Abstract
Furstenberg, using tools from topological dynamics, defined the notion of a
central subset of positive integers, and proved a powerful combinatorial
theorem about such sets. Using the algebraic structure of the Stone-Čech
compactification, this combinatorial theorem has been generalized and extended
to the Central Sets Theorem. The algebraic techniques also discovered many
sets, which are not central, that satisfy the conclusion of the Central Sets
Theorem. We call such sets $C$ sets. Since $C$ sets are defined
combinatorially, it is natural to ask if this notion admits a dynamical
characterization similar to Furstenberg’s original definition of a central
set? In this paper we give a positive answer to this question by proving a
dynamical characterization of $C$ sets.
## 1 Introduction
Furstenberg, in his book connecting dynamical systems with combinatorial
number theory, defined the concept of a central subset of positive integers
[5, Definition 8.3] and proved several important properties of such sets, all
using notions from topological dynamics. For instance, whenever a central set
is finitely partitioned, at least one cell of the partition contains a central
set [5, Theorem 8.8]. Many of the remaining important properties of central
sets follows from a powerful combinatorial theorem [5, Proposition 8.21] also
due to Furstenberg.
Inspired by the fruitful interaction between Ramsey Theory and ultrafilters on
semigroups, Bergelson and Hindman, with the assistance of B. Weiss, later
proved an algebraic characterization of central sets in $\mathbb{N}$ [1,
Section 6]. Using this algebraic characterization as a definition enabled them
to easily extend the notion of a central set to any semigroup. The definition
in [5] also extends naturally to arbitrary semigroups, and the algebraic and
dynamical characterizations were proved to be equivalent in general by H. Shi
and H. Yang in [11].
This algebraic definition turns out to have several advantages over the
original dynamical definition. For instance, the fact that central sets are
‘preserved under finite partitions’ (this is a concise way of stating [5,
Theorem 8.8]) easily follows from the algebraic definition. More importantly,
the combinatorial result [5, Proposition 8.21]—and stronger combinatorial
statements about central sets—follow from a relatively simple recursive
construction.
As an example, we state (currently) the strongest combinatorial theorem about
central sets commonly used. We first state this theorem for commutative
semigroups. (In Section 4 of this paper, we shall give the simplest statement
of the Central Sets Theorem currently known for arbitrary semigroups. The
statement of this ‘general’ version of the Central Sets Theorem is necessarily
complicated because of noncommutativity.)
In the statement of this theorem, and in the remainder of this paper, we let
$\mathcal{P}_{\\!f}(X)$ denote the collection of all nonempty finite subsets
of a set $X$; let ${}^{\hbox{$A$}}{\hskip-2.0ptB}$ denote the collection of
all functions with domain $A$ and codomain $B$; and, for typographical
convenience, we let
$\mathcal{T}=\hbox{${}^{\hbox{$\mathbb{N}$}}{\hskip-2.0ptS}$}$ for a given set
$S$. (Generally, the set $S$ in question will be clear from context.)
###### Theorem 1.1 (Central Sets Theorem).
Let $(S,+)$ be a commutative semigroup and $A\subseteq S$ central. Then there
exist functions $\alpha\colon\mathcal{P}_{\\!f}(\mathcal{T})\to S$ and
$H\colon\mathcal{P}_{\\!f}(\mathcal{T})\to\mathcal{P}_{\\!f}(\mathbb{N})$ that
satisfy the following two statements:
* (1)
If $F$, $G\in\mathcal{P}_{\\!f}(\mathcal{T})$ and $F\subsetneq G$, then $\max
H(F)<\min H(G)$.
* (2)
Whenever $m\in\mathbb{N}$, $G_{1}$, $G_{2}$, …, $G_{m}$ is a finite sequence
in $\mathcal{P}_{\\!f}(\mathcal{T})$ with $G_{1}\subsetneq
G_{2}\subsetneq\cdots\subsetneq G_{m}$ and for each $i\in\\{1,2,\ldots,m\\}$,
$f_{i}\in G_{i}$, then we have
$\sum_{i=1}^{m}\Bigl{(}\alpha(G_{i})+\sum_{t\in H(G_{i})}f_{i}\Bigr{)}\in A.$
###### Proof.
This was proved by De, Hindman, and Strauss in [4, Theorem 2.2]. ∎
###### Remark 1.2.
It’s an accident of history that [5, Proposition 8.21] is also known in the
literature as the Central Sets Theorem. Depending on how one counts, there are
about four different versions of ‘the’ Central Sets Theorem. (Happily each
newer version implies or easily reduces to the previous version.) From this
point on in this paper we shall only refer to Theorem 1.1 and Theorem 4.3 as
the Central Sets Theorem.
We shall call sets that satisfy the conclusion of the Central Sets Theorem $C$
sets. Despite the combinatorial power (and its name), the Central Sets Theorem
is not strong enough to combinatorially characterize central sets. In short
there are $C$ sets which are not central sets.
This somewhat surprising situation was first discovered in the context of an
important type of $C$ sets called the quasi-central sets. The quasi-central
sets were first defined algebraically [7, Definition 1.2] and given a
combinatorially characterization [7, Theorem 3.7] in a paper of Hindman,
Maleki, and Strauss. The fact that quasi-central sets also satisfy the
conclusion of the Central Sets Theorem follows from the proof of [4, Theorem
2.2].
Since quasi-central sets are defined algebraically, it is natural to wonder if
this notion admits a dynamical characterization similar to Furstenberg’s
original definition of central sets. In their recent paper, Burns and Hindman
prove such a dynamical characterization of quasi-central sets [2, Theorem
3.4].
However, their paper didn’t provide a dynamical characterization of $C$ sets.
(The fact that the notions of $C$ sets and quasi-central sets are distinct
follows from an example constructed in a recent paper [6] of Hindman.) In this
paper we fill this lacuna and prove a dynamical characterization of $C$ sets
in Theorem 4.8. This characterization will be a special case of a more general
result in Theorem 3.3 that gives a dynamical characterization of members of
idempotent ultrafilters in compact subsemigroups of the Stone-Čech
compactification.
### Acknowledgements
The characterization proved here is a generalization of part of the author’s
dissertation research conducted under the guidance of Neil Hindman. I want to
thank Dr. Hindman for his excellent advisement and helpful comments on this
paper.
## 2 Preliminaries on Compact Subsemigroups
In this section we state the basic definitions, conventions, and results we
need to prove our dynamical characterization of members of certain idempotent
ultrafilters. None of the results and definitions in this section are due to
the author. We also omit any proofs, but we do give references to where proofs
can be found.
We start by giving a brief review of the algebraic structure of the Stone-Čech
compactification of a discrete semigroup.
Given a discrete nonempty space $S$ we take the points of $\beta S$ to be the
collection of all ultrafilters on $S$. We identify points of $S$ with the
principal ultrafilters in $\beta S$. (Thus we pretend that $S\subseteq\beta
S$.) Given $A\subseteq S$, put $\overline{A}=\\{p\in\beta S:A\in p\\}$. Then
the collection $\\{\overline{A}:A\subseteq S\\}$ is a basis for a compact
Hausdorff topology on $\beta S$. This topology is the Stone-Čech
compactification of the discrete space $S$. The proofs for all of these
assertions can be found in [8, Sections 3.2 and 3.3].
Given a discrete semigroup $(S,\cdot)$, we can extend the semigroup operation
to $\beta S$ [8, Theorem 4.1] such that for $p$, $q\in\beta S$ and $A\subseteq
S$, we have $A\in p\cdot q$ if and only if $\\{x\in S:x^{-1}A\in q\\}\in p$
[8, Theorem 4.12] where $x^{-1}A=\\{y\in S:xy\in A\\}$. With this operation,
$(\beta S,\cdot)$ becomes a compact Hausdorff right-topological semigroup. The
word ‘right-topological’ means that for every $q\in\beta S$ the function
$\rho_{q}\colon\beta S\to\beta S$, defined by $\rho_{q}(p)=p\cdot q$, is
continuous.
###### Definition 2.1.
Let $S$ be a nonempty discrete space and $\mathcal{K}$ a filter on $S$.
* (a)
$\overline{\mathcal{K}}=\\{p\in\beta S:\mathcal{K}\subseteq p\\}$.
* (b)
$\mathcal{L}(\mathcal{K})=\\{A\subseteq S:S\setminus A\not\in\mathcal{K}\\}$.
As is well known, the function $\mathcal{K}\mapsto\overline{\mathcal{K}}$ is a
bijection from the collection of all filters on $S$ onto the collection of all
compact subspaces of $\beta S$ [8, Theorem 3.20]. We also have the following
important theorem relating the above two concepts.
###### Theorem 2.2.
Let $S$ be a nonempty discrete space and $\mathcal{K}$ a filter on $S$.
* (a)
$\overline{\mathcal{K}}=\\{p\in\beta S:\mbox{$A\in\mathcal{L}(\mathcal{K})$
for all $A\in p$}\\}$.
* (b)
Let $\mathcal{B}\subseteq\mathcal{L}(\mathcal{K})$ be closed under finite
intersections. Then there exists a $p\in\beta S$ with $\mathcal{B}\subseteq
p\subseteq\mathcal{L}(\mathcal{K})$.
###### Proof.
Both of these assertions follow from [8, Theorem 3.11]. ∎
If $(S,\cdot)$ is discrete semigroup and $\mathcal{K}$ a filter on $S$, then
there are precise conditions on $\mathcal{K}$ which guarantee that
$\overline{\mathcal{K}}$ is a compact subsemigroup of $\beta S$ [3, Theorem
2.6]. Hence $\overline{\mathcal{K}}$ is a compact Hausdorff right-topological
semigroup for a suitable filter $\mathcal{K}$.
###### Theorem 2.3.
Let $T$ be a compact Hausdorff right-topological semigroup.
* (a)
$T$ contains at least one idempotent, that is, there exists $x\in T$ such that
$x=x\cdot x$.
* (b)
$T$ contains an ideal, called the smallest ideal and denoted as $K(T)$, that
is contained in every ideal of $T$. Additionally, $K(T)$ also contains at
least one idempotent.
###### Proof.
The proofs of statements (a) and (b) are given in [8, Theorem 2.5] and [8,
Theorem 2.8], respectively. ∎
###### Remark 2.4.
It does follow that $c\ell_{T}K(T)$ is a compact subsemigroup of $T$, and
hence by (a) this subsemigroup also contains an idempotent. While the smallest
ideal $K(T)$ itself may not be closed, it is the union of all of the minimal
left ideals of $T$, which are closed, so the fact that it contains an
idempotent is also immediate.
We are now in a position to give the algebraic definitions of a central set
and quasi-central set in a semigroup.
###### Definition 2.5.
Let $(S,\cdot)$ be a semigroup and $A\subseteq S$.
* (a)
We call $A$ a central set if and only if there exists an idempotent $p\in
K(\beta S)$ such that $A\in p$.
* (b)
We call $A$ a quasi-central set if and only if there exists an idempotent
$p\in c\ell_{\beta S}K(\beta S)$ such that $A\in p$.
To finish this section, we give the definition of a dynamical system and
relate this notion to the algebraic structure of the Stone-Čech
compactification.
###### Definition 2.6.
A pair $(X,\langle T_{s}\rangle_{s\in S})$ is a dynamical system if and only
if it satisfies the following four conditions:
* (1)
$X$ is a compact Hausdorff space.
* (2)
$S$ is a semigroup.
* (3)
$T_{s}\colon X\to X$ is continuous for every $s\in S$.
* (4)
For every $s$, $t\in S$ we have $T_{st}=T_{s}\circ T_{t}$.
###### Theorem 2.7.
Let $(X,\langle T_{s}\rangle_{s\in S})$ be a dynamical system. Then we can
extend $(X,\langle T_{s}\rangle_{s\in S})$ to a semigroup action on $\beta S$.
More precisely, for each $p\in\beta S\setminus S$ we can define $T_{p}\colon
X\to X$ such that for every $q$, $r\in\beta S$, $T_{qr}=T_{q}\circ T_{r}$.
Furthermore, given $p\in\beta S$, $x$ and $y$ in $X$, we have $T_{p}(x)=y$ if
and only if for every neighborhood $U$ of $y$, $\\{s\in S:T_{s}(x)\in U\\}\in
p$.
###### Proof.
Both of these assertions follow from [8, Theorems 3.27 and Corollary 4.22]. ∎
Be warned that for $p\in\beta S\setminus S$, $T_{p}$ is usually not
continuous.
## 3 Dynamical Characterization of Members of Idempotent Ultrafilters
###### Definition 3.1.
Let $(X,\langle T_{s}\rangle_{s\in S})$ be a dynamical system, $x$ and $y$
points in $X$, and $\mathcal{K}$ a filter on $S$. The pair $(x,y)$ is called
jointly $\mathcal{K}$-recurrent if and only if for every neighborhood $U$ of
$y$ we have $\\{s\in S:\mbox{$T_{s}(x)\in U$ and $T_{s}(y)\in
U$}\\}\in\mathcal{L}(\mathcal{K})$.
###### Lemma 3.2.
Let $(X,\langle T_{s}\rangle_{s\in S})$ be a dynamical system, let $x$ and $y$
be points in $X$, and let $\mathcal{K}$ be a filter $S$ such that
$\overline{\mathcal{K}}$ is a compact subsemigroup of $\beta S$. The following
statements are equivalent.
* (a)
The pair $(x,y)$ is jointly $\mathcal{K}$-recurrent.
* (b)
There exists $p\in\overline{\mathcal{K}}$ such that $T_{p}(x)=y=T_{p}(y)$.
* (c)
There exists an idempotent $p\in\overline{\mathcal{K}}$ such that
$T_{p}(x)=y=T_{p}(y)$.
###### Proof.
(a) $\implies$ (b). For each neighborhood $U$ of $y$ put $B_{U}=\\{s\in
S:\mbox{$T_{s}(x)\in U$ and $T_{s}(y)\in U$}\\}$. Observe that since $B_{U\cap
V}=B_{U}\cap B_{V}$ for $U$ and $V$ neighborhoods of $y$, we have that the
collection $\mathcal{B}=\\{B_{U}:\mbox{$U$ is a neighborhood of $y$}\\}$ is
closed under finite intersections. Also, by assumption we have that
$\mathcal{B}\subseteq\mathcal{L}(\mathcal{K})$. Hence by Theorem 2.2 we can
pick $p\in\overline{\mathcal{K}}$ with $\mathcal{B}\subseteq p$.
For every neighborhood $U$ of $y$ we have $B_{U}\subseteq\\{s\in S:T_{s}(x)\in
U\\}$ and $B_{U}\subseteq\\{s\in S:T_{s}(y)\in U\\}$. Therefore $\\{s\in
S:T_{s}(x)\in U\\}\in p$ and $\\{s\in S:T_{s}(y)\in U\\}\in p$. It now follows
from Theorem 2.7 that $T_{p}(x)=y=T_{p}(y)$.
(b) $\implies$ (c). Put
$M=\\{p\in\overline{\mathcal{K}}:T_{p}(x)=y=T_{p}(y)\\}$. By Theorem 2.3 it
suffices to show that $M$ is a compact subsemigroup of $\beta S$.
$M$ is nonempty by assumption.
To see that $M$ is compact, we simply show that $M$ is closed. Let $p\not\in
M$, then either $T_{p}(x)\neq y$ or $T_{p}(y)\neq y$. First assume that
$T_{p}(x)\neq y$. By Theorem 2.7 pick $U$ a neighborhood of $y$ such that
$\\{s\in S:T_{s}(x)\in U\\}\not\in p$. Put $A=\\{s\in S:T_{s}(x)\in U\\}$ and
note that $S\setminus A\in p$. We have that $(\overline{S\setminus A})\cap
M=\emptyset$, that is, $\overline{S\setminus A}$ is a (basic) neighborhood of
$p$ that misses $M$. (If $q\in(\overline{S\setminus A})\cap M$, then it
follows that $A\in q$ and $S\setminus A\in q$, a contradiction.) The
construction of a (basic) neighborhood of $p$ that misses $M$ when
$T_{p}(y)\neq y$ is similar. Therefore $M$ is a closed subset of $\beta S$.
To see that $M$ is a subsemigroup, let $q$, $r\in M$. Then by Theorem 2.7 and
assumption we have $T_{qr}(x)=T_{q}\circ T_{r}(x)=T_{q}(y)=y=T_{q}\circ
T_{r}(y)=T_{qr}(y)$. ∎
###### Theorem 3.3 (Main Result).
Let $(S,\cdot)$ be a semigroup, let $\mathcal{K}$ be a filter on $S$ such that
$\overline{\mathcal{K}}$ is a compact subsemigroup of $\beta S$, and let
$A\subseteq S$. Then $A$ is a member of an idempotent in
$\overline{\mathcal{K}}$ if and only if there exists a dynamical system
$(X,\langle T_{s}\rangle_{s\in S})$ with points $x$ and $y$ in $X$ and there
exists a neighborhood $U$ of $y$ such that the pair $(x,y)$ is jointly
$\mathcal{K}$-recurrent and $A=\\{s\in S:T_{s}(x)\in U\\}$.
###### Proof.
($\Rightarrow$) Let $R=S\cup\\{e\\}$ be the semigroup with an identity $e$
adjoined to $S$. (For expository convenience, we add this new identity even if
$S$ already contains an identity.) Give $\\{0,1\\}$ the discrete topology and
give $X=\hbox{${}^{\hbox{$R$}}{\hskip-2.0pt\\{0,1\\}}$}$ the product topology.
Then $X$ is a compact Hausdorff space.
For each $s\in S$, define $T_{s}\colon X\to X$ by $T_{s}(f)=f\circ\rho_{s}$.
It’s a routine exercise, or see [8, Theorem 19.14], to show that $(X,\langle
T_{s}\rangle_{s\in S})$ is a dynamical system.
Now let $x=\mathbf{1}_{A}$ be the characteristic function of $A$, pick an
idempotent $r\in\overline{\mathcal{K}}$ with $A\in r$, and put $y=T_{r}(x)$.
Then we have that
$T_{r}(y)=T_{r}\bigl{(}T_{r}(x)\bigr{)}=T_{rr}(x)=T_{r}(x)=y$. Therefore by
Lemma 3.2 we have that the pair $(x,y)$ is jointly $\mathcal{K}$-recurrent.
Put $U=\\{w\in X:w(e)=y(e)\\}$ and observe that $U$ is a (subbasic) open
neighborhood of $y$. (The set $U$ is equal to the inverse image of
$\\{y(e)\\}$ under the projection map.) To help us show that $U$ is the
neighborhood of $y$ we are looking for, we first will show that $y(e)=1$.
Since $y=T_{r}(x)$ we have that $\\{s\in S:T_{s}(x)\in U\\}\in r$ by Theorem
2.7. We can pick $s\in A$ such that $T_{s}(x)\in U$. Then by definition of $U$
we have that $y(e)=T_{s}(x)(e)=x\bigl{(}\rho_{s}(e)\bigr{)}=x(es)=x(s)$. Also
by our choice of $s\in A$ we have $x(s)=\mathbf{1}_{A}(s)=1$.
To finish up this direction observe that for $s\in S$ the following logical
relation is true:
$\displaystyle s\in A$ $\displaystyle\iff\mathbf{1}_{A}(s)=1,$
$\displaystyle\iff x(s)=1,$ $\displaystyle\iff x(es)=1,$
$\displaystyle\iff(x\circ\rho_{s})(e)=1,$ $\displaystyle\iff
T_{s}(x)(e)=1=y(e),$ $\displaystyle\iff T_{s}(x)\in U.$
Hence $A=\\{s\in S:T_{s}(x)\in U\\}$.
($\Leftarrow$) Pick a dynamical system $(X,\langle T_{s}\rangle_{s\in S})$,
pick points $x$ and $y$ in $X$, and pick $U$ a neighborhood of $y$ as
guaranteed by assumption. By Lemma 3.2 pick an idempotent
$r\in\overline{\mathcal{K}}$ such that $T_{r}(x)=y=T_{r}(y)$. Since $U$ is a
neighborhood of $y$ and $T_{r}(x)=y$, we have $A=\\{s\in S:T_{s}(x)\in U\\}\in
r$ by Theorem 2.7. ∎
###### Remark 3.4.
One remarkable thing about Lemma 3.2 and Theorem 3.3 is that, while the
results are much more general, the proofs are essentially trivial
modifications of the proofs of [2, Lemma 3.3 and Theorem 3.4].
## 4 A Dynamical Characterization of $C$ sets
In this section we give an application of our main result in Section 3 to
prove a dynamical characterization of $C$ sets. We start by giving the
combinatorial definition of a $C$ set. As mentioned in Section 1 this
combinatorial definition is rather complicated; but we shall soon state an
algebraic characterization showing that $C$ sets are members of idempotents in
a certain compact subsemigroup.
In the following definition, given an indexed family $\langle A_{i}:i\in
I\rangle$ of sets, we let $\hbox{\bigmath\char 2\relax}_{i\in I}A_{i}$
represent its cartesian product. (We reserve the use of the symbol $\prod$ for
our semigroup operation.) Recall from Section 1 that given a set $S$ we let
$\mathcal{T}=\hbox{${}^{\hbox{$\mathbb{N}$}}{\hskip-2.0ptS}$}$ and
$\mathcal{P}_{\\!f}(X)$ is the collection of all nonempty finite subsets of
$X$.
###### Definition 4.1.
Let $(S,\cdot)$ be a semigroup.
* (a)
For each positive integer $m$ put
$\mathcal{J}_{m}=\\{(t_{1},t_{2},\ldots,t_{m})\in\mathbb{N}^{m}:t_{1}<t_{2}<\cdots<t_{m}\\}$.
* (b)
Given $m\in\mathbb{N}$, $a\in S^{m+1}$, $t\in\mathcal{J}_{m}$, and
$f\in\mathcal{T}$, put
$x(m,a,t,f)=\prod_{i=1}^{m}\bigr{(}a(i)f(t_{i})\bigl{)}a(m+1)$.
* (c)
We call a subset $A\subseteq S$ a $C$ set if and only if there exist functions
$m\colon\mathcal{P}_{\\!f}(\mathcal{T})\to\mathbb{N}$,
$\alpha\in\hbox{\bigmath\char
2\relax}_{F\in\mathcal{P}_{\\!f}(\mathcal{T})}S^{m(F)+1}$, and
$\tau\in\hbox{\bigmath\char
2\relax}_{F\in\mathcal{P}_{\\!f}(\mathcal{T})}\mathcal{J}_{m(F)}$ such that
the following two statements are satisfied:
* (1)
If $F$, $G\in\mathcal{P}_{\\!f}(\mathcal{T})$ and $F\subsetneq G$, then
$\tau(F)\bigl{(}m(F)\bigr{)}<\tau(G)(1)$.
* (2)
Whenever $m\in\mathbb{N}$, $G_{1}$, $G_{2}$, …, $G_{m}$ is a finite sequence
in $\mathcal{P}_{\\!f}(\mathcal{T})$ with $G_{1}\subsetneq
G_{2}\subsetneq\cdots\subsetneq G_{m}$, and for each $i\in\\{1,2,\ldots,m\\}$,
$f_{i}\in G_{i}$, then we have
$\prod_{i=1}^{m}x(m(G_{i}),\alpha(G_{i}),\tau(G_{i}),f_{i})\in A.$
###### Remark 4.2.
This definition of a $C$ set is different from the original (and more
complicated) definition given in [4, Definition 3.3(i)]. It is a result in the
author’s dissertation, to be included in a forthcoming paper [10], that our
simpler definition of a $C$ set is equivalent to the original definition.
Before giving an algebraic characterization of $C$ sets, we pause to state the
Central Sets Theorem.
###### Theorem 4.3 (Central Sets Theorem).
Central sets in a semigroup are $C$ sets.
###### Proof.
This is proved in [10] using the our definition of a $C$ set, and is proved in
[4, Corollary 3.10] under the original definition. ∎
To give the algebraic definition of a $C$ set we shall need the following
(curiously named) combinatorial notion closely related to $C$ sets.
###### Definition 4.4.
Let $(S,\cdot)$ be a semigroup.
* (a)
We call a subset $A\subseteq S$ a $J$ set if and only if for every
$F\in\mathcal{P}_{\\!f}(\mathcal{T})$, there exist $m\in\mathbb{N}$, $a\in
S^{m+1}$, and $t\in\mathcal{J}_{m}$ such that for all $f\in F$, $x(m,a,t,f)\in
A$.
* (b)
$J(S)=\\{p\in\beta S:\mbox{$A$ is a $J$ set for every $A\in p$}\\}$.
###### Remark 4.5.
I must point out that $J$ sets are not named after the author! The term $J$
set is derived from the term ‘$J_{Y}$ set’ introduced as [7, Definition
2.4(b)] in a different and earlier paper. This definition of a $J$ set is also
different from the original (and more complicated) definition given in [4,
Definition 3.3(e)]. The fact that these two definitions are equivalent is
proved by the author in his dissertation and in the forthcoming paper [10].
###### Theorem 4.6.
Let $(S,\cdot)$ be a semigroup and $\mathcal{K}=\\{A\subseteq
S:\mbox{$S\setminus A$ is not a $J$ set}\\}$. Then $\mathcal{K}$ is a filter
on $S$ with $J(S)=\overline{\mathcal{K}}$ and $J(S)$ is a compact subsemigroup
of $\beta S$.
###### Proof.
To show that $\mathcal{K}$ is nonempty, doesn’t contain the empty set, and is
closed under supersets is a routine exercise. The fact that $\mathcal{K}$ is
closed under finite intersections follows from [10] using the new definition
or [9, Theorem 2.14] using the old equivalent definition of $J$ sets.
Observe that, under the assumption that $\mathcal{K}$ is a filter,
$\mathcal{L}(\mathcal{K})=\\{A\subseteq S:\mbox{$A$ is a $J$ set}\\}$. Hence
it follows from [10] (new definition) or [8, Theorem 3.11] (old equivalent
definition) that $J(S)=\overline{\mathcal{K}}$.
Finally, the fact that $J(S)$ is a subsemigroup follows from [4, Theorem 3.5].
∎
Since $J(S)$ is a compact subsemigroup, in fact an ideal, of $\beta S$, by
Theorem 2.3 we have that $J(S)$ contains idempotents. The following theorem
connects these idempotent elements with $C$ sets.
###### Theorem 4.7.
Let $(S,\cdot)$ be a semigroup and $A\subseteq S$. Then $A$ is a $C$ set if
and only if there exists an idempotent $p\in J(S)$ such that $A\in p$.
###### Proof.
This is proved in [10] for the new definition or [4, Theorem 3.8] for the old
definition. ∎
Using these facts and the general results in Section 3 we end with the
following dynamical characterization of $C$ sets.
###### Theorem 4.8.
Let $(S,\cdot)$ be a semigroup and $A\subseteq S$. Then $A$ is a $C$ set if
and only if there exists a dynamical system $(X,\langle T_{s}\rangle_{s\in
S})$ with points $x$ and $y$ in $X$ and there exists a neighborhood $U$ of $y$
such that $\\{s\in S:\mbox{$T_{s}(x)\in U$ and $T_{s}(y)\in U$}\\}$ is a $J$
set and $A=\\{s\in S:T_{s}(x)\in U\\}$.
###### Proof.
Let $\mathcal{K}=\\{B\subseteq S:\mbox{$S\setminus B$ is not a $J$ set}\\}$
and note by Theorem 4.6 that $\overline{\mathcal{K}}=J(S)$ and
$\mathcal{L}(\mathcal{K})=\\{A\subseteq S:\mbox{$A$ is a $J$ set}\\}$. Since
Theorem 4.7 characterizes $C$ sets in terms of idempotents in
$\overline{\mathcal{K}}$, we can apply Theorem 3.3 to prove our statement. ∎
## References
* [1] Vitaly Bergelson and Neil Hindman, Nonmetrizable topological dynamics and Ramsey theory, Transactions of the American Mathematical Society 320 (1990), no. 1, 293–320.
* [2] Shea D. Burns and Neil Hindman, Quasi-central sets and their dynamical characterization, Topology Proceedings 31 (2007), no. 2, 445–455.
* [3] Dennis Davenport, The minimal ideal of compact subsemigroups of $\beta S$, Semigroup Forum 41 (1990), 201–213.
* [4] Dibyendu De, Neil Hindman, and Dona Strauss, A new and stronger central sets theorem, Fundamenta Mathematicae 199 (2008), 155–175.
* [5] H. Furstenberg, Recurrence in Ergodic Theory and Combinatorial Number Theory, Princeton University Press, 1981.
* [6] Neil Hindman, Small sets satisfying the Central Sets Theorem, Integers 9 Supplement (2009), article 5.
* [7] Neil Hindman, Amir Maleki, and Dona Strauss, Central sets and their combinatorial characterization, Journal of Combinatorial Theory (Series A) 74 (1996), 188–208.
* [8] Neil Hindman and Dona Strauss, Algebra in the Stone-Čech Compactification, De Gruyter expositions in mathematics, no. 27, Walter de Gruyter, 1998.
* [9] , Cartesian products of sets satisfying the Central Sets Theorem, Topology Proceedings 35 (2010), 203–223.
* [10] John H. Johnson, A new and simpler Central Sets Theorem, (in preparation), November 2011.
* [11] Hong ting Shi and Hong wei Yang, Nonmetrizable topological dynamical characterization of central sets, Fundamenta Mathematicae (1996), no. 150, 1–9.
|
arxiv-papers
| 2011-12-04T04:10:01 |
2024-09-04T02:49:24.924950
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "John H. Johnson",
"submitter": "John Johnson",
"url": "https://arxiv.org/abs/1112.0715"
}
|
1112.0722
|
arxiv-papers
| 2011-12-04T05:03:14 |
2024-09-04T02:49:24.931287
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Ion Olaru",
"submitter": "Ion Olaru",
"url": "https://arxiv.org/abs/1112.0722"
}
|
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1112.0724
|
# Trudinger-Moser inequalities on complete noncompact Riemannian manifolds
Yunyan Yang yunyanyang@ruc.edu.cn Department of Mathematics, Renmin
University of China, Beijing 100872, P. R. China
###### Abstract
Let $(M,g)$ be a complete noncompact Riemannian $n$-manifold ($n\geq 2$). If
there exist positive constants $\alpha$, $\tau$ and $\beta$ such that
$\sup_{u\in W^{1,n}(M),\,\|u\|_{1,\tau}\leq
1}\int_{M}\left(e^{\alpha|u|^{\frac{n}{n-1}}}-\sum_{k=0}^{n-2}\frac{\alpha^{k}|u|^{\frac{nk}{n-1}}}{k!}\right)dv_{g}\leq\beta,$
where $\|u\|_{1,\tau}=\|\nabla_{g}u\|_{L^{n}(M)}+\tau\|u\|_{L^{n}(M)}$, then
we say that Trudinger-Moser inequality holds. Suppose Trudinger-Moser
inequality holds, we prove that there exists some positive constant $\epsilon$
such that ${\rm Vol}_{g}(B_{x}(1))\geq\epsilon$ for all $x\in M$. Also we give
a sufficient condition under which Trudinger-Moser inequality holds, say the
Ricci curvature of $(M,g)$ has lower bound and its injectivity radius is
positive. Moreover, Adams inequality is discussed in this paper. For
application of Trudinger-Moser inequalities, we obtain existence results for
some quasilinear equations with nonlinearity of exponential growth.
###### keywords:
Trudinger-Moser inequality, Adams inequality, exponential growth
###### MSC:
58E35, 35J60
††journal: ***
## 1 Introduction
Let $\Omega$ be a smooth bounded domain in $\mathbb{R}^{n}$ ($n\geq 2)$ and
$C_{0}^{\infty}(\Omega)$ be a space of smooth functions with compact support
in $\Omega$. Let $W_{0}^{m,p}(\Omega)$ be the completion of
$C_{0}^{\infty}(\Omega)$ under the Sobolev norm
$\|u\|_{W_{0}^{m,p}(\Omega)}:=\left(\sum_{l=0}^{m}\int_{\Omega}|\nabla^{l}u|^{p}dx\right)^{1/p}.$
(1.1)
Assume that $m$ is an integer satisfying $1\leq m<n$. Then Sobolev embedding
theorem asserts that $W_{0}^{m,p}(\Omega)\hookrightarrow L^{q}(\Omega)$,
$1\leq q\leq{np}/{(n-mp)}$. Concerning the limiting case $mp=n$, one has
$W_{0}^{m,n/m}(\Omega)\hookrightarrow L^{q}(\Omega)$ for all $q\geq 1$. But
the embedding is not valid for $q=\infty$. To fill this gap, it is natural to
find the maximal growth function $g:\mathbb{R}\rightarrow\mathbb{R}^{+}$ such
that
$\sup_{u\in W_{0}^{m,n/m}(\Omega),\,\|u\|_{W_{0}^{m,n/m}(\Omega)}\leq
1}\int_{\Omega}g(u)dx<\infty.$
In the case $m=1$, Trudinger [38] and Pohozaev [33] found independently that
the maximal growth is of exponential type. More precisely, there exist two
positive constants $\alpha_{0}$ and $C$ depending only on $n$ such that
$\sup_{u\in W_{0}^{1,n}(\Omega),\,\|u\|_{W_{0}^{1,n}(\Omega)}\leq
1}\int_{\Omega}e^{\alpha_{0}|u|^{\frac{n}{n-1}}}dx\leq C|\Omega|,$ (1.2)
where $|\Omega|$ denotes the Lebesgue measure of $\Omega$. Moser [30] obtained
the best constant $\alpha_{n}=n\omega_{n-1}^{1/(n-1)}$ such that the above
supremum is finite when $\alpha_{0}$ is replaced by $\alpha_{n}$, where
$\omega_{n-1}$ is the area of the unit sphere in $\mathbb{R}^{n}$. Moser’s
work relies on a rearrangement argument [17]. In literature the kind of
inequalities like (1.2) are called Trudinger-Moser inequalities.
Adams [2] generalized inequality (1.2) to the case of general $m:1\leq m<n$ as
follows. For any $u\in W_{0}^{m,n/m}(\Omega)$, the $l$-th order gradient of
$u$ reads
$\nabla^{l}u=\left\\{\begin{array}[]{lll}\Delta^{\frac{l}{2}}u,&{\rm
if}\,\,\,l\,\,\,{\rm is}\,\,\,{\rm even,}\\\\[6.45831pt]
\nabla\Delta^{\frac{l-1}{2}}u,&{\rm if}\,\,\,l\,\,\,{\rm is}\,\,\,{\rm
odd,}\end{array}\right.$ (1.3)
there exits a positive constant $C_{m,n}$ such that
$\sup_{u\in W_{0}^{m,n/m}(\Omega),\,\|u\|_{W_{0}^{m,n/m}(\Omega)}\leq
1}\int_{\Omega}e^{\beta_{0}|u|^{\frac{n}{n-m}}}dx\leq C_{m,n}|\Omega|,$ (1.4)
where $\beta_{0}$ is the best constant depending only on $n$ and $m$, namely
$\beta_{0}=\beta_{0}(m.n):=\left\\{\begin{array}[]{lll}\frac{n}{\omega_{n-1}}\left[\frac{\pi^{{n}/{2}}2^{m}\Gamma\left({(m+1)}/{2}\right)}{\Gamma\left({(n-m+1)}/{2}\right)}\right]&{\rm
when}\,\,\,m\,\,\,{\rm is}\,\,\,odd\\\\[6.45831pt]
\frac{n}{\omega_{n-1}}\left[\frac{\pi^{{n}/{2}}2^{m}\Gamma\left({m}/{2}\right)}{\Gamma\left({(n-m)}/{2}\right)}\right]&{\rm
when}\,\,\,m\,\,\,{\rm is}\,\,\,even.\end{array}\right.$ (1.5)
The inequality (1.4) is known as Adams inequality. Adams first represented a
function $u$ in terms of its gradient function $\nabla^{m}u$ by using a
convolution operator. Then using the O’Neil’s idea [31] of rearrangement of
convolution of two functions and the idea which originally goes back to
Garcia, he obtained (1.4).
There are many types of extensions for Trudinger-Moser inequality and Adams
inequality. One is to establish such inequalities on the whole euclidian space
$\mathbb{R}^{n}$. Cao [8] employed the decreasing rearrangement argument to
prove that for all $\alpha<4\pi$ and $A>0$, there exists a constant $C$
depending only on $\alpha$ and $A$ such that for all $u\in
W^{1,2}(\mathbb{R}^{2})$ with $\int_{\mathbb{R}^{2}}|\nabla u|^{2}dx\leq
1,\,\int_{\mathbb{R}^{2}}u^{2}dx\leq A$, there holds
$\int_{\mathbb{R}^{2}}\left(e^{\alpha u^{2}}-1\right)dx\leq C.$ (1.6)
His argument was generalized to $n$-dimensional case by do Ó [12] and Panda
[32] independently. Later, Adachi-Tanaka [1] gave another type of
generalization. All these inequalities are subcritical ones since
$\alpha<\alpha_{n}$. It was Ruf [35] who first proved the critical Trudinger-
Moser inequality in the whole euclidian space $\mathbb{R}^{2}$ and gave out
extremal functions via more delicate analysis. This result was generalized to
$n$-dimensional case by Li-Ruf [25] through combining symmetrization and blow-
up analysis. Subsequently, using the decreasing rearrangement argument and
Young’s inequality, Adimurthi-Yang [4] derived an interpolation of Trudinger-
Moser inequality and Hardy inequality in $\mathbb{R}^{n}$, which can be viewed
as a singular Trudinger-Moser inequality. Another kind of singular Trudinger-
Moser inequality was recently established by Wang-Ye [39] through the method
of blow-up analysis.
Substantial progresses on Adams inequality in $\mathbb{R}^{n}$ was also made
recently. Following lines of Adams, Kozono et al. [19] obtained subcritical
Adams inequality in the whole euclidian space $\mathbb{R}^{n}$. Based on
rearrangement argument of Trombetti-Vazquez [37], Ruf-Sani [36] proved the
critical Adams inequalities in $\mathbb{R}^{n}$ as follows. Let $m$ be an even
integer less than $n$. Assume that $u\in W_{0}^{m,n/m}(\mathbb{R}^{n})$ and
$\|(-\Delta+I)^{m/2}u\|_{L^{n/m}(\mathbb{R}^{n})}\leq 1$. There exists a
constant $C>0$ depending only on $n$ and $m$ such that
$\int_{\mathbb{R}^{n}}\left(e^{\beta_{0}|u|^{\frac{n}{n-m}}}-\sum_{k=0}^{j-2}\frac{\beta_{0}^{k}|u|^{\frac{nk}{n-m}}}{k!}\right)dx<C,$
where $j$ is the smallest integer great than or equal to $n/m$.
Another extension is to establish Trudinger-Moser inequality and Adams
inequality on compact Riemannian manifolds. Let $(M,g)$ be a compact
Riemannian $n$-manifold. For $u\in W^{1,n}(M)$, it was shown by Aubin [5] that
$\exp(\alpha|u|^{n/(n-1)}\|u\|_{W^{1,n}(M)}^{-n/(n-1)})$ is integrable for
sufficiently small $\alpha>0$ which does not depend on $u$. In fact, this is
an easy consequence of Trudinger-Moser inequality and finite partition of
unity on $M$. Let $\tilde{\alpha}$ be the supremum of the above $\alpha$’s. It
was first found by Cherrier [9] that $\tilde{\alpha}=\alpha_{n}$. Cherrier
[10] obtained similar results for $u\in W_{m,n/m}(M)$. Following the lines of
Adams, Fontana [15] obtained critical Adams inequality on $(M,g)$. In 1997,
using the method of blow-up analysis, Ding et al. [11] established a nice
Trudinger-Moser inequality on compact Riemannian surface and successfully
applied it to deal with the prescribed Gaussian curvature problem. Adapting
the argument of Ding et al., Li [21, 22] and Li-Liu [23] proved the existence
of extremal functions for Trudinger-Moser inequalities. Their idea was also
employed by the author [40, 41, 42] to find extremal functions for various
Trudinger-Moser type inequalities. For vector bundles over a compact
Riemannian 2-manifold, Li-Liu-Yang obtained Trudinger-Moser inequalities in
[24].
Among other contributions, we mention the following results. Using the method
of blow-up analysis, Adimurthi-Druet [3] proved that when
$0\leq\alpha<\lambda_{1}(\Omega)$, there holds
$\sup_{u\in W_{0}^{1,2}(\Omega),\,\|\nabla u\|_{2}\leq 1}\int_{\Omega}e^{4\pi
u^{2}(1+\alpha\|u\|_{2}^{2})}dx<\infty,$
where $\lambda_{1}(\Omega)$ is the first eigenvalue of Laplacian on bounded
smooth domain $\Omega\subset\mathbb{R}^{2}$. Moreover, the supremum is
infinite when $\alpha\geq\lambda_{1}(\Omega)$. Later this result was
generalized by the author [43] and Lu-Yang [27, 28, 29].
Although there are fruitful results on euclidian space and compact Riemannian
manifolds, we know little about Trudinger-Moser inequalities on complete
noncompact Riemannian manifolds. In this paper, we concern this problem. Let
$(M,g)$ be any complete noncompact Riemannian $n$-manifold. Throughout this
paper, all the manifolds are assumed to be without boundary, and of dimension
$n\geq 2$. We say that Trudinger-Moser inequality holds on $(M,g)$ if there
exist positive constants $\alpha$, $\tau$ and $\beta$ such that
$\sup_{u\in W^{1,n}(M),\,\|u\|_{1,\tau}\leq
1}\int_{M}\left(e^{\alpha|u|^{\frac{n}{n-1}}}-\sum_{k=0}^{n-2}\frac{\alpha^{k}|u|^{\frac{nk}{n-1}}}{k!}\right)dv_{g}\leq\beta,$
(1.7)
where
$\|u\|_{1,\tau}=\left(\int_{M}|\nabla_{g}u|^{n}dv_{g}\right)^{1/n}+\tau\left(\int_{M}|u|^{n}dv_{g}\right)^{1/n}.$
(1.8)
If the above supremum is infinite for all $\alpha>0$ and $\tau>0$, then we say
that Trudinger-Moser inequality is not valid on $(M,g)$. Motivated by Sobolev
embedding (Hebey [18], Chapter 3), in this paper, we propose and answer the
following three questions.
$(Q_{1})$ Which kind of complete noncompact Riemannian manifolds can possibly
make Trudinger-Moser inequalities hold?
$(Q_{2})$ What geometric assumptions should we consider in order to obtain
Trudinger-Moser inequalities on complete noncompact Riemannian manifolds?
$(Q_{3})$ Are those geometric assumptions in $(Q_{2})$ necessary?
This paper is organized as follows: In Section 2, we state our main results.
From section 3 to section 5, we answer the questions $(Q_{1})$-$(Q_{3})$,
respectively. Adams inequalities are considered in section 6. Finally,
Trudinger-Moser inequalities are applied to nonlinear analysis in section 7.
## 2 Main results
In this section, we answer questions $(Q_{1})$-$(Q_{3})$, and give an
application of Trudinger-Moser inequality. Throughout this paper, we denote
for simplicity a function
$\zeta:\mathbb{N}\times[0,\infty)\rightarrow\mathbb{R}$ by
$\zeta(l,t)=e^{t}-\sum_{k=0}^{l-2}\frac{t^{k}}{k!},\quad\forall l\geq 2.$
(2.1)
From ([44], lemma 2.1 and lemma 2.2), we know that
$\left(\zeta(l,t)\right)^{q}\leq\zeta(l,qt)$ (2.2)
and
$\zeta(l,t)\leq\frac{1}{\mu}\zeta(l,\mu t)+\frac{1}{\nu}\zeta(l,\nu t).$ (2.3)
for all $l\geq 2$, $q\geq 1$, $t\in[0,\infty)$, and $\mu>0$, $\nu>0$
satisfying $1/\mu+1/\nu=1$.
The following proposition answers question $(Q_{1})$.
Proposition 2.1. Let $(M,g)$ be a complete Riemannian $n$-manifold. Suppose
that Trudinger-Moser inequality holds on $(M,g)$, i.e. there exist positive
constants $\alpha$, $\tau$ and $\beta$ such that (1.7) holds. Then the Sobolev
space $W^{1,n}(M)$ is embedded in $L^{q}(M)$ continuously for any $q\geq n$.
Furthermore, for any $r>0$ there exists a positive constant $\epsilon$
depending only on $n$, $\alpha$, $\tau$, $\beta$ and $r$ such that ${\rm
Vol}_{g}(B_{x}(r))\geq\epsilon$ for all $x\in M$, where $B_{x}(r)$ denotes the
geodesic ball centered at $x$ with radius $r$.
From proposition 2.1 we know that there are indeed complete noncompact
Riemannian manifolds such that Trudinger-Moser inequalities are not valid,
namely
Corollary 2.2. For any integer $n\geq 2$, there exists a complete noncompact
Riemannian $n$-manifold on which Trudinger-Moser inequality is not valid.
To answer question $(Q_{2})$, we have the following:
Theorem 2.3. Let $(M,g)$ be a complete noncompact Riemannian $n$-manifold.
Suppose that its Ricci curvature has lower bound, namely ${\rm Rc}_{(M,g)}\geq
Kg$ for some constant $K\in\mathbb{R}$, and its injectivity radius is strictly
positive, namely ${\rm inj}_{(M,g)}\geq i_{0}$ for some constant $i_{0}>0$.
Then we have
$(i)$ for any $0\leq\alpha<\alpha_{n}=n\omega_{n-1}^{1/(n-1)}$, there exists
positive constants $\tau$ and $\beta$ depending only on $n$, $\alpha$, $K$ and
$i_{0}$ such that (1.7) holds. As a consequence, $W^{1,n}(M)$ is embedded in
$L^{q}(M)$ continuously for any $q\geq n$;
$(ii)$ for any $\alpha>\alpha_{n}$ and any $\tau>0$, the supremum in (1.7) is
infinite;
$(iii)$ for any $\alpha>0$ and any $u\in W^{1,n}(M)$, there holds
$\zeta(n,\alpha|u|^{n/(n-1)})\in L^{1}(M)$.
Now we turn to question $(Q_{3})$. The following proposition implies that one
of the hypotheses of theorem 2.3, the injectivity radius is strictly positive,
can not be removed.
Proposition 2.4. For any integer $n\geq 2$, there exists a complete noncompact
Riemannian $n$-manifold, whose Ricci curvature has lower bound, such that
Trudinger-Moser inequality is not valid on it.
We shall construct complete noncompact Riemannian manifolds on which
Trudinger-Moser inequalities hold, but their Ricci curvatures are unbounded
from below. This implies that the other hypothesis of theorem 2.3, Ricci
curvature has lower bound, is not necessary. Namely
Proposition 2.5. For any integer $n\geq 2$, there exists a complete noncompact
Riemannian $n$-manifold on which Trudinger-Moser inequality holds, but its
Ricci curvature is unbounded from below.
Let us explain the idea of proving proposition 2.1 and theorem 2.3. The first
part of conclusions of proposition 2.1, $W^{1,n}(M)\hookrightarrow L^{q}(M)$
for all $q\geq n$, is based on an observation
$\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}=\sum_{k=n-1}^{\infty}\frac{\alpha^{k}}{k!}\int_{M}|u|^{\frac{nk}{n-1}}dv_{g}.$
To find some $\epsilon>0$ such that ${\rm Vol}_{g}(B_{x}(r))\geq\epsilon$ for
all $x\in M$, we employ the method of Carron ([18], lemma 3.2) who obtained
similar result for Sobolev embedding. For the proof of theorem 2.3, we first
derive a uniform local Trudinger-Moser inequality (lemma 4.2 below). Then
using harmonic coordinates and Gromov’s covering lemma, we get the desired
global Trudinger-Moser inequality. The proofs of corollary 2.2, proposition
2.4 and proposition 2.5 are all based on construction of Riemannian manifolds.
Concerning Adams inequalities on complete noncompact Riemannian manifolds, we
have the following:
Theorem 2.6. Let $(M,g)$ be a complete noncompact Riemannian $n$-manifold.
Suppose that there exist positive constants $C(k)$ and $i_{0}$ such that
$|\nabla_{g}^{k}{\rm Rc}_{(M,g)}|\leq C(k)$, $k=0,1,\cdots,m-1$, ${\rm
inj}_{(M,g)}\geq i_{0}>0$. Let $j=n/m$ when $n/m$ is an integer, and
$j=[n/m]+1$ when $n/m$ is not an integer, where $[n/m]$ denotes the integer
part of $n/m$. Then we conclude the following:
$(i)$ there exist positive constants $\alpha_{0}$ and $\beta$ depending only
on $n$, $m$, $C(k)$, $k=1,\cdots,m-1$, and $i_{0}$ such that
$\sup_{\|u\|_{W^{m,n/m}(M)}\leq
1}\int_{M}\zeta\left(j,\alpha_{0}|u|^{\frac{n}{n-m}}\right)dv_{g}\leq\beta.$
As a consequence, $W^{m,{n}/{m}}(M)$ is embedded in $L^{q}(M)$ continuously
for any $q\geq{n}/{m}$;
$(ii)$ for any $\alpha>0$ and any $u\in W^{m,n/m}(M)$, there holds
$\zeta(j,\alpha|u|^{n/(n-m)})\in L^{1}(M)$.
The proof of theorem 2.6 is similar to that of theorem 2.3. It should be
remarked that the existing proofs of Trudinger-Moser inequalities or Adams
inequalities for the euclidian space $\mathbb{R}^{n}$ are all based on
rearrangement argument, which is difficult to be applied to complete
noncompact Riemannian manifold case. Our method is from uniform local
estimates to global estimates. It does not depend on the rearrangement theory
directly.
Trudinger-Moser inequality plays an important role in nonlinear analysis. Let
$(M,g)$ be a complete noncompact Riemannian $n$-manifold. $\nabla_{g}$ denotes
its covariant derivative, and ${\rm div}_{g}$ denotes its divergence operator.
Assume the Ricci curvature of $(M,g)$ has lower bound and the injectivity
radius is strictly positive. We consider the existence results for the
following quasilinear equation.
$-{\rm div}_{g}(|\nabla_{g}u|^{n-2}\nabla_{g}u)+v(x)|u|^{n-2}u=\phi(x)f(x,u),$
(2.4)
where $v(x)$, $\phi(x)$ and $f(x,t)$ are all continuous functions, and
$f(x,t)$ behaves like $e^{\alpha t^{n/(n-1)}}$ as $t\rightarrow+\infty$. In
the case that $(M,g)$ is the standard euclidean space $\mathbb{R}^{n}$ and
$\phi(x)=|x|^{-\beta}$ $(0\leq\beta<n)$, problem (2.4) has been studied by do
Ó et. al. [13, 14], Adimurthi-Yang [4], Yang [44], Lam-Lu [20] and Zhao [45].
Let $O$ be a fixed point of $M$ and $d_{g}(\cdot,\cdot)$ be the geodesic
distance between two points of $(M,g)$. Assume that $\phi(x)$ satisfies the
following hypotheses.
$(\phi_{1})$ $\phi(x)\in L^{p}_{\rm loc}(M)$ for some $p>1$, i. e., for any
$R>0$ there holds $\phi(x)\in L^{p}(B_{O}(R))$;
$(\phi_{2})$ $\phi(x)>0$ for all $x\in M$ and there exist positive constants
$C_{0}$ and $R_{0}$ such that $\phi(x)\leq C_{0}$ for all $x\in M\setminus
B_{O}(R_{0})$.
The potential $v(x)$ is assumed to satisfy the following:
$(v_{1})$ there exists some constant $v_{0}>0$ such that $v(x)\geq v_{0}$ for
all $x\in M$;
$(v_{2})$ either $v(x)\in L^{1/{(n-1)}}(M)$ or $v(x)\rightarrow+\infty$ as
$d_{g}(O,x)\rightarrow+\infty$.
The nonlinearity $f(x,t)$ satisfies the following hypotheses.
$(f_{1})$ there exist constants $\alpha_{0}$, $b_{1}$, $b_{2}>0$ such that for
all $(x,t)\in M\times\mathbb{R}^{+}$,
$|f(x,t)|\leq b_{1}t^{n-1}+b_{2}\zeta\left(n,\alpha_{0}t^{n/(n-1)}\right);$
$(f_{2})$ there exists some constant $\mu>n$ such that for all $x\in M$ and
$t>0$,
$0<\mu F(x,t)\equiv\mu\int_{0}^{t}f(x,s)ds\leq tf(x,t);$
$(f_{3})$ there exist constants $R_{1}$, $A_{1}>0$ such that if $t\geq R_{1}$,
then for all $x\in M$ there holds
$F(x,t)\leq A_{1}f(x,t).$
Define a function space
$E=\left\\{u\in W^{1,n}(M):\int_{M}v(x)|u|^{n}dv_{g}<\infty\right\\}.$ (2.5)
We say that $u\in E$ is a weak solution of problem (2.4) if for all
$\varphi\in E$ we have
$\int_{M}\left(|\nabla_{g}u|^{n-2}\nabla_{g}u\nabla_{g}\varphi+v(x)|u|^{n-2}u\varphi\right)dv_{g}=\int_{M}\phi(x){f(x,u)}\varphi
dv_{g}.$
Define a weighted eigenvalue for the $n$-Laplace operator by
$\lambda_{\phi}=\inf_{u\in E,\,u\not\equiv
0}\frac{\int_{M}(|\nabla_{g}u|^{n}+v(x)|u|^{n})dv_{g}}{\int_{M}\phi(x)|u|^{n}dv_{g}}.$
(2.6)
Then we state the following:
Theorem 2.7. Let $(M,g)$ be a complete noncompact Riemannian $n$-manifold.
Suppose that ${\rm Rc}_{(M,g)}\geq Kg$ for some constant $K\in\mathbb{R}$, and
${\rm inj}_{(M,g)}\geq i_{0}$ for some positive constant $i_{0}$. Assume that
$v(x)$ is a continuous function satisfying $(v_{1})$ and $(v_{2})$, $\phi(x)$
is a continuous function satisfying $(\phi_{1})$ and $(\phi_{2})$,
$f:M\times\mathbb{R}\rightarrow\mathbb{R}$ is a continuous function and the
hypotheses $(f_{1})$, $(f_{2})$ and $(f_{3})$ are satisfied. Furthermore we
assume
$(f_{4})$ $\limsup_{t\rightarrow 0+}{nF(x,t)}/{t^{n}}<\lambda_{\phi}$
uniformly in $x\in M$;
$(f_{5})$ there exist constants $q>n$ and $C_{q}$ such that for all $(x,t)\in
M\times[0,\infty)$
$f(x,t)\geq C_{q}t^{q-1},$
where
$C_{q}>\left(\frac{q-n}{q}\right)^{{(q-n)}/{n}}\left(\frac{p\alpha_{0}}{(p-1)\alpha_{n}}\right)^{(q-n)(n-1)/n}S_{q}^{q}$
and
$S_{q}=\inf_{u\in
E\setminus\\{0\\}}\frac{\left(\int_{M}(|\nabla_{g}u|^{n}+v(x)|u|^{n})dv_{g}\right)^{1/n}}{\left(\int_{M}\phi(x)|u|^{q}dv_{g}\right)^{1/q}}.$
(2.7)
Then the problem (2.4) has a nontrivial nonnegative weak solution.
Remark 2.8. We shall prove that $S_{q}$ can be attained (lemma 7.2 below).
When $(M,g)$ is the standard euclidian space $\mathbb{R}^{n}$,
$\phi(x)=|x|^{-\beta}$ for $0\leq\beta<n$, $(f_{1})$-$(f_{4})$ and
$None$
uniformly in $x$, where $\mathcal{M}$ is some sufficiently large number, we
obtained similar existence result in [44]. The following proposition implies
that the set of functions satisfying $(f_{1})$-$(f_{5})$ is not empty and
assumptions $(f_{1})$-$(f_{5})$ do not imply $(H_{5})$.
Proposition 2.9. There exist continuous functions
$f:M\times\mathbb{R}\rightarrow\mathbb{R}$ such that $(f_{1})$-$(f_{5})$ are
satisfied, but $(H_{5})$ is not satisfied.
We also consider multiplicity results for a perturbation of the problem (2.4),
namely
$-{\rm
div}_{g}(|\nabla_{g}u|^{n-2}\nabla_{g}u)+v(x)|u|^{n-2}u=\phi(x)f(x,u)+\epsilon
h(x),$ (2.8)
where $h(x)\in E^{*}$, the dual space of $E$. If $h\not\equiv 0$ and
$\epsilon>0$ is sufficiently small, under some assumptions there exist at
least two distinct weak solutions to (2.8). Precisely, we have the following
theorem.
Theorem 2.10. Let $(M,g)$ be a complete noncompact Riemannian $n$-manifold.
Suppose that ${\rm Rc}_{(M,g)}\geq Kg$ for some constant $K\in\mathbb{R}$, and
${\rm inj}_{(M,g)}\geq i_{0}$ for some positive constant $i_{0}$. Assume
$f(x,t)$ is continuous in $M\times\mathbb{R}$ and $(f_{1})$-$(f_{5})$ are
satisfied. Both $v(x)$ and $\phi(x)$ are continuous in $M$ and $(v_{1})$,
$(v_{2})$, $(\phi_{1})$, $(\phi_{2})$ are satisfied, $h$ belongs to $E^{*}$,
the dual space of $E$, with $h\geq 0$ and $h\not\equiv 0$. Then there exists
$\epsilon_{0}>0$ such that if $0<\epsilon<\epsilon_{0}$, then the problem
(2.8) has at least two distinct nonnegative weak solutions.
The proofs of theorem 2.7 and theorem 2.10 are based on theorem 2.3, Mountain-
pass theorem and Ekeland’s variational principle. Though similar idea was used
in the case $(M,g)$ is the standard euclidian space $\mathbb{R}^{n}$ [4, 13,
14, 20, 44], technical difficulties caused by manifold structure must be
smoothed.
## 3 Necessary conditions
In this section, we consider the necessary conditions under which Trudinger-
Moser inequality holds. Precisely we shall prove proposition 2.1 and corollary
2.2. Firstly we have the following:
Lemma 3.1. Let $(M,g)$ be a complete Riemannian $n$-manifold. Suppose that
there exist constants $q>n$, $A>0$ and $\tau>0$ such that for all $u\in
W^{1,n}(M)$, there holds
$\left(\int_{M}|u|^{q}dv_{g}\right)^{1/q}\leq A\|u\|_{1,\tau},$ (3.1)
where $\|u\|_{1,\tau}$ is defined by (1.8). Then for any $r>0$ there exists
some positive constant $\epsilon$ depending only on $A$, $n$, $q$, $\tau$, and
$r$ such that for all $x\in M$, ${\rm Vol}_{g}(B_{x}(r))\geq\epsilon$.
Proof. Let $r>0$, $x\in M$, and $\phi\in W^{1,n}(M)$ be such that $\phi=0$ on
$M\setminus B_{x}(r)$. By Hölder’s inequality,
$\left(\int_{M}|\phi|^{n}dv_{g}\right)^{1/n}\leq{\rm
Vol}_{g}(B_{x}(r))^{\frac{1}{n}-\frac{1}{q}}\left(\int_{M}|\phi|^{q}dv_{g}\right)^{1/q}.$
This together with (3.1) gives
$\left(1-\tau A{\rm
Vol}_{g}(B_{x}(r))^{\frac{1}{n}-\frac{1}{q}}\right)\left(\int_{M}|\phi|^{q}dv_{g}\right)^{1/q}\leq
A\left(\int_{M}(|\nabla\phi|^{n}dv_{g}\right)^{1/n}.$ (3.2)
Fix $x\in M$ and $R>0$. Then either
${\rm Vol}_{g}(B_{x}(R))>\left(\frac{1}{2\tau A}\right)^{{nq}/{(q-n)}}$ (3.3)
or
${\rm Vol}_{g}(B_{x}(R))\leq\left(\frac{1}{2\tau A}\right)^{{nq}/{(q-n)}}.$
(3.4)
If (3.4) holds, then we have
$1-\tau A{\rm Vol}_{g}(B_{x}(R))^{\frac{1}{n}-\frac{1}{q}}\geq{1}/{2},$
and whence for all $r\in(0,R]$ and all $\phi\in W^{1,n}(M)$ with $\phi=0$ on
$M\setminus B_{x}(r)$,
$\left(\int_{M}|\phi|^{q}dv_{g}\right)^{1/q}\leq
2A\left(\int_{M}(|\nabla\phi|^{n}dv_{g}\right)^{1/n}.$ (3.5)
Now we set
$\phi(y)=\left\\{\begin{array}[]{lll}r-d_{g}(x,y)&{\rm when}\quad
d_{g}(x,y)\leq r\\\\[6.45831pt] 0&{\rm when}\quad
d_{g}(x,y)>r.\end{array}\right.$
Clearly $\phi\in W^{1,n}(M)$, $\phi=0$ on $M\setminus B_{x}(r)$, $\phi\geq
r/2$ on $B_{x}(r/2)$, and $|\nabla\phi|=1$ almost everywhere in $B_{x}(r)$. It
then follows from (3.5) that
$\frac{r}{2}{\rm Vol}_{g}(B_{x}(r/2))^{1/q}\leq 2A{\rm
Vol}_{g}(B_{x}(r))^{1/n}.$
Hence we have for all $r\leq R$,
${\rm Vol}_{g}(B_{x}(r))\geq\left(\frac{r}{4A}\right)^{n}{\rm
Vol}_{g}(B_{x}(r/2))^{n/q}.$
By induction we obtain for any positive integer $m$,
${\rm
Vol}_{g}(B_{x}(R))\geq\left(\frac{R}{2A}\right)^{n\alpha(m)}\left(\frac{1}{2}\right)^{n\beta(m)}{\rm
Vol}_{g}(B_{x}(R/2^{m}))^{(n/q)^{m}},$ (3.6)
where
$\alpha(m)=\sum_{j=1}^{m}({n}/{q})^{j-1},\quad\beta(m)=\sum_{j=1}^{m}j(n/q)^{j-1}.$
On one hand we know from ([7], Theorem 3.98) that ${\rm
Vol}_{g}(B_{x}(r))=\frac{\omega_{n-1}}{n}r^{n}(1+o(r))$, where $\omega_{n-1}$
is the area of the euclidean unit sphere in $\mathbb{R}^{n}$, and
$o(r)\rightarrow 0$ as $r\rightarrow 0$. One can see without any difficulty
that
$\lim_{m\rightarrow\infty}{\rm Vol}_{g}(B_{x}(R/2^{m}))^{(n/q)^{m}}=1.$
On the other hand we have
$\sum_{j=1}^{\infty}({n}/{q})^{j-1}=\frac{q}{q-n},\quad\sum_{j=1}^{\infty}j(n/q)^{j-1}=\frac{q^{2}}{(q-n)^{2}}.$
Hence, passing to the limit $m\rightarrow\infty$ in (3.6), one concludes that
${\rm
Vol}_{g}(B_{x}(R))\geq\left(\frac{R}{2^{(2q-n)/(q-n)}A}\right)^{nq/(q-n)}.$
This together with (3.3), (3.4) implies that
${\rm Vol}_{g}(B_{x}(R))\geq\min\left\\{\frac{1}{2\tau
A},\frac{R}{2^{(2q-n)/(q-n)}A}\right\\}^{nq/(q-n)}$
and completes the proof of the lemma. $\hfill\Box$
It should be pointed out that the above argument is a modification of that of
Carron ([18], lemma 3.2). Note that the condition (3.1) implies that
$W^{1,n}(M)$ is continuously embedded in $L^{q}(M)$ for some $q>n$. This is
different from the assumption of ([18], lemma 3.2).
To prove proposition 2.1, we also need the following interpolation inequality.
Lemma 3.2. Let $\tau$ be any positive real number. Suppose there exist
positive constants $q_{1}$, $q_{2}$, $A_{1}$ and $A_{2}$ such that
$q_{2}>q_{1}>0$ and
$\left(\int_{M}|u|^{q_{i}}dv_{g}\right)^{1/{q_{i}}}\leq A_{i}\|u\|_{1,\tau}$
(3.7)
for all $u\in W^{1,n}(M)$, $i=1,2$. Then for all $q:q_{1}<q<q_{2}$ there
exists a positive constant $A=A(A_{1},A_{2},q_{1},q_{2})$ such that
$\left(\int_{M}|u|^{q}dv_{g}\right)^{1/{q}}\leq A\|u\|_{1,\tau}$ (3.8)
for all $u\in W^{1,n}(M)$.
Proof. For any $u\in W^{1,n}(M)\setminus\\{0\\}$, we set
$\widetilde{u}=u/\|u\|_{1,\tau}$. It follows from (3.7) that
$\left(\int_{M}|\widetilde{u}|^{q_{i}}dv_{g}\right)^{1/{q_{i}}}\leq
A_{i},\,\,\,i=1,2.$
Assume $q_{1}<q<q_{2}$. Since
$|\widetilde{u}|^{q}\leq|\widetilde{u}|^{q_{1}}+|\widetilde{u}|^{q_{2}}$,
there holds
$\int_{M}|\widetilde{u}|^{q}dv_{g}\leq\int_{M}|\widetilde{u}|^{q_{1}}dv_{g}+\int_{M}|\widetilde{u}|^{q_{2}}dv_{g}\leq
A_{1}^{q_{1}}+A_{2}^{q_{2}}.$
Hence
$\left(\int_{M}|{u}|^{q}dv_{g}\right)^{1/q}\leq(A_{1}^{q_{1}}+A_{2}^{q_{2}})^{\frac{1}{q}}\|u\|_{1,\tau}.$
Take
$A=\max\\{(A_{1}^{q_{1}}+A_{2}^{q_{2}})^{1/q_{1}},(A_{1}^{q_{1}}+A_{2}^{q_{2}})^{1/q_{2}}\\}$.
Then (3.8) follows immediately. $\hfill\Box$
Proof of proposition 2.1. Assume there exist positive constants $\alpha$,
$\tau$ and $\beta$ such that (1.7) holds. For any $u\in W^{1,n}(M)$ we set
$\widetilde{u}=u/\|u\|_{1,\tau}$. It follows from (1.7) that
$\int_{M}\sum_{k=n-1}^{\infty}\frac{\alpha^{k}|\widetilde{u}|^{\frac{nk}{n-1}}}{k!}dv_{g}\leq\beta.$
Particularly for any integer $k\geq n-1$ there holds
$\int_{M}\frac{\alpha^{k}|\widetilde{u}|^{\frac{nk}{n-1}}}{k!}dv_{g}\leq\beta,$
and thus
$\left(\int_{M}|u|^{\frac{nk}{n-1}}dv_{g}\right)^{\frac{n-1}{nk}}\leq\left(\frac{k!\beta}{\alpha^{k}}\right)^{\frac{n-1}{nk}}\|u\|_{1,\tau}.$
For any $q\geq n$, there exists some $k\geq n-1$ such that
$\frac{nk}{n-1}\leq q<\frac{n(k+1)}{n-1}.$
In fact we can choose $k=[(n-1)p/n]$, the integer part of $(n-1)p/n$. By lemma
3.2, there exists a positive constant $A$ depending only on $n$, $q$,
$\alpha$, and $\beta$ such that
$\left(\int_{M}|{u}|^{q}dv_{g}\right)^{1/q}\leq A\|u\|_{1,\tau}.$
This implies that $W^{1,n}(M)\hookrightarrow L^{q}(M)$ continuously. Now we
fix some $q>n$, say $q=n+1$. Then by lemma 3.1, there exists some constant
$\epsilon>0$ depending only on $n$, $\alpha$, $\tau$, $\beta$ and $r$ such
that for all $x\in M$, ${\rm Vol}_{g}(B_{x}(r))\geq\epsilon$. $\hfill\Box$
Proof of corollary 2.2. For any complete noncompact Riemannian $n$-manifold
$(M,g)$, if Trudinger-Moser inequality holds, then by proposition 2.1, there
exists some constant $\epsilon>0$ such that ${\rm
Vol}_{g}(B_{x}(r))\geq\epsilon$ for all $x\in M$. Hence if there exists some
complete noncompact Riemannian $n$-manifold $(M,g)$ such that
$\inf_{x\in M}{\rm Vol}_{g}(B_{x}(r))=0,$
then we conclude that Trudinger-Moser inequality is not valid on it. Now we
construct such complete Riemannian manifolds. Consider the warped product
$M=\mathbb{R}\times N,\,\,\,g(t,\theta)=dt^{2}+f(t)ds_{N}^{2},$
where $(N,ds_{N}^{2})$ is a compact $(n-1)$-Riemannian manifold, $dt^{2}$ is
the euclidian metric of $\mathbb{R}$, and $f$ is a smooth function satisfying
$f(t)>0,\forall t\in\mathbb{R}$ and $\lim_{t\rightarrow+\infty}f(t)=0$. If
$y=(t_{1},m_{1})$ and $z=(t_{2},m_{2})$ are two points of ${M}$, then
$d_{g}(y,z)\geq|t_{2}-t_{1}|$. This together with the compactness of $N$
implies that $({M},g)$ is complete. In addition, for any $x=(t,m)\in M$, there
holds
$B_{x}(1)\subset(t-1,t+1)\times N.$
Therefore
$\displaystyle{\rm Vol}_{g}(B_{x}(1))$ $\displaystyle\leq$ $\displaystyle{\rm
Vol}_{g}\left((t-1,t+1)\times N\right){}$ (3.9) $\displaystyle\leq$
$\displaystyle{\rm Vol}_{ds_{N}^{2}}(N)\int_{t-1}^{t+1}f(t)dt{}$
$\displaystyle=$ $\displaystyle 2{\rm Vol}_{ds_{N}^{2}}(N)f(\xi)$
$\displaystyle\rightarrow$ $\displaystyle 0\,\,\,{\rm
as}\,\,\,t\rightarrow+\infty,$
where we used the integral mean value theorem, $\xi$ is some point in
$(t-1,t+1)$. This gives the desired result. $\hfill\Box$
## 4 Sufficient conditions
In this section, we investigate sufficient conditions under which Trudinger-
Moser inequality holds. Precisely we shall prove theorem 2.3 and proposition
2.4. Firstly we have the following key observation:
Lemma 4.1. Let $\mathbb{B}_{0}(\delta)\subset\mathbb{R}^{n}$ be a ball
centered at $0$ with radius $\delta$. If
$0\leq\alpha\leq\alpha_{n}=n\omega_{n-1}^{1/(n-1)}$, then there exists some
constant $C$ depending only on $n$ such that for all $u\in
W_{0}^{1,n}(\mathbb{B}_{0}(\delta))$ satisfying
$\int_{\mathbb{B}_{0}(\delta)}|\nabla u|^{n}dx\leq 1$, there holds
$\int_{\mathbb{B}_{0}(\delta)}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dx\leq
C\delta^{n}\left(\frac{\alpha}{\alpha_{n}}\right)^{n-1}\int_{\mathbb{B}_{0}(\delta)}|\nabla
u|^{n}dx.$ (4.1)
Proof. Let $\widetilde{u}=u/\|\nabla u\|_{L^{n}(\mathbb{B}_{0}(\delta))}$.
Since $\|\nabla u\|_{L^{n}(\mathbb{B}_{0}(\delta))}\leq 1$ and
$0\leq\alpha\leq\alpha_{n}$, we have
$\displaystyle\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)$ $\displaystyle=$
$\displaystyle\sum_{k=n-1}^{\infty}\frac{\alpha^{k}|u|^{\frac{nk}{n-1}}}{k!}{}$
(4.2) $\displaystyle=$
$\displaystyle\sum_{k=n-1}^{\infty}\left(\frac{\alpha}{\alpha_{n}}\right)^{k}\frac{\alpha_{n}^{k}\|\nabla
u\|_{L^{n}(\mathbb{B}_{0}(\delta))}^{\frac{nk}{n-1}}|\widetilde{u}|^{\frac{nk}{n-1}}}{k!}{}$
$\displaystyle\leq$ $\displaystyle\|\nabla
u\|_{L^{n}(\mathbb{B}_{0}(\delta))}^{n}\left(\frac{\alpha}{\alpha_{n}}\right)^{n-1}\zeta\left(n,\alpha_{n}|\widetilde{u}|^{\frac{n}{n-1}}\right).$
It follows from the classical Trudinger-Moser inequality ((1.2) with
$\alpha_{0}$ replaced by $\alpha_{n}$) that
$\int_{\mathbb{B}_{0}(\delta)}\zeta\left(n,\alpha_{n}|\widetilde{u}|^{\frac{n}{n-1}}\right)dx\leq
C\delta^{n}$ (4.3)
for some constant $C$ depending only on $n$. Integrating (4.2) on
$\mathbb{B}_{0}(\delta)$, we immediately obtain (4.1) by using (4.3). This
concludes the lemma. $\hfill\Box$
Let $(M,g)$ be a complete Riemannian $n$-manifold with ${\rm Ric}_{(M,g)}\geq
Kg$ for some $K\in\mathbb{R}$ and ${\rm inj}_{(M,g)}\geq i_{0}$ for some
$i_{0}>0$. Then we have the following local version of Trudinger-moser
inequality which is the key estimate for the proof of theorem 2.3:
Lemma 4.2. For any $\alpha:0<\alpha<\alpha_{n}$ there exists some constant
$\delta$ depending only on $n$, $\alpha$, $K$ and $i_{0}$ such that for all
$x\in M$ and all $u\in C_{0}^{\infty}(B_{x}(\delta))$ with
$\|\nabla_{g}u\|_{L^{n}(B_{x}(\delta))}\leq 1$, there holds
$\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}\leq
C\int_{M}|\nabla_{g}u|^{n}dv_{g}$
for some constant $C$ depending only on $n$, $\alpha$, $K$ and $i_{0}$.
Proof. By (Hebey [18], theorem 1.3), we know that for any $\epsilon>0$ there
exists a positive constant $\delta$ depending only on $\epsilon$, $n$, $K$ and
$i_{0}$ satisfying the following property: for any $x\in M$ there exists a
harmonic coordinate chart $\phi:B_{x}(\delta)\rightarrow\mathbb{R}^{n}$ such
that $\phi(x)=0$, and the components $(g_{jl})$ of $g$ in this chart satisfy
$e^{-\epsilon}\delta_{jl}\leq g_{jl}\leq e^{\epsilon}\delta_{jl}$
as bilinear forms. One then has that
$\phi(B_{x}(\delta))\subset\mathbb{B}_{0}(e^{\epsilon/2}\delta)$. Let $u$ be a
function in $C_{0}^{\infty}(B_{x}(\delta))$ and
$\|\nabla_{g}u\|_{L^{n}(B_{x}(\delta))}\leq 1$. It is not difficult to see
that
$\displaystyle\int_{B_{x}(\delta)}|\nabla_{g}u|^{n}dv_{g}$ $\displaystyle\geq$
$\displaystyle
e^{-n\epsilon}\int_{\mathbb{B}_{0}(e^{\epsilon/2}\delta)}|\nabla(u\circ\phi^{-1})(x)|^{n}dx,$
(4.4)
$\displaystyle\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}$
$\displaystyle\leq$ $\displaystyle
e^{n\epsilon/2}\int_{\mathbb{B}_{0}(e^{\epsilon/2}\delta)}\zeta\left(n,\alpha|(u\circ\phi^{-1})(x)|^{\frac{n}{n-1}}\right)dx.$
(4.5)
For any fixed $\alpha:0<\alpha<\alpha_{n}$, there exists some $\epsilon_{0}$
depending only on $n$ and $\alpha$ such that when
$0<\epsilon\leq\epsilon_{0}$, it follows from (4.4) and
$\|\nabla_{g}u\|_{L^{n}(B_{x}(\delta))}\leq 1$ that
$\alpha\left(\int_{\mathbb{B}_{0}(e^{\epsilon/2}\delta)}|\nabla(u\circ\phi^{-1})(x)|^{n}dx\right)^{1/(n-1)}\leq\alpha
e^{n\epsilon_{0}/(n-1)}<\alpha_{n}.$
Now let $\epsilon=\epsilon_{0}$ be fixed and $\delta$ depending only on
$\epsilon_{0}$, $n$, $K$ and $i_{0}$ be chosen as above. By lemma 4.1, there
exists a constant $C_{1}=C_{1}(n)$ depending only on $n$ such that
$\int_{\mathbb{B}_{0}(e^{\epsilon_{0}/2}\delta)}\zeta\left(n,\alpha|(u\circ\phi^{-1})(x)|^{\frac{n}{n-1}}\right)dx\leq
C_{1}e^{n\epsilon_{0}/2}\delta^{n}\int_{\mathbb{B}_{0}(e^{\epsilon_{0}/2}\delta)}|\nabla(u\circ\phi^{-1})(x)|^{n}dx.$
This together with (4.4) and (4.5) implies that
$\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}\leq
C_{1}e^{2n\epsilon_{0}}\delta^{n}\int_{M}|\nabla u|^{n}dv_{g}.$
Take $C=C_{1}e^{2n\epsilon_{0}}\delta^{n}$. We conclude that $C$ depends on
$n$, $\alpha$, $K$ and $i_{0}$. $\hfill\Box$
Proof of theorem 2.3. $(i)$ For any $\alpha:0<\alpha<\alpha_{n}$, let
$\delta=\delta(n,\alpha,K,i_{0})$ be chosen as in lemma 4.2. Independently, by
Gromov’s covering lemma (Hebey [18], lemma 1.6), we can select a sequence
$(x_{j})$ of points of $M$ such that
$(a)$ $M=\cup_{j}B_{x_{j}}(\delta/2)$, and for any $j\not=l$ there holds
$B_{x_{j}}(\delta/4)\cap B_{x_{l}}(\delta/4)=\varnothing$;
$(b)$ there exists $N$ depending only on $n$, $K$ and $\delta$ such that each
point of $M$ has a neighborhood which intersects at most $N$ of the
$B_{x_{j}}(\delta)$’s.
For any $j$, we take a cut-off function $\phi_{j}\in
C_{0}^{\infty}(B_{x_{j}}(\delta))$ satisfying $0\leq\phi_{j}\leq 1$,
$\phi_{j}\equiv 1$ on $B_{x_{j}}(\delta/2)$, and $|\nabla_{g}\phi_{j}|\leq
4/\delta$. It follows that for all $j$
$|\nabla_{g}\phi_{j}^{2}|=2\phi_{j}|\nabla_{g}\phi_{j}|\leq\frac{8}{\delta}\phi_{j}.$
(4.6)
By the covering properties $(a)$ and $(b)$, we have
$1\leq\sum_{j}\phi_{j}(x)\leq N\,\,\,{\rm for\,\,\,all}\,\,\,x\in M.$ (4.7)
Set $\tau=8/\delta$. Assume $u\in C_{0}^{\infty}(M)$ satisfies
$\|u\|_{1,\tau}=\left(\int_{M}|\nabla
u|^{n}dv_{g}\right)^{1/n}+\tau\left(\int_{M}|u|^{n}dv_{g}\right)^{1/n}\leq 1.$
It follows from (4.6) and the Minkowvsky inequality that
$\displaystyle\left(\int_{M}|\nabla_{g}(\phi_{j}^{2}u)|^{n}dv_{g}\right)^{1/n}\leq\left(\int_{M}\phi_{j}^{2n}|\nabla_{g}u|^{n}dv_{g}\right)^{1/n}+\left(\int_{M}|\nabla_{g}\phi_{j}^{2}|^{n}|u|^{n}dv_{g}\right)^{1/n}\leq\|u\|_{1,\tau}\leq
1.$
In view of lemma 4.2, this leads to
$\displaystyle\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}$
$\displaystyle\leq$
$\displaystyle\sum_{j}\int_{B_{\delta/2}(x_{j})}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}{}$
(4.8) $\displaystyle\leq$
$\displaystyle\sum_{j}\int_{B_{\delta}(x_{j})}\zeta\left(n,\alpha|\phi_{j}^{2}u|^{\frac{n}{n-1}}\right)dv_{g}$
$\displaystyle\leq$ $\displaystyle
C\sum_{j}\int_{M}|\nabla(\phi_{j}^{2}u)|^{n}dv_{g}$
for some constant $C$ depending only on $n$, $\alpha$, $K$ and $i_{0}$. In
addition we have by using (4.6) and $0\leq\phi_{j}\leq 1$ that
$\displaystyle\int_{M}|\nabla_{g}(\phi_{j}^{2}u)|^{n}dv_{g}$
$\displaystyle\leq$ $\displaystyle
2^{n}\int_{M}\left(\phi_{j}^{2n}|\nabla_{g}u|^{n}+|\nabla_{g}\phi_{j}^{2}|^{n}|u|^{n}\right)dv_{g}{}$
$\displaystyle\leq$ $\displaystyle
2^{n}\int_{M}\phi_{j}|\nabla_{g}u|^{n}dv_{g}+\frac{16^{n}}{\delta^{n}}\int_{M}\phi_{j}|u|^{n}dv_{g}.$
In view of (4.7), it follows that
$\displaystyle\sum_{j}\int_{M}|\nabla_{g}(\phi_{j}^{2}u)|^{n}dv_{g}$
$\displaystyle\leq$ $\displaystyle
2^{n}N\int_{M}|\nabla_{g}u|^{n}dv_{g}+\frac{16^{n}}{\delta^{n}}N\int_{M}|u|^{n}dv_{g}$
$\displaystyle\leq$ $\displaystyle 2^{n}N+\frac{16^{n}}{\tau\delta^{n}}N.$
This together with (4.8) implies
$\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}\leq{C}$
for some constant ${C}$ depending only on $n$, $\alpha$, $K$ and $i_{0}$. By
the density of $C_{0}^{\infty}(M)$ in $W^{1,n}(M)$, the inequality (1.7) holds
for the above $\alpha$, $\tau$ and $C$.
By proposition 2.1, we have that $W^{1,n}(M)$ is continuously embedded in
$L^{q}(M)$ for any $q\geq n$.
$(ii)$ Fix some point $z\in M$, let $r=r(x)=d_{g}(z,x)$ be the geodesic
distance between $x$ and $z$. Without loss of generality, we may assume the
injectivity radius of $(M,g)$ at $z$ is strictly larger than 1. Take a
function sequence
$\phi_{\epsilon}(x)=\left\\{\begin{array}[]{ll}1,&{\rm when}\quad
r<\epsilon\\\\[6.45831pt]
\left(\log\frac{1}{\epsilon}\right)^{-1}\log\frac{1}{r},&{\rm
when}\quad\epsilon\leq r\leq 1\\\\[6.45831pt] 0,&{\rm when}\quad
r>1.\end{array}\right.$
Then $\phi_{\epsilon}\in W^{1,n}(M)$ and for any constant $\tau>0$ there holds
$\displaystyle\|\phi_{\epsilon}\|_{1,\tau}=\left(\log\frac{1}{\epsilon}\right)^{(1-n)/n}\omega_{n-1}^{1/n}\left(1+O\left(\frac{1}{\log{\epsilon}}\right)\right).$
Set
$\widetilde{\phi}_{\epsilon}=\phi_{\epsilon}/\|\phi_{\epsilon}\|_{1,\tau}$.
Then we have on the geodesic ball $B_{z}(\epsilon)\subset M$,
$\zeta(n,\alpha\widetilde{\phi}_{\epsilon}^{\frac{n}{n-1}})=e^{\alpha\widetilde{\phi}_{\epsilon}^{\frac{n}{n-1}}}-\sum_{k=0}^{n-2}\frac{\alpha^{k}\widetilde{\phi}_{\epsilon}^{\frac{nk}{n-1}}}{k!}\geq\epsilon^{\alpha\omega_{n-1}^{-\frac{1}{n-1}}(1+O(1/\log\epsilon))}+O\left(\left(\log\frac{1}{\epsilon}\right)^{n-2}\right).$
Note that $\alpha\omega_{n-1}^{-\frac{1}{n-1}}>n$ for any $\alpha>\alpha_{n}$.
Hence, when $\alpha>\alpha_{n}$, we have
$\displaystyle\int_{M}\zeta(n,\alpha|\widetilde{\phi}_{\epsilon}|^{\frac{n}{n-1}})dv_{g}$
$\displaystyle\geq$
$\displaystyle\int_{B_{z}(\epsilon)}\zeta(n,\alpha|\widetilde{\phi}_{\epsilon}|^{\frac{n}{n-1}})dv_{g}$
$\displaystyle\geq$
$\displaystyle\frac{\omega_{n-1}}{n}(1+o_{\epsilon}(1))\epsilon^{n-\alpha\omega_{n-1}^{-1/(n-1)}(1+O(1/\log\epsilon))}+o_{\epsilon}(1).$
$\displaystyle\rightarrow$ $\displaystyle+\infty\quad{\rm
as}\quad\epsilon\rightarrow 0.$
This ends the proof of $(ii)$.
$(iii)$ Take $\alpha_{0}:0<\alpha_{0}<\alpha_{n}$. By $(i)$ there exists some
$\tau_{0}=\tau_{0}(n,\alpha_{0},K,i_{0})>0$ such that
$\Lambda_{\alpha_{0}}:=\sup_{\|u\|_{1,\tau_{0}}\leq
1}\int_{M}\zeta(n,\alpha_{0}|u|^{\frac{n}{n-1}})dv_{g}<\infty.$
Given any $\alpha>0$ and any $u\in W^{1,n}(M)$. Since $C_{0}^{\infty}(M)$ is
dense in $W^{1,n}(M)$ under the norm $\|\cdot\|_{W^{1,n}(M)}$, which is
equivalent to the norm $\|\cdot\|_{1,\tau_{0}}$, we can choose some $u_{0}\in
C_{0}^{\infty}(M)$ such that
$2^{\frac{n}{n-1}}\alpha\|u-u_{0}\|_{1,\tau_{0}}^{\frac{n}{n-1}}<\alpha_{0}.$
(4.9)
Since $\zeta(n,t)$ is increasing in $t$ for $t\geq 0$, we obtain by using
(2.3)
$\displaystyle\int_{M}\zeta(n,\alpha|u|^{\frac{n}{n-1}})dv_{g}$
$\displaystyle\leq$
$\displaystyle\int_{M}\zeta(n,2^{\frac{n}{n-1}}\alpha|u-u_{0}|^{\frac{n}{n-1}}+2^{\frac{n}{n-1}}\alpha|u_{0}|^{\frac{n}{n-1}})dv_{g}{}$
(4.10) $\displaystyle\leq$
$\displaystyle\frac{1}{\mu}\int_{M}\zeta(n,2^{\frac{n}{n-1}}\alpha\mu|u-u_{0}|^{\frac{n}{n-1}})dv_{g}$
$\displaystyle\quad+\frac{1}{\nu}\int_{M}\zeta(n,2^{\frac{n}{n-1}}\alpha\nu|u_{0}|^{\frac{n}{n-1}})dv_{g},$
where $1/\mu+1/\nu=1$. In view of (4.9), we can take $\mu>1$ sufficiently
close to $1$ such that
$2^{\frac{n}{n-1}}\alpha\mu\|u-u_{0}\|_{1,\tau_{0}}^{\frac{n}{n-1}}<\alpha_{0}.$
Hence
$\int_{M}\zeta(n,2^{\frac{n}{n-1}}\alpha\mu|u-u_{0}|^{\frac{n}{n-1}})dv_{g}\leq\Lambda_{\alpha_{0}}.$
(4.11)
Since $u_{0}\in C_{0}^{\infty}(M)$, particularly $u_{0}$ has compact support,
there holds
$\int_{M}\zeta(n,2^{\frac{n}{n-1}}\alpha\nu|u_{0}|^{\frac{n}{n-1}})dv_{g}<\infty.$
(4.12)
Combining (4.10), (4.11) and (4.12), we obtain
$\int_{M}\zeta(n,\alpha|u|^{\frac{n}{n-1}})dv_{g}<\infty.$
This completes the proof of $(iii)$. $\hfill\Box$
Now we shall prove proposition 2.4. Let us recall some notations from
Riemannian geometry. In any chart, the Christoffel symbols of the Levi-Civita
connection are given by
$\Gamma_{ij}^{k}=\frac{1}{2}g^{mk}\left(\partial_{i}g_{mj}+\partial_{j}g_{mi}-\partial_{m}g_{ij}\right),$
(4.13)
where $g_{ij}$’s are the components of $g$, $(g^{ij})$ denotes the inverse
matrix of $(g_{ij})$. Here and in the sequel the Einstein’s summation
convention is adopted. Denote the Riemannian curvature of $(M,g)$, a
$(4,0)$-type tensor field, by ${\rm Rm}_{(M,g)}$. The components of ${\rm
Rm}_{(M,g)}$ are given by the relation
$R_{ijkl}=g_{i\alpha}\left(\partial_{k}\Gamma_{jl}^{\alpha}-\partial_{l}\Gamma_{jk}^{\alpha}+\Gamma_{k\beta}^{\alpha}\Gamma_{jl}^{\beta}-\Gamma_{l\beta}^{\alpha}\Gamma_{jk}^{\beta}\right).$
(4.14)
Similarly, the components of the Ricci curvature ${\rm Rc}_{(M,g)}$ of $(M,g)$
are given by the relation
$R_{ij}=g^{\alpha\beta}R_{i\alpha j\beta}.$ (4.15)
Proof of proposition 2.4. In view of proposition 2.1, it suffices to construct
a complete noncompact Riemannian $n$-manifold $(M,g)$ such that its Ricci
curvature has lower bound and there holds
$\inf_{x\in M}{\rm Vol}_{g}(B_{x}(1))=0.$
Again we consider the warped product
${M}=\mathbb{R}\times N,\,\,\,g(x,\theta)=dx^{2}+f(x)ds_{N}^{2},$
where $(N,ds_{N}^{2})$ is a compact $(n-1)$-Riemannian manifold, $dx^{2}$ is
the euclidean metric of $\mathbb{R}$, and $f$ is a smooth function satisfying
$f(x)>0,\forall x\in\mathbb{R}$. In the following we calculate the Ricci
curvature of $({M},g)$. In some product chart $(\mathbb{R}\times
U,Id\times\phi)$ ($\\{x,y^{2},\cdots,y^{n}\\}$), $g_{11}=1$, $g_{1\alpha}=0$,
$g_{\alpha\beta}=fh_{\alpha\beta}$, $g^{11}=1$, $g^{1,\alpha}=0$, and
$g^{\alpha\beta}=f^{-1}h^{\alpha\beta}$. Equivalently
$g=dx^{2}+f(x)h_{\alpha\beta}dy^{\alpha}dy^{\beta},$
where $(h_{\alpha\beta})$ denote components of the metric $ds_{N}^{2}$. Here
and in the sequel, all indices $\alpha$, $\beta$, $\mu$, $\nu$ and $\lambda$
run from $2$ to $n$. In view of (4.13), the Christoffel symbols of the Levi-
Civita connection was calculated as follows:
$\displaystyle\Gamma_{11}^{1}=\Gamma_{11}^{\alpha}=\Gamma_{1\alpha}^{1}=0,\,\,\,\Gamma_{1\alpha}^{\beta}=\frac{1}{2}g^{\mu\beta}\partial_{1}g_{\mu\alpha}=\frac{f^{\prime}}{2f}\delta_{\alpha}^{\beta}$
$\displaystyle\Gamma_{\alpha\beta}^{1}=-\frac{1}{2}\partial_{1}g_{\alpha\beta}=-\frac{f^{\prime}}{2}h_{\alpha\beta},\quad\Gamma_{\alpha\beta}^{\gamma}=\widetilde{\Gamma}_{\alpha\beta}^{\gamma},$
where $\delta_{\alpha}^{\beta}$ is equal to $1$ when $\alpha=\beta$, and $0$
when $\alpha\not=\beta$, $\widetilde{\Gamma}_{\alpha\beta}^{\gamma}$’s are
components of the Christoffel symbols of Levi-Civita connection on
$(N,ds_{N}^{2})$. In view of (4.14), the components of the Riemannian
curvature reads
$\displaystyle R_{1\alpha 1\beta}$ $\displaystyle=$ $\displaystyle
g_{11}\partial_{1}\Gamma_{\alpha\beta}^{1}$ $\displaystyle=$
$\displaystyle\frac{{f^{\prime}}^{2}-2ff^{\prime\prime}}{4f}h_{\alpha\beta}$
$\displaystyle R_{1\alpha\beta\gamma}$ $\displaystyle=$ $\displaystyle
g_{11}\left(\partial_{\beta}\Gamma_{\alpha\gamma}^{1}-\partial_{\gamma}\Gamma_{\alpha\beta}^{1}+\Gamma_{\beta
k}^{1}\Gamma_{\alpha\gamma}^{k}-\Gamma_{\gamma
k}^{1}\Gamma_{\alpha\beta}^{k}\right)$ $\displaystyle=$
$\displaystyle\frac{f^{\prime}}{2}\left(-\partial_{\beta}h_{\alpha\gamma}+\partial_{\gamma}h_{\alpha\beta}-h_{\beta\mu}\widetilde{\Gamma}_{\alpha\gamma}^{\mu}+h_{\gamma\mu}\widetilde{\Gamma}_{\alpha\beta}^{\mu}\right)$
$\displaystyle R_{\alpha\beta\gamma\mu}$ $\displaystyle=$ $\displaystyle
g_{\alpha\lambda}\left(\partial_{\gamma}\Gamma_{\beta\mu}^{\lambda}-\partial_{\mu}\Gamma_{\beta\gamma}^{\lambda}+\Gamma_{\gamma
k}^{\lambda}\Gamma_{\beta\mu}^{k}-\Gamma_{\mu
k}^{\lambda}\Gamma_{\beta\gamma}^{k}\right)$ $\displaystyle=$ $\displaystyle
f\widetilde{R}_{\alpha\beta\gamma\mu}+g_{\alpha\lambda}\left(\Gamma_{\gamma
1}^{\lambda}\Gamma_{\beta\mu}^{1}-\Gamma_{\mu
1}^{\lambda}\Gamma_{\beta\gamma}^{1}\right)$ $\displaystyle=$ $\displaystyle
f\widetilde{R}_{\alpha\beta\gamma\mu}+\frac{{f^{\prime}}^{2}}{4}\left(h_{\alpha\mu}h_{\beta\gamma}-h_{\alpha\gamma}h_{\beta\mu}\right),$
where $\widetilde{R}_{\alpha\beta\gamma\mu}$’s denote the components of
Riemannian curvature of $(N,ds_{N}^{2})$. In view of (4.15), we get the
components of the Ricci curvature as follows.
$\displaystyle R_{11}$ $\displaystyle=$ $\displaystyle
g^{\alpha\beta}R_{1\alpha 1\beta}$ $\displaystyle=$
$\displaystyle(n-1)\frac{{f^{\prime}}^{2}-2ff^{\prime\prime}}{4f^{2}}$
$\displaystyle R_{1\alpha}$ $\displaystyle=$ $\displaystyle
g^{\beta\gamma}R_{1\beta\alpha\gamma}$ $\displaystyle=$
$\displaystyle\frac{f^{\prime}}{2f}h^{\beta\gamma}\left(-\partial_{\alpha}h_{\beta\gamma}+\partial_{\gamma}h_{\alpha\beta}-h_{\alpha\mu}\widetilde{\Gamma}_{\beta\gamma}^{\mu}+h_{\gamma\mu}\widetilde{\Gamma}_{\alpha\beta}^{\mu}\right)$
$\displaystyle R_{\alpha\beta}$ $\displaystyle=$ $\displaystyle
g^{11}R_{\alpha 1\beta 1}+g^{\mu\nu}R_{\alpha\mu\beta\nu}$ $\displaystyle=$
$\displaystyle\frac{{f^{\prime}}^{2}-2ff^{\prime\prime}}{4f}h_{\alpha\beta}+\widetilde{R}_{\alpha\beta}+\frac{{f^{\prime}}^{2}}{4f}h^{\mu\nu}\left(h_{\alpha\nu}h_{\mu\beta}-h_{\alpha\beta}h_{\mu\nu}\right)$
$\displaystyle=$
$\displaystyle\frac{(2-n){f^{\prime}}^{2}-2ff^{\prime\prime}}{4f}h_{\alpha\beta}+\widetilde{R}_{\alpha\beta},$
where $\widetilde{R}_{\alpha\beta}$’s are components of the Ricci curvature of
$(N,ds_{N}^{2})$. If we assume the functions $f$, $f^{\prime}/f$ and
$f^{\prime\prime}/f$ are all bounded, then in the chart $(\mathbb{R}\times
U,Id\times\phi)$, the eigenvalues of the matrix $(R_{jl})$ and the matrix
$(g_{jl})$ are uniformly bounded. Thus there exists some constant
$K_{1}\in\mathbb{R}$ such that $(R_{jl})\geq K_{1}(g_{jl})$. Note that
$(N,ds_{N}^{2})$ is compact. There exists some constant $K\in\mathbb{R}$ such
that $Ric_{(M,g)}\geq Kg$ as bilinear forms. If we further assume
$\lim_{x\rightarrow+\infty}f(x)=0$, then by (LABEL:volume), we have ${\rm
Vol}_{g}(B_{y}(1))\rightarrow 0$ as $x\rightarrow+\infty$, where
$y=(x,m)\in\mathbb{R}\times N$. One can check that the following functions
satisfy all the above assumptions on $f$.
$\bullet$ $f$ is a smooth positive function defined on $\mathbb{R}$ and
satisfies
$f(x)=\left\\{\begin{array}[]{lll}(1+x^{2})e^{-x+\sin x}&{\rm
when}&x>1\\\\[6.45831pt] 1,&{\rm when}&x<0\end{array}\right.$
$\bullet$ $f$ is a smooth positive function defined on $\mathbb{R}$ and
satisfies
$f(x)=\left\\{\begin{array}[]{lll}\frac{1}{\log x}&{\rm
when}&x>2\\\\[6.45831pt] 1,&{\rm when}&x<0.\end{array}\right.$
This gives the desired result.$\hfill\Box$
## 5 Proof of proposition 2.5
In this section, we shall construct complete noncompact Riemannain
$n$-manifolds to show that the condition Ricci curvature has lower bound in
theorem 2.3 is not necessarily needed.
Proof of proposition 2.5. It suffices to construct a complete noncompact
Riemannian $n$-manifold on which Trudinger-Moser embedding holds, but its
Ricci curvature has no lower bound. For this purpose, we consider the
Riemannian manifold $(\mathbb{R}^{n},g)$, where $\mathbb{R}^{n}$ is the
euclidian space and
$g=dx_{1}^{2}+f(x_{1})dx_{2}^{2}+\cdots+f(x_{1})dx_{n}^{2},$
and $f$ is a smooth function on $\mathbb{R}$ such that $a\leq f\leq b$ for two
positive constants $a$ and $b$. Clearly $(\mathbb{R}^{n},g)$ is complete and
noncompact. In view of Trudinger-Moser inequality on the standard euclidian
space $\mathbb{R}^{n}$ [8, 12, 32], one can easily see that if $\alpha$ is
chosen sufficiently small, then the supremum
$\sup_{u\in W^{1,n}(M),\,\|u\|_{W^{1,n}}\leq
1}\int_{\mathbb{R}^{n}}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}$
is finite, i.e. Trudinger-Moser inequality holds on the manifold
$(\mathbb{R}^{n},g)$, where
$\|u\|_{W^{1,n}}=\left(\int_{\mathbb{R}^{n}}(|\nabla_{g}u|^{n}+|u|^{n})dv_{g}\right).$
In the following, we shall further choose $f$ such that the Ricci curvature of
$(\mathbb{R}^{n},g)$ is unbounded from below. By (4.15),
$R_{11}=(n-1)\frac{{f^{\prime}}^{2}-2ff^{\prime\prime}}{4f^{2}}.$ (5.1)
It suffices to find a sequence of points $(x^{(m)})$ of $\mathbb{R}^{n}$ such
that $R_{11}(x^{(m)})\rightarrow-\infty$. One choice of $f$ is that
$f(t)=2+\sin t^{2}$. In this case, we have
$\displaystyle f^{\prime}(x_{1})=2+2x_{1}\cos x_{1}^{2},\quad
f^{\prime\prime}(x_{1})=2\cos x_{1}^{2}-4x_{1}^{2}\sin x_{1}^{2}.$
Thus (5.1) implies
$R_{11}(x)=(n-1)\frac{(2+2x_{1}\cos x_{1}^{2})^{2}-2(2+\sin x_{1}^{2})(2\cos
x_{1}^{2}-4x_{1}^{2}\sin x_{1}^{2})}{4(2+\sin x_{1}^{2})^{2}}.$
Choosing $x^{(m)}=\left(\sqrt{2m\pi+3\pi/2},0,\cdots,0\right)$, we obtain
$R_{11}({x^{(m)}})=-4m\pi-3\pi+n-1\rightarrow-\infty\quad{\rm as}\quad
m\rightarrow\infty.$
Another choice of $f$ is that $f(t)=e^{\sin t^{2}}$. In this case, we have
$\displaystyle f^{\prime}(x_{1})$ $\displaystyle=$ $\displaystyle
2x_{1}e^{\sin x_{1}^{2}}\cos x_{1}^{2},\quad f^{\prime\prime}(x_{1})=e^{\sin
x_{1}^{2}}\left(-4x_{1}^{2}\sin x_{1}^{2}+4x_{1}^{2}\cos^{2}x_{1}^{2}+2\cos
x_{1}^{2}\right).$
In view of (5.1), we obtain
$R_{11}(x)=(n-1)(2x_{1}^{2}\sin
x_{1}^{2}+x_{1}^{2}\cos^{2}x_{1}^{2}-2x_{1}^{2}\cos^{2}x_{1}^{2}-\cos
x_{1}^{2}).$
Again, we select $x^{(m)}=\left(\sqrt{2m\pi+3\pi/2},0,\cdots,0\right)$ and
conclude $R_{11}(x^{(m)})\rightarrow-\infty$ as $m\rightarrow\infty$.
$\hfill\Box$
## 6 Adams inequalities
In this section, we concern Adams inequalities on complete noncompact
Riemannian manifolds. Precisely we shall prove theorem 2.6. The method we
adopted here is similar to that of theorem 2.3.
Proof of theorem 2.6. $(i)$ Suppose that ${\rm inj}_{(M,g)}\geq i_{0}>0$ and
there exist constants $C(k)$ such that $|\nabla^{k}{\rm Rc}_{(M,g)}|\leq
C(k)$, $k=0,1,\cdots,m-1$. It follows from (Hebey [18], theorem 1.3) that for
any $Q>1$ and $\alpha\in(0,1)$, the harmonic radius $r_{H}=r_{H}(Q,m,\alpha)$
is positive. Namely, for any $Q>1$, $\alpha\in(0,1)$, and $x\in M$, there
exists a harmonic coordinate chart
$\psi:B_{x}(r_{H})\rightarrow\mathbb{R}^{n}$ such that
$\left\\{\begin{array}[]{lll}Q^{-1}\delta_{lq}\leq g_{lq}\leq
Q\delta_{lq}\quad{\rm as\,\,a\,\,bilinear\,\,form};\\\\[6.45831pt]
\sum_{1\leq|\beta|\leq
m}\|\partial^{\beta}g_{lq}\|_{C^{0}(B_{x}(r_{H}))}+\sum_{|\beta|=m}\|\partial^{\beta}g_{lq}\|_{C^{\alpha}(B_{x}(r_{H}))}\leq
Q-1.\end{array}\right.$ (6.1)
Now we fix $Q>1$ and $\alpha\in(0,1)$. Without loss of generality, we may
assume $\psi(x)=0$. Particularly we have that for any $r:0<r\leq r_{H}$
$\mathbb{B}_{0}(r/\sqrt{Q})\subset\psi(B_{x}(r))\subset\mathbb{B}_{0}({\sqrt{Q}}r).$
Let $\eta\in C_{0}^{\infty}(\mathbb{R}^{n})$ be such that $0\leq\eta\leq 1$,
and
$\eta=\left\\{\begin{array}[]{lll}1&\,\,\,{\rm
on}&\mathbb{B}_{0}(r_{H}/(4\sqrt{Q})),\\\\[6.45831pt] 0&\,\,\,{\rm
on}&\mathbb{R}^{n}\setminus\mathbb{B}_{0}(r_{H}/(2\sqrt{Q})).\end{array}\right.$
Then $\eta\circ\psi\in C_{0}^{\infty}(M)$ satisfies $0\leq\eta\circ\psi\leq
1$, $\eta\circ\psi\equiv 1$ on $B_{x}(r_{H}/(4Q))$, and $\eta\circ\psi\equiv
0$ on $M\setminus B_{x}(r_{H}/2)$. By Gromov’s covering lemma (Hebey [18],
lemma 1.6), there exists a sequence of points $(x_{k})$ of $M$ such that
$M=\cup_{k}B_{x_{k}}(r_{H}/(4Q))$ (6.2)
and there exists some integer $N$ such that for any $x\in M$, $x$ belongs to
at most $N$ balls in the covering. Let
$\psi_{k}:B_{x_{k}}(r_{H})\rightarrow\mathbb{R}^{n}$ be as the above $\psi$
and set $\eta_{k}=\eta\circ\psi_{k}$. By (6.1), the components of the metric
tensor are $C^{m}$-controlled in the charts $(B_{x_{k}}(r_{H}),\psi_{k})$. It
then follows that there exists some constant $C_{1}>0$ depending only on
$r_{H}$ and $Q$ such that $|\nabla_{g}^{l}\eta_{k}|\leq C_{1}$ for all all
$l:0\leq l\leq m$ and all $k\in\mathbb{N}$, where $\nabla_{g}^{l}$ is defined
by (1.3).
Assume $u\in C^{\infty}(M)$ satisfies $\|u\|_{W^{m,n/m}(M)}\leq 1$. Then we
get
$\eta_{k}^{m+1}u\in C_{0}^{\infty}(B_{x_{k}}(r_{H}/2))$
and
$\|\nabla_{g}^{m}(\eta_{k}^{m+1}u)\|_{L^{\frac{n}{m}}(B_{x_{k}}(r_{H}/2))}\leq
C_{2}$ (6.3)
for some constant $C_{2}$ depending only on $n$, $m$, and $C_{1}$. By the
standard elliptic estimates (Gilbarg-Trudinger [16], Chapter 9), one can see
that
$\|\nabla_{\mathbb{R}^{n}}^{m}((\eta_{k}^{m+1}u)\circ\psi_{k}^{-1})\|_{L^{\frac{n}{m}}(\mathbb{B}_{0}(\sqrt{Q}r_{H}))}\leq
C_{3}$ (6.4)
for some constant $C_{3}$ depending only on $n$, $m$, $Q$, $r_{H}$ and
$C_{1}$. Let $j$ be the smallest integer great than or equal to $n/m$.
Similarly as we derived (4.8), we calculate by using (6.2), (6.3) and the
relation ${(j-1)n}/{(n-m)}\geq{n}/{m}$
$\displaystyle\int_{M}\zeta\left(j,\alpha|u|^{\frac{n}{n-m}}\right)dv_{g}$
$\displaystyle\leq$
$\displaystyle\sum_{k}\int_{B_{x_{k}}(r_{H}/(4Q))}\zeta\left(j,\alpha|u|^{\frac{n}{n-m}}\right)dv_{g}{}$
(6.5) $\displaystyle\leq$
$\displaystyle\sum_{k}\int_{B_{x_{k}}(r_{H}/2)}\zeta\left(j,\alpha|\eta_{k}^{m+1}u|^{\frac{n}{n-m}}\right)dv_{g}{}$
$\displaystyle\leq$
$\displaystyle\sum_{k}\left(\frac{\|\nabla_{g}^{m}(\eta_{k}^{m+1}u)\|_{L^{\frac{n}{m}}(B_{x_{k}}(r_{H}/2))}}{C_{2}}\right)^{\frac{(j-1)n}{n-m}}\int_{B_{x_{k}}(r_{H}/2)}\zeta\left(j,\alpha
C_{2}^{\frac{n}{n-m}}|\eta_{k}^{m+1}u|^{\frac{n}{n-m}}\right)dv_{g}{}$
$\displaystyle\leq$
$\displaystyle\sum_{k}\frac{\|\nabla_{g}^{m}(\eta_{k}^{m+1}u)\|_{L^{\frac{n}{m}}(B_{x_{k}}(r_{H}/2))}^{\frac{n}{m}}}{C_{2}^{\frac{n}{m}}}\int_{B_{x_{k}}(r_{H}/2)}\zeta\left(j,\alpha
C_{2}^{\frac{n}{n-m}}|\eta_{k}^{m+1}u|^{\frac{n}{n-m}}\right)dv_{g}.$
Noting that $Q^{-1}\delta_{lq}\leq g_{lq}\leq Q\delta_{lq}$ as a bilinear
form, we have
$\int_{B_{x_{k}}(r_{H}/2)}\zeta\left(j,\alpha
C_{2}^{\frac{n}{n-m}}|\eta_{k}^{m+1}u|^{\frac{n}{n-m}}\right)dv_{g}\leq
Q^{\frac{n}{2}}\int_{\mathbb{B}_{0}(\sqrt{Q}r_{H})}\zeta\left(j,\alpha
C_{2}^{\frac{n}{n-m}}|(\eta_{k}^{m+1}u)\circ\psi_{k}^{-1}|^{\frac{n}{n-m}}\right)dx.$
(6.6)
In view of (6.4), we take
$\alpha_{0}=\beta_{0}/(C_{2}C_{3})^{\frac{n}{n-m}}.$ (6.7)
Then for any $\alpha:0<\alpha\leq\alpha_{0}$, it follows from Adams inequality
(1.4) that
$\int_{\mathbb{B}_{0}(\sqrt{Q}r_{H})}\zeta\left(j,\alpha
C_{2}^{\frac{n}{n-m}}|(\eta_{k}^{m+1}u)\circ\psi_{k}^{-1}|^{\frac{n}{n-m}}\right)dx\leq
C_{m,n}|\mathbb{B}_{0}(\sqrt{Q}r_{H})|.$ (6.8)
Clearly there exists some constant $C_{4}>0$ depending only on $n$, $m$, $Q$
and $r_{H}$ such that
$|\nabla_{g}^{l}\eta_{k}^{m+1}|^{\frac{n}{m}}\leq C_{4}\eta_{k},\quad\forall
l=0,1,\cdots,m.$ (6.9)
Since $1\leq\sum_{k}\eta_{k}(x)\leq N$ for all $x\in M$, we obtain by
combining (6.5)-(6.9) that
$\displaystyle\int_{M}\zeta\left(j,\alpha|u|^{\frac{n}{n-m}}\right)dv_{g}$
$\displaystyle\leq$ $\displaystyle
C_{5}\sum_{k}\int_{M}|\nabla_{g}^{m}(\eta_{k}^{m+1}u)|^{\frac{n}{m}}dv_{g}$
$\displaystyle\leq$ $\displaystyle
C_{5}\sum_{k}\sum_{l=0}^{m}(C_{m}^{l})^{\frac{n}{m}}\int_{M}|\nabla_{g}^{m-k}\eta_{k}^{m+1}\nabla_{g}^{l}u|^{\frac{n}{m}}dv_{g}$
$\displaystyle\leq$ $\displaystyle
C_{4}C_{5}\sum_{l=0}^{m}(C_{m}^{l})^{\frac{n}{m}}\int_{M}(\sum_{k}\eta_{k})|\nabla_{g}^{l}u|^{\frac{n}{m}}dv_{g}$
$\displaystyle\leq$ $\displaystyle
C_{4}C_{5}N\sum_{l=0}^{m}(C_{m}^{l})^{\frac{n}{m}}\int_{M}|\nabla_{g}^{l}u|^{\frac{n}{m}}dv_{g}$
$\displaystyle\leq$ $\displaystyle C_{6}$
for constants $C_{5}$ and $C_{6}$ depending only on $n$, $m$, $Q$ and $r_{H}$,
where $C_{m}^{l}=\frac{m!}{l!\,(m-l)!}$.
According to (Hebey [18], theorem 2.8), $C_{0}^{\infty}(M)$ is dense in
$W^{m,\frac{n}{m}}(M)$. Hence for any $u\in W^{m,\frac{n}{m}}(M)$, there
exists a sequence $(u_{k})$ in $C_{0}^{\infty}(M)$ such that
$\|u_{k}-u\|_{W^{m,\frac{n}{m}}(M)}\rightarrow 0$ as $k\rightarrow\infty$.
Assume $\|u\|_{W^{m,\frac{n}{m}}(M)}\leq 1$. Then for any
$\alpha:0<\alpha<\alpha_{0}$ there holds
$\displaystyle\int_{M}\zeta\left(j,\alpha|u|^{\frac{n}{n-m}}\right)dv_{g}$
$\displaystyle\leq$
$\displaystyle\lim_{k\rightarrow\infty}\int_{M}\zeta\left(j,\alpha|u_{k}|^{\frac{n}{n-m}}\right)dv_{g}\leq
C_{6}.$
Using the same method of deriving $W^{1,n}(M)\hookrightarrow L^{q}(M)$
continuously for all $q\geq n$ in theorem 2.3, we obtain the continuous
embedding $W^{m,{n}/{m}}(M)\hookrightarrow L^{q}(M)$ for any $q\geq{n}/{m}$.
$(ii)$ Let $\alpha>0$ be any real number and $u$ be any function belonging to
the space $W^{m,\frac{n}{m}}(M)$. Since $C_{0}^{\infty}(M)$ is dense in
$W^{m,\frac{n}{m}}(M)$, there exists some $u_{0}\in C_{0}^{\infty}(M)$ such
that
$\alpha\|u-u_{0}\|_{W^{m,\frac{n}{m}}(M)}^{\frac{n}{n-m}}<\alpha_{0}/2,$
(6.10)
where $\alpha_{0}$ is defined by (6.7). Using (2.3) and an elementary
inequality
$|a|^{p}\leq(1+\epsilon)|a-b|^{p}+c(\epsilon,p)|b|^{p},$
where $\epsilon>0$, $p>1$ and $c(\epsilon,p)$ is a constant depending only on
$\epsilon$ and $p$, we have
$\displaystyle\int_{M}\zeta\left(j,\alpha|u|^{\frac{n}{n-m}}\right)dv_{g}$
$\displaystyle\leq$
$\displaystyle\int_{M}\zeta\left(j,(1+\epsilon)\alpha|u-u_{0}|^{\frac{n}{n-m}}+c(\epsilon,n/(n-m))\alpha|u_{0}|^{\frac{n}{n-m}}\right)dv_{g}$
(6.11) $\displaystyle\leq$
$\displaystyle\frac{1}{\mu}\int_{M}\zeta\left(j,\mu(1+\epsilon)\alpha|u-u_{0}|^{\frac{n}{n-m}}\right)dv_{g}$
$\displaystyle+\frac{1}{\nu}\int_{M}\zeta\left(j,\nu
c(\epsilon,n/(n-m))\alpha|u_{0}|^{\frac{n}{n-m}}\right)dv_{g},$
where $\mu>1$, $\nu>1$ and $1/\mu+1/\nu=1$. Choosing $\epsilon$ sufficiently
small and $\mu$ sufficiently close to $1$ such that
$\mu(1+\epsilon)\alpha_{0}/2\leq\alpha_{0}$, in view of (6.10), we have by
part $(i)$
$\int_{M}\zeta\left(j,\mu(1+\epsilon)\alpha|u-u_{0}|^{\frac{n}{n-m}}\right)dv_{g}\leq
C_{6}.$ (6.12)
Note that $u_{0}\in C_{0}^{\infty}(M)$, particularly $u_{0}$ has compact
support. It follows that
$\int_{M}\zeta\left(j,\nu
c(\epsilon,n/(n-m))\alpha|u_{0}|^{\frac{n}{n-m}}\right)dv_{g}<\infty.$ (6.13)
Inserting (6.12) and (6.13) into (6.11), we complete the proof of part $(ii)$.
$\hfill\Box$
## 7 Applications of Trudinger-Moser inequalities
In this section, we consider applications of theorem 2.3, namely the existence
and multiplicity results for the problem (2.4) and its perturbation (2.8).
Specifically we shall prove theorem 2.7 and theorem 2.10. Throughout this
section, we use the notations introduced in section 2. Let $(M,g)$ be a
complete noncompact Riemannian $n$-manifold with ${\rm Rc}_{(M,g)}\geq Kg$ for
some $K\in\mathbb{R}$ and ${\rm inj}_{(M,g)}\geq i_{0}>0$. Assume $\phi(x)$
satisfies the hypotheses $(\phi_{1})$ and $(\phi_{2})$, $v(x)$ satisfies the
hypotheses $(v_{1})$ and $(v_{2})$. Let $E$ be a function space defined by
(2.5). If $u\in E$, then the $E$-norm of $u$ is defined by
$\|u\|_{E}=\left(\int_{M}(|\nabla_{g}u|^{n}+v|u|^{n})dv_{g}\right)^{1/n}.$
The following compact embedding result is very important in our analysis.
Proposition 7.1. For any $q\geq n$, the function space $E$ is compactly
embedded in $L^{q}(M)$.
Proof. Let $(u_{k})$ be a sequence of functions with $\|u_{k}\|_{E}\leq C$ for
some constant $C$. It suffices to prove that up to a subsequence, $(u_{k})$
converges in $L^{q}(M)$ for any $q\geq n$. Clearly $(u_{k})$ is bounded in
$W^{1,n}(M)$, and thus we can assume that for any $q>1$, up to a subsequence
$\displaystyle u_{k}\rightharpoonup u_{0}\quad{\rm weakly\,\,in}\quad E{}$
$\displaystyle u_{k}\rightarrow u_{0}\quad{\rm strongly\,\,in}\quad L^{q}_{\rm
loc}(M)$ (7.1) $\displaystyle u_{k}\rightarrow u_{0}\quad{\rm
a.\,e.}\,\,\,{\rm in}\quad M.$
If $v(x)\in L^{1/(n-1)}(M)$, using the same argument of ([44], Lemma 2.4), we
conclude that $E\hookrightarrow L^{q}(M)$ compactly for any $q>1$. So, in view
of $(v_{2})$, we may assume $v(x)\rightarrow\infty$ as
$d_{g}(O,x)\rightarrow\infty$, where $O$ is a fixed point of $M$. Given any
$\epsilon>0$, there exists some $R>0$ such that $v(x)>(2C)^{n}/\epsilon$ when
$d_{g}(O,x)\geq R$. Hence
$\frac{(2C)^{n}}{\epsilon}\int_{M\setminus
B_{O}(R)}|u_{k}-u_{0}|^{n}dv_{g}<\int_{M}v|u_{k}-u_{0}|^{n}dv_{g}\leq(2C)^{n}.$
This gives
$\int_{M\setminus B_{O}(R)}|u_{k}-u_{0}|^{n}dv_{g}<\epsilon.$
By (LABEL:Lploc), we have
$\lim_{k\rightarrow\infty}\int_{B_{O}(R)}|u_{k}-u_{0}|^{n}dv_{g}=0.$
Hence for the above $\epsilon$, there exists some $l\in\mathbb{N}$ such that
when $k>l$,
$\int_{M}|u_{k}-u_{0}|^{n}dv_{g}<2\epsilon.$
This implies $u_{k}\rightarrow u_{0}$ strongly in $L^{n}(M)$ as
$k\rightarrow\infty$.
It follows from $(i)$ of theorem 2.3 that $(u_{k})$ is bounded in $L^{q}(M)$
for any $q\geq n$. Now fixing $q>n$, we get by Hölder’s inequality
$\int_{M}|u_{k}-u_{0}|^{q}dv_{g}\leq\left(\int_{M}|u_{k}-u_{0}|^{n}dv_{g}\right)^{1/n}\left(\int_{M}|u_{k}-u_{0}|^{\frac{n(q-1)}{n-1}}dv_{g}\right)^{1-1/n}.$
This together with the fact that $u_{k}\rightarrow u_{0}$ in $L^{n}(M)$
implies $u_{k}\rightarrow u_{0}$ in $L^{q}(M)$. $\hfill\Box$
Let $S_{q}$ be defined by (2.7). Then we have the following:
Proposition 7.2. For any $q>n$, $S_{q}$ is attained by some nonnegative
function $u\in E\setminus\\{0\\}$.
Proof. Assume $q>n$. It is easy to see that
$S_{q}^{n}=\inf_{\int_{M}\phi|u|^{q}dv_{g}=1}\int_{M}\left(|\nabla
u|^{n}+v|u|^{n}\right)dv_{g}.$
Choosing a sequence of functions $(u_{k})\subset E$ such that
$\int_{M}\phi|u_{k}|^{q}dv_{g}=1$ and
$\lim_{k\rightarrow\infty}\int_{M}\left(|\nabla
u_{k}|^{n}+v|u_{k}|^{n}\right)dv_{g}=S_{q}^{n}.$
By proposition 7.1, there exists some $u\in E$ such that up to a subsequence,
$u_{k}\rightharpoonup u$ weakly in $E$, $u_{k}\rightarrow u$ strongly in
$L^{q}(M)$ for any $q\geq n$, and $u_{k}\rightarrow u$ almost everywhere in
$M$. Since $u_{k}\rightarrow u$ strongly in $L^{s}(B_{O}(R_{0}))$ for all
$s>1$ and $\phi\in L^{p}(B_{O}(R_{0}))$, we have by using Hölder’s inequality
that
$\lim_{k\rightarrow\infty}\int_{B_{O}(R_{0})}\phi|u_{k}|^{q}dv_{g}=\int_{B_{O}(R_{0})}\phi|u|^{q}dv_{g}.$
(7.2)
In view of $(v_{2})$, we have
$\displaystyle\int_{M\setminus B_{O}(R_{0})}\phi||u_{k}|^{q}-|u|^{q}|dv_{g}$
$\displaystyle\leq$ $\displaystyle
qC_{0}\int_{M}(|u_{k}|^{q-1}+|u|^{q-1})|u_{k}-u|dv_{g}$ $\displaystyle\leq$
$\displaystyle
qC_{0}\left\\{\left(\int_{M}|u_{k}|^{q}dv_{g}\right)^{1-1/q}+\left(\int_{M}|u|^{q}dv_{g}\right)^{1-1/q}\right\\}$
$\displaystyle\quad\times\left(\int_{M}|u_{k}-u|^{q}dv_{g}\right)^{1/q}$
$\displaystyle\rightarrow$ $\displaystyle 0\,\,\,{\rm
as}\,\,\,k\rightarrow\infty.$
This together with (7.2) implies
$\int_{M}\phi|u|^{q}dv_{g}=\lim_{k\rightarrow\infty}\int_{M}\phi|u_{k}|^{q}dv_{g}=1.$
(7.3)
Since $u_{k}\rightharpoonup u$ weakly in $E$, we have
$\displaystyle\int_{M}|\nabla
u|^{n}dv_{g}=\lim_{k\rightarrow\infty}\int_{M}|\nabla u|^{n-2}\nabla u\nabla
u_{k}dv_{g}\leq\limsup_{k\rightarrow\infty}\left(\int_{M}|\nabla
u_{k}|^{n}dv_{g}\right)^{\frac{1}{n}}\left(\int_{M}|\nabla
u|^{n}dv_{g}\right)^{1-\frac{1}{n}},$
from which we obtain
$\int_{M}|\nabla u|^{n}dv_{g}\leq\limsup_{k\rightarrow\infty}\int_{M}|\nabla
u_{k}|^{n}dv_{g}.$ (7.4)
In addition, we have by Fatou’s lemma
$\int_{M}v|u|^{n}dv_{g}\leq\limsup_{k\rightarrow\infty}\int_{M}v|u_{k}|^{n}dv_{g}.$
(7.5)
Combining (7.3), (7.4) and (7.5), we conclude that $S_{q}$ is attained by
$u\in E\setminus\\{0\\}$. Since $|u|\in E$, one can easily see that $S_{q}$ is
also attained by $|u|$. $\hfill\Box$
Now we get back to the problem (2.4). Since we are interested in nonnegative
weak solutions, without loss of generality we may assume $f(x,t)\equiv 0$ for
all $(x,t)\in M\times(-\infty,0]$. By $(f_{1})$, we have for all $(x,t)\in
M\times\mathbb{R}$,
$|F(x,t)|\leq\frac{b_{1}}{n}|t|^{n}+b_{2}t\zeta\left(n,|t|^{\frac{n}{n-1}}\right).$
This together with $(\phi_{1})$, $(\phi_{2})$ and (2.2) implies that for any
$u\in E$ there holds
$\displaystyle\int_{M}\phi F(x,u)dv_{g}$ $\displaystyle\leq$
$\displaystyle\|\phi\|_{L^{p}(B_{O}(R_{0}))}\|F(x,u)\|_{L^{q}(M)}+C_{0}\int_{M}F(x,u)dv_{g}$
$\displaystyle\leq$
$\displaystyle\|\phi\|_{L^{p}(B_{O}(R_{0}))}\left(\frac{b_{1}}{n}\|u\|_{L^{qn}(M)}^{n}+b_{2}\|u\zeta(n,|u|^{\frac{n}{n-1}})\|_{L^{q}(M)}\right)$
$\displaystyle+C_{0}\frac{b_{1}}{n}\|u\|_{L^{n}(M)}^{n}+C_{0}b_{2}\|u\zeta(n,|u|^{\frac{n}{n-1}})\|_{L^{1}(M)}$
$\displaystyle\leq$ $\displaystyle
C\left(\|u\|_{L^{qn}(M)}^{n}+\|u\|_{L^{qn}(M)}\|\zeta(n,\frac{qn}{n-1}|u|^{\frac{n}{n-1}})\|_{L^{1}(M)}^{1-\frac{1}{n}}\right.$
$\displaystyle\left.\|u\|_{L^{n}(M)}^{n}+\|u\|_{L^{n}(M)}\|\zeta(n,\frac{n}{n-1}|u|^{\frac{n}{n-1}})\|_{L^{1}(M)}\right),$
where $C$ is a constant depending only on $n$, $b_{1}$, $b_{2}$, $C_{0}$ and
$\|\phi\|_{L^{p}(B_{O}(R_{0}))}$, and $1/p+1/q=1$. By theorem 2.3, $u\in
L^{s}(M)$ for all $s\geq n$, and for any $\alpha>0$ there holds
$\zeta(n,\alpha|u|^{\frac{n}{n-1}})\in L^{1}(M)$. Hence
$\int_{M}\phi F(x,u)dv_{g}<+\infty,\quad\forall u\in E.$
Based on this, we can define a functional on $E$ by
$J(u)=\frac{1}{n}\|u\|_{E}^{n}-\int_{M}\phi F(x,u)dv_{g}.$ (7.6)
By ([13], proposition 1) and the standard argument [34], we have
$J\in\mathcal{C}^{1}(E,\mathbb{R})$. Clearly the critical point of $J$ is a
weak solution to (2.4). Concerning the geometry of $J$, the following two
lemmas imply that $J$ has a mountain pass structure.
Lemma 7.3. Assume that $(f_{1})$, $(f_{2})$, and $(f_{3})$ are satisfied. Then
for any nonnegative, compactly supported function $u\in E\setminus\\{0\\}$,
there holds $J(tu)\rightarrow-\infty$ as $t\rightarrow+\infty$.
Proof. By $(f_{2})$ and $(f_{3})$, there exist $c_{1}$, $c_{2}>0$ and $\mu>n$
such that $F(x,s)\geq c_{1}s^{\mu}-c_{2}$ for all $(x,s)\in
M\times[0,+\infty)$. Assume ${\rm supp}\,u\subset B_{O}(R_{1})$ for some
$R_{1}>0$. We have
$\displaystyle J(tu)$ $\displaystyle=$
$\displaystyle\frac{t^{n}}{n}\|u\|_{E}^{n}-\int_{B_{O}(R_{1})}\phi
F(x,tu)dv_{g}$ $\displaystyle\leq$
$\displaystyle\frac{t^{n}}{n}\|u\|_{E}^{n}-c_{1}t^{\mu}\int_{B_{O}(R_{1})}\phi
u^{\mu}dv_{g}-c_{2}\int_{B_{O}(R_{1})}\phi dv_{g}.$
This gives the desired result since $\phi(x)>0$ for all $x\in M$ and
$\mu>n$.$\hfill\Box$
Lemma 7.4. Assume that $(f_{1})$ and $(f_{4})$ are satisfied. Then there exist
sufficiently small constants $r>0$ and $\delta>0$ such that $J(u)\geq\delta$
for all $u$ with $\|u\|_{E}=r$.
Proof. By $(f_{1})$ and $(f_{4})$, there exists some constants
$\theta\in(0,1)$ and $C>0$ such that
$F(x,s)\leq\frac{(1-\theta)\lambda_{\phi}}{n}|s|^{n}+C|s|^{n+1}\zeta\left(n,\alpha_{0}|s|^{\frac{n}{n-1}}\right)$
for all $(x,s)\in M\times\mathbb{R}$. By definition of $\lambda_{\phi}$,
$\frac{(1-\theta)\lambda_{\phi}}{n}\int_{M}\phi|u|^{n}dv_{g}\leq\frac{1-\theta}{n}\|u\|_{E}^{n}.$
(7.7)
Note that $\phi$ satisfies $(\phi_{1})$ and $(\phi_{2})$. We have by Hölder’s
inequality and (2.2) that
$\displaystyle\int_{M}\phi|u|^{n+1}\zeta\left(n,\alpha_{0}|u|^{\frac{n}{n-1}}\right)dv_{g}$
$\displaystyle\leq$
$\displaystyle\|\phi\|_{L^{p}(B_{O}(R_{0}))}\left(\int_{M}|u|^{(n+1)q}dv_{g}\right)^{1/q}\left(\int_{M}\zeta\left(n,q^{\prime}\alpha_{0}|u|^{\frac{n}{n-1}}\right)dv_{g}\right)^{1/q^{\prime}}$
(7.8)
$\displaystyle+C_{0}\left(\int_{M}|u|^{(n+1)\beta}dv_{g}\right)^{1/\beta}\left(\int_{M}\zeta\left(n,\gamma\alpha_{0}|u|^{\frac{n}{n-1}}\right)dv_{g}\right)^{1/\gamma},$
where $1/p+1/q+1/q^{\prime}=1$ and $1/\beta+1/\gamma=1$. Fix
$\alpha=\beta_{0}/2$, where $\beta_{0}$ is defined by (1.5). It follows from
$(i)$ of theorem 2.3 that there exists some constant $\tau$ depending only on
$\alpha$, $n$, $K$ and $i_{0}$ such that
$\Lambda_{\alpha}:=\sup_{\|u\|_{1,\tau}\leq
1}\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}<+\infty.$ (7.9)
Let $r$ be a positive constant to be determined later. Now suppose
$\|u\|_{E}=r$. It is easy to see that $\|u\|_{1,\tau}\leq r+\tau
r/v_{0}^{1/n}$. Clearly one can select $r$ sufficiently small such that
$q^{\prime}\alpha_{0}\|u\|_{1,\tau}^{n/(n-1)}<\alpha$ and
$\gamma\alpha_{0}\|u\|_{1,\tau}^{n/(n-1)}<\alpha$. It follows from (7.9) that
$\sup_{\|u\|_{E}=r}\int_{M}\zeta\left(n,q^{\prime}\alpha_{0}|u|^{\frac{n}{n-1}}\right)dv_{g}\leq\Lambda_{\alpha}$
and
$\sup_{\|u\|_{E}=r}\int_{M}\zeta\left(n,\gamma\alpha_{0}|u|^{\frac{n}{n-1}}\right)dv_{g}\leq\Lambda_{\alpha},$
provided that $r$ is chosen sufficiently small. Inserting these two
inequalities into (7.8), then using the embedding $E\hookrightarrow L^{s}(M)$
for all $s\geq n$ (proposition 7.1) and (7.7), we obtain
$J(u)\geq\frac{\theta}{n}\|u\|_{E}^{n}-\tilde{C}\|u\|_{E}^{n+1}$
for some constant $\tilde{C}$ depending only on $\alpha$, $n$, $K$ and
$i_{0}$, provided that $\|u\|_{E}$ is sufficiently small. This gives the
desired result. $\hfill\Box$
To estimate the min-max level of $J$, we state the following:
Lemma 7.5. Assume $(f_{5})$. There exists some nonnegative function $u^{*}\in
E$ such that
$\sup_{t\geq
0}J(tu^{*})<\frac{1}{n}\left(\frac{(p-1)\alpha_{n}}{p\alpha_{0}}\right)^{n-1}.$
Proof. Let $u^{*}$ be given by proposition 7.2, namely $u^{*}\geq 0$,
$\|u^{*}\|_{E}=S_{q}$, and $\int_{M}\phi|u^{*}|^{q}dv_{g}=1$. Then for any
$t\geq 0$ there holds
$\displaystyle J(tu^{*})$ $\displaystyle=$
$\displaystyle\frac{1}{n}\|tu^{*}\|_{E}^{n}-\int_{M}\phi(x)F(x,tu^{*})dv_{g}$
$\displaystyle\leq$
$\displaystyle\frac{S_{q}^{n}}{n}t^{n}-\frac{C_{q}}{q}t^{q}$
$\displaystyle\leq$
$\displaystyle\frac{q-n}{nq}\frac{S_{q}^{{nq}/{(q-n)}}}{C_{q}^{{n}/{(q-n)}}}$
$\displaystyle<$
$\displaystyle\frac{1}{n}\left(\frac{(p-1)\alpha_{n}}{p\alpha_{0}}\right)^{n-1}.$
Here we have used the hypothesis $(f_{5})$. $\hfill\Box$
Adapting the proof of ([44], lemma 3.4), we obtain the following compactness
result.
Lemma 7.6. Assume $(f_{1})$, $(f_{2})$ and $(f_{3})$. Let $(u_{j})\subset E$
be an arbitrary Palais-Smale sequence of $J$, i.e.,
$J(u_{j})\rightarrow c,\,\,J^{\prime}(u_{j})\rightarrow 0\,\,{\rm
in}\,\,E^{*}\,\,{\rm as}\,\,j\rightarrow\infty,$ (7.10)
where $E^{*}$ denotes the dual space of $E$. Then there exist a subsequence of
$(u_{j})$ (still denoted by $(u_{j})$) and $u\in E$ such that
$u_{j}\rightharpoonup u$ weakly in $E$, $u_{j}\rightarrow u$ strongly in
$L^{q}(M)$ for all $q\geq n$, and
$\displaystyle\left\\{\begin{array}[]{lll}\nabla u_{j}(x)\rightarrow\nabla
u(x)\quad{\rm a.\,\,e.\,\,\,in}\quad M\\\\[6.45831pt]
{\phi(x)F(x,\,u_{j})}\rightarrow{\phi(x)F(x,\,u)}\,\,{\rm
strongly\,\,in}\,\,L^{1}(M).\end{array}\right.$
Furthermore $u$ is a weak solution of (2.4).
Proof. Assume $(u_{j})$ is a Palais-Smale sequence of $J$. By $(\ref{PS})$, we
have
$\displaystyle\frac{1}{n}\|u_{j}\|_{E}^{n}-\int_{M}\phi(x)F(x,u_{j})dv_{g}\rightarrow
c\,\,{\rm as}\,\,j\rightarrow\infty,$ (7.12)
$\displaystyle\left|\int_{M}\left(|\nabla_{g}u_{j}|^{n-2}\nabla_{g}u_{j}\nabla_{g}\psi+v|u_{j}|^{n-2}u_{j}\psi\right)dv_{g}-\int_{M}\phi(x)f(x,u_{j})\psi
dv_{g}\right|\leq\sigma_{j}\|\psi\|_{E}\qquad\quad$ (7.13)
for all $\psi\in E$, where $\sigma_{j}\rightarrow 0$ as $j\rightarrow\infty$.
Note that $f(x,s)\equiv 0$ for all $(x,s)\in M\times(-\infty,0]$. By
$(f_{2})$, we have $0\leq\mu F(x,u_{j})\leq u_{j}f(x,u_{j})$ for some $\mu>n$.
Taking $\psi=u_{j}$ in (7.13) and multiplying (7.12) by $\mu$, we have
$\displaystyle\left(\frac{\mu}{n}-1\right)\|u_{j}\|_{E}^{n}$
$\displaystyle\leq$ $\displaystyle\mu|c|+\int_{M}\phi(x)\left(\mu
F(x,u_{j})-f(x,u_{j})u_{j}\right)dv_{g}+\sigma_{j}\|u_{j}\|_{E}$
$\displaystyle\leq$ $\displaystyle\mu|c|+\sigma_{j}\|u_{j}\|_{E}.$
Therefore $\|u_{j}\|_{E}$ is bounded. It then follows from (7.12) and (7.13)
that
$\int_{M}\phi(x)f(x,u_{j})u_{j}dv_{g}\leq
C,\quad\int_{M}\phi(x)F(x,u_{j})dv_{g}\leq C$ (7.14)
for some constant $C$ depending only on $\mu$, $n$ and $c$. By proposition
7.1, there exists some $u\in E$ such that $u_{j}\rightharpoonup u$ weakly in
$E$, $u_{j}\rightarrow u$ strongly in $L^{q}(M)$ for any $q\geq n$, and
$u_{j}\rightarrow u$ almost everywhere in $M$. By $(f_{3})$, there exist
positive constants $A_{1}$ and $R_{1}$ such that $F(x,s)\leq A_{1}f(x,s)$ for
all $s\geq R_{1}$. Particularly for any $A>R_{1}$ there holds
$F(x,s)\leq A_{1}f(x,s),\quad\forall s\geq A.$ (7.15)
Now we prove that $\phi(x)F(x,u_{j})\rightarrow\phi F(x,u)$ strongly in
$L^{1}(M)$. To this end, for any $\epsilon>0$, we take
$A>\max\\{A_{1}C/\epsilon,R_{1}\\}$, where $C$ is given by (7.14). Then we
have by (7.15)
$\int_{|u_{j}|>A}\phi(x)F(x,u_{j})dv_{g}\leq\frac{A_{1}}{A}\int_{M}\phi(x)f(x,u_{j})u_{j}dv_{g}<\epsilon.$
(7.16)
In the same way
$\int_{|u|>A}\phi(x)F(x,u)dv_{g}<\epsilon.$ (7.17)
By $(f_{1})$, we have for $(x,s)\in M\times[0,\infty)$
$\displaystyle f(x,s)$ $\displaystyle\leq$ $\displaystyle
b_{1}s^{n-1}+b_{2}\zeta\left(n,\alpha_{0}s^{\frac{n}{n-1}}\right)$
$\displaystyle=$ $\displaystyle
b_{1}s^{n-1}+b_{2}s^{n}\sum_{k=n-1}^{\infty}\frac{\alpha_{0}^{k}s^{\frac{n}{n-1}(k-n+1)}}{k!}$
$\displaystyle\leq$ $\displaystyle
b_{1}s^{n-1}+b_{2}s^{n}\alpha_{0}^{n-1}e^{\alpha_{0}s^{\frac{n}{n-1}}}.$
Hence for all $(x,s)\in M\times[0,A]$ there holds
$f(x,s)\leq\left(b_{1}+b_{2}\alpha_{0}^{n-1}Ae^{\alpha_{0}A^{\frac{n}{n-1}}}\right)s^{n-1}.$
It follows that
$F(x,s)\leq\frac{b_{1}+b_{2}\alpha_{0}^{n-1}Ae^{\alpha_{0}A^{\frac{n}{n-1}}}}{n}s^{n},\quad\forall
s\in[0,A].$
for all $(x,s)\in M\times[0,A]$, which implies
$|\phi(x)\chi_{\\{|u_{j}|\leq A\\}}(x)F(x,u_{j})|\leq
C_{1}\phi(x)|u_{j}|^{n},$ (7.18)
where
$C_{1}={(b_{1}+b_{2}\alpha_{0}^{n-1}Ae^{\alpha_{0}A^{{n}/{(n-1)}}})}/{n}$ and
$\chi_{\\{|u_{j}|\leq A\\}}(x)$ denotes the characteristic function of the set
$\\{x\in M:|u_{j}(x)|\leq A\\}$. By an inequality $||a|^{n}-|b|^{n}|\leq
n|a-b|(|a|^{n-1}+|b|^{n-1})\,(\forall a,b\in\mathbb{R})$ and Hölder’s
inequality, we get
$\displaystyle\int_{M}\phi||u_{j}|^{n}-|u|^{n}|dv_{g}$ $\displaystyle\leq$
$\displaystyle n\int_{M}\phi|u_{j}-u|(|u_{j}|^{n-1}+|u|^{n-1})dv_{g}$
$\displaystyle\leq$ $\displaystyle
n\left(\int_{M}\phi|u_{j}-u|^{n}dv_{g}\right)^{\frac{1}{n}}\left\\{\left(\int_{M}\phi|u_{j}|^{n}dv_{g}\right)^{1-\frac{1}{n}}+\left(\int_{M}\phi|u|^{n}dv_{g}\right)^{1-\frac{1}{n}}\right\\}.$
Hence $\phi|u_{j}|^{n}\rightarrow\phi|u|^{n}$ in $L^{1}(M)$ since
$u_{j}\rightarrow u$ strongly in $L^{n}(M)$. In view of (7.18), we conclude
from the generalized Lebesgue’s dominated convergence theorem
$\lim_{j\rightarrow\infty}\int_{M}\phi(x)\chi_{\\{|u_{j}|\leq
A\\}}(x)F(x,u_{j})dv_{g}=\int_{M}\phi(x)\chi_{\\{|u|\leq
A\\}}(x)F(x,u)dv_{g}.$
This together with (7.16) and (7.17) implies that there exists some
$m\in\mathbb{N}$ such that when $j>m$ there holds
$\left|\int_{M}\phi F(x,u_{j})dv_{g}-\int_{M}\phi
F(x,u)dv_{g}\right|<3\epsilon.$
Therefore
$\lim_{j\rightarrow\infty}\int_{M}\phi F(x,u_{j})dv_{g}=\int_{M}\phi
F(x,u)dv_{g}.$
Using the same method as that of proving ([4], (4.26)), we have
$\nabla_{g}u_{j}(x)\rightarrow\nabla_{g}u(x)$ for almost every $x\in M$ and
$|\nabla_{g}u_{j}|^{n-2}\nabla_{g}u_{j}\rightharpoonup|\nabla_{g}u|^{n-2}\nabla_{g}u\quad{\rm
weakly\,\,in}\quad\left(L^{\frac{n}{n-1}}(M)\right)^{n}.$
Passing to the limit $j\rightarrow\infty$ in $(\ref{2})$, we obtain
$\int_{M}\left(|\nabla_{g}u|^{n-2}\nabla_{g}u\nabla\psi+v|u|^{n-2}u\psi\right)dv_{g}-\int_{M}\phi(x)f(x,u)\psi
dv_{g}=0$
for all $\psi\in C_{0}^{\infty}(M)$. Since $C_{0}^{\infty}(M)$ is dense in $E$
under the norm $\|\cdot\|_{E}$, $u$ is a weak solution of $(\ref{equa})$.
$\hfill\Box$
We say more words on lemma 7.6. Suppose $(M,g)$ is the standard euclidian
space $\mathbb{R}^{n}$ and $\phi(x)=|x|^{-\beta}$, $0\leq\beta<n$. The author
[44] proved that $\phi F(x,u_{j})\rightarrow\phi F(x,u)$ in
$L^{1}(\mathbb{R}^{n})$ under the assumption $E\hookrightarrow
L^{q}(\mathbb{R}^{n})$ compactly for all $q\geq 1$. While Lam-Lu [20] observed
that the convergence still holds under the assumption $E\hookrightarrow
L^{q}(\mathbb{R}^{n})$ for all $q\geq n$. Here we generalized these two
situations.
The following lemma is a nontrivial consequence of theorem 2.3. It is
sufficient for our use when we consider the existence and multiplicity results
for problems (2.4) and (2.8).
Lemma 7.7. Let $(u_{j})\subset E$ be any sequence of functions satisfying
$\|u_{j}\|_{E}\leq 1$, $u_{j}\rightharpoonup u_{0}$ weakly in $E$,
$\nabla_{g}u_{j}\rightarrow\nabla_{g}u_{0}$ almost everywhere in M, and
$u_{j}\rightarrow u_{0}$ strongly in $L^{n}(M)$ as $j\rightarrow\infty$. Then
$(i)$ for any $\alpha:0<\alpha<\alpha_{n}$, there holds
$\sup_{j}\int_{M}\zeta\left(n,\alpha|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}<\infty;$
(7.19)
$(ii)$ for any $\alpha:0<\alpha<\alpha_{n}$ and
$q:0<q<(1-\|u_{0}\|_{E}^{n})^{-1/(n-1)}$, there holds
$\sup_{j}\int_{M}\zeta\left(n,q\alpha|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}<\infty.$
(7.20)
Proof. $(i)$ For any fixed $\alpha:0<\alpha<\alpha_{n}$, it follows from part
$(i)$ of theorem 2.3 that there exists a positive constant $\tau_{\alpha}$
depending only on $\alpha$, $n$, $K$ and $i_{0}$ such that
$\mathcal{B}_{\alpha}=\sup_{u\in W^{1,n}(M),\,\|u\|_{1,\tau_{\alpha}}\leq
1}\int_{M}\zeta\left(n,\alpha|u|^{\frac{n}{n-1}}\right)dv_{g}<\infty.$ (7.21)
Note that $v\geq v_{0}$ in $M$. Since $\|u_{j}\|_{E}\leq 1$, we get
$\displaystyle\|u_{j}\|_{1,\tau_{\alpha}}=\left(\int_{M}|\nabla_{g}u_{j}|^{n}dv_{g}\right)^{\frac{1}{n}}+\tau_{\alpha}\left(\int_{M}|u_{j}|^{n}dv_{g}\right)^{\frac{1}{n}}\leq
1+\frac{\tau_{\alpha}}{v_{0}^{1/n}}.$
There exists some small positive number $\alpha_{0}$ such that
$\alpha_{0}\|u_{j}\|_{1,\tau_{\alpha}}^{\frac{n}{n-1}}\leq\alpha$. Hence by
(7.21), there holds
$\sup_{j}\int_{M}\zeta\left(n,\alpha_{0}|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}\leq\sup_{j}\int_{M}\zeta\left(n,\alpha\left|\frac{u_{j}}{\|u_{j}\|_{1,\tau_{\alpha}}}\right|^{\frac{n}{n-1}}\right)dv_{g}\leq\mathcal{B}_{\alpha}.$
This allows us to define
$\alpha^{*}=\sup\left\\{\alpha:\sup_{j}\int_{M}\zeta\left(n,\alpha|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}<\infty\right\\}.$
To prove (7.19), it suffices to prove that $\alpha^{*}\geq\alpha_{n}$. Suppose
not, we have $\alpha^{*}<\alpha_{n}$. Take two constants $\alpha^{\prime}$ and
$\alpha^{\prime\prime}$ such that
$\alpha^{*}<\alpha^{\prime}<\alpha^{\prime\prime}<\alpha_{n}$. By part $(i)$
of theorem 2.3 again, there exists some constant
$\tau_{\alpha^{\prime\prime}}$ depending only on $\alpha^{\prime\prime}$, $n$,
$K$ and $i_{0}$ such that
$\mathcal{B}_{\alpha^{\prime\prime}}=\sup_{u\in
W^{1,n}(M),\,\|u\|_{1,\tau_{\alpha^{\prime\prime}}}\leq
1}\int_{M}\zeta\left(n,\alpha^{\prime\prime}|u|^{\frac{n}{n-1}}\right)dv_{g}<\infty.$
(7.22)
Since $u_{j}\rightarrow u_{0}$ strongly in $L^{n}(M)$ and
$\nabla_{g}u_{j}\rightarrow\nabla_{g}u_{0}$ a. e. in M, we obtain by using
Brezis-Lieb’s lemma [6]
$\|u_{j}-u_{0}\|_{1,\tau_{\alpha\,^{\prime\prime}}}=\left(\int_{M}|\nabla_{g}u_{j}|^{n}dv_{g}-\int_{M}|\nabla_{g}u_{0}|^{n}dv_{g}\right)^{{1}/{n}}+o_{j}(1),$
where $o_{j}(1)\rightarrow 0$ as $j\rightarrow\infty$. Since
$u_{j}\rightharpoonup u_{0}$ weakly in $E$, there holds
$\lim_{j\rightarrow+\infty}\int_{M}|\nabla_{g}u_{0}|^{n-2}\nabla_{g}u_{0}\nabla_{g}u_{j}\,dv_{g}=\int_{M}|\nabla_{g}u_{0}|^{n}dv_{g}.$
This immediately implies that
$\int_{M}|\nabla_{g}u_{0}|^{n}dv_{g}\leq\limsup_{j\rightarrow+\infty}\int_{M}|\nabla_{g}u_{j}|^{n}dv_{g}\leq
1.$
Hence
$\|u_{j}-u_{0}\|_{1,\tau_{\alpha\,^{\prime\prime}}}\leq 1+o_{j}(1).$
It follows from (2.3) that for any $\epsilon>0$ there exists some constant
$\tilde{c}$ depending only on $\epsilon$ and $n$ such that
$\displaystyle\zeta\left(n,\alpha^{\prime}|u_{j}|^{\frac{n}{n-1}}\right)\leq\frac{1}{\mu}\zeta\left(n,\alpha^{\prime}(1+\epsilon)\mu|u_{j}-u_{0}|^{\frac{n}{n-1}}\right)+\frac{1}{\nu}\zeta\left(n,\alpha^{\prime}\tilde{c}\nu|u_{0}|^{\frac{n}{n-1}}\right),$
(7.23)
where $1/\mu+1/\nu=1$. Choosing $\epsilon$ sufficiently small and $\mu$
sufficiently close to $1$ such that
$\alpha^{\prime}(1+\epsilon)\mu\|u_{j}-u_{0}\|_{1,\tau_{\alpha^{\prime\prime}}}^{\frac{n}{n-1}}<\alpha^{\prime\prime},$
provided that $j$ is sufficiently large. This together with (7.22) implies
that
$\sup_{j}\int_{M}\zeta\left(n,\alpha^{\prime}(1+\epsilon)\mu|u_{j}-u_{0}|^{\frac{n}{n-1}}\right)dv_{g}\leq\mathcal{B}_{\alpha^{\prime\prime}}.$
(7.24)
In addition, we have by part $(iii)$ of theorem 2.3 that
$\int_{M}\zeta\left(n,\alpha^{\prime}\tilde{c}\nu|u_{0}|^{\frac{n}{n-1}}\right)dv_{g}<+\infty.$
(7.25)
Inserting (7.24) and (7.25) into (7.23), we get
$\sup_{j}\int_{M}\zeta\left(n,\alpha^{\prime}|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}<+\infty,$
which contradicts the definition of $\alpha^{*}$ and thus ends the proof of
part $(i)$.
$(ii)$ Given any $\alpha:0<\alpha<\alpha_{n}$ and any
$q:0<q<(1-\|u_{0}\|_{E}^{n})^{-1/(n-1)}$. By (2.3), $\forall\epsilon>0$, there
exist constants $\tilde{c}>0$, $\mu>1$ and $\nu>1$ $(1/\mu+1/\nu=1)$ such that
$\displaystyle\int_{M}\zeta\left(n,q\alpha|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}\leq\frac{1}{\mu}\int_{M}\zeta\left(n,q\alpha(1+\epsilon)\mu|u_{j}-u_{0}|^{\frac{n}{n-1}}\right)dv_{g}+\frac{1}{\nu}\int_{M}\zeta\left(n,q\alpha\tilde{c}\nu|u_{0}|^{\frac{n}{n-1}}\right)dv_{g}.$
By Brezis-Lieb’s lemma [6],
$\|u_{j}-u_{0}\|_{E}^{\frac{n}{n-1}}\leq(1-\|u_{0}\|_{E}^{n})^{\frac{1}{n-1}}+o_{j}(1).$
If we choose $\epsilon$ sufficiently small and $\mu$ sufficiently close to $1$
such that
$q\alpha(1+\epsilon)\mu\|u_{j}-u_{0}\|_{E}^{\frac{n}{n-1}}\leq(\alpha+\alpha_{n})/2,$
provided that $j$ is sufficiently large. It then follows from part $(i)$ that
$\sup_{j}\int_{M}\zeta\left(n,q\alpha(1+\epsilon)\mu|u_{j}-u_{0}|^{\frac{n}{n-1}}\right)dv_{g}<+\infty.$
By part $(iii)$ of theorem 2.3, we have
$\int_{M}\zeta\left(n,q\alpha\tilde{c}\nu|u_{0}|^{\frac{n}{n-1}}\right)dv_{g}<+\infty.$
Therefore (7.20) holds. $\hfill\Box$
Remark 7.8. In lemma 7.7, if $u_{0}\equiv 0$, then the conclusion of $(ii)$ is
weaker than that of $(i)$. If $u_{0}\not\equiv 0$, then the conclusion of
$(i)$ is a special case of that of $(ii)$. If $(M,g)$ has dimension two, the
assumption $\nabla_{g}u_{j}\rightarrow\nabla_{g}u_{0}$ almost everywhere in
$M$ can be removed.
Proof of theorem 2.7. It follows from lemma 7.3 and lemma 7.4 that $J$
satisfies all the hypothesis of the mountain-pass theorem except for the
Palais-Smale condition: $J\in\mathcal{C}^{1}(E,\mathbb{R})$; $J(0)=0$;
$J(u)\geq\delta>0$ when $\|u\|_{E}=r$; $J(e)<0$ for some $e\in E$ with
$\|e\|_{E}>r$. Then using the mountain-pass theorem without the Palais-Smale
condition [34], we can find a sequence $(u_{j})$ in $E$ such that
$J(u_{j})\rightarrow c>0,\quad J^{\prime}(u_{j})\rightarrow 0\,\,{\rm
in}\,\,E^{*},$
where
$c=\min_{\gamma\in\Gamma}\max_{u\in\gamma}J(u)\geq\delta$
is the min-max value of $J$, where
$\Gamma=\\{\gamma\in\mathcal{C}([0,1],E):\gamma(0)=0,\gamma(1)=e\\}$. This is
equivalent to (7.12) and $(\ref{2})$. By lemma 7.6, up to a subsequence, there
holds
$\left\\{\begin{array}[]{lll}u_{j}\rightharpoonup u\,\,{\rm
weakly\,\,in}\,\,E\\\\[6.45831pt] u_{j}\rightarrow u\,\,{\rm
strongly\,\,in}\,\,L^{q}(M),\,\,\forall q\geq n\\\\[6.45831pt]
\lim\limits_{j\rightarrow\infty}\int_{M}\phi(x){F(x,u_{j})}dv_{g}=\int_{M}\phi(x){F(x,u)}dv_{g}\\\\[6.45831pt]
u\,\,{\rm
is\,\,a\,\,weak\,\,solution\,\,of}\,\,(\ref{equa}).\end{array}\right.$ (7.26)
Now suppose by contradiction $u\equiv 0$. Since $F(x,0)=0$ for all $x\in M$,
it follows from (7.12) and (7.26) that
$\lim_{j\rightarrow\infty}\|u_{j}\|_{E}^{n}=nc>0.$ (7.27)
By lemma 7.5,
$0<c<\frac{1}{n}\left(\frac{(p-1)\alpha_{n}}{p\alpha_{0}}\right)^{n-1}$. Thus
there exists some $\eta_{0}>0$ and $m>0$ such that
$\|u_{j}\|_{E}^{n}\leq\left(\frac{p-1}{p}\frac{\alpha_{n}}{\alpha_{0}}-\eta_{0}\right)^{n-1}$
for all $j>m$. Choose $q>1$ sufficiently close to $1$ such that
$q\alpha_{0}\|u_{j}\|_{E}^{\frac{n}{n-1}}\leq(1-1/p)\alpha_{n}-\alpha_{0}\eta_{0}/2$
for all $j>m$. By $(f_{1})$,
$|f(x,u_{j})u_{j}|\leq
b_{1}|u_{j}|^{n}+b_{2}|u_{j}|\zeta\left(n,\alpha_{0}|u_{j}|^{\frac{n}{n-1}}\right).$
It follows from (2.2), Hölder’s inequality, and part $(i)$ of lemma 7.7 that
$\displaystyle\int_{M}\phi{|f(x,u_{j})u_{j}|}dv_{g}$ $\displaystyle\leq$
$\displaystyle
b_{1}\int_{M}\phi|u_{j}|^{n}dv_{g}+b_{2}\int_{M}\phi|u_{j}|\zeta\left(n,\alpha_{0}|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}$
$\displaystyle\leq$ $\displaystyle
b_{1}\int_{M}\phi{|u_{j}|^{n}}dv_{g}+b_{2}\left(\int_{M}\phi|u_{j}|^{q^{\prime}}dv_{g}\right)^{1/{q^{\prime}}}\left(\int_{M}\phi\zeta\left(n,q\alpha_{0}|u_{j}|^{\frac{n}{n-1}}\right)dv_{g}\right)^{1/{q}}$
$\displaystyle\leq$ $\displaystyle
b_{1}\int_{M}\phi{|u_{j}|^{n}}dv_{g}+C\left(\int_{M}\phi{|u_{j}|^{q^{\prime}}}dv_{g}\right)^{1/{q^{\prime}}}\rightarrow
0\quad{\rm as}\quad j\rightarrow\infty,$
where $1/q+1/q^{\prime}=1$ and $C$ is some constant which is independent of
$j$. Here we have used (7.26) again (precisely $u_{j}\rightarrow u$ in
$L^{r}(\mathbb{R}^{N})$ for all $r\geq n$) in the above estimates. Inserting
this into (7.13) with $\psi=u_{j}$, we have
$\|u_{j}\|_{E}\rightarrow 0\quad{\rm as}\quad j\rightarrow\infty,$
which contradicts (7.27). Therefore $u\not\equiv 0$ and we obtain a nontrivial
weak solution of (2.4). Finally $u$ is nonnegative since $f(x,s)\equiv 0$ for
all $(x,s)\in M\times(-\infty,0]$. $\hfill\Box$
Proof of theorem 2.10. Since the proof is very similar to that of ([44],
theorem 1.2), we only give its sketch and emphasize the difference between
these two situations. Instead of $J:E\rightarrow\mathbb{R}$ defined by (7.6),
we consider functionals for all $u\in E$ and $\epsilon>0$
$J_{\epsilon}(u)=\frac{1}{n}\|u\|_{E}^{n}-\int_{M}\phi(x){F(x,u)}dv_{g}-\epsilon\int_{M}hudv_{g}.$
Firstly, lemma 7.6 still holds if we replace $J$ by $J_{\epsilon}$. Namely for
any Palais-Smale sequence $(u_{j})\subset E$ of $J_{\epsilon}$, there exist a
subsequence of $(u_{j})$ (still denoted by $(u_{j})$) and $u\in E$ such that
$u_{j}\rightharpoonup u$ weakly in $E$, $u_{j}\rightarrow u$ strongly in
$L^{q}(M)$ for all $q\geq n$, and
$\left\\{\begin{array}[]{lll}\nabla_{g}u_{j}(x)\rightarrow\nabla_{g}u(x)\quad{\rm
a.\,\,e.\,\,\,in}\quad M\\\\[6.45831pt]
{\phi(x)F(x,\,u_{j})}\rightarrow{\phi(x)F(x,\,u)}\,\,{\rm
strongly\,\,in}\,\,L^{1}(M)\\\\[6.45831pt] u\,\,{\rm
is\,\,a\,\,weak\,\,solution\,\,of}\,\,(\ref{equa1}).\end{array}\right.$ (7.28)
Secondly, using the same method in the first two steps of the proof of ([44],
theorem 1.2), we have the following:
$(a)$ there exist constants $\epsilon_{1}>0$, $\delta>0$, and a sequence of
functions $(v_{j})\subset E$ such that $J_{\epsilon}(v_{j})\rightarrow c_{M}$
and $J^{\prime}_{\epsilon}(v_{j})\rightarrow 0$ as $j\rightarrow\infty$,
provided that $0<\epsilon<\epsilon_{1}$. In addition, $v_{j}$ is bounded in
$E$, $v_{j}\rightharpoonup u_{M}$ weakly in $E$ and $u_{M}$ is a weak solution
of (2.8). Here $c_{M}$ is the min-max value of $J_{\epsilon}$ and satisfies
$0<c_{M}<\frac{1}{n}\left(1-\frac{1}{p}\right)^{n-1}\left(\frac{\alpha_{n}}{\alpha_{0}}\right)^{n-1}-\delta;$
(7.29)
$(b)$ there exists a constant $\epsilon_{2}:0<\epsilon_{2}<\epsilon_{1}$ such
that for any $\epsilon:0<\epsilon<\epsilon_{2}$, there exist positive constant
$r_{\epsilon}$ with $r_{\epsilon}\rightarrow 0$ as $\epsilon\rightarrow 0$ and
sequence $(u_{j})\subset E$ such that
$J_{\epsilon}(u_{j})\rightarrow c_{\epsilon}=\inf_{\|u\|_{E}\leq
r_{\epsilon}}J_{\epsilon}(u)<0,\quad J_{\epsilon}^{\prime}(u_{j})\rightarrow
0\quad{\rm in}\quad E^{*}\quad{\rm as}\quad j\rightarrow\infty.$
In addition, $u_{j}\rightarrow u_{0}$ strongly in $E$, where $u_{0}$ is a weak
solution of (2.8) with $J_{\epsilon}(u_{0})=c_{\epsilon}$.
Thirdly, there exists $\epsilon_{0}:0<\epsilon_{0}<\epsilon_{2}$ such that if
$0<\epsilon<\epsilon_{0}$, then $u_{M}\not\equiv u_{0}$. Suppose by
contradiction that $u_{M}\equiv u_{0}$. Then $v_{j}\rightharpoonup u_{0}$
weakly in $E$. By $(a)$,
$J_{\epsilon}(v_{j})\rightarrow c_{M}>0,\quad|\langle
J_{\epsilon}^{\prime}(v_{j}),\varphi\rangle|\leq\gamma_{j}\|\varphi\|_{E}$
(7.30)
with $\gamma_{j}\rightarrow 0$ as $j\rightarrow\infty$. On one hand we have by
(7.28),
$\int_{M}\phi(x)F(x,v_{j})dv_{g}\rightarrow\int_{M}\phi(x)F(x,u_{0})dv_{g}\quad{\rm
as}\quad j\rightarrow\infty.$ (7.31)
On the other hand, since $v_{j}\rightharpoonup u_{0}$ weakly in $E$ and $h\in
E^{*}$, it follows that
$\int_{M}hv_{j}dv_{g}\rightarrow\int_{M}hu_{0}dv_{g}\quad{\rm as}\quad
j\rightarrow\infty.$ (7.32)
Inserting (7.31) and (7.32) into the first equality of (7.30), we obtain
$\frac{1}{n}\|v_{j}\|_{E}^{n}=c_{M}+\int_{M}\phi(x)F(x,u_{0})dv_{g}+\epsilon\int_{M}hu_{0}dv_{g}+o_{j}(1).$
(7.33)
In the same way, one can derive
$\frac{1}{n}\|u_{j}\|_{E}^{n}=c_{\epsilon}+\int_{M}\phi(x)F(x,u_{0})dv_{g}+\epsilon\int_{M}hu_{0}dv_{g}+o_{j}(1).$
(7.34)
Combining (7.33) and (7.34), we have
$\|v_{j}\|_{E}^{n}-\|u_{0}\|_{E}^{n}=n\left(c_{M}-c_{\epsilon}+o_{j}(1)\right).$
(7.35)
From $(b)$, we know that $c_{\epsilon}\rightarrow 0$ as $\epsilon\rightarrow
0$. This together with (7.29) leads to the existence of
$\epsilon_{0}:0<\epsilon_{0}<\epsilon_{2}$ such that if
$0<\epsilon<\epsilon_{0}$, then
$0<c_{M}-c_{\epsilon}<\frac{1}{n}\left(\frac{p-1}{p}\frac{\alpha_{n}}{\alpha_{0}}\right)^{n-1}.$
(7.36)
Write
$w_{j}=\frac{v_{j}}{\|v_{j}\|_{E}},\quad
w_{0}=\frac{u_{0}}{\left(\|u_{0}\|_{E}^{n}+n(c_{M}-c_{\epsilon})\right)^{1/n}}.$
It follows from (7.35) and $v_{j}\rightharpoonup u_{0}$ weakly in $E$ that
$w_{j}\rightharpoonup w_{0}$ weakly in $E$. Note that
$\int_{M}\phi(x)\zeta\left(n,\alpha_{0}|v_{j}|^{n/(n-1)}\right)dv_{g}=\int_{M}\phi(x)\zeta\left(n,\alpha_{0}\|v_{j}\|_{E}^{{n}/{(n-1)}}|w_{j}|^{n/(n-1)}\right)dv_{g}.$
By (7.35) and (7.36), a straightforward calculation shows
$\lim_{j\rightarrow\infty}\alpha_{0}\|v_{j}\|_{E}^{\frac{n}{n-1}}\left(1-\|w_{0}\|_{E}^{n}\right)^{\frac{1}{n-1}}<\left(1-\frac{1}{p}\right)\alpha_{n}.$
Hence lemma 7.7 together with (2.3) implies that
$\phi(x)\zeta\left(n,\alpha_{0}|v_{j}|^{n/(n-1)}\right)$ is bounded in
$L^{q}(M)$ for some $q:1<q<n/(n-1)$. By $(f_{1})$,
$|f(x,v_{j})|\leq
b_{1}|v_{j}|^{n-1}+b_{2}\zeta(n,\alpha_{0}|v_{j}|^{\frac{n}{n-1}}).$
By the definition of $\zeta$ there exists a constant $c>0$ such that
$|f(x,v_{j})\chi_{\\{|v_{j}|\leq 1\\}}(x)|\leq
c|v_{j}|^{n-1},\quad|f(x,v_{j})\chi_{\\{|v_{j}|>1\\}}(x)|\leq
c\zeta(n,\alpha_{0}|v_{j}|^{\frac{n}{n-1}}),$
where $\chi_{B}$ denotes the characteristic function of $B\subset M$. Hence
$\displaystyle\left|\int_{M}\phi(x)f(x,v_{j})(v_{j}-u_{0})dv_{g}\right|$
$\displaystyle\leq$ $\displaystyle
c\int_{M}\phi(x)\left(|v_{j}|^{n-1}+\zeta(n,\alpha_{0}|v_{j}|^{\frac{n}{n-1}})\right)|v_{j}-u_{0}|dv_{g}$
$\displaystyle\leq$ $\displaystyle
c\left\|\phi|v_{j}|^{n-1}\right\|_{L^{\frac{n}{n-1}}(M)}\|v_{j}-u_{0}\|_{L^{n}(M)}$
$\displaystyle+c\left\|\phi\zeta(n,\alpha_{0}|v_{j}|^{\frac{n}{n-1}})\right\|_{L^{q}(M)}\|v_{j}-u_{0}\|_{L^{q^{\prime}}(M)}.$
Since $1<q<n/(n-1)$, we have $q^{\prime}>n$. Then it follows from the compact
embedding $E\hookrightarrow L^{r}(M)$ for all $r\geq n$ that
$\lim_{j\rightarrow\infty}\int_{M}\phi(x)f(x,v_{j})(v_{j}-u_{0})dv_{g}=0.$
(7.37)
Taking $\varphi=v_{j}-u_{0}$ in (7.30), we have by using (7.32) and (7.37)
that
$\int_{M}\left(|\nabla_{g}v_{j}|^{n-2}\nabla_{g}v_{j}\nabla_{g}(v_{j}-u_{0})+v(x)|v_{j}|^{n-2}v_{j}(v_{j}-u_{0})\right)dv_{g}\rightarrow
0.$ (7.38)
However the fact $v_{n}\rightharpoonup u_{0}$ weakly in $E$ leads to
$\int_{M}\left(|\nabla_{g}u_{0}|^{n-2}\nabla_{g}u_{0}\nabla_{g}(v_{j}-u_{0})+v(x)|u_{0}|^{n-2}u_{0}(v_{j}-u_{0})\right)dv_{g}\rightarrow
0.$ (7.39)
Subtracting (7.39) from (7.38), using the well known inequality (see [26],
chapter 10)
$2^{n-1}|b-a|^{n}\leq\langle|b|^{n-2}b-|a|^{n-2}a,b-a\rangle,\quad\forall
a,b\in\mathbb{R}^{n},$
we have $\|v_{j}-u_{0}\|_{E}^{n}\rightarrow 0$ as $j\rightarrow\infty$. This
together with (7.35) implies that $c_{M}=c_{\epsilon}$, which is absurd since
$c_{M}>0$ and $c_{\epsilon}<0$. Therefore $u_{M}\not\equiv u_{0}$. Since
$f(x,s)\equiv 0$ for all $(x,s)\in M\times(-\infty,0]$, one can easily see
that $u_{M}\geq 0$ and $u_{0}\geq 0$. This completes the proof of the theorem.
$\hfill\Box$
Finally we shall construct examples of $f$’s to show that $(f_{1})$-$(f_{5})$
do not imply $(H_{5})$.
Proof of proposition 2.9. Let $\phi$ satisfies the hypotheses $(\phi_{1})$ and
$(\phi_{2})$, $p>1$ be given in $(\phi_{1})$, $l$ be an integer satisfying
$l\geq n$, $q=nl/(n-1)+1$ and $S_{q}$ be defined by (2.7). In view of lemma
7.2, $S_{q}$ is attained by some nonnegative function $u\in E$. Let $C_{q}$ be
a positive number such that
$C_{q}>\left(\frac{q-n}{q}\right)^{{(q-n)}/{n}}\left(\frac{p\alpha_{0}}{(p-1)\alpha_{n}}\right)^{(q-n)(n-1)/n}S_{q}^{q}.$
Let $\chi:[0,\infty)\rightarrow\mathbb{R}$ be a smooth function such that
$0\leq\chi\leq 1$, $\chi\equiv 0$ on $[0,A]$, $\chi\equiv 1$ on $[2A,\infty)$,
and $|\chi\,^{\prime}|\leq 2/A$, where $A$ is a positive constant to be
determined later. We set
$\displaystyle
f(t)=\left\\{\begin{array}[]{lll}2^{l}l!C_{q}\sum_{k=l}^{\infty}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!},&t\geq
0\\\\[6.45831pt] 0,&t<0.\end{array}\right.$
Now we check $(f_{1})$-$(f_{5})$ for appropriate choice of $A$ as follows.
$(f_{1})$: If $A>1$, then $0\leq t^{n/(n-1)}-\chi(t)t^{1/(n-1)}\leq
t^{n/(n-1)}$ for all $t\geq 0$. Thus
$\displaystyle f(t)$ $\displaystyle=$ $\displaystyle
2^{l}l!C_{q}\left(e^{t^{n/(n-1)}-\chi(t)t^{1/(n-1)}}-\sum_{k=0}^{l-1}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}\right)$
$\displaystyle\leq$ $\displaystyle
2^{l}l!C_{q}\left(e^{t^{n/(n-1)}}-\sum_{k=0}^{l-1}\frac{t^{\frac{nk}{n-1}}}{k!}\right)$
$\displaystyle\leq$ $\displaystyle 2^{l}l!C_{q}\zeta(n,t^{n/(n-1)})$
for all $t\geq 0$. So $(f_{1})$ is satisfied when $A>1$.
$(f_{2})$: When $t\in[0,A]$, we have $\chi(t)=0$ and
$\int_{0}^{t}f(t)dt=2^{l}l!C_{q}\sum_{k=l}^{\infty}\int_{0}^{t}\frac{t^{\frac{nk}{n-1}}}{k!}dt\leq{2^{l}l!C_{q}t}\sum_{k=l}^{\infty}\frac{t^{\frac{nk}{n-1}}}{k!}=tf(t).$
(7.41)
When $t\geq A$, we claim that if $A$ is chosen sufficiently large, say $A\geq
4^{n-1}$, then
$\int_{A}^{t}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}dt\leq\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k+1}}{(k+1)!}-\frac{A^{\frac{n(k+1)}{n-1}}}{(k+1)!}.$
(7.42)
In fact, if we set
$\gamma(t)=\int_{A}^{t}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}dt-\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k+1}}{(k+1)!}+\frac{A^{\frac{n(k+1)}{n-1}}}{(k+1)!},$
then $\gamma(A)=0$ and
$\gamma\,^{\prime}(t)=\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}-\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}\left(\frac{n}{n-1}t^{\frac{1}{n-1}}-\chi\,^{\prime}(t)t^{\frac{1}{n-1}}-\frac{1}{n-1}\chi(t)t^{\frac{1}{n-1}-1}\right).$
Let $A\geq 4^{n-1}$. Then for $t\in[A,\infty)$ there holds
$\displaystyle\frac{n}{n-1}t^{\frac{1}{n-1}}-\chi\,^{\prime}(t)t^{\frac{1}{n-1}}-\frac{1}{n-1}\chi(t)t^{\frac{1}{n-1}-1}$
$\displaystyle\geq$
$\displaystyle\left(\frac{n}{n-1}-\frac{2}{A}\right)A^{\frac{1}{n-1}}-\frac{1}{n-1}A^{\frac{1}{n-1}-1}$
$\displaystyle\geq$ $\displaystyle
4\left(\frac{n}{n-1}-\frac{2}{4(n-1)}-\frac{1}{4(n-1)^{2}}\right)$
$\displaystyle>$ $\displaystyle 1.$
Hence $\gamma\,^{\prime}(t)\leq 0$ and thus our claim (7.42) holds. Note that
$\int_{0}^{A}\frac{t^{\frac{nk}{n-1}}}{k!}dt=\frac{A^{\frac{n(k+1)}{n-1}}}{(k+1)!}\frac{k+1}{\frac{nk}{n-1}+1}A^{-\frac{1}{n-1}}\leq\frac{A^{\frac{n(k+1)}{n-1}}}{(k+1)!}.$
(7.43)
It follows from (7.42) and (7.43) that when $t\geq A$,
$\int_{0}^{t}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k}}{k!}dt\leq\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{k+1}}{(k+1)!},$
and whence
$\int_{0}^{t}f(t)dt\leq f(t)\leq\frac{1}{\mu}tf(t)$ (7.44)
for some $\mu>n$. This together with (7.41) implies that $(f_{2})$ holds.
$(f_{3})$: Let $A\geq 4^{n-1}$. In view of (7.44), when $t\geq A$,
$F(t)=\int_{0}^{t}f(t)dt\leq f(t).$
Hence $(f_{3})$ is satisfied.
$(f_{4})$: Since $l>n$, we get $F(t)/t^{n}\rightarrow 0$ as $t\rightarrow 0+$.
Hence $(f_{4})$ holds.
$(f_{5})$: Note that $t^{n/(n-1)}-t^{1/(n-1)}\geq t^{n/(n-1)}/2$ for all
$t\geq 2$. Let $A\geq 2$. Then for all $t\geq A$ there holds
$f(t)\geq
2^{l}l!C_{q}\frac{(t^{\frac{n}{n-1}}-\chi(t)t^{\frac{1}{n-1}})^{l}}{l!}\geq
2^{l}C_{q}(t^{\frac{n}{n-1}}/2)^{l}=C_{q}t^{q-1}.$
When $t\in[0,A]$, we get
$f(t)\geq 2^{l}l!C_{q}\frac{t^{\frac{nl}{n-1}}}{l!}=2^{l}C_{q}t^{q-1}.$
Hence $(f_{5})$ is satisfied. In short, $f(t)$ satisfies $(f_{1})$-$(f_{5})$
if $A\geq 4^{n-1}$.
Finally we check that $(H_{5})$ does not hold. When $t\geq 2A$, we have
$f(t)=2^{l}l!C_{q}\left(e^{t^{\frac{n}{n-1}}-t^{\frac{1}{n-1}}}-\sum_{k=0}^{l-1}\frac{(t^{\frac{n}{n-1}}-t^{\frac{1}{n-1}})^{k}}{k!}\right).$
It follows that
$\lim_{t\rightarrow+\infty}tf(t)e^{-t^{\frac{n}{n-1}}}=0.$
Thus $f(t)$ does not satisfy $(H_{5})$. $\hfill\Box$
Acknowledgements. This work was partly supported by the NSFC 11171347 and the
NCET program 2008-2011.
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|
arxiv-papers
| 2011-12-04T06:17:05 |
2024-09-04T02:49:24.937014
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Yunyan Yang",
"submitter": "Yunyan Yang",
"url": "https://arxiv.org/abs/1112.0724"
}
|
1112.0739
|
# On the UMD constants for a class of iterated $L_{p}(L_{q})$ spaces
Yanqi QIU Inst. Math. Jussieu, Équipe d’Analyse Fonctionnelle Université
Paris VI, 75252 Paris Cedex 05, France yanqi-qiu@math.jussieu.fr
###### Abstract.
Let $1<p\neq q<\infty$ and $(D,\mu)=(\\{\pm
1\\},\frac{1}{2}\delta_{-1}+\frac{1}{2}\delta_{1})$. Define by recursion:
$X_{0}=\mathbb{C}$ and $X_{n+1}=L_{p}(\mu;L_{q}(\mu;X_{n}))$. In this paper,
we show that there exist $c_{1}=c_{1}(p,q)>1$ depending only on $p,q$ and
$c_{2}=c_{2}(p,q,s)$ depending on $p,q,s$, such that the $\text{UMD}_{s}$
constants of $X_{n}$’s satisfy $c_{1}^{n}\leq C_{s}(X_{n})\leq c_{2}^{n}$ for
all $1<s<\infty$. Similar results will be showed for the analytic UMD
constants. We mention that the first super-reflexive non-UMD Banach lattices
were constructed by Bourgain. Our results yield another elementary
construction of super-reflexive non-UMD Banach lattices, i.e. the inductive
limit of $X_{n}$, which can be viewed as iterating infinitely many times
$L_{p}(L_{q})$.
###### Key words and phrases:
UMD property, analytic UMD property, iterated $L_{p}(L_{q})$ spaces, super-
reflexive non-UMD Banach lattices
The author was partially supported by ANR grant 2011 BS01 008 01
## 1\. Introduction
A Banach space $X$ is UMD if for all (or equivalently, for some) $1<s<\infty$
there is a constant $C>0$ depending only on $s$ and $X$ such that
(1)
$\displaystyle\sup_{\varepsilon_{k}\in\\{-1,1\\}}\|\sum_{k=0}^{n}\varepsilon_{k}df_{k}\|_{L_{s}(X)}\leq
C\|\sum_{k=0}^{n}df_{k}\|_{L_{s}(X)}$
for all $n\geq 0$ and all $X$-valued martingale difference sequences
$(df_{k})_{k=0}^{n}$. The best such $C$ is called the $\text{UMD}_{s}$
constant of $X$ and will be denoted by $C_{s}(X)$ in the sequel. It is well-
known that in the above definition, we can restrict to the dyadic martingale
differences and the best constant remains the same. The UMD property for
Banach spaces was introduced by Maurey and Pisier. The reader is refered to
Burkholder’s papers [5, 7] for the details of the UMD property.
Let $\mathbb{T}=\\{z\in\mathbb{C}:|z|=1\\}$ be the one dimensional torus
equipped with the normalised Haar measure $m$. Consider the canonical
filtration on the probability space
$(\mathbb{T}^{\mathbb{N}},m^{\otimes\mathbb{N}})$ defined by
$\sigma(z_{0})\subset\sigma(z_{0},z_{1})\subset\cdots\subset\sigma(z_{0},z_{1},\cdots,z_{n})\subset\cdots.$
By definition, a Hardy martingale in $L_{s}(\mathbb{T}^{\mathbb{N}};X)$ is a
martingale $f=(f_{n})_{n\geq 0}$ with respect to the canonical filtration such
that $\sup_{n}\|f_{n}\|_{L_{s}}<\infty$, and such that the martingale
difference $df_{n}=f_{n}-f_{n-1}$ (by convention, $df_{0}:=f_{0}$) is analytic
in the last variable $z_{n}$, i.e., $df_{n}$ has the form:
$df_{n}(z_{0},\cdots,z_{n-1},z_{n})=\sum_{k\geq
1}\phi_{n,k}(z_{0},\cdots,z_{n-1})z_{n}^{k}.$
In the above definition of UMD spaces, if the Banach space is over the complex
field $\mathbb{C}$, and if we restrict to the Hardy martingales, then a
different class of Banach spaces is defined, i.e. the analytic UMD class (AUMD
by abreviation). The best constant is called the $\text{AUMD}_{s}$ constant of
$X$ and will be denoted by $C_{s}^{a}(X)$. Note that UMD implies AUMD but not
conversely, for instance, $L_{1}(\mathbb{T},m)$ is an AUMD space which is not
UMD (cf. [9]).
It is well-known that UMD implies super-reflexivity but not conversely. The
first super-reflexive non-UMD Banach space was constructed by Pisier in [11].
Super-reflexive non-UMD Banach lattices were later constructed by Bourgain in
[2, 3]. We refer to Rubio de Francia’s paper [13] for some open problems
related to the super-reflexive non-UMD Banach lattices.
The main topic of this paper is the investigation of the UMD constants of a
family of iterated $L_{p}(L_{q})$-spaces. As a consequence of our results, we
give an elementary construction of super-reflexive non-UMD Banach lattices.
## 2\. Some elementary inequalities
We will use the following lemma.
###### Lemma 2.1.
Let $(\Omega,\nu)$ be a measure space such that $\nu$ is finite. Suppose that
$\alpha\neq 1$ and $0<\alpha<\infty$. If $F,f\in L_{\alpha}(\Omega,\nu)\bigcap
L_{1}(\Omega,\nu)$ satisfy
$\int(|F|+|g|)^{\alpha}d\nu\leq\int(|f|+|g|)^{\alpha}d\nu$
for all $g\in L_{\infty}(\Omega,\nu)$. Then $|F|\leq|f|$ a.e..
###### Proof.
Consider first those $g\in L_{\infty}(\Omega,\nu)$ such that there exists
$\delta>0$ and $|g|\geq\delta$ a.e.. If $F,f$ satisfy the condition in the
statement, then for all $\varepsilon>0$, we have
(2)
$\displaystyle\int(\varepsilon|F|+|g|)^{\alpha}d\nu\leq\int(\varepsilon|f|+|g|)^{\alpha}d\nu.$
By the mean value theorem, there exists $\theta=\theta_{\varepsilon}\in(0,1)$,
such that
$\frac{(\varepsilon|f|+|g|)^{\alpha}-|g|^{\alpha}}{\varepsilon}=\alpha(\theta\varepsilon|f|+|g|)^{\alpha-1}|f|.$
If $\alpha<1$, then
$(\theta\varepsilon|f|+|g|)^{\alpha-1}|f|\leq|g|^{\alpha-1}|f|\in
L_{1}(\Omega,\nu)$ and if $\alpha>1$, then for $0<\varepsilon<1$, we have
$0<\theta\varepsilon<1$ and hence
$(\theta\varepsilon|f|+|g|)^{\alpha-1}|f|\leq
2^{\alpha-1}(|f|^{\alpha}+|g|^{\alpha-1}|f|)\in L_{1}(\Omega,\nu)$. By the
dominated convergence theorem, we have
$\lim_{\varepsilon\to
0^{+}}\frac{\int(\varepsilon|f|+|g|)^{\alpha}d\nu-\int|g|^{\alpha}d\nu}{\varepsilon}=\alpha\int|f||g|^{\alpha-1}d\nu.$
The same equality holds for $F$. Combining this with (2), we get
$\int|F||g|^{\alpha-1}d\nu\leq\int|f||g|^{\alpha-1}d\nu.$
Replacing $g$ by $|g|^{\frac{1}{\alpha-1}}$ yields
$\int|F||g|d\nu\leq\int|f||g|d\nu.$
By approximation, the above inequality holds for all $g\in
L_{\infty}(\Omega,\nu)$. Hence $|F|\leq|f|$ a.e., as announced. ∎
###### Proposition 2.2.
Let $(\Omega,\nu)$ be a measure space such that $\nu$ is finite. Suppose that
$1\leq p\neq q<\infty$. If $F,f\in L_{p}(\Omega,\nu)\bigcap L_{q}(\Omega,\nu)$
satisfy
$\int(|F|^{q}+|g|^{q})^{p/q}d\nu\leq\int(|f|^{q}+|g|^{q})^{p/q}d\nu$
for all $g\in L_{\infty}(\Omega,\nu)$. Then $|F|\leq|f|$ a.e..
###### Proof.
This is just a reformulation of Lemma 2.1. ∎
Let $D=\\{-1,1\\}$ be the Bernoulli probability space equipped with the
measure $\mu=\frac{1}{2}\delta_{-1}+\frac{1}{2}\delta_{1}$. For any $1\leq
q\leq\infty$, the 2-dimensional $\ell_{q}$-space will be denoted by
$\ell_{q}^{(2)}$.
###### Proposition 2.3.
Suppose that $1\leq p\neq q\leq\infty$. Let $P$ be the projection on
$L_{p}(\mu;\ell_{q}^{(2)})$ defined by
$\begin{array}[]{cccc}P:&L_{p}(\mu;\ell_{q}^{(2)})&\rightarrow&L_{p}(\mu;\ell_{q}^{(2)})\\\
&(f,g)&\mapsto&(\mathbb{E}f,g)\end{array},$
where $\mathbb{E}$ is the expectation. Then $P$ is not contractive.
###### Proof.
Assume first that both $p,q$ are finite. If $P$ is contractive, then for any
two functions $f$ and $g$, we have
$\int(|\mathbb{E}f|^{q}+|g|^{q})^{p/q}d\mu\leq\int(|f|^{q}+|g|^{q})^{p/q}d\mu.$
By Proposition 2.2, it follows that $|\mathbb{E}(f)|\leq|f|$, which is a
contradiction, hence $P$ is not contractive.
If $p=\infty$ and $1<q<\infty$, then $p^{\prime}=1$ and $1<q^{\prime}<\infty$.
Since the adjoint map $P^{*}$ on $L_{1}(\mu;\ell_{q^{\prime}}^{(2)})$ has the
same form as $P$, the preceding argument shows that $P^{*}$ and hence $P$ is
not contractive.
If $p=\infty$ and $q=1$. Assume $P$ is contractive, then we have
(3)
$\displaystyle\big{\|}|\mathbb{E}f|+|g|\big{\|}_{\infty}\leq\big{\|}|f|+|g|\big{\|}_{\infty}.$
Consider $f=1+\varepsilon,g=1-\varepsilon$, where $\varepsilon:D\rightarrow D$
is the identity function. Then the left hand side of (3) equals to 3 while the
right hand side equals to 2. This contradiction shows that $P$ is not
contractive.
If $1\leq p<\infty$ and $q=\infty$, then $1<p^{\prime}\leq\infty$ and
$q^{\prime}=1$, hence $P^{*}$ is not contractive. It follows that $P$ is not
contractive. ∎
The norm of $P$ on $L_{p}(\mu;\ell_{q}^{(2)})$ will be denoted by $c(p,q)$ in
the sequel. If $p=q$, then $c(p,p)=1$. If $1\leq p\neq q\leq\infty$, then
(4) $\displaystyle c(p,q)>1.$
###### Remark 2.4.
It is not difficult to check that $c(\infty,1)=c(1,\infty)=\frac{3}{2}$. But
we do not know the exact value of $c(p,q)$ for general $p\neq q$.
As usual, we set
$H_{p}(\mathbb{T})=\\{f\in L_{p}(\mathbb{T},m):\hat{f}(k)=0,\forall
k\in\mathbb{Z}_{<0}\\}.$
We will say that a measurable function $f:\mathbb{T}\rightarrow\mathbb{C}$ is
bounded from below, if there exists $\delta>0$, such that $|f|\geq\delta$ a.e.
on $\mathbb{T}$. If $f\in L_{p}(\mathbb{T})$ is bounded from below, then the
geometric mean $M(|f|)$ of $|f|$ is defined by
$\log M(|f|)=\int_{\mathbb{T}}\log|f(z)|dm(z).$
In particular, if $f:\mathbb{D}\rightarrow\mathbb{C}$ is an outer function,
then
(5) $\displaystyle M(|f|)=|f(0)|=|\mathbb{E}f|.$
The following elementary proposition will be used in §4 when we treat the
analytic UMD property.
###### Proposition 2.5.
Suppose that $1\leq p\neq q<\infty$. Define $\kappa(p,q)$ to be the best
constant $C$ satisfying the property: For any measurable partition
$\mathbb{T}=A\dot{\cup}B$ with $m(A)=m(B)=\frac{1}{2}$, for any function
$f=f_{1}\chi_{A}+f_{2}\chi_{B}$ with $f_{1}>0,f_{2}>0$ and any function
$g=g_{1}\chi_{A}+g_{2}\chi_{B}$, we have
$\int_{\mathbb{T}}(M(|f|)^{q}+|g|^{q})^{p/q}dm\leq
C^{p}\int_{\mathbb{T}}(|f|^{q}+|g|^{q})^{p/q}dm.$
Then $\kappa(p,q)>1$.
###### Proof.
Assume $k(p,q)\leq 1$. Fix any measurable partition $\mathbb{T}=A\dot{\cup}B$
such that $m(A)=m(B)=\frac{1}{2}$. Consider the 2-valued functions
$f=f_{1}\chi_{A}+f_{2}\chi_{B}$ and $g=g_{1}\chi_{A}+g_{2}\chi_{B}$ with
$f_{1},f_{2}$ positive scalars. By Proposition 2.2, $M(f)\leq f$. However, one
can easily check that $M(f)=f_{1}^{1/2}f_{2}^{1/2}$. If $f_{1}>f_{2}$, then
$M(f)>f_{2}^{1/2}f_{2}^{1/2}=f_{2}$, which contradicts to $M(f)\leq f$. Whence
the announced statement. ∎
## 3\. UMD constants of iterated $L_{p}(L_{q})$ spaces
The following definition is essential in the sequel.
###### Definition 3.1.
Consider a Banach space $X$ with a fixed family of vectors $\\{x_{i}\\}_{i\in
I}$. We define $S(X;\\{x_{i}\\})$ to be the best constant $C$ such that
(6)
$\displaystyle\Big{\|}\sum_{k=0}^{N}\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{1}(\Omega,\mathbb{P};X)}\leq
C\Big{\|}\sum_{k=0}^{N}\theta_{k}x_{i_{k}}\Big{\|}_{L_{\infty}(\Omega,\mathbb{P};X)}$
holds for any $N\in\mathbb{N}$, any probability space
$(\Omega,\mathcal{F},\mathbb{P})$ equipped with a filtration
$\mathcal{A}_{0}\subset\mathcal{A}_{1}\subset\cdots\subset\mathcal{A}_{n}\subset\cdots\subset\mathcal{F}$,
any $N+1$ distinct indices $\\{i_{0},i_{1},\cdots,i_{N}\\}\subset I$ and any
$N+1$ functions $\theta_{0},\theta_{1},\cdots,\theta_{N}$ in
$L_{\infty}(\Omega,\mathcal{F},\mathbb{P})$.
If there does not exist such constant, we set $S(X;\\{x_{i}\\})=\infty$.
In what follows, we are mostly interested in the special case when
$\\{x_{i}\\}$ is a 1-unconditional basic sequence, since in this case we can
relate $S(X;\\{x_{i}\\})$ to the UMD constants of $X$. If $\\{x_{i}\\}$ is
clear from the context and there is no confusion, we will use the simplified
notation $S(X)$ for $S(X;\\{x_{i}\\})$. In particular, if $X$ has a natural
basis, then $S(X)$ will always mean to be calculated with this basis.
We will need the following well-known Stein inequality in UMD spaces, which
was originally proved by Bourgain [4]. For the sake of completeness, we
include the proof.
###### Theorem 3.2.
Let $X$ be a UMD space. Then for any $1<s<\infty$, any finite sequences of
functions $(F_{k})_{k\geq 0}$ in $L_{s}(\Omega,\mathbb{P};X)$ and any
filtration
$\mathcal{A}_{0}\subset\mathcal{A}_{1}\subset\cdots\subset\mathcal{A}_{n}\subset\cdots$
on $(\Omega,\mathbb{P})$, we have
(7)
$\displaystyle\Big{\|}\sum_{k}\varepsilon_{k}\mathbb{E}_{k}(F_{k})\Big{\|}_{L_{s}(\mu_{\infty}\times\mathbb{P};X)}\leq
C_{s}(X)\Big{\|}\sum_{k}\varepsilon_{k}F_{k}\Big{\|}_{L_{s}(\mu_{\infty}\times\mathbb{P};X)},$
where $\mathbb{E}_{k}=\mathbb{E}^{\mathcal{A}_{k}}$ and
$(\varepsilon_{k})_{k\geq 0}$ is the usual Rademacher sequence on
$(D^{\mathbb{N}},\mu_{\infty})$, $\mu_{\infty}=\mu^{\otimes\mathbb{N}}$.
###### Proof.
Let $f=\sum_{k}\varepsilon_{k}F_{k}$ and
$f^{\prime}=\sum_{k}\varepsilon_{k}\mathbb{E}_{k}(F_{k})$. Then if
$\mathcal{C}_{2j}=\mathcal{A}_{j}\otimes\sigma(\varepsilon_{0},\cdots,\varepsilon_{j})$
and
$\mathcal{C}_{2j-1}=\mathcal{A}_{j}\otimes\sigma(\varepsilon_{0},\cdots,\varepsilon_{j-1})$,
we have
$f^{\prime}=\sum_{j}(\mathbb{E}^{\mathcal{C}_{2j}}-\mathbb{E}^{\mathcal{C}_{2j-1}})(f).$
Indeed,
$\mathbb{E}^{\mathcal{C}_{2j}}(f)=\sum_{0}^{j}\varepsilon_{k}\mathbb{E}_{j}(F_{k})$
and
$\mathbb{E}^{\mathcal{C}_{2j-1}}(f)=\sum_{0}^{j-1}\varepsilon_{k}\mathbb{E}_{j}(F_{k})$.
Hence
$(\mathbb{E}^{\mathcal{C}_{2j}}-\mathbb{E}^{\mathcal{C}_{2j-1}})(f)=\varepsilon_{j}\mathbb{E}_{j}(F_{j})$.
It follows (see the next remark) that
$\|f^{\prime}\|_{L_{s}(\mu_{\infty}\times\mathbb{P};X)}\leq
C_{s}(X)\|f\|_{L_{s}(\mu_{\infty}\times\mathbb{P};X)},$
whence (7). ∎
###### Remark 3.3.
By an extreme point argument, we have
$\sup_{-1\leq\alpha_{k}\leq
1}\|\sum_{k=0}^{n}\alpha_{k}df_{k}\|_{L_{s}(X)}=\sup_{\varepsilon_{k}\in\\{-1,1\\}}\|\sum_{k=0}^{n}\varepsilon_{k}df_{k}\|_{L_{s}(X)}.$
Hence we have
$\sup_{-1\leq\alpha_{k}\leq
1}\|\sum_{k=0}^{n}\alpha_{k}df_{k}\|_{L_{s}(X)}\leq
C_{s}(X)\|\sum_{k=0}^{n}df_{k}\|_{L_{s}(X)}.$
###### Proposition 3.4.
Let $X$ be a UMD space. Assume that $\\{x_{i}\\}_{i\in I}$ is a
1-unconditional basic sequence in $X$. Then for any $1<s<\infty$, any finite
sequence of functions $(\theta_{k})_{k\geq 0}$ in $L_{s}(\Omega,\mathbb{P})$
and any filtration
$\mathcal{A}_{0}\subset\mathcal{A}_{1}\subset\cdots\subset\mathcal{A}_{n}\subset\cdots$
on $(\Omega,\mathbb{P})$, we have
(8)
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}_{k}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{s}(\Omega,\mathbb{P};X)}\leq
C_{s}(X)\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{s}(\Omega,\mathbb{P};X)}.$
###### Proof.
For any $i_{k}$’s, consider the sequence $(F_{k})_{k\geq 0}$ in
$L_{s}(\Omega,\mathbb{P};X)$ defined by $F_{k}(w)=\theta_{k}(w)x_{i_{k}}$.
Then $\mathbb{E}_{k}(F_{k})=\mathbb{E}_{k}(\theta_{k})x_{i_{k}}$. By the
1-unconditionality of $\\{x_{i}\\}_{i\in I}$, for any fixed choice of signs
$\varepsilon_{k}\in\\{-1,1\\}$ and $w\in\Omega$, we have
$\Big{\|}\sum_{k}\varepsilon_{k}F_{k}(w)\Big{\|}_{X}=\Big{\|}\sum_{k}\varepsilon_{k}\theta_{k}(w)x_{i_{k}}\Big{\|}_{X}=\Big{\|}\sum_{k}\theta_{k}(w)x_{i_{k}}\Big{\|}_{X}.$
It follows that
$\Big{\|}\sum_{k}\varepsilon_{k}F_{k}\Big{\|}_{L_{s}(\mu_{\infty}\times\mathbb{P};X)}=\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{s}(\Omega,\mathbb{P};X)}.$
Similarly, we have
$\Big{\|}\sum_{k}\varepsilon_{k}\mathbb{E}_{k}(F_{k})\Big{\|}_{L_{s}(\mu_{\infty}\times\mathbb{P};X)}=\Big{\|}\sum_{k}\mathbb{E}_{k}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{s}(\Omega,\mathbb{P};X)}.$
By these equalities, (8) follows from (7). ∎
Let $X$ be as in Proposition 3.4, $\\{x_{i}\\}_{i\in I}$ is a 1-unconditional
basic sequence in $X$. Assume that $\theta_{k}\in
L_{\infty}(\Omega,\mathbb{P})$. By an application of the contractive
inclusions $L_{\infty}(\Omega,\mathbb{P};X)\subset
L_{s}(\Omega,\mathbb{P};X)\subset L_{1}(\Omega,\mathbb{P};X)$, we have
(9)
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}_{k}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{1}(\Omega,\mathbb{P};X)}\leq
C_{s}(X)\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{\infty}(\Omega,\mathbb{P};X)}.$
Hence
(10) $\displaystyle S(X;\\{x_{i}\\})\leq C_{s}(X)$
for all $1<s<\infty$.
###### Theorem 3.5.
Let $E$ be a Banach space with a 1-unconditional basis $\\{e_{i}:i\in I\\}$,
let $F$ be another Banach space. By definition, $E(F)$ is the completion of
the algebraic tensor product $E\otimes F$ under the norm defined as follows:
if $x=\sum_{i}e_{i}\otimes x_{i}\in E\otimes F$, where $(x_{i})$ is a finite
supported sequence in $F$, then
$\|x\|_{E(F)}:=\Big{\|}\sum_{i}e_{i}\big{\|}x_{i}\big{\|}_{F}\Big{\|}_{E}.$
For any fixed family of vectors $\\{f_{j}:j\in J\\}$ in $F$, consider the
family of vectors $\\{e_{i}\otimes f_{j}:i\in I,j\in J\\}$. Then we have
$S(E(F))\geq S(E)S(F),$
where $S(E(F))$, $S(E)$ and $S(F)$ are defined with respect to the mentioned
families of vectors respectively.
###### Proof.
From the definition, for any $\varepsilon>0$, there exist finite number of
distinct indices $\\{i_{k}:1\leq k\leq N_{1}\\}\subset I$ and $\\{j_{n}:1\leq
n\leq N_{2}\\}\subset J$, and there exist functions $\theta_{k}\in
L_{\infty}(\Omega^{\prime},\mathbb{P}^{\prime}),1\leq k\leq N_{1}$ and
functions $\xi_{n}\in L_{\infty}(\Omega_{0},\mathbb{P}_{0}),1\leq n\leq N_{2}$
satisfying
$\|\sum_{k}\theta_{k}e_{i_{k}}\|_{L_{\infty}(\Omega^{\prime},\mathbb{P}^{\prime};E)}\leq
1$
and
$\|\sum_{n}\xi_{n}f_{j_{n}}\|_{L_{\infty}(\Omega_{0},\mathbb{P}_{0};F)}\leq 1$
such that
$\Big{\|}\sum_{k}\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})e_{i_{k}}\Big{\|}_{L_{1}(\Omega^{\prime},\mathbb{P}^{\prime};E)}\geq
S(E)-\varepsilon$
and
$\Big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})f_{j_{n}}\Big{\|}_{L_{1}(\Omega_{0},\mathbb{P}_{0};F)}\geq
S(F)-\varepsilon.$
Let
$(\Omega,\mathbb{P})=(\Omega^{\prime}\times\Omega_{0}^{\mathbb{N}},\mathbb{P}^{\prime}\otimes\mathbb{P}_{0}^{\otimes\mathbb{N}})$,
the general element in $\Omega$ will be denoted by
$w=(w^{\prime},(w_{l})_{l\geq 0})$. Consider the $\sigma$-algebras
$\mathcal{F}_{k,n}$ defined on $(\Omega,\mathbb{P})$ by
$\mathcal{F}_{k,n}:=\mathcal{A}_{k}\otimes\underbrace{\mathcal{B}_{\infty}\otimes\cdots\otimes\mathcal{B}_{\infty}}_{k-1\text{
times }}\otimes\mathcal{B}_{n}\otimes\mathcal{C}_{\geq k+1},$
where $\mathcal{B}_{\infty}=\sigma(\mathcal{B}_{n}:n\geq 0)$ is a
$\sigma$-algebra on $(\Omega_{0},\mathbb{P}_{0})$, $\mathcal{B}_{0}$ is
assumed to be trivial and $\mathcal{C}_{\geq k+1}$ is the trivial
$\sigma$-algebra on
$(\Omega_{0}^{\mathbb{N}_{\geq{k+1}}},\mathbb{P}_{0}^{\mathbb{N}_{\geq
k+1}})$. It is easy to check that $\mathcal{F}_{k,n}$ is a filtration with
respect to the lexigraphic order, i.e. if $(k,n)<(k^{\prime},n^{\prime})$
(that is $k<k^{\prime}$ or $k=k^{\prime}$ but $n<n^{\prime}$), then
$\mathcal{F}_{k,n}\subset\mathcal{F}_{k^{\prime},n^{\prime}}$.
Now let us define $h:\Omega\rightarrow E(F)$ by
$h(w)=h(w^{\prime},(w_{l}))=\sum_{k,n}\theta_{k}(w^{\prime})\xi_{n}(w_{k})e_{i_{k}}\otimes
f_{j_{n}}.$
Let $h_{k,n}(w)=\theta_{k}(w^{\prime})\xi_{n}(w_{k})$, then
$h=\sum_{k,n}h_{k,n}e_{i_{k}}\otimes f_{j_{n}}$. Clearly, we have
(11)
$\displaystyle\mathbb{E}^{\mathcal{F}_{k,n}}(h_{k,n})(w)=\Big{[}\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})\Big{]}(w^{\prime})\Big{[}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})\Big{]}(w_{k})\quad
a.e..$
By the 1-unconditionality of $\\{e_{i}:i\in I\\}$, for a.e. $w\in\Omega$, we
have
$\displaystyle\|h(w)\|_{E(F)}$ $\displaystyle=$
$\displaystyle\Big{\|}\sum_{k,n}\theta_{k}(w^{\prime})\xi_{n}(w_{k})e_{i_{k}}\otimes
f_{j_{n}}\Big{\|}_{E(F)}$ $\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}\big{\|}\sum_{n}\theta_{k}(w^{\prime})\xi_{n}(w_{k})f_{j_{n}}\big{\|}_{F}\Big{\|}_{E}$
$\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}|\theta_{k}(w^{\prime})|\big{\|}\sum_{n}\xi_{n}(w_{k})f_{j_{n}}\big{\|}_{F}\Big{\|}_{E}$
$\displaystyle\leq$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}|\theta_{k}(w^{\prime})|\Big{\|}_{E}=\Big{\|}\sum_{k}e_{i_{k}}\theta_{k}(w^{\prime})\Big{\|}_{E}\leq
1.$
Hence $\|h\|_{L_{\infty}(\Omega,\mathbb{P};E(F))}\leq 1$. If we denote
$\widetilde{h}=\sum_{k,n}\mathbb{E}^{\mathcal{F}_{k,n}}(h_{k,n})e_{i_{k}}\otimes
f_{j_{n}},$
then by (11),
$\displaystyle\|\widetilde{h}(w)\|_{E(F)}$ $\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}|\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})(w^{\prime})|\big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})(w_{k})f_{j_{n}}\big{\|}_{F}\Big{\|}_{E}.$
By Jensen’s inequality, we have
$\displaystyle\int\Big{\|}\sum_{k}e_{i_{k}}|\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})(w^{\prime})|\big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})(w_{k})f_{j_{n}}\big{\|}_{F}\Big{\|}_{E}d\mathbb{P}_{0}^{\otimes\mathbb{N}}((w_{l}))$
$\displaystyle\geq$
$\displaystyle\Big{\|}\int\sum_{k}e_{i_{k}}|\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})(w^{\prime})|\big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})(w_{k})f_{j_{n}}\big{\|}_{F}d\mathbb{P}_{0}^{\otimes\mathbb{N}}((w_{l}))\Big{\|}_{E}$
$\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}|\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})(w^{\prime})|\Big{\|}_{E}\cdot\Big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})f_{j_{n}}\Big{\|}_{L_{1}(\Omega_{0},\mathbb{P}_{0};F)}$
$\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}e_{i_{k}}\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})(w^{\prime})\Big{\|}_{E}\cdot\Big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})f_{j_{n}}\Big{\|}_{L_{1}(\Omega_{0},\mathbb{P}_{0};F)}.$
Note that in the last equality, we used the 1-unconditionality assumption on
$\\{e_{i}:i\in I\\}$. By integrating both sides with respect to $\int
d\mathbb{P}^{\prime}(w^{\prime})$, we get
$\displaystyle\Big{\|}\sum_{k,n}\mathbb{E}^{\mathcal{F}_{k,n}}(h_{k,n})e_{i_{k}}\otimes
f_{j_{n}}\Big{\|}_{L_{1}(\Omega,\mathbb{P};E(F))}$ $\displaystyle\geq$
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}^{\mathcal{A}_{k}}(\theta_{k})e_{i_{k}}\Big{\|}_{L_{1}(\Omega^{\prime},\mathbb{P}^{\prime};E)}\cdot\Big{\|}\sum_{n}\mathbb{E}^{\mathcal{B}_{n}}(\xi_{n})f_{j_{n}}\Big{\|}_{L_{1}(\Omega_{0},\mathbb{P}_{0};F)}$
$\displaystyle\geq$ $\displaystyle(S(E)-\varepsilon)(S(F)-\varepsilon).$
Therefore $S(E(F))\geq(S(E)-\varepsilon)(S(F)-\varepsilon)$. Since
$\varepsilon>0$ is arbitrary, it follows that $S(E(F))\geq S(E)S(F)$ as
desired. ∎
###### Remark 3.6.
If $E$ is a Banach lattice which is $p$-convex and $q$-concave $($see [10] for
the details$)$ with $1\leq p\leq q\leq\infty$ and $F$ is a Banach space. Then
the preceding proof is valid with $S_{q,p}(E)$ and $S_{q,p}(F)$ defined using
(6) with $L_{p}$-norm on the left hand side and $L_{q}$-norm on the right hand
side.
###### Remark 3.7.
Let $1\leq p<q\leq\infty$. If we define $C_{q,p}(X)$ as the best constant $C$
in (1) with $L_{p}$-norm on the left hand side and $L_{q}$-norm on the right
hand side, it is well-known that $X$ is in the UMD class if and only if
$C_{q,p}(X)<\infty$. The preceding argument shows that under the same
assumption of Theorem 3.5, we have $C_{\infty,1}(E(F))\geq
S(E)C_{\infty,1}(F)$. Moreover, if $E$ is $p$-convex and $q$-concave we have
$C_{q,p}(E(F))\geq S_{q,p}(E)C_{q,p}(F)$.
###### Lemma 3.8.
Suppose that $1\leq p\neq q\leq\infty$. If
$E_{1}=\ell_{p}^{(2)}(\ell_{q}^{(2)})$, then
$S(E_{1})\geq c(p,q)>1.$
###### Proof.
Denote by $\\{e_{1}^{p},e_{2}^{p}\\}$, $\\{e_{1}^{q},e_{2}^{q}\\}$ the
canonical basis of $\ell_{p}^{(2)}$ and $\ell_{q}^{(2)}$ respectively .Then
$\\{e_{1}^{p}\otimes e_{1}^{q},e_{1}^{p}\otimes e_{2}^{q},e_{2}^{p}\otimes
e_{1}^{q},e_{2}^{p}\otimes e_{2}^{q}\\}$ is the canonical 1-unconditional
basis of $\ell_{p}^{(2)}(\ell_{q}^{(2)})$. Consider the probability space
$(D,\mu)$ equipped with the filtration
$\\{\phi,D\\}\subset\sigma(\varepsilon)$, where $\varepsilon$ is the identity
function on $D$. Define a linear map $T:L_{\infty}(D;E_{1})\rightarrow
L_{1}(D;E_{1})$ by setting
$T\Big{[}a_{ij}(\varepsilon)e_{i}^{p}\otimes
e_{j}^{q}\Big{]}=\left\\{\begin{array}[]{lc}\mathbb{E}(a_{ij})e_{i}^{p}\otimes
e_{j}^{q},\text{ if }j=1\\\ a_{ij}(\varepsilon)e_{i}^{p}\otimes
e_{j}^{q},\text{ if }j=2\end{array}.\right.$
By definition of $S(E_{1})$ we have
$S(E_{1})\geq\|T\|_{L_{\infty}(D;E_{1})\rightarrow L_{1}(D;E_{1})}$. Now for
any $a,b$ two scalar functions on $D$ , consider
$f(\varepsilon)=e_{1}^{p}\otimes\Big{[}a(\varepsilon)e_{1}^{q}+b(\varepsilon)e_{2}^{q}\Big{]}+e_{2}^{p}\otimes\Big{[}a(-\varepsilon)e_{1}^{q}+b(-\varepsilon)e_{2}^{q}\Big{]}.$
Then
$(Tf)(\varepsilon)=e_{1}^{p}\otimes\Big{[}\mathbb{E}(a)e_{1}^{q}+b(\varepsilon)e_{2}^{q}\Big{]}+e_{2}^{p}\otimes\Big{[}\mathbb{E}(a)e_{1}^{q}+b(-\varepsilon)e_{2}^{q}\Big{]}.$
If $p,q$ are both finite, then for any fixed $\varepsilon\in D$, we have
$\displaystyle\|f(\varepsilon)\|_{E_{1}}\\!\\!\\!$ $\displaystyle=$
$\displaystyle\\!\\!\Big{\\{}(|a(\varepsilon)|^{q}+|b(\varepsilon)|^{q})^{p/q}+(|a(-\varepsilon)|^{q}+|b(-\varepsilon)|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$
$\displaystyle\\!\\!\Big{\\{}(|a(1)|^{q}+|b(1)|^{q})^{p/q}+(|a(-1)|^{q}+|b(-1)|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$
$\displaystyle\\!\\!2^{1/p}\Big{\\{}\frac{1}{2}(|a(1)|^{q}+|b(1)|^{q})^{p/q}+\frac{1}{2}(|a(-1)|^{q}+|b(-1)|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$
$\displaystyle\\!\\!2^{1/p}\Big{\\{}\int(|a(\varepsilon)|^{q}+|b(\varepsilon)|^{q})^{p/q}d\mu(\varepsilon)\Big{\\}}^{1/p}$
$\displaystyle=$
$\displaystyle\\!\\!2^{1/p}\big{\|}(a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}.$
Similarly,
$\|(Tf)(\varepsilon)\|_{E_{1}}=2^{1/p}\big{\|}(\mathbb{E}a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}.$
It follows that
$\|f\|_{L_{\infty}(D;E_{1})}=2^{1/p}\big{\|}(a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}$
and
$\|Tf\|_{L_{1}(D;E_{1})}=2^{1/p}\big{\|}(\mathbb{E}a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}.$
Hence
(12) $\displaystyle\|T\|_{L_{\infty}(D;E_{1})\rightarrow
L_{1}(D;E_{1})}\geq\frac{\|Tf\|_{L_{1}(D;E_{1})}}{\|f\|_{L_{\infty}(D;E_{1})}}=\frac{\big{\|}(\mathbb{E}a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}}{\big{\|}(a,b)\big{\|}_{L_{p}(\mu;\ell_{q}^{(2)})}}.$
Similarly, if $q=\infty$ and $p$ is finite, then
$\|f\|_{L_{\infty}(D;E_{1})}=2^{1/p}\|(a,b)\|_{L_{p}(\mu;\ell_{\infty}^{(2)})}$
and
$\|Tf\|_{L_{1}(D;E_{1})}=2^{1/p}\big{\|}(\mathbb{E}a,b)\big{\|}_{L_{p}(\mu;\ell_{\infty}^{(2)})}.$
If $p=\infty$ and $q$ is finite, then
$\|f\|_{L_{\infty}(D;E_{1})}=\|(a,b)\|_{L_{\infty}(\mu;\ell_{q}^{(2)})}$ and
$\|Tf\|_{L_{1}(D;E_{1})}=\big{\|}(\mathbb{E}a,b)\big{\|}_{L_{\infty}(\mu;\ell_{q}^{(2)})}$.
Therefore, (12) holds in full generality. By Proposition 2.3, we have
$\|T\|_{L_{\infty}(D;E_{1})\rightarrow L_{1}(D;E_{1})}\geq\|P\|=c(p,q).$
Hence $S(E_{1})\geq c(p,q)>1$, as announced. ∎
###### Remark 3.9.
Let $(e_{k})_{k\geq 0}$ be the canonical basis of
$\ell_{p}=\ell_{p}(\mathbb{N})$, then $S(\ell_{p})=1$. Indeed, if
$(\theta_{k})_{k\geq 0}$ is a finite sequence of functions, then
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}_{k}(\theta_{k})e_{k}\Big{\|}_{L_{1}(\ell_{p})}$
$\displaystyle\leq$
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}_{k}(\theta_{k})e_{k}\Big{\|}_{L_{p}(\ell_{p})}=\Big{\|}(\sum_{k}|\mathbb{E}_{k}(\theta_{k})|^{p})^{1/p}\Big{\|}_{L_{p}}$
$\displaystyle=$
$\displaystyle\Big{\|}\sum_{k}|\mathbb{E}_{k}(\theta_{k})|^{p}\Big{\|}_{L_{1}}^{1/p}=(\sum_{k}\big{\|}\mathbb{E}_{k}(\theta_{k})\big{\|}_{p}^{p})^{1/p}$
$\displaystyle\leq$
$\displaystyle(\sum_{k}\big{\|}\theta_{k}\big{\|}_{p}^{p})^{1/p}=\Big{\|}\sum_{k}\theta_{k}e_{k}\Big{\|}_{L_{p}(\ell_{p})}$
$\displaystyle\leq$
$\displaystyle\Big{\|}\sum_{k}\theta_{k}e_{k}\Big{\|}_{L_{\infty}(\ell_{p})}.$
###### Theorem 3.10.
Suppose that $1\leq p,q\leq\infty$. Let $E_{1}=\ell_{p}^{(2)}(\ell_{q}^{(2)})$
and define by recursion: $E_{n+1}=\ell_{p}^{(2)}(\ell_{q}^{(2)}(E_{n}))$. Then
for any $1<s<\infty$, we have
$C_{s}(E_{n})\geq S(E_{n})\geq c(p,q)^{n},$
where $S(E_{n})$ is computed with respect to the canonical basis of $E_{n}$.
In particular, if $p\neq q$, then $C_{s}(E_{n})$ has at least an exponential
growth with respect to $n$.
###### Proof.
By Theorem 3.5,
$S(E_{n+1})\geq S(\ell_{p}^{(2)}(\ell_{q}^{(2)}))S(E_{n}).$
By Lemma 3.8, we have $S(E_{n+1})\geq c(p,q)S(E_{n})$. It follows that
$S(E_{n})\geq c(p,q)^{n}$. Since the canonical basis of $E_{n}$ is
1-unconditional, by (10), for any $1<s<\infty$, we have $C_{s}(E_{n})\geq
S(E_{n})$. ∎
The following simple observation shows that the exponential growth of
$C_{s}(E_{n})$ is optimal.
###### Proposition 3.11.
Suppose $1<p\neq q<\infty$. Let $X$ be a Banach space. Define by recursion:
$Y_{0}=X$ and $Y_{n+1}=L_{p}(\mathbb{T};L_{q}(\mathbb{T};Y_{n}))$. Then for
all $1<s<\infty$, there exists $\chi=\chi(p,q,s)$, such that
$C_{s}(Y_{n})\leq\chi^{n}C_{s}(X).$
###### Proof.
We will use the following well-known fact (see e.g. [5, 6]) about UMD
constants: for any $1<r,s<\infty$, there exist $\alpha(r,s)$ and $\beta(r,s)$
such that for all Banach space $X$,
(13) $\displaystyle\alpha(r,s)C_{s}(X)\leq C_{r}(X)\leq\beta(r,s)C_{s}(X).$
We will also use the elementary identity $C_{s}(L_{s}(X))=C_{s}(X)$. Combining
these, we have
$\displaystyle C_{s}(Y_{n+1})$ $\displaystyle=$ $\displaystyle
C_{s}(L_{p}(L_{q}(Y_{n})))\leq\beta(s,p)C_{p}(L_{p}(L_{q}(Y_{n})))$
$\displaystyle=$
$\displaystyle\beta(s,p)C_{p}(L_{q}(Y_{n}))\leq\beta(s,p)\beta(p,q)C_{q}(L_{q}(Y_{n}))$
$\displaystyle=$
$\displaystyle\beta(s,p)\beta(p,q)C_{q}(Y_{n})\leq\beta(s,p)\beta(p,q)\beta(q,s)C_{s}(Y_{n}).$
Let $\chi=\beta(s,p)\beta(p,q)\beta(q,s)$, then
$C_{s}(E_{n})\leq\chi^{n}C_{s}(X).$ ∎
###### Remark 3.12.
Even if one of $p,q$ is infinite or equals to 1, then since
$\dim(E_{n})=4^{n}$, we have $C_{s}(E_{n})\lesssim\sqrt{\dim E_{n}}=2^{n}$.
Indeed, the Banach-Mazur distance between $E_{n}$ and $\ell_{2}^{\dim E_{n}}$
is $\leq\sqrt{\dim E_{n}}$ $($cf. e.g. [14]$)$.
## 4\. Analytic UMD constants
The main idea in §3 can be easily adapted for treating the analytic UMD
property. In this section, all spaces are over $\mathbb{C}$.
Denote the general element in $\mathbb{T}^{\mathbb{N}}$ be $z=(z_{n})_{n\geq
0}$ and let $m_{\infty}=m^{\otimes\mathbb{N}}$ be the Haar measure on
$\mathbb{T}^{\mathbb{N}}$. Recall the canonical filtration on
$(\mathbb{T}^{\mathbb{N}},m_{\infty})$ defined by
$\sigma(z_{0})\subset\sigma(z_{0},z_{1})\subset\cdots\subset\sigma(z_{0},z_{1},\cdots,z_{n})\subset\cdots.$
From now on, we will denote
$\mathcal{G}_{n}=\sigma(z_{0},z_{1},\cdots,z_{n})$. Recall that
$H_{s}(\mathbb{T}^{\mathbb{N}})$ is the subspace of
$L_{s}(\mathbb{T}^{\mathbb{N}},m_{\infty})$ consisting of limit values of
Hardy martingales, i.e. $f\in H_{s}(\mathbb{T}^{\mathbb{N}})$ if and only if
$f\in L_{s}(\mathbb{T}^{\mathbb{N}},m_{\infty})$ and the associated martingale
$(\mathbb{E}^{\mathcal{G}_{n}}f)_{n\geq 0}$ is a Hardy martingale. For
convenience, we always assume $z_{0}\equiv 1$ such that $\mathcal{G}_{0}$ is a
trivial $\sigma$-algebra.
###### Definition 4.1.
Let $X$ be a Banach space and let $\\{x_{i}\\}_{i\in I}$ be a family of
vectors in $X$. The number $S^{a}(X;\\{x_{i}\\})$ is defined to be the best
constant $C$ such that for any $N\in\mathbb{N}$ and any finite sequence of
functions $(\theta_{k})_{k=0}^{N}$ in $H_{\infty}(\mathbb{T}^{\mathbb{N}})$,
we have
$\Big{\|}\sum_{k}\mathbb{E}^{\mathcal{G}_{k}}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{1}(m_{\infty};X)}\leq
C\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{\infty}(m_{\infty};X)}$
If there does not exist such constant, we set $S^{a}(X;\\{x_{i}\\})=\infty$.
If $\\{x_{i}\\}$ is clear from the context, then $S^{a}(X;\\{x_{i}\\})$ will
be simplified as $S^{a}(X)$.
The Stein type inequality still holds in this setting, more precisely, we have
###### Proposition 4.2.
Let $X$ be an AUMD space. For any $1\leq s<\infty$, let $(F_{k})_{k\geq 0}$ be
an arbitrary finite sequence in $H_{s}(\mathbb{T}^{\mathbb{N}};X)$. Then we
have
(14)
$\displaystyle\Big{\|}\sum_{k}\zeta_{k}\mathbb{E}^{\mathcal{G}_{k}}(F_{k})(z)\Big{\|}_{L_{s}(X)}\leq
C_{s}^{a}(X)\Big{\|}\sum_{k}\zeta_{k}F_{k}(z)\Big{\|}_{L_{s}(X)},$
where $\zeta=(\zeta_{k})_{k\geq 0}$ is an independent copy of
$z=(z_{k})_{k\geq 0}$ and
$L_{s}(X)=L_{s}(\mathbb{T}_{z}^{\mathbb{N}}\times\mathbb{T}_{\zeta}^{\mathbb{N}},m_{\infty}\times
m_{\infty};X)$.
###### Proof.
Consider the filtration on
$\mathbb{T}_{z}^{\mathbb{N}}\times\mathbb{T}_{\zeta}^{\mathbb{N}}$ defined by
$\mathcal{B}_{2j}=\sigma(z_{0},\cdots,z_{j})\otimes\sigma(\zeta_{0},\cdots,\zeta_{j})$
and
$\mathcal{B}_{2j-1}=\sigma(z_{0},\cdots,z_{j})\otimes\sigma(\zeta_{0},\cdots,\zeta_{j-1})$.
Then $f=\sum_{k}\zeta_{k}F_{k}(z)$ is a Hardy martingale with respect to the
above filtration. Let
$f^{\prime}=\sum_{k}\zeta_{k}\mathbb{E}^{\mathcal{G}_{k}}(F_{k})$. Then we
have
$f^{\prime}=\sum_{j}(\mathbb{E}^{\mathcal{B}_{2j}}-\mathbb{E}^{\mathcal{B}_{2j-1}})(f)$.
It follows (see Remark 3.3) that $\|f^{\prime}\|_{L_{s}(X)}\leq
C_{s}^{a}(X)\|f\|_{L_{s}(X)}$, whence (14). ∎
###### Proposition 4.3.
Let $X$ be an AUMD space. Assume that $\\{x_{i}\\}_{i\in I}$ is a
1-unconditional basic sequence in $X$. Then for any $1\leq s<\infty$ and any
finite sequence of functions $(\theta_{k})_{k\geq 0}$ in
$H_{s}(\mathbb{T}^{\mathbb{N}})$,
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}^{\mathcal{G}_{k}}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{s}(m_{\infty};X)}\leq
C_{s}^{a}(X)\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{s}(m_{\infty};X)}.$
###### Proof.
It follows verbatim the proof of Proposition 3.4. ∎
Let $X$ be as in Proposition 4.3, $\\{x_{i}\\}$ is a 1-unconditional basic
sequence in $X$. Then for all $1\leq s<\infty$, we have
$\displaystyle\Big{\|}\sum_{k}\mathbb{E}^{\mathcal{G}_{k}}(\theta_{k})x_{i_{k}}\Big{\|}_{L_{1}(m_{\infty};X)}\leq
C_{s}^{a}(X)\Big{\|}\sum_{k}\theta_{k}x_{i_{k}}\Big{\|}_{L_{\infty}(m_{\infty};X)}.$
Hence
$S^{a}(X;\\{x_{i}\\})\leq C_{s}^{a}(X)$
for all $1\leq s<\infty$.
###### Theorem 4.4.
Let $E$ be a Banach space with a 1-unconditional basis $\\{e_{i}:i\in I\\}$,
let $F$ be another Banach space. Let $E(F)$ be defined as in Theorem 3.5. For
any fixed family of vectors $\\{f_{j}:j\in J\\}$ in $F$, consider the family
of vectors $\\{e_{i}\otimes f_{j}:i\in I,j\in J\\}$ in $E(F)$, then we have
$S^{a}(E(F))\geq S^{a}(E)S^{a}(F),$
where $S^{a}(E(F))$, $S^{a}(E)$ and $S^{a}(F)$ are defined with respect to the
mentioned families of vectors respectively.
###### Proof.
The proof is similar to the proof of Theorem 3.5. We mention the slight
difference concerning the filtration. Consider the infinite tensor product
$L_{\infty}(\mathbb{T}^{\mathbb{N}})\otimes
L_{\infty}(\mathbb{T}^{\mathbb{N}})\otimes\cdots$, define
$z_{k,n}=\underbrace{1\otimes\cdots\otimes 1}_{k\quad\text{times}}\otimes
z_{n}\otimes 1\otimes\cdots,\text{if }n\geq 1$
and
$z_{k,0}=z_{k}\otimes 1\otimes 1\otimes\cdots.$
Then the filtration defined by
$\mathcal{F}^{a}_{k,n}:=\sigma\left(z_{j}:j\leq(k,n)\right)$ is an analytic
filtration, where the order on $\mathbb{N}\times\mathbb{N}$ is the lexigraphic
order as defined in the proof of Theorem 3.5. This filtration plays the role
similar to that of $(\mathcal{F}_{k,n})_{k,n}$ in the proof of Theorem 3.5.
Note that we may restrict to the functions $\theta_{k},\xi_{n}$ depending only
on finitely many variables. Thus only a finite subset of
$\mathbb{N}\times\mathbb{N}$ is used. ∎
The following lemma requires slightly more efforts than Lemma 3.8.
###### Lemma 4.5.
Suppose that $1\leq p\neq q<\infty$. If
$E_{1}=\ell_{p}^{(2)}(\ell_{q}^{(2)})$, then
$S^{a}(E_{1})\geq\kappa(p,q)>1.$
###### Proof.
We will use the notations in the proof of Lemma 3.8. Define a linear map
$U:H_{\infty}(\mathbb{T},m;E_{1})\rightarrow H_{1}(\mathbb{T},m;E_{1})$ by
$U\Big{[}a_{ij}(z)e_{i}^{p}\otimes
e_{j}^{q}\Big{]}=\left\\{\begin{array}[]{lc}\mathbb{E}(a_{ij})e_{i}^{p}\otimes
e_{j}^{q},\text{ if }j=1\\\ a_{ij}(z)e_{i}^{p}\otimes e_{j}^{q},\text{ if
}j=2\end{array}.\right.$
If $C=\|U\|_{H_{\infty}(E_{1})\rightarrow H_{1}(E_{1})}$, then
$S^{a}(E_{1})\geq C$. By definition, for any $a,b,c,d$ functions in
$H_{\infty}(\mathbb{T})$, we have
$\displaystyle\int_{\mathbb{T}}\Big{\\{}(|\mathbb{E}a|^{q}+|b(z)|^{q})^{p/q}+(|\mathbb{E}c|^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}dm(z)$
$\displaystyle\leq$ $\displaystyle
C\sup_{z\in\mathbb{T}}\Big{\\{}(|a(z)|^{q}+|b(z)|^{q})^{p/q}+(|c(z)|^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}.$
Note that if $a,c$ are outer functions, then by (5), we have
$|\mathbb{E}a|=M(|a|)$ and $|\mathbb{E}c|=M(|c|)$. So for any functions
$a,b,c,d\in H_{\infty}(\mathbb{T})$ such that $a,c$ are outer, we have
$\displaystyle\int_{\mathbb{T}}\Big{\\{}(M(|a|)^{q}+|b(z)|^{q})^{p/q}+(M(|c|)^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}dm(z)$
$\displaystyle\leq$ $\displaystyle
C\sup_{z\in\mathbb{T}}\Big{\\{}(|a(z)|^{q}+|b(z)|^{q})^{p/q}+(|c(z)|^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}.$
By the classical Szegö’s condition, if
$a^{\prime},b^{\prime},c^{\prime},d^{\prime}$ are functions in
$L_{\infty}(\mathbb{T})$ which are bounded from below, then there are outer
functions $a,b,c,d\in H_{\infty}(\mathbb{T})$, such that
$|a^{\prime}|=|a|,|b^{\prime}|=|b|,|c^{\prime}|=|c|,|d^{\prime}|=|d|$ a.e..
Hence (4) still holds for any 2-valued non-vanishing functions $a,b,c,d\in
L_{\infty}(\mathbb{T})$ (note that for a function taking only two values, non-
vanishing is the same as bounded from below). By approximation, we can further
relax the non-vanishing condition on $b,d$. Now consider any measurable
partition $\mathbb{T}=A\dot{\cup}B$, such that $m(A)=m(B)=\frac{1}{2}$. If
$a=u\chi_{A}+v\chi_{B}$, $c=v\chi_{A}+u\chi_{B}$, $b=w\chi_{A}+t\chi_{B}$ and
$d=t\chi_{A}+w\chi_{B}$, then it is easy to check that for all
$z\in\mathbb{T}$, we have
$\displaystyle\Big{\\{}(|a(z)|^{q}+|b(z)|^{q})^{p/q}+(|c(z)|^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$
$\displaystyle\Big{\\{}(|u|^{q}+|w|^{q})^{p/q}+(|v|^{q}+|t|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$ $\displaystyle
2^{1/p}\Big{\\{}\int_{\mathbb{T}}(|a|^{q}+|b|^{q})^{p/q}dm\Big{\\}}^{1/p}.$
Similarly for all $z\in\mathbb{T}$, we have
$\displaystyle\Big{\\{}(M(|a|)^{q}+|b(z)|^{q})^{p/q}+(M(|c|)^{q}+|d(z)|^{q})^{p/q}\Big{\\}}^{1/p}$
$\displaystyle=$ $\displaystyle
2^{1/p}\Big{\\{}\int_{\mathbb{T}}(M(|a|)^{q}+|b|^{q})^{p/q}dm\Big{\\}}^{1/p}.$
Substituting these equalities to (4), we get
$\Big{\\{}\int_{\mathbb{T}}(M(|a|)^{q}+|b|^{q})^{p/q}dm\Big{\\}}^{1/p}\leq
C\Big{\\{}\int_{\mathbb{T}}(|a|^{q}+|b|^{q})^{p/q}dm\Big{\\}}^{1/p}.$
By Proposition 2.5, we have $C\geq\kappa(p,q)$. This completes the proof. ∎
###### Theorem 4.6.
Suppose that $1\leq p\neq q<\infty$. If $E_{n}$’s are defined as in Theorem
3.10, then for any $1\leq s<\infty$, we have
$C_{s}^{a}(E_{n})\geq S^{a}(E_{n})\geq\kappa(p,q)^{n}.$
Moreover, if $1<p,q<\infty$, then there exists $\kappa_{2}=\kappa_{2}(p,q,s)$,
such that
$C_{s}^{a}(E_{n})\leq\kappa_{2}^{n}.$
###### Proof.
The first part of proof is identical to the proof of Theorem 3.10. The second
part follows from the fact that $C_{s}^{a}(E_{n})\leq C_{s}(E_{n})$ and
Proposition 3.11. ∎
## 5\. Construction and further discussions
For the sake of clearness, we introduce the family $X_{n}(p,q)$, which is
defined as follows: Let $X_{0}(p,q)=\mathbb{R}$, and define by recursion that
$X_{n+1}(p,q)=L_{p}(D,\mu;L_{q}(D,\mu;X_{n}(p,q))).$
In the complex case, $X^{\mathbb{C}}_{n}(p,q)$ is defined similarly.
Obviously, $X_{n}(p,q)$ is isometric to $E_{n}$ defined in the previous
sections using $p,q$. Our main purpose for introducing $X_{n}$’s is the
existence of canonical isometric inclusion $X_{n}(p,q)\subset X_{n+1}(p,q)$.
By these inclusions, the union $\cup_{n}X_{n}(p,q)$ is a normed space and its
completion will be denoted by $X(p,q)$. We have
$X(p,q):=\overline{\cup_{n}X_{n}(p,q)}\simeq\lim_{\longrightarrow}X_{n}(p,q),$
where the last term is the inductive limit of $X_{n}(p,q)$’s associated to the
canonical inclusions. In the complex case, $X^{\mathbb{C}}(p,q)$ is defined
similarly.
###### Remark 5.1.
If $1\leq p=q<\infty$, then $X(p,p)$ is the real space
$L^{p}_{\mathbb{R}}(D^{\mathbb{N}},\mu^{\otimes\mathbb{N}})$ and
$X^{\mathbb{C}}(p,p)$ is the complex space
$L_{\mathbb{C}}^{p}(D^{\mathbb{N}},\mu^{\otimes\mathbb{N}})$.
We have the following complex interpolation result.
###### Proposition 5.2.
Let $1<p_{0},p_{1},q_{0},q_{1}<\infty$ and $0<\theta<1$. Then we have the
following isometric isomorphism:
$X^{\mathbb{C}}(p_{\theta},q_{\theta})=[X^{\mathbb{C}}(p_{0},q_{0}),X^{\mathbb{C}}(p_{1},q_{1})]_{\theta},$
with $\frac{1}{p}=\frac{\theta}{p_{1}}+\frac{1-\theta}{p_{0}}$ and
$\frac{1}{q}=\frac{\theta}{q_{1}}+\frac{1-\theta}{q_{0}}$.
###### Proof.
Note that $X(p,q)$ is a Banach lattice of functions on
$(D^{\mathbb{N}},\mu^{\otimes\mathbb{N}})$. Clearly, $X(p,q)$ is
$\min(p,q)$-convex and $\max(p,q)$-concave in the sense of §1.d in [10], and
hence by Theorem 1.f.1 (p. 80) and Proposition 1.e.3 (p. 61) in [10] it is
reflexive. Then the above result is a particular case of a classical formula
going back to Calderón ([8], p. 125). ∎
Recall that a Banach space $X$ over the complex field is $\theta$-Hilbertian
($0\leq\theta\leq 1$) if there exists an interpolation pair $(X_{0},X_{1})$ of
Banach spaces such that $X$ is isometric with $[X_{0},X_{1}]_{\theta}$ and
$X_{1}$ is a Hilbert space.
###### Corollary 5.3.
Let $1<p\neq q<\infty$. Then $X(p,q)$ is non-UMD and $X^{\mathbb{C}}(p,q)$ is
non-AUMD. Moreover, there exists $0<\theta<1$ such that $X^{\mathbb{C}}(p,q)$
is $\theta$-Hilbertian. In particular, $X^{\mathbb{C}}(p,q)$ and a fortiori
$X(p,q)$ is super-reflexive.
###### Proof.
It follows easily from Theorem 3.10 and Theorem 4.6 that $X(p,q)$ is non-UMD
and $X^{\mathbb{C}}(p,q)$ is non-AUMD.
For $0<\theta<1$ small enough, such that
$\max(\frac{1/p-\theta/2}{1-\theta},\frac{1/q-\theta/2}{1-\theta})<1$, we can
find $1<\tilde{p},\tilde{q}<\infty$ satisfying the equalities:
$\frac{1}{p}=\frac{\theta}{2}+\frac{1-\theta}{\tilde{p}},\quad\frac{1}{q}=\frac{\theta}{2}+\frac{1-\theta}{\tilde{q}}.$
By Proposition 5.2, we have
$X^{\mathbb{C}}(p,q)=[X^{\mathbb{C}}(\tilde{p},\tilde{q}),X^{\mathbb{C}}(2,2)]_{\theta}.$
Since
$X^{\mathbb{C}}(2,2)=L^{2}_{\mathbb{C}}(D^{\mathbb{N}},\mu^{\otimes\mathbb{N}})$
is Hilbertian, $X^{\mathbb{C}}(p,q)$ is $\theta$-Hilbertian. The super-
reflexivity of $X^{\mathbb{C}}(p,q)$ follows from the well-known fact that any
$\theta$-Hilbertian space is super-reflexive for $\theta>0$ (cf.[12]). ∎
###### Remark 5.4.
Let $1<p\neq q<\infty$. For any $0<\eta<1$, let
$\frac{1}{p_{\eta}}=\frac{1-\eta}{p}+\frac{\eta}{q}$ and
$\frac{1}{q_{\eta}}=\frac{1-\eta}{q}+\frac{\eta}{p}$. By Proposition 5.2, we
have
$X^{\mathbb{C}}(p_{\eta},q_{\eta})=[X^{\mathbb{C}}(p,q),X^{\mathbb{C}}(q,p)]_{\eta}.$
Note that in this interpolation scale, there is only one UMD space
corresponding to $\eta=\frac{1}{2}$.
For futher discussions, let us now turn to the non-atomic case and modify
slightly the definitions. For any $1<p,q<\infty$, consider the family of
spaces $Z_{n}=Z_{n}(p,q)$ defined by recursion: $Z_{0}=\mathbb{C}$ and
$Z_{n+1}=Z_{n}(L_{p}(\mathbb{T},m;L_{q}(\mathbb{T},m))$. From the definition,
we have
$Z_{n}(p,q)\subset Z_{n+1}(p,q).$
Thus we can define
$Z(p,q)=\lim_{\longrightarrow}Z_{n}(p,q).$
To avoid ambiguity, let us emphasize the inclusions $Z_{n}(p,q)\subset
Z_{n+1}(p,q)$ used to define the inductive limit. For simplicity of notations,
we will write $L_{p_{1}}L_{p_{2}}=L_{p_{1}}(L_{p_{2}})$,
$L_{p_{1}}L_{p_{2}}L_{p_{3}}=L_{p_{1}}(L_{p_{2}}(L_{p_{3}}))$, etc. With these
notations, one can easily see the difference between $X_{n}$ and $Z_{n}$ as
follows:
$X_{n+1}=L_{p}(L_{q}(X_{n}))=L_{p}L_{q}\underbrace{L_{p}L_{q}\cdots
L_{p}L_{q}}_{X_{n}},$
where $L_{p}=L_{p}(D,\mu)$ and $L_{q}=L_{q}(D,\mu)$ are two dimensional. And
$Z_{n+1}=Z_{n}(L_{p}(L_{q}))=\underbrace{L_{p}L_{q}\cdots
L_{p}L_{q}}_{Z_{n}}L_{p}L_{q},$
where $L_{p}=L_{p}(\mathbb{T},m)$ and $L_{q}=L_{q}(\mathbb{T},m)$.
###### Remark 5.5.
The main purpose of introducing the spaces $Z_{n}(p,q)$ is that we have
lattice isometric isomorphisms $L_{p}(Z_{n}(p,q))\simeq Z_{n}(p,q)$ for all
$n$ and moreover, these isomorphisms are compatible with the inclusion of
$Z_{n}(p,q)\subset Z_{n+1}(p,q)$ (the word “compatible” will be explained by a
commutative diagram in the sequel) and this will be used to show some
additional properties for $Z(p,q)$. The family of $X_{n}(p,q)$’s shares the
property of having lattice isometric isomorphisms $L_{p}(X_{n}(p,q))\simeq
X_{n}(p,q)$ for all $n$, but the isomorphisms are not compatible with the
inclusions $X_{n}(p,q)\subset X_{n+1}(p,q)$.
The $Z(p,q)$’s are Banach lattices of functions on the infinite torus
$\mathbb{T}^{\mathbb{N}}$, they have the following properties.
###### Proposition 5.6.
Let $1<p,q<\infty$. We have isomorphisms
$Z(p,q)\simeq Z(q,p)$
and
$L_{p}(Z(p,q))\simeq L_{q}(Z(p,q)).$
If $p\neq q$, then $Z(p,q)$ does not have unconditional basis.
###### Proof.
Since $L_{p}(\mathbb{T})$ and $L_{p}(\mathbb{T}\times\mathbb{T})$ are
isometric as Banach lattices, we have isometric isomorphisms which are
compatible with the inclusions $Z_{n}\subset Z_{n+1}$, that is we have the
commutative diagram
$\begin{CD}Z_{n}(p,q)@>{\text{ inclusion }}>{}>Z_{n+1}(p,q)\\\ @V{\text{
isometric }}V{\simeq}V@V{\simeq}V{\text{ isometric}}V\\\ L_{p}(Z_{n}(p,q))@
>\text{ inclusion }>>L_{p}(Z_{n+1}(p,q)).\end{CD}$
By taking Banach space inductive limit, we have
$Z(p,q)\xrightarrow[\text{isometric}]{\simeq}L_{p}(Z(p,q)).$
If $p\neq q$, then $Z(p,q)$ and hence $L_{p}(Z(p,q))$ is non-UMD. By a result
of D.J. Aldous (see [1], Proposition 4), $Z(p,q)$ has no unconditional basis.
It is easy to see that $Z(p,q)$ and $Z(q,p)$ complementably embed into each
other. Since $\ell_{p}^{(2)}(L_{p})=L_{p}$ as Banach lattices, we have
$\ell_{p}^{(2)}(L_{p}(Z(p,q)))=L_{p}(Z(p,q)).$
Moreover, since $L_{p}(Z(p,q))=Z(p,q)$, the above isometry implies that as
Banach space $Z(p,q)=Z(p,q)\oplus Z(p,q)$. Similarly, $Z(q,p)=Z(q,p)\oplus
Z(q,p)$. By the classical Pełcyński decomposition method, we have
$Z(p,q)\simeq Z(q,p)$. Hence
$\qquad\quad L_{p}(Z(p,q))=Z(p,q)\simeq Z(q,p)=L_{q}(Z(q,p))\simeq
L_{q}(Z(p,q)).$
∎
Let $(p_{i})_{i\geq 1}$ be a sequence of real numbers such that
$1<p_{i}<\infty$. Define
$X[(p_{i})]=\lim_{\longrightarrow}L_{p_{n}}\cdots L_{p_{2}}L_{p_{1}}$
and
$Z[(p_{i})]=\lim_{\longrightarrow}L_{p_{1}}L_{p_{2}}\cdots L_{p_{n}}.$
###### Problem.
Under which condition is $X[(p_{i})]$ or $Z[(p_{i})]$ in the UMD class ?
We have the following observations on the necessary condition:
* (i)
A trivial necessary condition is that there exist $1<p_{0},p_{\infty}<\infty$,
such that $p_{0}\leq p_{i}\leq p_{\infty}$ for all $i\geq 1$.
* (ii)
If the above condition is satisfied, then the sequence $(p_{i})$ has at least
one cluster point $1<p<\infty$. Then a necessary condition is that the
sequence has only one cluster point, i.e. $\lim_{i\to\infty}p_{i}=p$. Indeed,
assume that the sequence $(p_{i})$ has two cluster points $1<p\neq q<\infty$,
so that there exist two subsequences of $(p_{i})$ which tend to $p,q$
respectively. Then one can easily show that by Theorem 3.10, both $X[(p_{i})]$
and $Z[(p_{i})]$ are non-UMD (they are in fact non-AUMD).
* (iii)
Now the speed of convergence of $(p_{i})$ will play a role. Since
$\ell_{p_{1}}^{(2)}(\ell_{p_{2}}^{(2)}(\cdots(\ell_{p_{n}}^{(2)})\cdots))$
embeds isometrically into $L_{p_{1}}L_{p_{2}}\cdots L_{p_{n}}$. A necessary
condition for $Z[(p_{i})]$ to be UMD is $\prod_{i}c(p_{2i},p_{2i+1})<\infty$.
Similarly, it is necessary that $\prod_{i}c(p_{2i+1},p_{2i+2})<\infty$.
Combining these, a necessary condition for $Z[(p_{i})]$ to be in the UMD class
is
$\prod_{i}c(p_{i},p_{i+1})<\infty.$
The same statement remains true for $X[(p_{i})]$. Note that by (4),
$c(p_{i},p_{i+1})>1$ if $p_{i}\neq p_{i+1}$.
Intuitively, if $p_{i}$ tends to $p$ sufficiently fast, then both $X[(p_{i})]$
and $Z[(p_{i})]$ are in the UMD class. The author obtained some partial
results in this direction, which will be treated elsewhere.
###### Remark 5.7.
Let $1<p<q<\infty$. We have the following Banach lattices isometries
$L_{p}L_{q}=L_{p}L_{p}L_{q},\quad L_{p}L_{q}=L_{p}L_{q}L_{q}.$
Since $L_{p}L_{r}L_{q}$ is an interpolation space between $L_{p}L_{p}L_{q}$
and $L_{p}L_{q}L_{q}$ for any $p\leq r\leq q$, the $\text{UMD}_{s}$ constant
of $L_{p}L_{r}L_{q}$ is actually the same as that of $L_{p}(L_{q})$. The same
argument shows that $L_{p}L_{u}L_{r}L_{v}L_{q}$ has the same $\text{UMD}_{s}$
constant with $L_{p}L_{q}$, provided $p\leq u\leq r\leq v\leq q$. More
generally, if $(p_{i})_{i=1}^{n}$ is a finite sequence, assume that
$(p_{i})_{i=k}^{l}$ is consecutive monotone (non-increasing or non-decreasing)
subsequence, then $L_{p_{1}}\cdots L_{p_{k}}\cdots L_{p_{l}}\cdots L_{p_{n}}$
and $L_{p_{1}}\cdots L_{p_{k}}L_{p_{l}}\cdots L_{p_{n}}$ have the same
$\text{UMD}_{s}$ constant for all $1<s<\infty$.
Our results have some applications in the non-commutative setting, i.e. on the
operator space UMD property, which will appear in a future publication.
## Acknowledgements
This work was carried out while the author was visiting at Texas A&M
University. The author would like to acknowledge the hospitality provided by
Department of Mathematics of Texas A&M. He would like to thank his advisor G.
Pisier for suggesting this problem and for the constant and stimulating
discussions. The author also appreciates the careful review of the paper by
the referee who suggested many changes to enhance the readability of the
paper.
## References
* [1] D. J. Aldous. Unconditional bases and martingales in $L_{p}(F)$. Math. Proc. Cambridge Philos. Soc., 85(1):117–123, 1979.
* [2] J. Bourgain. Some remarks on Banach spaces in which martingale difference sequences are unconditional. Ark. Mat., 21(2):163–168, 1983.
* [3] J. Bourgain. On martingales transforms in finite-dimensional lattices with an appendix on the $K$-convexity constant. Math. Nachr., 119:41–53, 1984.
* [4] J. Bourgain. Vector-valued singular integrals and the $H^{1}$-BMO duality. In Probability theory and harmonic analysis (Cleveland, Ohio, 1983), volume 98 of Monogr. Textbooks Pure Appl. Math., pages 1–19. Dekker, New York, 1986.
* [5] D. L. Burkholder. A geometrical characterization of Banach spaces in which martingale difference sequences are unconditional. Ann. Probab., 9(6):997–1011, 1981.
* [6] D. L. Burkholder. A geometric condition that implies the existence of certain singular integrals of Banach-space-valued functions. In Conference on harmonic analysis in honor of Antoni Zygmund, Vol. I, II (Chicago, Ill., 1981), Wadsworth Math. Ser., pages 270–286. Wadsworth, Belmont, CA, 1983.
* [7] Donald L. Burkholder. Martingales and singular integrals in Banach spaces. In Handbook of the geometry of Banach spaces, Vol. I, pages 233–269. North-Holland, Amsterdam, 2001.
* [8] A.-P. Calderón. Intermediate spaces and interpolation, the complex method. Studia Math., 24:113–190, 1964.
* [9] D. J. H. Garling. On martingales with values in a complex Banach space. Math. Proc. Cambridge Philos. Soc., 104(2):399–406, 1988.
* [10] Joram Lindenstrauss and Lior Tzafriri. Classical Banach spaces. II, volume 97 of Ergebnisse der Mathematik und ihrer Grenzgebiete. Springer-Verlag, Berlin, 1979. Function spaces.
* [11] G. Pisier. Un exemple concernant la super-réflexivité. In Séminaire Maurey-Schwartz 1974–1975: Espaces $L^{p}$ applications radonifiantes et géométrie des espaces de Banach, Annexe No. 2, page 12. Centre Math. École Polytech., Paris, 1975\.
* [12] G. Pisier. Some applications of the complex interpolation method to Banach lattices. J. Analyse Math., 35:264–281, 1979.
* [13] José L. Rubio de Francia. Martingale and integral transforms of Banach space valued functions. In Probability and Banach spaces (Zaragoza, 1985), volume 1221 of Lecture Notes in Math., pages 195–222. Springer, Berlin, 1986.
* [14] Nicole Tomczak-Jaegermann. Banach-Mazur distances and finite-dimensional operator ideals, volume 38 of Pitman Monographs and Surveys in Pure and Applied Mathematics. Longman Scientific & Technical, Harlow, 1989.
|
arxiv-papers
| 2011-12-04T09:46:28 |
2024-09-04T02:49:24.951273
|
{
"license": "Public Domain",
"authors": "Yanqi Qiu",
"submitter": "Yanqi Qiu",
"url": "https://arxiv.org/abs/1112.0739"
}
|
1112.0854
|
# New approach for normalization and photon-number distributions of photon-
added (-subtracted) squeezed thermal states
Li-Yun Hu1,2† and Zhi-Ming Zhang2‡ 1College of Physics & Communication
Electronics, Jiangxi Normal University, Nanchang 330022, China
2Key Laboratory of Photonic Information Technology of Guangdong Higher
Education Institutes,
SIPSE & LQIT, South China Normal University, Guangzhou 510006, China
${\dagger}$E-mail: hlyun2008@126.com; ‡ E-mail: zmzhang@scnu.edu.cn.
###### Abstract
Using the thermal field dynamics theory to convert the thermal state to a
“pure” state in doubled Fock space, it is found that the average value of
$e^{fa^{{\dagger}}a}$ under squeezed thermal state (STS) is just the
generating function of Legendre polynomials, a remarkable result. Based on
this point, the normalization and photon-number distributions of m-photon
added (or subtracted) STS are conviently obtained as the Legendre polynomials.
This new concise method can be expanded to the entangled case.
## I Introduction
Nonclassicality of optical fields is helpful in understanding fundamentals of
quantum optics and have many applications in quantum information processing 1
. To generate and manipulate various nonclassical optical fields, subtracting
or adding photons from/to traditional quantum states or Gaussian states is
proposed 2 ; 2a ; 3 ; 4 ; 5 ; 6 ; 7 ; 8 . For example, the photon addition and
subtraction have been successfully demonstrated experimentally for probing
quantum commutation rules by Parigi et al. 6 . Recently, photon-added
(-subtracted) Gaussian states have received more attention from both
experimentalists and theoreticians 9 ; 10 ; 11 ; 11a ; 12 ; 13 ; 14 ; 15 ; 16
; 17 , since these states exhibit an abundant of nonclassical properties and
may give access to a complete engineering of quantum states and to fundamental
quantum phenomena.
Theoretically, the normalization factors of such quantum states are essential
for studying their nonclassical properties. Very recent, Fan and Jiang 18
present a new concise approach for normalizing m-photon-added (-subtracted)
squeezed vacuum state (pure state) by constructing generating function.
However, most systems are not isolated, but are immersed in a thermal
reservoir, thus it is often the case that we have no enough information to
specify completely the state of a system. In such a situation, the system only
can be described by mixed states, such as thermal states. In addition, the
squeezed thermal states (STSs) can be considered as the generalized Gaussian
states.
In this paper, we shall extend this case to the mixed state, i.e., by using
the thermal field dynamics (TFD) theory to convert the thermal state to a
“pure” state in doubled Fock space, we present a new concise method for
normalizing photon-added (-subtracted) squeezed thermal states (PASTSs,
PSSTSs) and for deriving their photon-number distributions (PNDs) which have
been a major topic of studies on quantum optics and quantum statistics. It is
found that the normalization factors and PNDs are related to the Legendre
polynomials in a compact form.
Our paper is arranged as follows. In section 2, based on Takahashi-Umezawa
TFD, we convert the thermal state $\rho_{th}$ to a pure state in doubled Fock
space by the partial trace,
$\rho_{th}=\mathtt{t\tilde{r}}\left[\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\right]$ (see Eq.(4) below). In section 3, we introduce the STS
$S_{1}\left|0(\beta)\right\rangle$ with $S_{1}$ being the single-mode
squeezing operator for the real mode. It is shown that the average value of
$e^{fa^{{\dagger}}a}$ under STS is just the generating function of Legendre
polynomials, a remarkable result. Based on this point, in sections 4 and 5,
the normalization factors and PNDs of m-photon added (or subtracted) STS are
obtained as the Legendre polynomials, respectively. The last section is
devoted to drawing a conclusion.
## II Representation of thermal state in doubled Fock space
We begin with briefly reviewing the properties of thermal state. For a single
field mode with frequency $\omega$ in a thermal equilibrium state
corresponding to absolute temperature $T$, the density operator is
$\rho_{th}=\sum_{n=0}^{\infty}\frac{n_{c}^{n}}{\left(n_{c}+1\right)^{n+1}}\left|n\right\rangle\left\langle
n\right|,$ (1)
where $n_{c}=[\exp(\hbar\omega/(kT))-1]^{-1}$ being the average photon number
of the thermal state $\rho_{th}$ and $k$ being Bltzmann’s constant. Note
$\left|n\right\rangle=a^{\dagger n}/\sqrt{n!}\left|0\right\rangle$ and the
normally ordering form of vacuum projector $\left|0\right\rangle\left\langle
0\right|=\colon\exp(-a^{\dagger}a)\colon$(the symbol $\colon\colon$ denotes
the normal ordering), one can put Eq.(1) into the following form
$\rho_{th}=\colon\frac{1}{n_{c}+1}e^{-\frac{1}{n_{c}+1}a^{{\dagger}}a}\colon=\frac{1}{n_{c}+1}e^{a^{{\dagger}}a\ln\frac{n_{c}}{n_{c}+1}},$
(2)
where in the last step, the operator identity $\exp\left(\lambda
a^{{\dagger}}a\right)=\colon\exp\left[\left(e^{\lambda}-1\right)a^{{\dagger}}a\right]\colon$is
used.
Recalling the thermal field dynamics (TFD) introduced by Takahashi and Umezawa
19 ; 20 ; 21 , its elemental spirit is to convert the calculations of ensemble
averages for a mixed state $\rho$, $\left\langle
A\right\rangle=\mathtt{tr}\left(A\rho\right)/\mathtt{tr}\left(\rho\right),$ to
equivalent expectation values with a pure state $\left|0(\beta)\right\rangle$,
i.e.,
$\left\langle A\right\rangle=\left\langle
0(\beta)\right|A\left|0(\beta)\right\rangle,$ (3)
where $\beta=1/kT$, $k$ is the Boltzmann constant. Thus, for the density
operator $\rho_{th}$, by using the partial trace method 22 , i.e.,
$\rho_{th}=\widetilde{\mathtt{tr}}\left[\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\right]$ where $\widetilde{\mathtt{tr}}$ denotes the trace
operation over the environment freedom (denoted as operator
$\tilde{a}^{{\dagger}}$), one can obtain the explicit expression of
$\left|0(\beta)\right\rangle$ in doubled Fock space,
$\left|0(\beta)\right\rangle=\text{sech}\theta\exp\left[a^{\dagger}\tilde{a}^{\dagger}\tanh\theta\right]\left|0\tilde{0}\right\rangle=S\left(\theta\right)\left|0\tilde{0}\right\rangle,$
(4)
where $\left|0\tilde{0}\right\rangle$ is annihilated by $\tilde{a}$ and $a$,
$[\tilde{a},\tilde{a}^{{\dagger}}]=1$, and $S\left(\theta\right)$ is the
thermal operator,
$S\left(\theta\right)\equiv\exp\left[\theta\left(a^{\dagger}\tilde{a}^{\dagger}-a\tilde{a}\right)\right]\
$with a similar form to the a two-mode squeezing operator except for the tilde
mode, and $\theta$ is a parameter related to the temperature by
$\tanh\theta=\exp\left(-\frac{\hbar\omega}{2kT}\right)$.
$\left|0(\beta)\right\rangle$ is named thermal vacuum state.
Let $\mathtt{Tr}$ denote the trace operation over both the system freedom
(expressed by $\mathtt{tr}$) and the environment freedom by
$\widetilde{\mathtt{tr}}$, i.e.,
$\mathtt{Tr}=\mathtt{tr}\widetilde{\mathtt{tr}}$, then we have
$\displaystyle\mathtt{tr}\left(A\rho_{th}\right)$
$\displaystyle=\mathtt{Tr}\left[A\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\right]$
$\displaystyle=\mathtt{tr}\left[A\widetilde{\mathtt{tr}}\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\right],$ (5)
and the average photon number of the thermal state $\rho_{th}$ is
$n_{c}=\mathtt{Tr}\left[a^{\dagger}a\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\right]=\sinh^{2}\theta.$ (6)
Here we should emphasize that
$\widetilde{\mathtt{tr}}\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|\neq\left\langle 0(\beta)\right|\left.0(\beta)\right\rangle,$
since $\left|0(\beta)\right\rangle$ involves both real mode $a$ and fictitious
mode $\tilde{a}$. From Eqs.(3) and (4) one can see that the worthwhile
convenience in Eq.(4) is at the expense of introducing a fictitious field (or
called a tilde-conjugate field) in the extended Hilbert space, i.e., the
original optical field state $\left|n\right\rangle$ in the Hilbert space
$\mathcal{H}$ is accompanied by a tilde state $\left|\tilde{n}\right\rangle$
in $\mathcal{\tilde{H}}$. A similar rule holds for operators: every Bose
annihilation operator $a$ acting on $\mathcal{H}$ has an image $\tilde{a}$
acting on $\mathcal{\tilde{H}}$. These operators in $\mathcal{H}$ are
commutative with those in $\mathcal{\tilde{H}}$.
## III Suqeezed thermal vacuum state
To realize our purpose, we first introduce the squeezed thermal vacuum state,
defined as $S_{1}\left(r\right)\left|0(\beta)\right\rangle$, where
$S_{1}\left(r\right)=\exp[r(a^{2}-a^{\dagger 2})/2]$ is the single-mode
squeezing operator for the real mode with $r$ being squeezing parameter. Note
that Eq.(4) and the Baker-Hausdorf lemma
$S_{1}\left(r\right)a^{{\dagger}}S_{1}^{{\dagger}}\left(r\right)=a^{{\dagger}}\cosh
r+a\sinh r,$ (7)
then we get
$\displaystyle S_{1}\left(r\right)\left|0(\beta)\right\rangle$
$\displaystyle=S_{1}\left(r\right)\text{sech}\theta\exp\left[a^{\dagger}\tilde{a}^{\dagger}\tanh\theta\right]\left|0\tilde{0}\right\rangle$
$\displaystyle=\text{sech}\theta\text{sech}^{1/2}r\exp\left[\left(a^{{\dagger}}\cosh
r+a\sinh r\right)\tilde{a}^{\dagger}\tanh\theta\right]$
$\displaystyle\times\exp\left[-\frac{a^{\dagger 2}}{2}\tanh
r\right]\left|0\tilde{0}\right\rangle,$ (8)
where we have used
$S_{1}\left(\lambda\right)\left|0\right\rangle=$sech${}^{1/2}\lambda\exp[-a^{\dagger
2}/2\tanh\lambda]\left|0\right\rangle.$ Further, note $e^{\tau
a\tilde{a}^{\dagger}}a^{\dagger}e^{-\tau
a\tilde{a}^{\dagger}}=a^{\dagger}+\tau\tilde{a}^{\dagger},$ and for operators
$A,B$ satisfying the conditions $\left[A,[A,B]\right]=\left[B,[A,B]\right]=0,$
we have $e^{A+B}=e^{A}e^{B}e^{-[A,B]/2},$ thus Eq.(8) can be put into the
following form
$\displaystyle S_{1}\left(r\right)\left|0(\beta)\right\rangle$
$\displaystyle=\text{sech}\theta\text{sech}^{1/2}r\exp\left\\{\frac{\tanh\theta}{\cosh
r}a^{{\dagger}}\tilde{a}^{\dagger}\right.$ $\displaystyle+\left.\frac{\tanh
r}{2}\left(\tilde{a}^{\dagger 2}\tanh^{2}\theta-a^{\dagger
2}\right)\right\\}\left|0\tilde{0}\right\rangle.$ (9)
Next, we shall use Eq.(9) to derive the average of operator
$e^{fa^{\dagger}a}$ under the suqeezed thermal vacuum state
$S_{1}\left(r\right)\left|0(\beta)\right\rangle$, which is a bridge for our
whole calculations. Notice
$e^{f/2a^{\dagger}a}a^{\dagger}e^{-f/2a^{\dagger}a}=a^{\dagger}e^{f/2}$ and
$e^{-f/2a^{\dagger}a}ae^{f/2a^{\dagger}a}=ae^{f/2}$, so we have
$\displaystyle
e^{f/2a^{\dagger}a}S_{1}\left(r\right)\left|0(\beta)\right\rangle$
$\displaystyle=\text{sech}\theta\text{sech}^{1/2}r\exp\left\\{\frac{\tanh\theta}{\cosh
r}a^{{\dagger}}\tilde{a}^{\dagger}e^{f/2}\right.$
$\displaystyle+\left.\frac{\tanh r}{2}\left(\tilde{a}^{\dagger
2}\tanh^{2}\theta-a^{\dagger
2}e^{f}\right)\right\\}\left|0\tilde{0}\right\rangle,$ (10)
which leads to
$\displaystyle\left\langle e^{fa^{\dagger}a}\right\rangle$
$\displaystyle\equiv\left\langle
0(\beta)\right|S_{1}^{{\dagger}}\left(r\right)e^{fa^{\dagger}a}S_{1}\left(r\right)\left|0(\beta)\right\rangle$
$\displaystyle=\text{sech}^{2}\theta\text{sech}r\left\langle
0\tilde{0}\right|\exp\left\\{\frac{e^{f/2}\tanh\theta}{\cosh
r}a\tilde{a}\right.$ $\displaystyle+\left.\frac{\tanh
r}{2}\left(\tilde{a}^{2}\tanh^{2}\theta-a^{2}e^{f}\right)\right\\}$
$\displaystyle\times\exp\left\\{\frac{\tanh\theta}{\cosh
r}a^{{\dagger}}\tilde{a}^{\dagger}e^{f/2}\right.$
$\displaystyle+\left.\frac{\tanh r}{2}\left(\tilde{a}^{\dagger
2}\tanh^{2}\theta-a^{\dagger
2}e^{f}\right)\right\\}\left|0\tilde{0}\right\rangle$
$\displaystyle=\left[Ce^{2f}-2Be^{f}+A\right]^{-1/2},$ (11)
where we have set
$\displaystyle A$ $\displaystyle=\cosh^{4}\theta+\cosh 2\theta\sinh^{2}r$
$\displaystyle=n_{c}^{2}+\left(2n_{c}+1\right)\cosh^{2}r,$ $\displaystyle B$
$\displaystyle=\sinh^{2}\theta\cosh^{2}\allowbreak\theta=n_{c}\left(n_{c}+1\right),$
$\displaystyle C$ $\displaystyle=\cosh^{4}\theta-\cosh 2\theta\cosh^{2}r$
$\displaystyle=n_{c}^{2}-\left(2n_{c}+1\right)\sinh^{2}r,$ (12)
and we have used the completness relation of coherent state $\int
d^{2}zd^{2}\tilde{z}\left|z\tilde{z}\right\rangle\left\langle
z\tilde{z}\right|/\pi^{2}=1,$ here $\left|z\right\rangle$ and
$\left|\tilde{z}\right\rangle$ is the coherent state in real and fictitious
modes, respectively, and the formula 23
$\displaystyle\int\frac{d^{2}z}{\pi}\exp\left(\zeta\left|z\right|^{2}+\xi
z+\eta z^{\ast}+fz^{2}+gz^{\ast 2}\right)$
$\displaystyle=\frac{1}{\sqrt{\zeta^{2}-4fg}}\exp\left[\frac{-\zeta\xi\eta+\xi^{2}g+\eta^{2}f}{\zeta^{2}-4fg}\right],$
(13)
whose convergent condition is Re$\left(\zeta\pm f\pm g\right)<0,\
$Re($\zeta^{2}-4fg)/(\zeta\pm f\pm g)<0.$ Eq.(11) is very important for later
calculation of photon-number distribution (PND) and normalization of photon-
added (-subtracted) squeezed thermal states (PASTS, PSSTS).
It is interesting to notice that the standard generating function of Legendre
polynomials 25 is given by
$\frac{1}{\sqrt{1-2xt+t^{2}}}=\sum_{m=0}^{\infty}P_{m}\left(x\right)t^{m},$
(14)
thus comparing Eq.(11) with Eq.(14) we find
$\displaystyle\left\langle e^{fa^{\dagger}a}\right\rangle$
$\displaystyle=A^{-1/2}\left[\frac{C}{A}e^{2f}-2\frac{B}{A}e^{f}+1\right]^{-1/2}$
$\displaystyle=A^{-1/2}\sum_{m=0}^{\infty}P_{m}\left(B/\sqrt{AC}\right)\left(\sqrt{C/A}e^{f}\right)^{m},$
(15)
which indicates that the average value of $e^{fa^{{\dagger}}a}$ under squeezed
thermal state (STS) is just the generating function of Legendre polynomials, a
remarkable result. Next, we shall examine the normalizations and PNDs of PASTS
and PSSTS by using Eqs.(11) and (15).
## IV Normalization and PND of PASTS
The $m$-photon-added scheme, denoted by the mapping $\rho\rightarrow
a^{{\dagger}m}\rho a^{m},$ was first proposed by Agarwal and Tara 2 . Here, we
introduce the PASTS. Theoretically, the PASTS can be obtained by repeatedly
operating the photon creation operator $a^{\dagger}$ on a STS, so its density
operator is given by
$\rho_{ad}=C_{a,m}^{-1}a^{{\dagger}m}S_{1}\rho_{th}S_{1}^{\dagger}a^{m},$ (16)
where $m$ is the added photon number (a non-negative integer), $C_{a,m}^{-1}$
is the normalization constant to be determined.
### IV.1 Normalization of PASTS
To fully describe a quantum state, its normalization is usually necessary.
Next, we shall employ the fact (5) and (11), (15) to realize our aim.
According to the normalization condition tr$\rho_{ad}=1$ and the TFD, we have
$\displaystyle C_{a,m}$
$\displaystyle=\mathtt{tr}\left[a^{{\dagger}m}S_{1}\rho_{th}S_{1}^{\dagger}a^{m}\right]$
$\displaystyle=\mathtt{Tr}\left[a^{{\dagger}m}S_{1}\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{m}\right]$ $\displaystyle=\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{m}a^{{\dagger}m}S_{1}\left|0(\beta)\right\rangle,$
(17)
which implies that the calculation of normaliation factor $C_{a,m}$ is
converted to a matrix element after introducing the thermal vacuum state
$\left|0(\beta)\right\rangle$.
Note the operator identity 24 $e^{\tau
a^{{\dagger}}a}=e^{-\tau}\vdots\exp[(1-e^{-\tau})a^{{\dagger}}a]\vdots$, we
see
$\sum_{m=0}^{\infty}\frac{\tau^{m}}{m!}a^{m}a^{{\dagger}m}=\vdots e^{\tau
a^{{\dagger}}a}\vdots=\left(\frac{1}{1-\tau}\right)^{a^{{\dagger}}a+1},$ (18)
where the symbol $\vdots$ $\vdots$ denotes antinormally ordering. Thus using
Eqs.(11), (17) and (18), ($e^{f}\rightarrow\frac{1}{1-\tau}$) we have
$\displaystyle\sum_{m=0}^{\infty}\frac{\tau^{m}}{m!}C_{a,m}$
$\displaystyle=\frac{1}{1-\tau}\left\langle
0(\beta)\right|S_{1}^{\dagger}e^{a^{{\dagger}}a\ln\frac{1}{1-\tau}}S_{1}\left|0(\beta)\right\rangle$
$\displaystyle=\left[A\tau^{2}-2D\tau+1\right]^{-1/2},$ (19)
where we have set
$\displaystyle D$ $\displaystyle=\cosh^{2}\theta\cosh 2r-\sinh^{2}r$
$\displaystyle=n_{c}\cosh 2r+\cosh^{2}r.$ (20)
Comparing Eq.(19) with Eq.(14), and taking
$\tau^{\prime}\rightarrow\sqrt{A}\tau$, we obtain
$\displaystyle\sum_{m=0}^{\infty}\tau^{\prime m}\frac{C_{a,m}}{m!A^{m/2}}$
$\displaystyle=\left[\tau^{\prime 2}-2D/\sqrt{A}\tau^{\prime}+1\right]^{-1/2}$
$\displaystyle=\sum_{m=0}^{\infty}P_{m}\left(D/\sqrt{A}\right)\tau^{\prime
m},$ (21)
thus the normalization constant of PASTS is given by
$C_{a,m}=m!A^{m/2}P_{m}\left(D/\sqrt{A}\right),$ (22)
which is identical with the result in Ref.26 . It is noted that, for the case
of no-photon-addition with $m=0$, $C_{a,0}=1$ as expected. Under the case of
$m$-photon-addition thermal state (no squeezing) with $D=\allowbreak n_{c}+1$,
$A=\allowbreak\left(n_{c}+1\right)^{2},$ and $P_{m}\left(1\right)=1$, then
$C_{a,m}=m!\left(n_{c}+1\right)^{m}.$ The same result as Eq.(32) found in
Ref.27 . In addition, when $r=0$ corresponding to photon-added thermal state,
Eq.(22) just reduces to $C_{a,m}=m!\cosh^{2m}\theta$ 27 .
### IV.2 PND of PASTS
The photon-number distribution (PND) is a key characteristic of every optical
field. The PND, i.e., the probability of finding $n$ photons in a quantum
state described by the density operator $\rho$, is
$\mathcal{P}(n)=\mathtt{tr}\left[\left|n\right\rangle\left\langle
n\right|\rho\right]$. In a similar spirit of deriving Eq.(22), noting
$a^{m}\left|n\right\rangle=\sqrt{n!/(n-m)!}\left|n-m\right\rangle$ and
$\left|n\right\rangle=a^{\dagger n}/\sqrt{n!}\left|0\right\rangle,$ the PND of
the PASTS can be calculated as
$\displaystyle\mathcal{P}_{a}(n)$
$\displaystyle=C_{a,m}^{-1}\mathtt{tr}\left[\left|n\right\rangle\left\langle
n\right|a^{{\dagger}m}S_{1}\rho_{th}S_{1}^{\dagger}a^{m}\right]$
$\displaystyle=C_{a,m}^{-1}\mathtt{Tr}\left[\left|n\right\rangle\left\langle
n\right|a^{{\dagger}m}S_{1}\left|0(\beta)\right\rangle\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{m}\right]$
$\displaystyle=\frac{n!C_{a,m}^{-1}}{l!}\left\langle
0(\beta)\right|S_{1}^{\dagger}\left|l\right\rangle\left\langle
l\right|S_{1}\left|0(\beta)\right\rangle$
$\displaystyle=\frac{n!C_{a,m}^{-1}}{\left(l!\right)^{2}}\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{\dagger l}\left|0\right\rangle\left\langle
0\right|a^{l}S_{1}\left|0(\beta)\right\rangle$ (23)
which leads to
$\displaystyle\sum_{l=0}^{\infty}\tau^{l}\frac{l!}{n!}C_{a,m}\mathcal{P}_{a}(n)$
$\displaystyle=\sum_{l=0}^{\infty}\frac{\tau^{l}}{l!}\left\langle
0(\beta)\right|S_{1}^{\dagger}\colon\left(a^{\dagger}a\right)^{l}e^{-a^{\dagger}a}\colon
S_{1}\left|0(\beta)\right\rangle$ $\displaystyle=\left\langle
0(\beta)\right|S_{1}^{\dagger}e^{a^{\dagger}a\ln\tau}S_{1}\left|0(\beta)\right\rangle,$
(24)
where $l=n-m$ and the vacuum projector operator
$\left|0\right\rangle\left\langle 0\right|=\colon e^{-a^{\dagger}a}\colon$ and
the operator identity $e^{\lambda
a^{\dagger}a}=\colon\exp[(e^{\lambda}-1)a^{\dagger}a]\colon$ are used.
Using Eq.(11) again ($e^{f}\rightarrow\tau$) and comparing Eq.(24) with Eq.
(14) we see
$\displaystyle\sum_{l=0}^{\infty}\tau^{l}\frac{l!}{n!}C_{a,m}\mathcal{P}_{a}(n)$
$\displaystyle=A^{-1/2}\left[\frac{C}{A}\tau^{2}-2\frac{B}{A}\tau+1\right]^{-1/2}$
$\displaystyle=A^{-1/2}\sum_{l=0}^{\infty}P_{l}\left(B/\sqrt{AC}\right)\left(\sqrt{C/A}\tau\right)^{l},$
(25)
which leads to the PND of PASTS
$\mathcal{P}_{a}(n)=\frac{n!C_{a,m}^{-1}\left(C/A\right)^{\left(n-m\right)/2}}{\left(n-m\right)!\sqrt{A}}P_{n-m}\left(B/\sqrt{AC}\right),$
(26)
a Legendre polynomial with a condition $n\geqslant m$ which implies that the
photon-number ($n$) involved in PASTS is always no-less than the photon-number
($m$) operated on the STS, and there is no photon distribution when $n<m$). It
is obvious that when $m=0$ corresponding to the STS, then the PND of STS is
also a Legendre distribution 28 .
## V Normalization and PND of PSSTS
Next, we turn our attention to discussing the photon-subtracted squeezed
thermal state (PSSTS), defined as
$\rho_{sb}=C_{s,m}^{-1}a^{m}S_{1}\rho_{th}S_{1}^{\dagger}a^{{\dagger}m},$ (27)
where $m$ is the subtracted photon number (a non-negative integer), and
$C_{s,m}$ is a normalized constant.
In a similar way to deriving Eq.(22), we have
$C_{s,m}=\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{{\dagger}m}a^{m}S_{1}\left|0(\beta)\right\rangle,$
(28)
so employing $e^{\lambda
a^{\dagger}a}=\colon\exp[(e^{\lambda}-1)a^{\dagger}a]\colon$ and Eq.(11)
($e^{f}\rightarrow 1+\tau$) we see
$\displaystyle\sum_{m=0}^{\infty}\frac{\tau^{m}}{m!}C_{s,m}$
$\displaystyle=\left\langle
0(\beta)\right|S_{1}^{\dagger}e^{a^{{\dagger}}a\ln(1+\tau)}S_{1}\left|0(\beta)\right\rangle$
$\displaystyle=\left[C\tau^{2}-2E\tau+1\right]^{-1/2},$ (29)
where
$\displaystyle E$ $\displaystyle=\cosh 2r\cosh^{2}\theta-\cosh^{2}r$
$\displaystyle=\frac{1}{2}\left[(2n_{c}+1)\cosh 2r-1\right].$ (30)
Comparing Eq.(29) with Eq.(14) yields
$C_{s,m}=m!C^{m/2}P_{m}\left(E/\sqrt{C}\right),$ (31)
which is the normalization factor of PSSTS. When $r=0$ corresponding to
photon-subtracted thermal state, Eq.(31) just reduces to
$C_{s,m}=m!\sinh^{2m}\theta$ 27 .
Using the same procession as obtaining Eq.(26), the PND of PSSTS is given by
$\displaystyle\mathcal{P}_{s}(n)$ $\displaystyle=C_{s,m}^{-1}\left\langle
0(\beta)\right|S_{1}^{\dagger}a^{{\dagger}m}\left|n\right\rangle\left\langle
n\right|a^{m}S_{1}\left|0(\beta)\right\rangle$
$\displaystyle=\frac{1}{n!}C_{s,m}^{-1}\left\langle
0(\beta)\right|S_{1}^{\dagger}\colon
a^{{\dagger}m+n}a^{m+n}e^{-a^{{\dagger}}a}\colon
S_{1}\left|0(\beta)\right\rangle,$ (32)
so $\left(k=m+n\right)$
$\displaystyle\sum_{k=0}^{\infty}\frac{\tau^{k}}{k!}n!C_{s,m}\mathcal{P}_{s}(n)$
$\displaystyle=\left\langle
0(\beta)\right|S_{1}^{\dagger}e^{a^{{\dagger}}a\ln\tau}S_{1}\left|0(\beta)\right\rangle$
$\displaystyle=\text{R.H.S of (\ref{1.21}),}$ (33)
whcih leads to the PND of PSSTS
$\mathcal{P}_{s}(n)=\frac{\left(m+n\right)!}{n!C_{s,m}\sqrt{A}}\left(C/A\right)^{m+n/2}P_{m+n}\left(B/\sqrt{AC}\right),$
(34)
a Legendre polynomial, which is same as the result of Ref.28 .
## VI Conclusion
In this paper, we present a new concise approach for normalizing m-photon-
added (-subtracted) STS and for deriving the PNDs, which improve the method
used in Refs. 26 ; 28 . That is to say, using the thermal field dynamics
theory, we convert the thermal state to a pure state in doubled Fock space in
which the calculations of ensemble averages under a mixed state $\rho$,
$\left\langle
A\right\rangle=\mathtt{tr}\left(A\rho\right)/\mathtt{tr}\left(\rho\right)\ $is
replaced by an equivalent expectation values with a pure state
$\left|0(\beta)\right\rangle$, i.e., $\left\langle A\right\rangle=\left\langle
0(\beta)\right|A\left|0(\beta)\right\rangle$. It is shown that the average
value of $e^{fa^{{\dagger}}a}$ under STS is just the generating function of
Legendre polynomials, a remarkable result. Based on this point, the
normalization and PNDs of m-photon added (or subtracted) STS are easily
obtained as the Legendre polynomials. The generating function of the Legendre
polynomials and the average value of $e^{fa^{{\dagger}}a}$ under STS are used
in the whole calculation.
ACKNOWLEDGEMENTS: Work supported by the National Natural Science Foundation of
China (Grant Nos. 11047133, 60978009), the Major Research Plan of the National
Natural Science Foundation of China (Grant No. 91121023), and the “973”
Project (Grant No. 2011CBA00200), as well as the Natural Science Foundation of
Jiangxi Province of China (No. 2010GQW0027).
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|
arxiv-papers
| 2011-12-05T07:55:27 |
2024-09-04T02:49:24.962879
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Li-Yun Hu and Zhi-Ming Zhang",
"submitter": "Liyun Hu",
"url": "https://arxiv.org/abs/1112.0854"
}
|
1112.0938
|
The LHCb collaboration
# Evidence for $C\\!P$ violation in time-integrated $\bm{D^{0}\rightarrow
h^{-}h^{+}}$ decay rates
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Kerzel37, T. Ketel24, A. Keune38, B. Khanji6, Y.M. Kim46, M. Knecht38, R.
Koopman24, P. Koppenburg23, A. Kozlinskiy23, L. Kravchuk32, K. Kreplin11, M.
Kreps44, G. Krocker11, P. Krokovny11, F. Kruse9, K. Kruzelecki37, M.
Kucharczyk20,25,37,j, T. Kvaratskheliya30,37, V.N. La Thi38, D. Lacarrere37,
G. Lafferty50, A. Lai15, D. Lambert46, R.W. Lambert24, E. Lanciotti37, G.
Lanfranchi18, C. Langenbruch11, T. Latham44, C. Lazzeroni55, R. Le Gac6, J.
van Leerdam23, J.-P. Lees4, R. Lefèvre5, A. Leflat31,37, J. Lefrançois7, O.
Leroy6, T. Lesiak25, L. Li3, L. Li Gioi5, M. Lieng9, M. Liles48, R. Lindner37,
C. Linn11, B. Liu3, G. Liu37, J. von Loeben20, J.H. Lopes2, E. Lopez Asamar35,
N. Lopez-March38, H. Lu38,3, J. Luisier38, A. Mac Raighne47, F. Machefert7,
I.V. Machikhiliyan4,30, F. Maciuc10, O. Maev29,37, J. Magnin1, S. Malde51,
R.M.D. Mamunur37, G. Manca15,d, G. Mancinelli6, N. Mangiafave43, U. Marconi14,
R. Märki38, J. Marks11, G. Martellotti22, A. Martens8, L. Martin51, A. Martín
Sánchez7, D. Martinez Santos37, A. Massafferri1, Z. Mathe12, C. Matteuzzi20,
M. Matveev29, E. Maurice6, B. Maynard52, A. Mazurov16,32,37, G. McGregor50, R.
McNulty12, M. Meissner11, M. Merk23, J. Merkel9, R. Messi21,k, S.
Miglioranzi37, D.A. Milanes13,37, M.-N. Minard4, J. Molina Rodriguez54, S.
Monteil5, D. Moran12, P. Morawski25, R. Mountain52, I. Mous23, F. Muheim46, K.
Müller39, R. Muresan28,38, B. Muryn26, B. Muster38, M. Musy35, J. Mylroie-
Smith48, P. Naik42, T. Nakada38, R. Nandakumar45, I. Nasteva1, M. Nedos9, M.
Needham46, N. Neufeld37, C. Nguyen-Mau38,o, M. Nicol7, V. Niess5, N.
Nikitin31, A. Nomerotski51, A. Novoselov34, A. Oblakowska-Mucha26, V.
Obraztsov34, S. Oggero23, S. Ogilvy47, O. Okhrimenko41, R. Oldeman15,d, M.
Orlandea28, J.M. Otalora Goicochea2, P. Owen49, K. Pal52, J. Palacios39, A.
Palano13,b, M. Palutan18, J. Panman37, A. Papanestis45, M. Pappagallo47, C.
Parkes50,37, C.J. Parkinson49, G. Passaleva17, G.D. Patel48, M. Patel49, S.K.
Paterson49, G.N. Patrick45, C. Patrignani19,i, C. Pavel-Nicorescu28, A. Pazos
Alvarez36, A. Pellegrino23, G. Penso22,l, M. Pepe Altarelli37, S.
Perazzini14,c, D.L. Perego20,j, E. Perez Trigo36, A. Pérez-Calero Yzquierdo35,
P. Perret5, M. Perrin-Terrin6, G. Pessina20, A. Petrella16,37, A.
Petrolini19,i, A. Phan52, E. Picatoste Olloqui35, B. Pie Valls35, B.
Pietrzyk4, T. Pilař44, D. Pinci22, R. Plackett47, S. Playfer46, M. Plo
Casasus36, G. Polok25, A. Poluektov44,33, E. Polycarpo2, D. Popov10, B.
Popovici28, C. Potterat35, A. Powell51, J. Prisciandaro38, V. Pugatch41, A.
Puig Navarro35, W. Qian52, J.H. Rademacker42, B. Rakotomiaramanana38, M.S.
Rangel2, I. Raniuk40, G. Raven24, S. Redford51, M.M. Reid44, A.C. dos Reis1,
S. Ricciardi45, K. Rinnert48, D.A. Roa Romero5, P. Robbe7, E. Rodrigues47,50,
F. Rodrigues2, P. Rodriguez Perez36, G.J. Rogers43, S. Roiser37, V.
Romanovsky34, M. Rosello35,n, J. Rouvinet38, T. Ruf37, H. Ruiz35, G.
Sabatino21,k, J.J. Saborido Silva36, N. Sagidova29, P. Sail47, B. Saitta15,d,
C. Salzmann39, M. Sannino19,i, R. Santacesaria22, C. Santamarina Rios36, R.
Santinelli37, E. Santovetti21,k, M. Sapunov6, A. Sarti18,l, C. Satriano22,m,
A. Satta21, M. Savrie16,e, D. Savrina30, P. Schaack49, M. Schiller24, S.
Schleich9, M. Schlupp9, M. Schmelling10, B. Schmidt37, O. Schneider38, A.
Schopper37, M.-H. Schune7, R. Schwemmer37, B. Sciascia18, A. Sciubba18,l, M.
Seco36, A. Semennikov30, K. Senderowska26, I. Sepp49, N. Serra39, J. Serrano6,
P. Seyfert11, M. Shapkin34, I. Shapoval40,37, P. Shatalov30, Y. Shcheglov29,
T. Shears48, L. Shekhtman33, O. Shevchenko40, V. Shevchenko30, A. Shires49, R.
Silva Coutinho44, T. Skwarnicki52, A.C. Smith37, N.A. Smith48, E. Smith51,45,
K. Sobczak5, F.J.P. Soler47, A. Solomin42, F. Soomro18, B. Souza De Paula2, B.
Spaan9, A. Sparkes46, P. Spradlin47, F. Stagni37, S. Stahl11, O. Steinkamp39,
S. Stoica28, S. Stone52,37, B. Storaci23, M. Straticiuc28, U. Straumann39,
V.K. Subbiah37, S. Swientek9, M. Szczekowski27, P. Szczypka38, T. Szumlak26,
S. T’Jampens4, E. Teodorescu28, F. Teubert37, C. Thomas51, E. Thomas37, J. van
Tilburg11, V. Tisserand4, M. Tobin39, S. Topp-Joergensen51, N. Torr51, E.
Tournefier4,49, M.T. Tran38, A. Tsaregorodtsev6, N. Tuning23, M. Ubeda
Garcia37, A. Ukleja27, P. Urquijo52, U. Uwer11, V. Vagnoni14, G. Valenti14, R.
Vazquez Gomez35, P. Vazquez Regueiro36, S. Vecchi16, J.J. Velthuis42, M.
Veltri17,g, B. Viaud7, I. Videau7, X. Vilasis-Cardona35,n, J. Visniakov36, A.
Vollhardt39, D. Volyanskyy10, D. Voong42, A. Vorobyev29, H. Voss10, S.
Wandernoth11, J. Wang52, D.R. Ward43, N.K. Watson55, A.D. Webber50, D.
Websdale49, M. Whitehead44, D. Wiedner11, L. Wiggers23, G. Wilkinson51, M.P.
Williams44,45, M. Williams49, F.F. Wilson45, J. Wishahi9, M. Witek25, W.
Witzeling37, S.A. Wotton43, K. Wyllie37, Y. Xie46, F. Xing51, Z. Xing52, Z.
Yang3, R. Young46, O. Yushchenko34, M. Zavertyaev10,a, F. Zhang3, L. Zhang52,
W.C. Zhang12, Y. Zhang3, A. Zhelezov11, L. Zhong3, E. Zverev31, A. Zvyagin37.
1Centro Brasileiro de Pesquisas Físicas (CBPF), Rio de Janeiro, Brazil
2Universidade Federal do Rio de Janeiro (UFRJ), Rio de Janeiro, Brazil
3Center for High Energy Physics, Tsinghua University, Beijing, China
4LAPP, Université de Savoie, CNRS/IN2P3, Annecy-Le-Vieux, France
5Clermont Université, Université Blaise Pascal, CNRS/IN2P3, LPC, Clermont-
Ferrand, France
6CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
7LAL, Université Paris-Sud, CNRS/IN2P3, Orsay, France
8LPNHE, Université Pierre et Marie Curie, Université Paris Diderot,
CNRS/IN2P3, Paris, France
9Fakultät Physik, Technische Universität Dortmund, Dortmund, Germany
10Max-Planck-Institut für Kernphysik (MPIK), Heidelberg, Germany
11Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg, Heidelberg,
Germany
12School of Physics, University College Dublin, Dublin, Ireland
13Sezione INFN di Bari, Bari, Italy
14Sezione INFN di Bologna, Bologna, Italy
15Sezione INFN di Cagliari, Cagliari, Italy
16Sezione INFN di Ferrara, Ferrara, Italy
17Sezione INFN di Firenze, Firenze, Italy
18Laboratori Nazionali dell’INFN di Frascati, Frascati, Italy
19Sezione INFN di Genova, Genova, Italy
20Sezione INFN di Milano Bicocca, Milano, Italy
21Sezione INFN di Roma Tor Vergata, Roma, Italy
22Sezione INFN di Roma La Sapienza, Roma, Italy
23Nikhef National Institute for Subatomic Physics, Amsterdam, The Netherlands
24Nikhef National Institute for Subatomic Physics and Vrije Universiteit,
Amsterdam, The Netherlands
25Henryk Niewodniczanski Institute of Nuclear Physics Polish Academy of
Sciences, Kraców, Poland
26AGH University of Science and Technology, Kraców, Poland
27Soltan Institute for Nuclear Studies, Warsaw, Poland
28Horia Hulubei National Institute of Physics and Nuclear Engineering,
Bucharest-Magurele, Romania
29Petersburg Nuclear Physics Institute (PNPI), Gatchina, Russia
30Institute of Theoretical and Experimental Physics (ITEP), Moscow, Russia
31Institute of Nuclear Physics, Moscow State University (SINP MSU), Moscow,
Russia
32Institute for Nuclear Research of the Russian Academy of Sciences (INR RAN),
Moscow, Russia
33Budker Institute of Nuclear Physics (SB RAS) and Novosibirsk State
University, Novosibirsk, Russia
34Institute for High Energy Physics (IHEP), Protvino, Russia
35Universitat de Barcelona, Barcelona, Spain
36Universidad de Santiago de Compostela, Santiago de Compostela, Spain
37European Organization for Nuclear Research (CERN), Geneva, Switzerland
38Ecole Polytechnique Fédérale de Lausanne (EPFL), Lausanne, Switzerland
39Physik-Institut, Universität Zürich, Zürich, Switzerland
40NSC Kharkiv Institute of Physics and Technology (NSC KIPT), Kharkiv, Ukraine
41Institute for Nuclear Research of the National Academy of Sciences (KINR),
Kyiv, Ukraine
42H.H. Wills Physics Laboratory, University of Bristol, Bristol, United
Kingdom
43Cavendish Laboratory, University of Cambridge, Cambridge, United Kingdom
44Department of Physics, University of Warwick, Coventry, United Kingdom
45STFC Rutherford Appleton Laboratory, Didcot, United Kingdom
46School of Physics and Astronomy, University of Edinburgh, Edinburgh, United
Kingdom
47School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom
48Oliver Lodge Laboratory, University of Liverpool, Liverpool, United Kingdom
49Imperial College London, London, United Kingdom
50School of Physics and Astronomy, University of Manchester, Manchester,
United Kingdom
51Department of Physics, University of Oxford, Oxford, United Kingdom
52Syracuse University, Syracuse, NY, United States
53CC-IN2P3, CNRS/IN2P3, Lyon-Villeurbanne, France, associated member
54Pontifícia Universidade Católica do Rio de Janeiro (PUC-Rio), Rio de
Janeiro, Brazil, associated to 2
55University of Birmingham, Birmingham, United Kingdom
56Physikalisches Institut, Universität Rostock, Rostock, Germany, associated
to 11
aP.N. Lebedev Physical Institute, Russian Academy of Science (LPI RAS),
Moscow, Russia
bUniversità di Bari, Bari, Italy
cUniversità di Bologna, Bologna, Italy
dUniversità di Cagliari, Cagliari, Italy
eUniversità di Ferrara, Ferrara, Italy
fUniversità di Firenze, Firenze, Italy
gUniversità di Urbino, Urbino, Italy
hUniversità di Modena e Reggio Emilia, Modena, Italy
iUniversità di Genova, Genova, Italy
jUniversità di Milano Bicocca, Milano, Italy
kUniversità di Roma Tor Vergata, Roma, Italy
lUniversità di Roma La Sapienza, Roma, Italy
mUniversità della Basilicata, Potenza, Italy
nLIFAELS, La Salle, Universitat Ramon Llull, Barcelona, Spain
oHanoi University of Science, Hanoi, Viet Nam
###### Abstract
A search for time-integrated $C\\!P$ violation in $D^{0}\rightarrow
h^{-}h^{+}$ ($h=K$, $\pi$) decays is presented using 0.62 $\mbox{\,fb}^{-1}$
of data collected by LHCb in 2011. The flavor of the charm meson is determined
by the charge of the slow pion in the $D^{*+}\rightarrow D^{0}\pi^{+}$ and
$D^{*-}\rightarrow\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}\pi^{-}$
decay chains. The difference in $C\\!P$ asymmetry between $D^{0}\rightarrow
K^{-}K^{+}$ and $D^{0}\rightarrow\pi^{-}\pi^{+}$, $\Delta A_{C\\!P}\equiv
A_{C\\!P}(K^{-}K^{+})\,-\,A_{C\\!P}(\pi^{-}\pi^{+})$, is measured to be
$\left[-0.82\pm 0.21(\mathrm{stat.})\pm 0.11(\mathrm{syst.})\right]\%$. This
differs from the hypothesis of $C\\!P$ conservation by $3.5$ standard
deviations.
###### pacs:
11.30.Er, 13.25.Ft
EUROPEAN ORGANIZATION FOR NUCLEAR RESEARCH (CERN)
| |
---|---|---
| | LHCb-PAPER-2011-023
| | CERN-PH-EP-2011-208
The charm sector is a promising place to probe for the effects of physics
beyond the Standard Model (SM). There has been a resurgence of interest in the
past few years since evidence for $D^{0}$ mixing was first seen
bib:babar_mixing_moriond ; bib:belle_mixing_moriond . Mixing is now well-
established bib:hfag at a level which is consistent with, but at the upper
end of, SM expectations falk_grossman_ligeti_nir_petrov . By contrast, no
evidence for $C\\!P$ violation in charm decays has yet been found.
The time-dependent $C\\!P$ asymmetry $A_{C\\!P}(f;\,t)$ for $D^{0}$ decays to
a CP eigenstate $f$ (with $f=\bar{f}$) is defined as
$A_{C\\!P}(f;\,t)\equiv\frac{\Gamma(D^{0}(t)\rightarrow f)-\Gamma(\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}(t)\rightarrow
f)}{\Gamma(D^{0}(t)\rightarrow f)+\Gamma(\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}(t)\rightarrow f)},$ (1)
where $\Gamma$ is the decay rate for the process indicated. In general
$A_{C\\!P}(f;\,t)$ depends on $f$. For $f=K^{-}K^{+}$ and $f=\pi^{-}\pi^{+}$,
$A_{C\\!P}(f;\,t)$ can be expressed in terms of two contributions: a direct
component associated with $C\\!P$ violation in the decay amplitudes, and an
indirect component associated with $C\\!P$ violation in the mixing or in the
interference between mixing and decay. In the limit of U-spin symmetry, the
direct component is equal in magnitude and opposite in sign for $K^{-}K^{+}$
and $\pi^{-}\pi^{+}$, though the size of U-spin breaking effects remains to be
quantified precisely bib:grossman_kagan_nir . The magnitudes of $C\\!P$
asymmetries in decays to these final states are expected to be small in the SM
bib:cicerone ; bib:lenz ; bib:grossman_kagan_nir ; bib:petrov , with
predictions of up to $\mathcal{O}(10^{-3})$. However, beyond the SM the rate
of $C\\!P$ violation could be enhanced bib:grossman_kagan_nir ;
bib:littlest_higgs .
The asymmetry $A_{C\\!P}(f;\,t)$ may be written to first order as
bib:cdf_paper ; bib:bigi_d2hh
$A_{C\\!P}(f;\,t)=a^{\mathrm{dir}}_{C\\!P}(f)\,+\,\frac{t}{\tau}a^{\mathrm{ind}}_{C\\!P},$
(2)
where $a^{\mathrm{dir}}_{C\\!P}(f)$ is the direct $C\\!P$ asymmetry, $\tau$ is
the $D^{0}$ lifetime, and $a^{\mathrm{ind}}_{C\\!P}$ is the indirect $C\\!P$
asymmetry. To a good approximation this latter quantity is universal
bib:grossman_kagan_nir ; bib:kagan_sokoloff . The time-integrated asymmetry
measured by an experiment, $A_{C\\!P}(f)$, depends upon the time-acceptance of
that experiment. It can be written as
$A_{C\\!P}(f)=a^{\mathrm{dir}}_{C\\!P}(f)\,+\,\frac{\langle
t\rangle}{\tau}a^{\mathrm{ind}}_{C\\!P},$ (3)
where $\langle t\rangle$ is the average decay time in the reconstructed
sample. Denoting by $\Delta$ the differences between quantities for
$D^{0}\rightarrow K^{-}K^{+}$ and $D^{0}\rightarrow\pi^{-}\pi^{+}$ it is then
possible to write
$\displaystyle\Delta A_{C\\!P}$ $\displaystyle\equiv$ $\displaystyle
A_{C\\!P}(K^{-}K^{+})\,-\,A_{C\\!P}(\pi^{-}\pi^{+})$ $\displaystyle=$
$\displaystyle\left[a^{\mathrm{dir}}_{C\\!P}(K^{-}K^{+})\,-\,a^{\mathrm{dir}}_{C\\!P}(\pi^{-}\pi^{+})\right]\,+\,\frac{\Delta\langle
t\rangle}{\tau}a^{\mathrm{ind}}_{C\\!P}.$
In the limit that $\Delta\langle t\rangle$ vanishes, $\Delta A_{C\\!P}$ is
equal to the difference in the direct $C\\!P$ asymmetry between the two
decays. However, if the time-acceptance is different for the $K^{-}K^{+}$ and
$\pi^{-}\pi^{+}$ final states, a contribution from indirect $C\\!P$ violation
remains.
The most precise measurements to date of the time-integrated $C\\!P$
asymmetries in $D^{0}\rightarrow K^{-}K^{+}$ and
$D^{0}\rightarrow\pi^{-}\pi^{+}$ were made by the CDF, BaBar, and Belle
collaborations bib:cdf_paper ; bib:babar_paper2008 ; bib:belle_paper2008 . The
Heavy Flavor Averaging Group (HFAG) has combined time-integrated and time-
dependent measurements of $C\\!P$ asymmetries, taking account of the different
decay time acceptances, to obtain world average values for the indirect
$C\\!P$ asymmetry of $a_{C\\!P}^{\mathrm{ind}}=(-0.03\pm 0.23)\%$ and the
difference in direct $C\\!P$ asymmetry between the final states of $\Delta
a_{C\\!P}^{\mathrm{dir}}=(-0.42\pm 0.27)\%$ bib:hfag .
In this Letter, we present a measurement of the difference in time-integrated
$C\\!P$ asymmetries between $D^{0}\rightarrow K^{-}K^{+}$ and
$D^{0}\rightarrow\pi^{-}\pi^{+}$, performed with 0.62 $\mbox{\,fb}^{-1}$ of
data collected at LHCb between March and June 2011. The flavor of the initial
state ($D^{0}$ or $\kern 1.99997pt\overline{\kern-1.99997ptD}{}^{0}$) is
tagged by requiring a $D^{*+}\rightarrow D^{0}\pi_{\mathrm{s}}^{+}$ decay,
with the flavor determined by the charge of the slow pion
($\pi_{\mathrm{s}}^{+}$). The inclusion of charge-conjugate modes is implied
throughout, except in the definition of asymmetries.
The raw asymmetry for tagged $D^{0}$ decays to a final state $f$ is given by
$A_{\mathrm{raw}}(f)$, defined as
$A_{\mathrm{raw}}(f)\equiv\frac{N(D^{*+}\rightarrow
D^{0}(f)\pi_{s}^{+})\,-\,N(D^{*-}\rightarrow\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}(f)\pi_{s}^{-})}{N(D^{*+}\rightarrow
D^{0}(f)\pi_{s}^{+})\,+\,N(D^{*-}\rightarrow\kern
1.99997pt\overline{\kern-1.99997ptD}{}^{0}(f)\pi_{s}^{-})},$ (5)
where $N(X)$ refers to the number of reconstructed events of decay $X$ after
background subtraction.
To first order the raw asymmetries may be written as a sum of four components,
due to physics and detector effects:
$A_{\mathrm{raw}}(f)=A_{C\\!P}(f)\,+\,A_{\mathrm{D}}(f)\,+\,A_{\mathrm{D}}(\pi_{\mathrm{s}}^{+})\,+\,A_{\mathrm{P}}(D^{*+}).$
(6)
Here, $A_{\mathrm{D}}(f)$ is the asymmetry in selecting the $D^{0}$ decay into
the final state $f$, $A_{\mathrm{D}}(\pi_{\mathrm{s}}^{+})$ is the asymmetry
in selecting the slow pion from the $D^{*+}$ decay chain, and
$A_{\mathrm{P}}(D^{*+})$ is the production asymmetry for $D^{*+}$ mesons. The
asymmetries $A_{\mathrm{D}}$ and $A_{\mathrm{P}}$ are defined in the same
fashion as $A_{\mathrm{raw}}$. The first-order expansion is valid since the
individual asymmetries are small.
For a two-body decay of a spin-0 particle to a self-conjugate final state
there can be no $D^{0}$ detection asymmetry, i.e.
$A_{\mathrm{D}}(K^{-}K^{+})=A_{\mathrm{D}}(\pi^{-}\pi^{+})=0.$ Moreover,
$A_{\mathrm{D}}(\pi_{\mathrm{s}}^{+})$ and $A_{\mathrm{P}}(D^{*+})$ are
independent of $f$ and thus in the first-order expansion of equation 5 those
terms cancel in the difference
$A_{\mathrm{raw}}(K^{-}K^{+})\,-\,A_{\mathrm{raw}}(\pi^{-}\pi^{+})$, resulting
in
$\Delta
A_{C\\!P}=A_{\mathrm{raw}}(K^{-}K^{+})\,-\,A_{\mathrm{raw}}(\pi^{-}\pi^{+}).$
(7)
To minimize second-order effects that are related to the slightly different
kinematic properties of the two decay modes and that do not cancel in $\Delta
A_{C\\!P}$, the analysis is performed in bins of the relevant kinematic
variables, as discussed later.
The LHCb detector is a forward spectrometer covering the pseudorapidity range
$2<\eta<5$, and is described in detail in Ref. LHCb . The Ring Imaging
Cherenkov (RICH) detectors are of particular importance to this analysis,
providing kaon-pion discrimination for the full range of track momenta used.
The nominal downstream beam direction is aligned with the $+z$ axis, and the
field direction in the LHCb dipole is such that charged particles are
deflected in the horizontal ($xz$) plane. The field polarity was changed
several times during data taking: about 60% of the data were taken with the
down polarity and 40% with the other.
Selections are applied to provide samples of $D^{*+}\rightarrow
D^{0}\pi_{\mathrm{s}}^{+}$ candidates, with $D^{0}\rightarrow K^{-}K^{+}$ or
$\pi^{-}\pi^{+}$. Events are required to pass both hardware and software
trigger levels. A loose $D^{0}$ selection is applied in the final state of the
software trigger, and in the offline analysis only candidates that are
accepted by this trigger algorithm are considered. Both the trigger and
offline selections impose a variety of requirements on kinematics and decay
time to isolate the decays of interest, including requirements on the track
fit quality, on the $D^{0}$ and $D^{*+}$ vertex fit quality, on the transverse
momentum ($\mbox{$p_{\mathrm{T}}$}>2$ GeV$/c$) and decay time
($ct>100\,\,\upmu\rm m$) of the $D^{0}$ candidate, on the angle between the
$D^{0}$ momentum in the lab frame and its daughter momenta in the $D^{0}$ rest
frame ($|\cos\theta|<0.9$), that the $D^{0}$ trajectory points back to a
primary vertex, and that the $D^{0}$ daughter tracks do not. In addition, the
offline analysis exploits the capabilities of the RICH system to distinguish
between pions and kaons when reconstructing the $D^{0}$ meson, with no tracks
appearing as both pion and kaon candidates.
A fiducial region is implemented by imposing the requirement that the slow
pion lies within the central part of the detector acceptance. This is
necessary because the magnetic field bends pions of one charge to the left and
those of the other charge to the right. For soft tracks at large angles in the
$xz$ plane this implies that one charge is much more likely to remain within
the 300 mrad horizontal detector acceptance, thus making
$A_{\mathrm{D}}(\pi_{\mathrm{s}}^{+})$ large. Although this asymmetry is
formally independent of the $D^{0}$ decay mode, it breaks the assumption that
the raw asymmetries are small and it carries a risk of second-order systematic
effects if the ratio of efficiencies of $D^{0}\rightarrow K^{-}K^{+}$ and
$D^{0}\rightarrow\pi^{-}\pi^{+}$ varies in the affected region. The fiducial
requirements therefore exclude edge regions in the slow pion $(p_{x},p)$
plane. Similarly, a small region of phase space in which one charge of slow
pion is more likely to be swept into the beampipe region in the downstream
tracking stations, and hence has reduced efficiency, is also excluded. After
the implementation of the fiducial requirements about 70% of the events are
retained.
The invariant mass spectra of selected $K^{-}K^{+}$ and $\pi^{-}\pi^{+}$ pairs
are shown in Fig. 1. The half-width at half-maximum of the signal lineshape is
8.6 MeV$/c^{2}$ for $K^{-}K^{+}$ and 11.2 MeV$/c^{2}$ for $\pi^{-}\pi^{+}$,
where the difference is due to the kinematics of the decays and has no
relevance for the subsequent analysis. The mass difference ($\delta m$)
spectra of selected candidates, where $\delta m\equiv
m(h^{-}h^{+}\pi_{\mathrm{s}}^{+})-m(h^{-}h^{+})-m(\pi^{+})$ for $h=K,\pi$, are
shown in Fig. 2.
Figure 1: Fits to the (a) $m(K^{-}K^{+})$ and (b) $m(\pi^{-}\pi^{+})$ spectra
of $D^{*+}$ candidates passing the selection and satisfying $0<\delta m<15$
MeV$/c^{2}$. The dashed line corresponds to the background component in the
fit, and the vertical lines indicate the signal window of 1844–1884
MeV$/c^{2}$.
Candidates are required to lie inside a wide $\delta m$ window of 0–15
MeV$/c^{2}$, and in Fig. 2 and for all subsequent results candidates are in
addition required to lie in a mass signal window of 1844–1884 MeV$/c^{2}$. The
$D^{*+}$ signal yields are approximately $1.44\times 10^{6}$ in the
$K^{-}K^{+}$ sample, and $0.38\times 10^{6}$ in the $\pi^{-}\pi^{+}$ sample.
Charm from $b$-hadron decays is strongly suppressed by the requirement that
the $D^{0}$ originate from a primary vertex, and accounts for only 3% of the
total yield. Of the events that contain at least one $D^{*+}$ candidate, 12%
contain more than one candidate; this is expected due to background soft pions
from the primary vertex and all candidates are accepted. The background-
subtracted average decay time of $D^{0}$ candidates passing the selection is
measured for each final state, and the fractional difference $\Delta\langle
t\rangle/\tau$ is obtained. Systematic uncertainties on this quantity are
assigned for the uncertainty on the world average $D^{0}$ lifetime $\tau$
$(0.04\%)$, charm from $b$-hadron decays $(0.18\%$), and the background-
subtraction procedure $(0.04\%)$. Combining the systematic uncertainties in
quadrature, we obtain $\Delta\langle t\rangle/\tau=\left[9.83\pm
0.22(\mathrm{stat.})\pm 0.19(\mathrm{syst.})\right]\%$. The $\pi^{-}\pi^{+}$
and $K^{-}K^{+}$ average decay time is $\langle t\rangle=\left(0.8539\pm
0.0005\right)$ ps, where the error is statistical only.
Figure 2: Fits to the $\delta m$ spectra, where the $D^{0}$ is reconstructed
in the final states (a) $K^{-}K^{+}$ and (b) $\pi^{-}\pi^{+}$, with mass lying
in the window of 1844–1884 MeV$/c^{2}$. The dashed line corresponds to the
background component in the fit.
Fits are performed on the samples in order to determine
$A_{\mathrm{raw}}(K^{-}K^{+})$ and $A_{\mathrm{raw}}(\pi^{-}\pi^{+})$. The
production and detection asymmetries can vary with $p_{\mathrm{T}}$ and
pseudorapidity $\eta$, and so can the detection efficiency of the two
different $D^{0}$ decays, in particular through the effects of the particle
identification requirements. The analysis is performed in 54 kinematic bins
defined by the $p_{\mathrm{T}}$ and $\eta$ of the $D^{*+}$ candidates, the
momentum of the slow pion, and the sign of $p_{x}$ of the slow pion at the
$D^{*+}$ vertex. The events are further partitioned in two ways. First, the
data are divided between the two dipole magnet polarities. Second, the first
60% of data are processed separately from the remainder, with the division
aligned with a break in data taking due to an LHC technical stop. In total,
216 statistically independent measurements are considered for each decay mode.
In each bin, one-dimensional unbinned maximum likelihood fits to the $\delta
m$ spectra are performed. The signal is described as the sum of two Gaussian
functions with a common mean $\mu$ but different widths $\sigma_{i}$,
convolved with a function $B(\delta m;s)=\Theta(\delta m)\,\delta m^{s}$
taking account of the asymmetric shape of the measured $\delta m$
distribution. Here, $s\simeq-0.975$ is a shape parameter fixed to the value
determined from the global fits shown in Fig. 2, $\Theta$ is the Heaviside
step function, and the convolution runs over $\delta m$. The background is
described by an empirical function of the form $1-e^{-(\delta m-\delta
m_{0})/\alpha}$, where $\delta m_{0}$ and $\alpha$ are free parameters
describing the threshold and shape of the function, respectively. The $D^{*+}$
and $D^{*-}$ samples in a given bin are fitted simultaneously and share all
shape parameters, except for a charge-dependent offset in the central value
$\mu$ and an overall scale factor in the mass resolution. The raw asymmetry in
the signal yields is extracted directly from this simultaneous fit. No fit
parameters are shared between the 216 subsamples of data, nor between the
$K^{-}K^{+}$ and $\pi^{-}\pi^{+}$ final states.
The fits do not distinguish between the signal and backgrounds that peak in
$\delta m$. Such backgrounds can arise from $D^{*+}$ decays in which the
correct slow pion is found but the $D^{0}$ is partially mis-reconstructed.
These backgrounds are suppressed by the use of tight particle identification
requirements and a narrow $D^{0}$ mass window. From studies of the $D^{0}$
mass sidebands (1820–1840 and 1890–1910 MeV$/c^{2}$), this contamination is
found to be approximately 1% of the signal yield and to have small raw
asymmetry (consistent with zero asymmetry difference between the $K^{-}K^{+}$
and $\pi^{-}\pi^{+}$ final states). Its effect on the measurement is estimated
in an ensemble of simulated experiments and found to be negligible; a
systematic uncertainty is assigned below based on the statistical precision of
the estimate.
A value of $\Delta A_{C\\!P}$ is determined in each measurement bin as the
difference between $A_{\mathrm{raw}}(K^{-}K^{+})$ and
$A_{\mathrm{raw}}(\pi^{-}\pi^{+})$. Testing these 216 measurements for mutual
consistency, we obtain $\chi^{2}/\mathrm{ndf}=211/215$ ($\chi^{2}$ probability
of 56%). A weighted average is performed to yield the result $\Delta
A_{C\\!P}=(-0.82\pm 0.21)\%$, where the uncertainty is statistical only.
Numerous robustness checks are made. The value of $\Delta A_{C\\!P}$ is
studied as a function of the time at which the data were taken (Fig. 3) and
found to be consistent with a constant value ($\chi^{2}$ probability of 57%).
The measurement is repeated with progressively more restrictive RICH particle
identification requirements, finding values of $(-0.88\pm 0.26)\%$ and
$(-1.03\pm 0.31)\%$; both of these values are consistent with the baseline
result when correlations are taken into account. Table 1 lists $\Delta
A_{C\\!P}$ for eight disjoint subsamples of data split according to magnet
polarity, the sign of $p_{x}$ of the slow pion, and whether the data were
taken before or after the technical stop. The $\chi^{2}$ probability for
consistency among the subsamples is 45%. The significances of the differences
between data taken before and after the technical stop, between the magnet
polarities, and between $p_{x}>0$ and $p_{x}<0$ are $0.4$, $0.6$, and $0.7$
standard deviations, respectively. Other checks include applying electron and
muon vetoes to the slow pion and to the $D^{0}$ daughters, use of different
kinematic binnings, validation of the size of the statistical uncertainties
with Monte Carlo pseudo-experiments, tightening of kinematic requirements,
testing for variation of the result with the multiplicity of tracks and of
primary vertices in the event, use of other signal and background
parameterizations in the fit, and imposing a full set of common shape
parameters between $D^{*+}$ and $D^{*-}$ candidates. Potential biases due to
the inclusive hardware trigger selection are investigated with the subsample
of data in which one of the signal final-state tracks is directly responsible
for the hardware trigger decision. In all cases good stability is observed.
For several of these checks, a reduced number of kinematic bins are used for
simplicity. No systematic dependence of $\Delta A_{C\\!P}$ is observed with
respect to the kinematic variables.
Figure 3: Time-dependence of the measurement. The data are divided into 19 disjoint, contiguous, time-ordered blocks and the value of $\Delta A_{C\\!P}$ measured in each block. The horizontal red dashed line shows the result for the combined sample. The vertical dashed line indicates the technical stop referred to in Table 1. Table 1: Values of $\Delta A_{C\\!P}$ measured in subsamples of the data, and the $\chi^{2}/\mathrm{ndf}$ and corresponding $\chi^{2}$ probabilities for internal consistency among the 27 bins in each subsample. The data are divided before and after a technical stop (TS), by magnet polarity (up, down), and by the sign of $p_{x}$ for the slow pion (left, right). The consistency among the eight subsamples is $\chi^{2}/\mathrm{ndf}=6.8/7$ (45%). Subsample | $\Delta A_{C\\!P}~{}[\%]$ | $\chi^{2}/\mathrm{ndf}$
---|---|---
Pre-TS, up, left | $-1.22\pm 0.59$ | $13/26~{}(98\%)$
Pre-TS, up, right | $-1.43\pm 0.59$ | $27/26~{}(39\%)$
Pre-TS, down, left | $-0.59\pm 0.52$ | $19/26~{}(84\%)$
Pre-TS, down, right | $-0.51\pm 0.52$ | $29/26~{}(30\%)$
Post-TS, up, left | $-0.79\pm 0.90$ | $26/26~{}(44\%)$
Post-TS, up, right | $+0.42\pm 0.93$ | $21/26~{}(77\%)$
Post-TS, down, left | $-0.24\pm 0.56$ | $34/26~{}(15\%)$
Post-TS, down, right | $-1.59\pm 0.57$ | $35/26~{}(12\%)$
All data | $-0.82\pm 0.21$ | $211/215~{}(56\%)$
Systematic uncertainties are assigned by: loosening the fiducial requirement
on the slow pion; assessing the effect of potential peaking backgrounds in
Monte Carlo pseudo-experiments; repeating the analysis with the asymmetry
extracted through sideband subtraction in $\delta m$ instead of a fit;
removing all candidates but one (chosen at random) in events with multiple
candidates; and comparing with the result obtained without kinematic binning.
In each case the full value of the change in result is taken as the systematic
uncertainty. These uncertainties are listed in Table 2. The sum in quadrature
is $0.11\%$. Combining statistical and systematic uncertainties in quadrature,
this result is consistent at the $1\sigma$ level with the current HFAG world
average bib:hfag .
Table 2: Summary of absolute systematic uncertainties for $\Delta A_{C\\!P}$. Source | Uncertainty
---|---
Fiducial requirement | 0.01%
Peaking background asymmetry | 0.04%
Fit procedure | 0.08%
Multiple candidates | 0.06%
Kinematic binning | 0.02%
Total | 0.11%
In conclusion, the time-integrated difference in $C\\!P$ asymmetry between
$D^{0}\rightarrow K^{-}K^{+}$ and $D^{0}\rightarrow\pi^{-}\pi^{+}$ decays has
been measured to be
$\Delta A_{C\\!P}=\left[-0.82\pm 0.21(\mathrm{stat.})\pm
0.11(\mathrm{syst.})\right]\%$
with 0.62 $\mbox{\,fb}^{-1}$ of 2011 data. Given the dependence of $\Delta
A_{C\\!P}$ on the direct and indirect $C\\!P$ asymmetries, shown in Eq.
(Evidence for $C\\!P$ violation in time-integrated $\bm{D^{0}\rightarrow
h^{-}h^{+}}$ decay rates), and the measured value $\Delta\langle
t\rangle/\tau=\left[9.83\pm 0.22(\mathrm{stat.})\pm
0.19(\mathrm{syst.})\right]\%$, the contribution from indirect $C\\!P$
violation is suppressed and $\Delta A_{C\\!P}$ is primarily sensitive to
direct $C\\!P$ violation. Dividing the central value by the sum in quadrature
of the statistical and systematic uncertainties, the significance of the
measured deviation from zero is $3.5\sigma$. This is the first evidence for
$C\\!P$ violation in the charm sector. To establish whether this result is
consistent with the SM will require the analysis of more data, as well as
improved theoretical understanding.
## Acknowledgements
We express our gratitude to our colleagues in the CERN accelerator departments
for the excellent performance of the LHC. We thank the technical and
administrative staff at CERN and at the LHCb institutes, and acknowledge
support from the National Agencies: CAPES, CNPq, FAPERJ and FINEP (Brazil);
CERN; NSFC (China); CNRS/IN2P3 (France); BMBF, DFG, HGF and MPG (Germany); SFI
(Ireland); INFN (Italy); FOM and NWO (The Netherlands); SCSR (Poland); ANCS
(Romania); MinES of Russia and Rosatom (Russia); MICINN, XuntaGal and GENCAT
(Spain); SNSF and SER (Switzerland); NAS Ukraine (Ukraine); STFC (United
Kingdom); NSF (USA). We also acknowledge the support received from the ERC
under FP7 and the Region Auvergne.
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|
arxiv-papers
| 2011-12-05T14:30:17 |
2024-09-04T02:49:24.971618
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "LHCb Collaboration: R. Aaij, C. Abellan Beteta, B. Adeva, M. Adinolfi,\n C. Adrover, A. Affolder, Z. Ajaltouni, J. Albrecht, F. Alessio, M. Alexander,\n G. Alkhazov, P. Alvarez Cartelle, A. A. Alves Jr, S. Amato, Y. Amhis, J.\n Anderson, R. B. Appleby, O. Aquines Gutierrez, F. Archilli, L. Arrabito, A.\n Artamonov, M. Artuso, E. Aslanides, G. Auriemma, S. Bachmann, J. J. Back, D.\n S. Bailey, V. Balagura, W. Baldini, R. J. Barlow, C. Barschel, S. Barsuk, W.\n Barter, A. Bates, C. Bauer, Th. Bauer, A. Bay, I. Bediaga, S. Belogurov, K.\n Belous, I. Belyaev, E. Ben-Haim, M. Benayoun, G. Bencivenni, S. Benson, J.\n Benton, R. Bernet, M.-O. Bettler, M. van Beuzekom, A. Bien, S. Bifani, T.\n Bird, A. Bizzeti, P. M. Bj{\\o}rnstad, T. Blake, F. Blanc, C. Blanks, J.\n Blouw, S. Blusk, A. Bobrov, V. Bocci, A. Bondar, N. Bondar, W. Bonivento, S.\n Borghi, A. Borgia, T. J. V. Bowcock, C. Bozzi, T. Brambach, J. van den Brand,\n J. Bressieux, D. Brett, M. Britsch, T. Britton, N. H. Brook, H. Brown, A.\n B\\\"uchler-Germann, I. Burducea, A. Bursche, J. Buytaert, S. Cadeddu, O.\n Callot, M. Calvi, M. Calvo Gomez, A. Camboni, P. Campana, A. Carbone, G.\n Carboni, R. Cardinale, A. Cardini, L. Carson, K. Carvalho Akiba, G. Casse, M.\n Cattaneo, Ch. Cauet, M. Charles, Ph. Charpentier, N. Chiapolini, K. Ciba, X.\n Cid Vidal, G. Ciezarek, P. E. L. Clarke, M. Clemencic, H. V. Cliff, J.\n Closier, C. Coca, V. Coco, J. Cogan, P. Collins, A. Comerma-Montells, F.\n Constantin, A. Contu, A. Cook, M. Coombes, G. Corti, G. A. Cowan, R. Currie,\n C. D'Ambrosio, P. David, P. N. Y. David, I. De Bonis, S. De Capua, M. De\n Cian, F. De Lorenzi, J. M. De Miranda, L. De Paula, P. De Simone, D. Decamp,\n M. Deckenhoff, H. Degaudenzi, L. Del Buono, C. Deplano, D. Derkach, O.\n Deschamps, F. Dettori, J. Dickens, H. Dijkstra, P. Diniz Batista, F. Domingo\n Bonal, S. Donleavy, F. Dordei, A. Dosil Su\\'arez, D. Dossett, A. Dovbnya, F.\n Dupertuis, R. Dzhelyadin, A. Dziurda, S. Easo, U. Egede, V. Egorychev, S.\n Eidelman, D. van Eijk, F. Eisele, S. Eisenhardt, R. Ekelhof, L. Eklund, Ch.\n Elsasser, D. Elsby, D. Esperante Pereira, L. Est\\`eve, A. Falabella, E.\n Fanchini, C. F\\\"arber, G. Fardell, C. Farinelli, S. Farry, V. Fave, V.\n Fernandez Albor, M. Ferro-Luzzi, S. Filippov, C. Fitzpatrick, M. Fontana, F.\n Fontanelli, R. Forty, M. Frank, C. Frei, M. Frosini, S. Furcas, A. Gallas\n Torreira, D. Galli, M. Gandelman, P. Gandini, Y. Gao, J-C. Garnier, J.\n Garofoli, J. Garra Tico, L. Garrido, D. Gascon, C. Gaspar, N. Gauvin, M.\n Gersabeck, T. Gershon, Ph. Ghez, V. Gibson, V.V. Gligorov, C. G\\\"obel, D.\n Golubkov, A. Golutvin, A. Gomes, H. Gordon, M. Grabalosa G\\'andara, R.\n Graciani Diaz, L. A. Granado Cardoso, E. Graug\\'es, G. Graziani, A. Grecu, E.\n Greening, S. Gregson, B. Gui, E. Gushchin, Yu. Guz, T. Gys, G. Haefeli, C.\n Haen, S. C. Haines, T. Hampson, S. Hansmann-Menzemer, R. Harji, N. Harnew, J.\n Harrison, P. F. Harrison, T. Hartmann, J. He, V. Heijne, K. Hennessy, P.\n Henrard, J. A. Hernando Morata, E. van Herwijnen, E. Hicks, K. Holubyev, P.\n Hopchev, W. Hulsbergen, P. Hunt, T. Huse, R. S. Huston, D. Hutchcroft, D.\n Hynds, V. Iakovenko, P. Ilten, J. Imong, R. Jacobsson, A. Jaeger, M. Jahjah\n Hussein, E. Jans, F. Jansen, P. Jaton, B. Jean-Marie, F. Jing, M. John, D.\n Johnson, C. R. Jones, B. Jost, M. Kaballo, S. Kandybei, M. Karacson, T. M.\n Karbach, J. Keaveney, I. R. Kenyon, U. Kerzel, T. Ketel, A. Keune, B. Khanji,\n Y. M. Kim, M. Knecht, R. Koopman, P. Koppenburg, A. Kozlinskiy, L. Kravchuk,\n K. Kreplin, M. Kreps, G. Krocker, P. Krokovny, F. Kruse, K. Kruzelecki, M.\n Kucharczyk, T. Kvaratskheliya, V. N. La Thi, D. Lacarrere, G. Lafferty, A.\n Lai, D. Lambert, R. W. Lambert, E. Lanciotti, G. Lanfranchi, C. Langenbruch,\n T. Latham, C. Lazzeroni, R. Le Gac, J. van Leerdam, J.-P. Lees, R. Lef\\`evre,\n A. Leflat, J. Lefran\\c{c}ois, O. Leroy, T. Lesiak, L. Li, L. Li Gioi, M.\n Lieng, M. Liles, R. Lindner, C. Linn, B. Liu, G. Liu, J. von Loeben, J. H.\n Lopes, E. Lopez Asamar, N. Lopez-March, H. Lu, J. Luisier, A. Mac Raighne, F.\n Machefert, I. V. Machikhiliyan, F. Maciuc, O. Maev, J. Magnin, S. Malde, R.\n M. D. Mamunur, G. Manca, G. Mancinelli, N. Mangiafave, U. Marconi, R.\n M\\\"arki, J. Marks, G. Martellotti, A. Martens, L. Martin, A. Mart\\'in\n S\\'anchez, D. Martinez Santos, A. Massafferri, Z. Mathe, C. Matteuzzi, M.\n Matveev, E. Maurice, B. Maynard, A. Mazurov, G. McGregor, R. McNulty, M.\n Meissner, M. Merk, J. Merkel, R. Messi, S. Miglioranzi, D. A. Milanes, M.-N.\n Minard, J. Molina Rodriguez, S. Monteil, D. Moran, P. Morawski, R. Mountain,\n I. Mous, F. Muheim, K. M\\\"uller, R. Muresan, B. Muryn, B. Muster, M. Musy, J.\n Mylroie-Smith, P. Naik, T. Nakada, R. Nandakumar, I. Nasteva, M. Nedos, M.\n Needham, N. Neufeld, C. Nguyen-Mau, M. Nicol, V. Niess, N. Nikitin, A.\n Nomerotski, A. Novoselov, A. Oblakowska-Mucha, V. Obraztsov, S. Oggero, S.\n Ogilvy, O. Okhrimenko, R. Oldeman, M. Orlandea, J. M. Otalora Goicochea, P.\n Owen, K. Pal, J. Palacios, A. Palano, M. Palutan, J. Panman, A. Papanestis,\n M. Pappagallo, C. Parkes, C. J. Parkinson, G. Passaleva, G. D. Patel, M.\n Patel, S. K. Paterson, G. N. Patrick, C. Patrignani, C. Pavel-Nicorescu, A.\n Pazos Alvarez, A. Pellegrino, G. Penso, M. Pepe Altarelli, S. Perazzini, D.\n L. Perego, E. Perez Trigo, A. P\\'erez-Calero Yzquierdo, P. Perret, M.\n Perrin-Terrin, G. Pessina, A. Petrella, A. Petrolini, A. Phan, E. Picatoste\n Olloqui, B. Pie Valls, B. Pietrzyk, T. Pila\\v{r}, D. Pinci, R. Plackett, S.\n Playfer, M. Plo Casasus, G. Polok, A. Poluektov, E. Polycarpo, D. Popov, B.\n Popovici, C. Potterat, A. Powell, J. Prisciandaro, V. Pugatch, A. Puig\n Navarro, W. Qian, J. H. Rademacker, B. Rakotomiaramanana, M. S. Rangel, I.\n Raniuk, G. Raven, S. Redford, M. M. Reid, A. C. dos Reis, S. Ricciardi, K.\n Rinnert, D. A. Roa Romero, P. Robbe, E. Rodrigues, F. Rodrigues, P. Rodriguez\n Perez, G. J. Rogers, S. Roiser, V. Romanovsky, M. Rosello, J. Rouvinet, T.\n Ruf, H. Ruiz, G. Sabatino, J. J. Saborido Silva, N. Sagidova, P. Sail, B.\n Saitta, C. Salzmann, M. Sannino, R. Santacesaria, C. Santamarina Rios, R.\n Santinelli, E. Santovetti, M. Sapunov, A. Sarti, C. Satriano, A. Satta, M.\n Savrie, D. Savrina, P. Schaack, M. Schiller, S. Schleich, M. Schlupp, M.\n Schmelling, B. Schmidt, O. Schneider, A. Schopper, M.-H. Schune, R.\n Schwemmer, B. Sciascia, A. Sciubba, M. Seco, A. Semennikov, K. Senderowska,\n I. Sepp, N. Serra, J. Serrano, P. Seyfert, M. Shapkin, I. Shapoval, P.\n Shatalov, Y. Shcheglov, T. Shears, L. Shekhtman, O. Shevchenko, V.\n Shevchenko, A. Shires, R. Silva Coutinho, T. Skwarnicki, A. C. Smith, N. A.\n Smith, E. Smith, K. Sobczak, F. J. P. Soler, A. Solomin, F. Soomro, B. Souza\n De Paula, B. Spaan, A. Sparkes, P. Spradlin, F. Stagni, S. Stahl, O.\n Steinkamp, S. Stoica, S. Stone, B. Storaci, M. Straticiuc, U. Straumann, V.\n K. Subbiah, S. Swientek, M. Szczekowski, P. Szczypka, T. Szumlak, S.\n T'Jampens, E. Teodorescu, F. Teubert, C. Thomas, E. Thomas, J. van Tilburg,\n V. Tisserand, M. Tobin, S. Topp-Joergensen, N. Torr, E. Tournefier, M. T.\n Tran, A. Tsaregorodtsev, N. Tuning, M. Ubeda Garcia, A. Ukleja, P. Urquijo,\n U. Uwer, V. Vagnoni, G. Valenti, R. Vazquez Gomez, P. Vazquez Regueiro, S.\n Vecchi, J. J. Velthuis, M. Veltri, B. Viaud, I. Videau, X. Vilasis-Cardona,\n J. Visniakov, A. Vollhardt, D. Volyanskyy, D. Voong, A. Vorobyev, H. Voss, S.\n Wandernoth, J. Wang, D. R. Ward, N. K. Watson, A. D. Webber, D. Websdale, M.\n Whitehead, D. Wiedner, L. Wiggers, G. Wilkinson, M. P. Williams, M. Williams,\n F. F. Wilson, J. Wishahi, M. Witek, W. Witzeling, S. A. Wotton, K. Wyllie, Y.\n Xie, F. Xing, Z. Xing, Z. Yang, R. Young, O. Yushchenko, M. Zavertyaev, F.\n Zhang, L. Zhang, W. C. Zhang, Y. Zhang, A. Zhelezov, L. Zhong, E. Zverev, A.\n Zvyagin",
"submitter": "Vincenzo Maria Vagnoni",
"url": "https://arxiv.org/abs/1112.0938"
}
|
1112.0960
|
# Antisymmetric tensor unparticle and the radiative lepton flavor violating
decays
E. O. Iltan,
Physics Department, Middle East Technical University
Ankara, Turkey
E-mail address: eiltan@metu.edu.tr
###### Abstract
We study the contribution of the tensor unparticle mediation to the branching
ratios of the radiative lepton flavor violating decays and predict a
restriction region for free parameters of the scenario by using experimental
upper limits. We observe that the branching ratios of the radiative lepton
flavor violating decays are sensitive to the fundamental mass scales of the
scenario and to the scale dimension of antisymmetric tensor unparticle. We
obtain a more restricted set for the free parameters in the case of the
$\mu\rightarrow e\gamma$ decay.
The radiative lepton flavor violating (LFV) decays $l_{i}\rightarrow
l_{j}\gamma$ reach great interest since their branching ratios (BRs) in the
framework of standard model (SM) are much below the experimental upper limits,
and, therefore, they are candidates to search and to test more fundamental
models beyond. The current experimental upper limits of the BRs read BR
$(\mu\rightarrow e\gamma)=2.4\,(1.2)\times 10^{-12}\,(10^{-11})\,(\,\,90\%CL)$
[1] ([2]), BR $(\tau\rightarrow e\gamma)=3.3\times 10^{-8}\,(\,\,90\%CL)$ [3]
and BR $(\tau\rightarrow\mu\gamma)=4.4\times 10^{-8}\,(\,\,90\%CL)$ [3]. There
is an extensive theoretical work in the literature in order to enhance BRs of
these decays. They were studied in the SM with the extended Higgs sector, so
called two Higgs doublet model (2HDM) [4]-[10], in supersymmetric models
[11]-[18], in a model independent way [19], in the framework of the 2HDM and
the supersymmetric model [20], in the SM including effective operators coming
from the possible unparticle effects [21]-[22], in little Higgs models [23]\-
[26], in seesaw models [27] -[30], in models with A(4) and S(4) flavor
symmetries [31], using the effective field theory with Higgs mediation [32],
in the Higgs triplet model [33], in the framework of Higgs-induced lepton
flavor violation [34].
In the present work, we consider the contribution of the antisymmetric tensor
unparticle mediation to the BRs of the radiative LFV decays (see [35] for the
contribution of the antisymmetric tensor unparticle mediation to the muon
anomalous magnetic dipole moment, to the electroweak precision observable $S$,
its effects in $Z$ invisible decays and see [36] for the contribution of the
antisymmetric tensor unparticle mediation to the lepton electric dipole
moment). Unparticles [37, 38], being massless due to the scale invariance and
having non integral scaling dimension $d_{U}$ around the scale
$\Lambda_{U}\sim 1.0\,TeV$, come out with the interaction of SM-ultraviolet
sector at some scale $M_{U}$:
${\cal{L}}_{eff}=\frac{C_{n}}{M_{U}^{d_{UV}+n-4}}\,O_{SM}\,O_{UV}\,,$ (1)
where $d_{UV}$ is the scaling dimension of the UV operator [39]. Around the
scale $\Lambda_{U}$ the effective interaction becomes [40]
${\cal{L}}_{eff}=\frac{C^{i}_{n}}{\Lambda_{n}^{d_{U}+n-4}}\,O_{SM,i}\,O_{U}\,.$
(2)
Here $O_{SM,i}$ is type $i$ SM operator, $n$ is its scaling dimension and
$\Lambda_{n}$ is the mass scale (see [35, 40] for details) which reads
$\Lambda_{n}=\Bigg{(}\frac{M_{U}^{d_{UV}+n-4}}{\Lambda_{U}^{d_{UV}-d_{U}}}\Bigg{)}^{\frac{1}{d_{U}+n-4}}\,.$
(3)
The antisymmetric tensor unparticle mediation induces these LFV decays at tree
level and we consider the case that the scale invariance is broken at some
scale $\mu$ after the electroweak symmetry breaking (see for example [41, 42]
for a possible interaction which causes that scale invariance is broken). The
effective lagrangian [35] which can drive the radiative LFV decays is
$\displaystyle{\cal{L}}_{eff}$ $\displaystyle=$
$\displaystyle\frac{g^{\prime}\,\lambda_{B}}{\Lambda_{2}^{d_{U}-2}}\,B_{\mu\nu}\,O^{\mu\nu}_{U}+\frac{g\,\lambda_{W}}{\Lambda_{4}^{d_{U}}}\,(H^{\dagger}\,\tau_{a}\,H)\,W^{a}_{\mu\nu}\,O^{\mu\nu}_{U}$
(4) $\displaystyle+$
$\displaystyle\frac{\lambda_{ij}}{\Lambda_{4}^{d_{U}}}\,y_{ij}\,\bar{l}_{i}\,H\,\sigma_{\mu\nu}\,l_{j}\,O^{\mu\nu}_{U}\,,$
where $l_{i(j)}$ is the lepton field, $H$ is the Higgs doublet, $g$ and
$g^{\prime}$ are weak couplings, $\lambda_{B}$, $\lambda_{W}$ and
$\lambda_{ij}$ are the unparticle-field tensor and unparticle-lepton-lepton
couplings, $B_{\mu\nu}$ is the field strength tensor of the $U(1)_{Y}$ gauge
boson with $B_{\mu}=c_{W}\,A_{\mu}+s_{W}\,Z_{\mu}$ and $W^{a}_{\mu\nu}$,
$a=1,2,3$, are field strength tensors of the $SU(2)_{L}$ gauge bosons with
$W^{3}_{\mu}=s_{W}\,A_{\mu}-c_{W}\,Z_{\mu}$ and $A_{\mu}$ ($Z_{\mu}$) is
photon (Z boson) field. The couplings $y_{ij}$ are responsible for the LF
violation and after the electroweak symmetry breaking we introduce modified
couplings $\xi_{ij}=\frac{v}{\sqrt{2}}\,y_{ij}$ which respect the mass
hierarchy of charged leptons. The process $l_{i}\rightarrow l_{j}\gamma$
appears in the tree level with the communication of two vertices111The first
vertex arises from the last term of the effective lagrangian and leads to the
$l_{i}\rightarrow l_{j}$ transition. The second one arises from the first and
second terms of the effective lagrangian and results in the
$O^{\mu\nu}_{U}\rightarrow A_{\nu}$ transition
$\lambda_{ij}\,\frac{\xi_{ij}}{\Lambda_{4}^{d_{U}}}\,\bar{l}_{i}\,\sigma_{\mu\nu}\,l_{j}\,O^{\mu\nu}_{U}$
and
$\Big{(}\,2\,i\,\frac{g^{\prime}\,c_{W}\,\lambda_{B}}{\Lambda_{2}^{d_{U}-2}}-i\,\frac{g\,v^{2}\,s_{W}\,\lambda_{W}}{2\,\Lambda_{4}^{d_{U}}}\,\Big{)}\,k_{\mu}\,\epsilon_{\nu}\,O^{\mu\nu}_{U}$,
by the antisymmetric tensor unparticle propagator (see Appendix and eq.(12))
and the matrix element of this process reads
$\displaystyle
M=a_{ij}\,\bar{l}_{i}\,\sigma_{\mu\nu}\,l_{j}\,k_{\mu}\,\epsilon_{\nu}\,,$ (5)
where
$\displaystyle
a_{ij}=\frac{i\,e\,\mu^{2\,(d_{U}-2)}\,\,A_{d_{U}}\,\lambda_{ij}\,\xi_{ij}}{sin\,(d_{U}\pi)\,\Lambda_{4}^{d_{U}}}\,\Bigg{(}\frac{\lambda_{B}}{\Lambda_{2}^{d_{U}-2}}-\frac{v^{2}\,\lambda_{W}}{4\,\Lambda_{4}^{d_{U}}}\Bigg{)}\,.$
(6)
Finally the decay width $\Gamma(l_{i}\rightarrow l_{j}\gamma)$ becomes
$\displaystyle\Gamma(l_{i}\rightarrow
l_{j}\gamma)=\frac{1}{8\,\pi}\,m_{i}^{3}\,|a_{ij}|^{2}\,,$ (7)
where $m_{i}$ is the mass of incoming lepton. Notice that, in this expression,
we ignore the mass of outgoing one.
Discussion
In this section we study the intermediate antisymmetric tensor unparticle
contribution to the radiative LFV decays $l_{i}\rightarrow l_{j}\gamma$ which
exist at tree level (see Fig.1). There are various free parameters in this
scenario and we restrict them by using the current experimental upper limits
of BRs of LFV decays. Now, we would like to present the free parameters and
discuss the restrictions predicted. The SM sector interacts with the UV one
and it appears as unparticle sector at a lower scale. The corresponding UV
(unparticle) operator $O_{UV}$ ($O_{U}$) has the scaling dimension $d_{UV}$
($d_{U}$) which is among the free parameters. We choose the scale dimension
$d_{U}$ in the range $1<d_{U}<2$. Notice that the scale dimension must satisfy
$d_{U}>2$ for antisymmetric tensor unparticle in order not to violate the
unitarity [43]. Our assumption is based on the fact that the scale invariance
is broken at some scale $\mu$ and one reaches to the particle sector. This
results in a relaxation on the values of $d_{U}$ and we choose $d_{U}$ in the
range $1<d_{U}<2$ so that the propagator for particle sector is obtained when
$d_{U}$ tends to one. Furthermore we choose the numerical value of $d_{UV}$ as
$d_{UV}=3$ which satisfies $d_{UV}>d_{U}$ (see [40]). The SM-ultraviolet
sector interaction scale $M_{U}$, the SM-unparticle sector interaction scale
$\Lambda_{U}$ and the scale $\mu$ which is the one that scale invariance is
broken belong to the free parameter set of the present scenario. Here we
choose $\mu\sim 1.0\,GeV$ and predict the restrictions for the others by using
the experimental upper limits of LFV decays. Finally, for the couplings
$\lambda_{B}$, $\lambda_{W}$ and $\lambda_{ij}$ we consider
$\lambda_{B}=\lambda_{W}=\lambda_{ij}=1$ and for $\xi_{ij}$ we respect the
mass hierarchy of charged leptons, namely we choose $\xi_{\tau\mu}>\xi_{\tau
e}>\xi_{\mu e}$ and we take $\xi_{\tau\mu}=0.1\,GeV$, $\xi_{\tau e}=0.01\,GeV$
and $\xi_{\mu e}=0.001\,GeV$ in our numerical calculations.
In Fig.2, we present the BR$(\mu\rightarrow e\,\gamma$) with respect to the
mass scale $M_{U}$ for $r_{U}=\frac{\Lambda_{U}}{M_{U}}=0.1$. Here, the solid
(long dashed-short dashed) line represents the BR for $d_{U}=1.7\,(1.8-1.9)$.
We observe that the BR is sensitive to the mass scale $M_{U}$ especially for
the large values of the scale dimension $d_{U}$ and it decreases almost three
orders in the range of $2000\,GeV<M_{U}<10000\,GeV$. The experimental upper
limit is reached for $d_{U}\sim 1.8\,(1.9)$ and $M_{U}\sim 8000\,(4500)\,GeV$.
Fig.3 is devoted to the BR$(\mu\rightarrow e\,\gamma$) with respect to the
scale parameter $d_{U}$ for $r_{U}=0.1$. Here the solid (long dashed-short
dashed-dotted) line represents the BR for
$M_{U}=3000\,(5000-8000-10000)\,GeV$. The BR strongly depends on $d_{U}$ and
decreases with the increasing values of $d_{U}$. The experimental upper limit
is observed in the range of $d_{U}\sim 1.78-1.88$ for $M_{U}\sim
5000-10000\,GeV$.
In Fig.4 we show the parameter $r_{U}$ with respect to $d_{U}$ for
BR$(\mu\rightarrow e\,\gamma)=2.4\times 10^{-12}$. Here the solid (long
dashed-short dashed) line represents $r_{U}$ for
$M_{U}=5000\,(8000-10000)\,GeV$. We see that the scale dimension $d_{U}$ and
$r_{U}$ can take values in the range $1.73-1.90$ and $0.05-0.12$, respectively
for $M_{U}=5000\,GeV$. For $M_{U}=10000\,GeV$ we have the range $1.65-1.9$ for
$d_{U}$ and $0.05-0.20$ for $r_{U}$.
Fig.5 represents the BR$(\tau\rightarrow e\,\gamma$) with respect to the mass
scale $M_{U}$. Here, the solid (long dashed-short dashed-dotted) line
represents the BR for $r_{U}=0.1$, $d_{U}=1.3$ ($r_{U}=0.1$,
$d_{U}=1.4$-$r_{U}=0.5$, $d_{U}=1.6$-$r_{U}=0.5$, $d_{U}=1.7$). It is observed
that the sensitivity of the BR to the mass scale $M_{U}$ increases with the
increasing values of the ratio $r_{U}$. The experimental upper limit is
reached for $r_{U}=0.1$ and $d_{U}\sim 1.3$ when the mass scale $M_{U}$ reads
$M_{U}\sim 4000\,GeV$. For $r_{U}=0.5$ one reaches the experimental limit in
the case of $d_{U}\sim 1.6$ and $M_{U}\sim 4000\,GeV$. Fig.5 shows the
BR$(\tau\rightarrow e\,\gamma$) with respect to the scale parameter $d_{U}$
for $r_{U}=0.1$. Here the solid (long dashed-short dashed) line represents the
BR for $M_{U}=2000\,(5000-10000)\,GeV$. The BR strongly depends on $d_{U}$ and
decreases with the increasing values of $d_{U}$ similar to the $\mu\rightarrow
e\,\gamma$ decay. One reaches the experimental upper limit in the range of
$d_{U}\sim 1.26-1.32$ for $M_{U}\sim 2000-10000\,GeV$.
Fig.7 is devoted to the parameter $r_{U}$ with respect to $d_{U}$ for
BR$(\tau\rightarrow e\,\gamma)=3.3\times 10^{-8}$. Here the solid (long
dashed-short dashed) line represents $r_{U}$ for
$M_{U}=5000\,(8000-10000)\,GeV$. This figure shows that the experimental upper
limit is reached for $M_{U}=5000\,GeV$ if the scale dimension $d_{U}$ and
$r_{U}$ can take values in the range $1.10-1.68$ and $0.05-1.00$,
respectively. For $M_{U}=10000\,GeV$ we have the range $1.10-1.64$ for $d_{U}$
and $0.05-1.00$ for $r_{U}$.
In Fig.8, we present the BR$(\tau\rightarrow\mu\,\gamma$) with respect to the
mass scale $M_{U}$ for $r_{U}=0.1$. Here, the solid (long dashed-short dashed-
dotted) line represents the BR for $d_{U}=1.4\,(1.5-1.6-1.7)$. We observe that
the BR decreases more than one order in the range of
$2000\,GeV<M_{U}<10000\,GeV$. The experimental upper limit is reached for
$d_{U}\sim 1.4$ and $M_{U}\sim 9000\,GeV$. Fig.9 represents the
BR$(\tau\rightarrow\mu\,\gamma)$ with respect to the scale parameter $d_{U}$
for $r_{U}=0.1$. Here the solid (long dashed-short dashed) line represents the
BR for $M_{U}=2000\,(5000-10000)\,GeV$. The experimental upper limit of the BR
is observed in the range of $d_{U}\sim 1.38-1.47$ for $M_{U}\sim
2000-10000\,GeV$.
In Fig. 10 we show the parameter $r_{U}$ with respect to $d_{U}$ for
BR$(\tau\rightarrow\mu\,\gamma)=4.4\times 10^{-8}$. Here the solid (long
dashed-short dashed) line represents $r_{U}$ for
$M_{U}=2000\,(5000-10000)\,GeV$. We see that the scale dimension $d_{U}$ and
$r_{U}$ can take values in the range $1.48-1.80\,(1.3-1.8)$ and
$0.10-0.55\,(0.05-1.00)$, respectively for $M_{U}=2000\,(5000)\,GeV$. For
$M_{U}=10000\,GeV$ we have the range $1.30-1.75$ for $d_{U}$ and $0.05-1.00$
for $r_{U}$.
As a summary, the BRs of radiative LFV decays are sensitive to the mass scale
$M_{U}$ especially for the large values of the scale dimension $d_{U}$ and
this sensitivity increases with the increasing values of the ratio $r_{U}$.
The experimental upper limit of the BR$(\mu\rightarrow e\,\gamma$) can be
reached for $r_{U}\sim 0.1$, $d_{U}>1.7$ and for larger values of $M_{U}$,
namely $M_{U}\sim 7000-9000\,GeV$. For BR$(\tau\rightarrow e\,\gamma$) one
reaches the experimental upper limit for $r_{U}\sim 0.1$, $d_{U}\sim 1.3$ and
$M_{U}>2000\,GeV$ . If we consider the $\tau\rightarrow\mu\,\gamma$ decay the
experimental upper limit of BR is obtained for $r_{U}\sim 0.1$, when $d_{U}$
is in the range $d_{U}\sim 1.4-1.5$ and $M_{U}>2000\,GeV$. We see that the
free parameters of this scenario are more restricted if the $\mu\rightarrow
e\,\gamma$ decay is considered. However the future more accurate measurements
of the upper limits of the LFV decays make it possible to obtain a more
restricted range for the free parameters of this scenario and they stimulate
to search the role and the nature of unparticle physics which is a candidate
to drive the lepton flavor violation.
Appendix
The scalar unparticle propagator reads [38, 44]
$\displaystyle\int\,d^{4}x\,e^{ip.x}<0|T\Big{(}O_{U}(x)\,O_{U}(0)\Big{)}0>=i\frac{A_{d_{U}}}{2\,\pi}\,\int_{0}^{\infty}\\!\\!\\!ds\,\frac{s^{d_{U}-2}}{p^{2}-s+i\epsilon}\\!=i\,\frac{A_{d_{U}}}{2\,sin\,(d_{U}\pi)}\,(-p^{2}-i\epsilon)^{d_{U}-2}$
(8)
where the factor $A_{d_{U}}$ is
$\displaystyle
A_{d_{U}}=\frac{16\,\pi^{5/2}}{(2\,\pi)^{2\,d_{U}}}\,\frac{\Gamma(d_{U}+\frac{1}{2})}{\Gamma(d_{U}-1)\,\Gamma(2\,d_{U})}\,.$
(9)
Now the tensor unparticle propagator is obtained by using the projection
operator $\Pi^{\mu\nu\alpha\beta}$
$\displaystyle\Pi_{\mu\nu\alpha\beta}=\frac{1}{2}(g_{\mu\alpha}\,g_{\nu\beta}-g_{\nu\alpha}\,g_{\mu\beta})\,,$
(10)
with the transverse and the longitudinal parts
$\displaystyle\Pi^{T}_{\mu\nu\alpha\beta}=\frac{1}{2}(P^{T}_{\mu\alpha}\,P^{T}_{\nu\beta}-P^{T}_{\nu\alpha}\,P^{T}_{\mu\beta})\,,\,\,\,\,\,\,\Pi^{L}_{\mu\nu\alpha\beta}=\Pi_{\mu\nu\alpha\beta}-\Pi^{T}_{\mu\nu\alpha\beta}\,,$
(11)
where $P^{T}_{\mu\nu}=g_{\mu\nu}-p_{\mu}\,p_{\nu}/{p^{2}}$ (see for example
[35] and references therein) and the propagator of antisymmetric tensor
unparticle becomes
$\displaystyle\int\,d^{4}x\,e^{ipx}\,<0|T\Big{(}O^{\mu\nu}_{U}(x)\,O^{\alpha\beta}_{U}(0)\Big{)}0>=i\,\frac{A_{d_{U}}}{2\,sin\,(d_{U}\pi)}\,\Pi^{\mu\nu\alpha\beta}(-p^{2}-i\epsilon)^{d_{U}-2}\,.$
On the other hand the propagator is modified if the scale invariance broken at
a certain scale and this modification is model dependent. Following the the
simple model [41, 45, 46] we take
$\displaystyle\int\,d^{4}x\,e^{ipx}\,<0|T\Big{(}O^{\mu\nu}_{U}(x)\,O^{\alpha\beta}_{U}(0)\Big{)}0>=i\,\frac{A_{d_{U}}}{2\,sin\,(d_{U}\pi)}\,\Pi^{\mu\nu\alpha\beta}(-(p^{2}-\mu^{2})-i\epsilon)^{d_{U}-2}\,.$
(12)
where $\mu$ is the scale that the scale invariance broken and the particle
sector comes out.
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Figure 1: Tree level diagram contributing to the $l_{i}\rightarrow
l_{j}\,\gamma$ decay due to antisymmetric tensor unparticle. Wavy (solid) line
represents the electromagnetic field (lepton field) and double dashed line the
antisymmetric tensor unparticle field.
Figure 2: $M_{U}$ dependence of the BR $(\mu\rightarrow e\,\gamma)$ for
$r_{U}=0.1$. Here, the solid (long dashed-short dashed) line represents the BR
for $d_{U}=1.7\,(1.8-1.9)$. Figure 3: $d_{U}$ dependence of the BR
$(\mu\rightarrow e\,\gamma)$ for $r_{U}=0.1$. Here the solid (long dashed-
short dashed-dotted) line represents the BR for
$M_{U}=3000\,(5000-8000-10000)\,GeV$. Figure 4: $r_{U}$ with respect to
$d_{U}$ for BR$(\mu\rightarrow e\,\gamma)=2.4\times 10^{-12}$. Here the solid
(long dashed-short dashed) line represents $r_{U}$ for
$M_{U}=5000\,(8000-10000)\,GeV$ Figure 5: $M_{U}$ dependence of the BR
$(\tau\rightarrow e\,\gamma)$. Here, the solid (long dashed-short dashed-
dotted) line represents the BR for $r_{U}=0.1$, $d_{U}=1.3$ ($r_{U}=0.1$,
$d_{U}=1.4$-$r_{U}=0.5$, $d_{U}=1.6$-$r_{U}=0.5$, $d_{U}=1.7$). Figure 6:
$d_{U}$ dependence of the BR $(\tau\rightarrow e\,\gamma)$ for $r_{U}=0.1$.
Here the solid (long dashed-short dashed) line represents the BR for
$M_{U}=2000\,(5000-10000)\,GeV$. Figure 7: $r_{U}$ with respect to $d_{U}$ for
BR$(\tau\rightarrow e\,\gamma)=3.3\times 10^{-8}$. Here the solid (long
dashed-short dashed) line represents $r_{U}$ for
$M_{U}=5000\,(8000-10000)\,GeV$. Figure 8: $M_{U}$ dependence of the
BR$(\tau\rightarrow\mu\,\gamma)$. Here, the solid (long dashed-short dashed-
dotted) line represents the BR for $d_{U}=1.4\,(1.5-1.6-1.7)$. Figure 9:
$d_{U}$ dependence of the BR$(\tau\rightarrow\mu\,\gamma)$ for $r_{U}=0.1$.
Here the solid (long dashed-short dashed) line represents the BR for
$M_{U}=2000\,(5000-10000)\,GeV$. Figure 10: $r_{U}$ with respect to $d_{U}$
for BR$(\tau\rightarrow\mu\,\gamma)=4.4\times 10^{-8}$. Here the solid (long
dashed-short dashed) line represents $r_{U}$ for
$M_{U}=2000\,(5000-10000)\,GeV$.
|
arxiv-papers
| 2011-12-05T15:10:44 |
2024-09-04T02:49:24.980859
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "E. Iltan",
"submitter": "Erhan Iltan",
"url": "https://arxiv.org/abs/1112.0960"
}
|
1112.1070
|
# Mechanism for puddle formation in graphene
S. Adam Center for Nanoscale Science and Technology, National Institute of
Standards and Technology, Gaithersburg, MD 20899, USA Suyong Jung Center for
Nanoscale Science and Technology, National Institute of Standards and
Technology, Gaithersburg, MD 20899, USA Maryland NanoCenter, University of
Maryland, College Park, MD 20472, USA Nikolai N. Klimov Center for Nanoscale
Science and Technology, National Institute of Standards and Technology,
Gaithersburg, MD 20899, USA Maryland NanoCenter, University of Maryland,
College Park, MD 20472, USA Physical Measurement Laboratory, National
Institute of Standards and Technology, Gaithersburg, MD 20899, USA Nikolai B.
Zhitenev Center for Nanoscale Science and Technology, National Institute of
Standards and Technology, Gaithersburg, MD 20899, USA Joseph A. Stroscio
Center for Nanoscale Science and Technology, National Institute of Standards
and Technology, Gaithersburg, MD 20899, USA M. D. Stiles Center for
Nanoscale Science and Technology, National Institute of Standards and
Technology, Gaithersburg, MD 20899, USA
###### Abstract
When graphene is close to charge neutrality, its energy landscape is highly
inhomogeneous, forming a sea of electron-like and hole-like puddles, which
determine the properties of graphene at low carrier density. However, the
details of the puddle formation have remained elusive. We demonstrate
numerically that in sharp contrast to monolayer graphene, the normalized
autocorrelation function for the puddle landscape in bilayer graphene depends
only on the distance between the graphene and the source of the long-ranged
impurity potential. By comparing with available experimental data, we find
quantitative evidence for the implied differences in scanning tunneling
microscopy measurements of electron and hole puddles for monolayer and bilayer
graphene in nominally the same disorder potential.
###### pacs:
73.22.Pr,68.37.Ef,81.05.ue
## I Introduction
In monolayer graphene, the hexagonal arrangement of carbon atoms dictates that
in the absence of atomic-scale disorder, graphene is a gapless
semiconductorDas Sarma et al. (2011); Castro Neto et al. (2009) that is always
metallic at low temperature.Fuhrer and Adam (2009) This metallic behavior
holds even in the presence of quantum interference and strong
disorder,Bardarson et al. (2007) in stark contrast to most other materials,
which undergo a metal-to-insulator transition at low carrier density.Anderson
(1958); Tanatar and Ceperley (1989) The physical origin for this robust
metallic state is that the ground-state of graphene at vanishing mean carrier
density becomes spatially inhomogeneous, breaking up into electron-rich and
hole-rich metallic regions connected by highly conducting p-n
junctions.Katsnelson et al. (2006) Bilayer graphene comprising two sheets of
graphene that become strongly coupled due to the $AB$ stacking
arrangementMcCann and Fal’ko (2006) shares some properties with regular
semiconductors (such as the parabolic band dispersion) and in other ways
behaves like monolayer graphene, including having chiral wavefunctions and
forming electron and hole puddles at low density.
These electron and hole puddles have now been observed in several experiments
of exfoliated graphene on an insulating SiO2 substrate including Refs. Martin
et al., 2008; Zhang et al., 2009; Deshpande et al., 2009, 2011; Jung et al.,
2011; Rutter et al., 2011. While these authors suggest that long-range charged
impurities in the substrate could be responsible for the spatial
inhomogeneity, detailed comparisons to microscopic models have not been made.
In this paper, we demonstrate that differences between the spatial properties
of puddles in monolayer and bilayer graphene can be quantitatively explained
by the differences in the screening properties of the two systems (that
ultimately arises from the differences in their band-structure). Numerical
results show that the correlation length for bilayer graphene is relatively
independent of density and significantly smaller than that of monolayer
graphene for a typical range of impurity densities. Finally, we find good
quantitative agreement when comparing our results with available experimental
data.
The rest of the paper is organized as follows. In Sec. II, we outline the the
theoretical model, providing a heuristic understanding of our results using
the Thomas-Fermi (TF) screening theory. However, the TF significantly
underestimates the effect of electronic screening in both monolayer and
bilayer graphene. It is therefore necessary to use the Random Phase
Approximation (RPA) screening theory, which we discuss in Sec. III. Our main
finding is that the puddle correlation length in bilayer graphene $(\xi\approx
3.5~{}{\rm nm})$ is relatively insensitive to the impurity concentration and
carrier doping. This is in contrast to monolayer graphene, where the puddle
correlation length varies from $3~{}{\rm nm}$ in dirty samples to more than
$35~{}{\rm nm}$ in clean samples.
The comparison with experiment is done in Sec. IV, where we examine three
different experimental results: (1) We consider first the experimentally
determined normalized correlation function $A(r)$ (see definition below)
obtained from the scanning tunneling microscopy (STM) data reported for
exfoliated bilayer graphene in Ref. Rutter et al., 2011. The full functional
form of $A(r)$ agrees with the theory where the only adjusted parameter in the
theory is the distance $d$ of the impurities from the graphene sheet. In
particular the experimentally determined correlation length $\xi=(3.68\pm
0.03)~{}{\rm nm}$, defined here as the half-width at half-maximum (HWHM) decay
length of $A(r)$ agrees well with the $d=1~{}{\rm nm}$ RPA theory value of
$\xi=3.4~{}{\rm nm}$. This value of $d$ is both reasonable and consistent with
those determined from other transport measurements on bilayer graphene.Adam
and Das Sarma (2008) (2) From the monolayer graphene STM experimental data
reported in Ref. Jung et al., 2011, we extract a correlation length $6~{}{\rm
nm}<\xi<11~{}{\rm nm}$. Since the measurements Jung et al. (2011); Rutter et
al. (2011) were made on the same exfoliated graphene sample containing both
single layer and bilayer graphene regions, we expect that the extrinsic
disorder potential is statistically identical for the two samples. Therefore,
using the value of $d=1~{}{\rm nm}$ (discussed above) and the disorder induced
Dirac point shift reported in Ref. Rutter et al., 2011, we calculate
theoretically (without any adjustable parameters) that a monolayer graphene
sample in the same disorder environment would have a puddle correlation length
$\xi=8~{}{\rm nm}$, in reasonable agreement with the experiment. (3) We then
compare $A(r)$ obtained using scanning Coulomb blockade spectroscopy reported
in Ref. Deshpande et al., 2011 with our theoretical results for monolayer
graphene at the Dirac point. The parameters used in the theory were obtained
from separate transport measurements on the same experimental sample.Deshpande
et al. (2011) The agreement between theory and experiment is remarkable since
it involves no adjustable parameters. Finally in Sec. V, we conclude by making
predictions for future experiments involving monolayer and bilayer graphene on
BN substrates.
## II Formalism
The doping level of graphene can be measured in a variety of different ways.
In transport measurements, the gate voltage potential that yields the
resistivity maximum identifies the extrinsic doping level due to extraneous
sources, such as charged impurities in the substrate impurities with density,
$n_{imp}$. While the width of the resistivity maximum is a measure of the
homogeneity of the sample, Adam et al. (2007) or the electron-hole puddle
distribution.
In local probe measurements, such as STM, the Dirac point energy relative to
the Fermi-level can be observed as a minimum in the tunneling differential
conductance, $dI/dV$ as a function of tunneling bias. Knowing the electronic
dispersion relation (see details below), for a particular gate voltage
$V_{g}$, this spatial map $V(r)$ of the Dirac point variation can then be used
to extract the spatial distribution of the local carrier density
(characterized by a width $n_{\rm rms}$). We can also characterize the puddles
through the radially averaged autocorrelation function
$C(r)=\frac{1}{2\pi}\int\limits_{0}^{2\pi}d\phi\langle\langle V({\bf
r})V(0)\rangle\rangle,$ (1)
where the angular brackets denote an average over the image area and the
$\phi$-integration averages over orientations.
Notice that while $C(0)=V_{\rm rms}^{2}$ (which is related to $n_{\rm rms}$)
characterizes the fluctuations in the puddle depth, $C(r)$ describes the
spatial profile of the electron and hole puddles. We find it useful to
consider the normalized correlation function $A(r)=C(r)/C(0)$. We will argue
below that $A(r)$ and $C(0)$ are quite different physical quantities that
depends quite differently on the parameters of the extrinsic impurity
potential. (In addition, for a typical STM experiment, where the shifts in the
Dirac point are determined Zhang et al. (2009) from the shifts in $dI/dV$ at
fixed $V_{g}$, the determination of $C(0)$ is complicated by the experimental
uncertainty in converting spatial maps of $dI/dV$ to Dirac point energy shifts
$V({\bf r})$. By contrast, the much smaller uncertainty in $A(r)$ is mostly
determined by the spatial resolution and image area.)
To theoretically compute the correlation functions for puddles in graphene, we
make two assumptions. First, the impurity potential comes from a random two
dimensional distribution of charged impurities displaced by a distance $d$
from the plane with density $n_{\rm imp}$. This model has been highly
successful in describing the effect of disorder in semiconductor
heterojunctionsAndo et al. (1982) and in graphene.Das Sarma et al. (2011) This
spatially varying potential gives rise to a varying charge density and local
variations in the screening of the potential. Second, we assume that it is
possible to find a global screening function $\epsilon(q,n_{\rm eff})$ that
adequately describes the effects of these local screening variations. Here,
the screening depends on the disorder potential only through an effective
carrier density $n_{\rm eff}$. This self-consistent screening model has been
used previously to understand the minimum conductivity problem in both
monolayerAdam et al. (2007) and bilayer graphene.Adam and Das Sarma (2008) If
these two assumptions are satisfied, the correlation function, aside from a
prefactor of $n_{\rm imp}$, then depends only on the screened impurity
potential
$\displaystyle C(r)$ $\displaystyle=$ $\displaystyle 2\pi n_{\rm
imp}\left(\frac{e^{2}}{\kappa}\right)^{2}\int_{0}^{\infty}dq\frac{q\exp(-2qd)}{[q\epsilon(q,n_{\rm
eff})]^{2}}J_{0}(qr),$ (2)
where $\kappa$ is the bulk (3D) dielectric constant, $\epsilon(q)$ is the
surface (2D) screening function in the plane, $-e$ the electron charge, and
$J_{0}(x)$ is a Bessel function.
As an illustration, consider the Thomas-Fermi (TF) screening for which the
surface dielectric function is given by $\epsilon(q,n_{\rm eff})=1+q_{\rm
TF}(n_{\rm eff})/q$, where $q_{\rm TF}(n_{\rm eff})$ (discussed below) is the
Thomas-Fermi screening wavevector. Shown in Fig. 1 is a calculation of $C(r)$
for different values of $q_{\rm TF}$ and $d$. Notice that for fixed $q_{\rm
TF}$, the spatial dependence of $C(r)$ depends on $d$ and $q_{\rm TF}$, but
not on $n_{\rm imp}$. On the other-hand, the function $C(0)$ (which, within
the TF can be calculated analytically) depends on $n_{\rm imp}$, $\kappa$, $d$
and $q_{\rm TF}$. This is why we find it useful to use the normalized
correlation function $A(r)=C(r)/C(0)$ that describes the spatial profile of
the screened impurity potential. We also emphasize that $A(r)$ contains
different information than the typical puddle size – for example, the puddle
correlation length $\xi$ (recall that $\xi$ is defined as the HWHM of $A(r)$)
describes the width of the screened impurity potential, and not the mean
impurity separation. For example, in the low impurity density limit, spatial
maps of the puddles would show isolated impurities, but $A(r)$ would not
change (for fixed $q_{\rm TF}$). Moreover, since $A(r)$ scales differently
with $q_{\rm TF}$ and $d$, in principle, both of these length scales can be
extracted from a measurement of $A(r)$.
Figure 1: Theoretical calculations for the correlation function $C(r)$ using
the Thomas-Fermi approximation. This autocorrelation function depends
separately on the typical distance $d$ of long-ranged impurities from the
graphene sheet, and $q_{\rm TF}$ the inverse effective screening length,
allowing them to be determined independently. Symbols show the correlation
length $\xi$, defined as the HWHM length.
Within the TF screening theory, any differences between monolayer and bilayer
graphene can only arise from differences in $q_{\rm TF}(n_{\rm eff})$. The
linear dispersion in monolayer graphene and the hyperbolic dispersion in
bilayer graphene gives rise to these differences. The low energy linearly
dispersing bands of monolayer graphene can be modeled by a single parameter,
the Fermi velocity $v_{\rm F}$, or equivalently, the effective fine-structure
constant $r_{s}=e^{2}/(\kappa\hbar v_{\rm F})\approx 0.8$. $r_{s}$
characterizes the strength of the electron-electron interaction for graphene
on a SiO2 substrateDas Sarma et al. (2011) and is useful because we are
interested in the screening properties of graphene. The Thomas-Fermi screening
wavevector is related to the density of states, and for monolayer graphene is
given by $q_{\rm TF}(n_{\rm eff})=4r_{s}\sqrt{\pi n_{\rm eff}}$, where $n_{\rm
eff}$ is the effective carrier density.
Bilayer graphene can be modeled with a hyperbolic dispersion with two
parameters $v_{\rm F}$ (throughout this manuscript, $v_{\rm F}$ is the Fermi
velocity of a single decoupled graphene sheet), and the low-energy effective
mass $m_{\rm eff}$. For simplicity we use for the two parameters $r_{s}$
(defined above) and $n_{0}=m_{\rm eff}^{2}v_{\rm F}^{2}/(\hbar^{2}\pi)\approx
2.3\times 10^{12}~{}{\rm cm}^{-2}$ which is the characteristic density scale
for the crossover from a (low density) parabolic to a (high density) linear
dispersion. The Thomas-Fermi screening wavevector for bilayer graphene is
given by $q_{\rm TF}(n_{\rm eff})=4r_{s}\sqrt{\pi n_{0}}\sqrt{1+n_{\rm
eff}/n_{0}}$.
As the system approaches the Dirac point, the fluctuations in carrier density
become larger than the average density. In this case, screening varies
spatially with the density fluctuations. We assume that it is possible to
describe the effect of this screening by using the screening for an ideal
system and an effective carrier density $n_{\rm eff}$ obtained self-
consistently.Adam et al. (2007) This is done by equating the squared Fermi
level shift with respect to the Dirac point with the square of the potential
fluctuations, $E^{2}[n=n_{\rm eff}]=C(0)$ where $C(0)$ is defined in Eq. 2 and
$E[n]=\hbar v_{\rm F}\sqrt{\pi n}$ for monolayer graphene, and $E[n]=v_{\rm
F}^{2}m_{\rm eff}\left[\sqrt{1+n/n_{0}}-1\right]$ for bilayer graphene. The
result of this procedure are shown in Fig. 2.
Figure 2: Effective carrier density as a function of impurity density
assuming $d=1~{}{\rm nm}$, $r_{s}=0.8$, and $n_{0}=2.3\times 10^{12}~{}{\rm
cm}^{-2}$. For bilayer graphene, the blue circles show the Thomas-Fermi
approximation and the red squares are RPA results. The empirical relation
$n_{\rm eff}=\sqrt{n_{\rm imp}~{}n_{1}}$ adequately captures the RPA results,
with $n_{1}=6.8\times 10^{11}~{}{\rm cm}^{-2}$.
Within the TF theory, we can now qualitatively discuss the main differences
between monolayer and bilayer graphene. For bilayer graphene, the inverse
screening length changes only slightly from the low density value of $q_{\rm
BLG}\approx 4r_{s}\sqrt{\pi n_{0}}$ that is set entirely by the band
parameters. As a consequence, the puddle correlation length does not change
with the impurity concentration or carrier density, and depends only on the
distance $d$ of the bilayer graphene sheet from the source of the long-ranged
impurity potential. In contrast, for monolayer graphene, $q_{\rm
MLG}=4r_{s}\sqrt{\pi n_{\rm eff}}$ depends essentially on $C(0)$ (and
therefore on $n_{\rm imp}$). This heuristic description (which we make more
quantitative below) implies that depending on the sample quality, choice of
substrate, or doping, the puddle correlation length in monolayer graphene (but
not bilayer graphene) could vary by more than an order of magnitude.
## III Random Phase Approximation
While the TF screening theory discussed in the previous section is useful to
obtain a qualitative picture, we find that it significantly underestimates the
effect of electronic screening. In both monolayer and bilayer graphene it
gives larger values for $C(0)$ and smaller values for $\xi$. In what follows
we use the the random phase approximation (RPA) where the screening function
is obtained using $\epsilon(q)=1+q_{\rm TF}{\tilde{\Pi}}(q)/q$.
The normalized polarizability ${\tilde{\Pi}}(q)$ for monolayerHwang and Das
Sarma (2007) and bilayerGamayun (2011) graphene are both available in the
literature. We note that for monolayer graphene ${\tilde{\Pi}}(q)$ depends
only on the dimensionless variable $x=q/(2\sqrt{\pi n_{\rm eff}})$
$\displaystyle{\tilde{\Pi}}(x)$ $\displaystyle=$ $\displaystyle
1+\theta(x-1)\left[\frac{\pi x}{4}-\frac{x}{2}{\rm
arccsc}(x)-\frac{\sqrt{1-x^{-2}}}{2}\right],$ (3) $\displaystyle\approx$
$\displaystyle\theta(1-x)+\theta(x-1)\frac{\pi x}{4},$
where $\theta(x)$ is a step-function and $D_{0}$ is the density of states with
$q_{\rm TF}=2\pi(e^{2}/\kappa)D_{0}=4\sqrt{\pi n_{\rm eff}}r_{s}$.
In contrast, for bilayer graphene, the polarizability depends both on the
scaled momentum transfer $x$, and on $\eta=n_{\rm eff}/n_{0}<8$, which
parameterizes the bilayer hyperbolic dispersion relation. For $\eta\ll 1$, the
bilayer graphene dispersion is quadratic, while for $1\ll\eta\leq 8$, the
dispersion is linear. For $\eta>8$, one must consider the effects of a second
higher-energy band that provides additional screeningMin et al. (2011) and is
not considered here. By restricting the density to $n\leq 8n_{0}$, we can
simplify the expression for the bilayer polarizability reported in Ref.
Gamayun, 2011. Using the bilayer density of states $D=D_{0}\sqrt{1+\eta}$,
with $D_{0}=2m/(\pi\hbar^{2})$, the normalized polarizability function
${\tilde{\Pi}}(x,\eta)=\Pi(q)/D_{0}$ is given by
$\displaystyle{\tilde{\Pi}}(x,\eta)$ $\displaystyle=$ $\displaystyle
f(x,\eta)+\theta(x-1)g(x,\eta),$ $\displaystyle f(x,\eta)$ $\displaystyle=$
$\displaystyle\left[1-x^{2}\eta+x^{4}(2+\eta+2\sqrt{1+\eta})\right]^{1/2}-\ln\left[\frac{2\sqrt{1+\eta
x^{2}}}{-1+\sqrt{1+\eta}}\right]-\frac{1}{2}+\frac{3x^{2}\eta-1}{2x\sqrt{\eta}}\arctan(x\sqrt{\eta})$
$\displaystyle\mbox{}+\sqrt{1-x^{2}\eta}\left(2{\rm arctanh}(\sqrt{1-\eta
x^{2}})-{\rm
arcsinh}\left[\frac{\sqrt{1-x^{2}\eta}(-1+\sqrt{1+\eta})}{x^{2}\eta}\right]\right)+\sqrt{1+\eta}$
$\displaystyle g(x,\eta)$ $\displaystyle=$
$\displaystyle\frac{-\sqrt{x^{2}-1}(1+\eta+2x^{2}\eta-\sqrt{1+\eta})}{2x(\sqrt{1+\eta}-1)}+\frac{(3x^{2}\eta-1)}{2x\sqrt{\eta}}\arccos\left[\sqrt{\frac{1+\eta}{1+x^{2}\eta}}\right]$
(4) $\displaystyle\mbox{}+{\rm
arctanh}\left[\frac{x\sqrt{x^{2}-1}\eta}{1+x^{2}\eta-\sqrt{1+\eta}}\right].$
The polarizability functions for monolayer and bilayer graphene are shown in
Fig. 3. What is left is to calculate the effective residual density $n_{\rm
eff}$ as a function of impurity concentration. As discussed earlier, this is
obtained by first calculating the autocorrelation function $C(0)$ from Eq. 2,
using the RPA results shown in Fig. 3. Figure 4 shows the autocorrelation
function $C(0)$ obtained using the RPA results (solid lines) as well the
Thomas-Fermi results (dashed lines). We note that except at very high density
(where both monolayer and bilayer graphene approach the “complete screening”
limit, with $C(0)=(4k_{\rm F}r_{s}d)^{-2}$), the Thomas-Fermi approximation
grossly underestimates the effect of screening. Moreover, for typical
densities in bilayer graphene $C(0)$ is approximately constant and independent
of carrier density consistent with the heuristic picture discussed at the end
of Sec. II.
Figure 3: The Random Phase Approximation polarizability function $\Pi(q)$
normalized by the density of states for monolayer graphene and for bilayer
graphene with a hyperbolic dispersion. Also shown is the parabolic
approximation for the bilayer, which can be obtained from the hyperbolic
dispersion when $\eta\ll 1$. The Thomas-Fermi approximation discussed in the
text corresponds to the assumption that the normalized $\Pi(q)=1$ for all $q$.
Figure 4: Potential autocorrelation function $C[0]$ for monolayer and bilayer
graphene. At the Dirac point, $k_{\rm F}$ is the Fermi wavevector arising from
the effective carrier density i.e. $k_{\rm F}=\sqrt{\pi n_{\rm eff}}$. For
large density $4k_{\rm F}r_{s}d\gg 1$, both the monolayer and bilayer results
approach the “complete screening” limit, defined here as $C(0)=(4k_{\rm
F}r_{s}d)^{-2}$. Notice that the Thomas-Fermi approximation shown as dashed
lines captures the correct qualitative behavior, but can give significantly
larger values for $C[0]$, and is therefore unsuitable for quantitative
comparisons.
The effective density calculated within the RPA is shown in Fig. 2. For
bilayer graphene, we find that the following empirical relationship adequately
describes the numerical results
$n_{\rm eff}=\sqrt{n_{\rm imp}~{}n_{1}},$ (5)
where $n_{1}=11.5\times 10^{11}~{}{\rm cm}^{-2}$ for the Thomas-Fermi
approximation, and $n_{1}=6.8\times 10^{11}~{}{\rm cm}^{-2}$ for the RPA
results. The scaling of the bilayer effective density $n_{\rm
eff}\sim\sqrt{n_{\rm imp}}$ can be anticipated for the TF approximation in the
limit $n_{\rm eff}\ll n_{0}$. However, it is surprising that the simple
empirical relation continues to hold both for the RPA screening theory, and
for larger values of $n_{\rm eff}$.
This dependence of $n_{\rm eff}\sim\sqrt{n_{\rm imp}}$ for bilayer graphene
should be contrasted with similar results obtained previously for monolayer
graphene,Adam et al. (2007) where $n_{\rm eff}=2r_{s}^{2}C[0]n_{\rm imp}$
cannot be captured by a similar empirical fit. For comparison, these earlier
results are also shown in Fig. 2, where we emphasize that for a given impurity
concentration ($n_{\rm imp}$), bilayer graphene exhibits larger density
fluctuations ($n_{\rm eff}$) than monolayer graphene. Finally, using Eq. 2, we
can also calculate the puddle correlation function, and the corresponding HWHM
correlation length, $\xi$, that is shown in Fig. 5
Figure 5: Theoretical results for the puddle correlation length at the Dirac
point as a function of impurity concentration. While the puddle size in
bilayer graphene ($\xi\approx 3.5~{}{\rm nm}$) is relatively insensitive to
the disorder concentration, the size of the puddles in monolayer graphene
varies from $3~{}{\rm nm}$ in dirty samples to more than $35~{}{\rm nm}$ in
clean samples.
## IV Comparison With Experiments
We now compare the calculated correlation functions with experiment. Figure
6(a) shows that for bilayer graphene, $A(r)$ extracted from the data reported
in Ref. Rutter et al., 2011 agrees with the calculation for $d=1~{}{\rm
nm}$.kn: (a) The circles show the experimental data and the RPA theory for
bilayer graphene is shown for $d=1~{}{\rm nm}$ (solid curve) and $d=0.5~{}{\rm
nm}$ (dashed curve). The theoretical results are insensitive to the impurity
concentration $n_{\rm imp}$ and to how far the doping is away from the Dirac
point. Consequently, the only free parameter in the theory is the distance $d$
of the impurities from the graphene sheet.
Figure 6: Comparison of theoretical results with experimental data. Top panel
shows the normalized correlation function $A(r)=C(r)/C(0)$ for bilayer
graphene. The circles are from the experimental data and solid curve is the
theory for bilayer graphene with $d=1~{}{\rm nm}$. The theory curve is
insensitive to impurity concentration and doping away from the Dirac point.
The error bars indicate single standard deviation uncertainties.kn: (a) The
small oscillation in the data over the monotonic decrease is a result of the
finite size of the experimental image. Bottom panel is the normalized puddle
correlation function in monolayer graphene at the Dirac point. Note the change
in $x$-axis scale from bilayer graphene in top panel. The solid curve is
obtained from the self-consistent screening theory. The black squares are the
results of a numerical mesoscopic density functional theory calculation for
the ground-state properties of monolayer graphene,Rossi and Das Sarma (2008)
while the circles are experimental data taken from Deshpande et al.(Ref.
Deshpande et al., 2011). Transport measurements on that same device set
$n_{\rm imp}=10^{11}{\rm cm}^{-2}$ which is the value used for the theory
curves. The theory also uses $d=1~{}{\rm nm}$, which is the typical distance
of the impurities from the graphene sheet extracted from transport
measurements of graphene on SiO2.Tan et al. (2007)
In Ref. Rutter et al., 2011, we reported a maximum peak-to-peak carrier
density fluctuation of $3.6\times 10^{11}~{}{\rm cm}^{-2}$. To extract an
impurity density (using the results in Fig. 2), we need to estimate $n_{\rm
rms}$ from this peak-to-peak value. By assuming that the carrier density has a
Gaussian distribution, and estimating that the peak-to-peak corresponds to a
measurement of $4\sigma$, we roughly estimate that $n_{\rm eff}=n_{\rm
rms}\approx 10^{11}~{}{\rm cm}^{-2}$ and that $n_{\rm imp}=n_{\rm
rms}^{2}/n_{1}\approx 1.2\times 10^{10}~{}{\rm cm}^{-2}$ for the substrate
induced impurities in that experiment.
We can use this value to predict theoretically the corresponding density
fluctuations in the adjacent monolayer sample reported in Ref. Jung et al.,
2011. However, one complication is that the theory discussed in Sec. III was
developed for monolayer graphene at the Dirac point, while the experimental
data was taken at a backgate induced density $n_{g}=1.4~{}\times
10^{12}~{}{\rm cm}^{-2}$. Very far from the Dirac point, i.e. when $n_{\rm
imp}/n_{g}\rightarrow 0$, the potential fluctuations $V_{\rm rms}$ can be
obtained from Eq. 2 by setting $n_{\rm eff}=n_{g}$ on the right-hand side. In
this case, the density fluctuations are
$n_{\rm rms}=\frac{2VV_{\rm rms}}{\pi\hbar^{2}v_{\rm F}^{2}}.$ (6)
We note that when $z\sim d\sqrt{n_{g}}\gg 1$, we can use the result
$C_{0}(z)=z^{-2}$ (see Fig. 4) to obtain $n_{\rm rms}\approx\sqrt{n_{\rm
imp}/(8\pi d^{2})}$. However, these constraints are not fully satisfied in the
experimental data. Calculating $n_{\rm rms}$ in the crossover between the
limits $n_{g}=0$ and $n_{g}\gg n_{\rm imp}$ is more complicated. For our
purposes, it is sufficient to extrapolate between the low-density and high-
density limits by adding the two contributions in quadrature, and solving for
$n_{\rm rms}$ self-consistently. This procedure gives
$\displaystyle n_{\rm rms}=2r_{s}\sqrt{n_{\rm imp}C^{\rm
RPA}(0)}\left[2n_{g}+3r_{s}^{2}n_{\rm imp}C^{\rm RPA}(0)\right]^{1/2},$ (7)
where the superscript indicates that the RPA screening approximation has been
used. In the limit that $n_{g}\gg n_{\rm imp}$, Eq. 7 reduces to Eq. 6, while
in the opposite limit $n_{g}\rightarrow 0$, Eq. 7 reduces to results shown in
Fig. 2.
Using the values for $n_{\rm imp}$ and $d$ determined from the bilayer data
discussed above, and using Eq. 7 for $n_{\rm rms}$ and Eq. 2 to calculate
$\xi$, we find theoretically (without any adjustable parameter) that $n_{\rm
rms}\approx 6~{}\times 10^{10}~{}{\rm cm}^{-2}$ and $\xi=8~{}{\rm nm}$, which
should be compared to the experimental values extracted from the data reported
in Ref. Jung et al., 2011. The area surveyed in Ref. Jung et al., 2011 was not
large enough to obtain $\xi$ accurately. However, by looking at different
real-space cuts of the autocorrelation function, we conclude that the
experimental data is consistent with a correlation length $6~{}{\rm
nm}<\xi<11~{}{\rm nm}$. This is in qualitative agreement with our theoretical
calculations. This result should be contrasted with bilayer graphene shown in
Fig. 6(a) where the experimentally determined $\xi=(3.68\pm 0.03)~{}{\rm nm}$,
and the $d=1~{}{\rm nm}$ RPA theory gives $\xi=3.4~{}{\rm nm}$.
To further confirm our results, we compare our calculations to measurements
made using scanning Coulomb blockade spectroscopy on a sample of monolayer
graphene.Deshpande et al. (2011) The circles in Fig. 6(b) are experimental
data for the normalized correlation and the solid line is the self-consistent
theory discussed above using $n_{\rm imp}=10^{11}{\rm cm}^{-2}$ and
$d=1~{}{\rm nm}$ for monolayer graphene at the Dirac point. The impurity
concentration and impurity distance were determined from transport
measurementsDeshpande et al. (2011); Adam et al. (2007) and as such, no
adjustable parameters were used in the calculation.
In Fig. 6 we also show (black squares) the results extracted from a numerical
mesoscopic density functional theoryRossi and Das Sarma (2008) using the same
parameters. The agreement between the two calculations provides a posteriori
justification for our assumption of a global screening function characterized
by the density $n_{\rm eff}$.
## V Conclusions
We conclude with the observation that our results require only that the source
of the disorder potential be uncorrelated charged impurities, and as such
should apply to graphene on other substrates. For example, recently graphene
devices with hexagonal BN gate insulators have been fabricated showing
transport properties similar to suspended grapheneDean et al. (2010) and
larger puddles than on SiO2 substrates.Xue et al. (2011); Decker et al. (2011)
These observations are consistent with both a much smaller charged impurity
density $n_{\rm imp}$ on the BN substrate and with a larger distance $d$ of
the impurities from the graphene layer. Both these scenarios are possible
because the BN substrate is typically placed on top of the usual SiO2 wafer
which would have similar charged disorder to the samples we study here. We
argue that an analysis similar to what we have performed here would be able to
uniquely determine both $d$ and $n_{\rm imp}$. Moreover, if a similar
experiment is done with bilayer graphene on BN substrates, we predict that the
puddle characteristics will not change much from what we find here with
bilayer graphene on SiO2.
## Acknowledgements
This work is supported in part by the NIST-CNST/UMD-NanoCenter Cooperative
Agreement. It is a pleasure to thank W. G. Cullen, M. S. Fuhrer, and M. Polini
for discussions, and P. W. Brouwer, E. Cockayne, G. Gallatin, J. McClelland
and R. McMichael for comments on the manuscript.
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* Jung et al. (2011) S. Jung, G. M. Rutter, N. N. Klimov, D. B. Newell, I. Calizo, A. R. Hight-Walker, N. B. Zhitenev, and J. A. Stroscio, Nature Phys. 7, 245 (2011).
* Rutter et al. (2011) G. R. Rutter, S. Jung, N. N. Klimov, D. B. Newell, N. B. Zhitenev, and J. A. Stroscio, Nature Phys. 7, 649 (2011).
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* Gamayun (2011) O. V. Gamayun, Phys. Rev. B 84, 085112 (2011).
* Min et al. (2011) H. Min, P. Jain, S. Adam, and M. D. Stiles, Phys. Rev. B p. 195117 (2011).
* kn: (a) The dominant uncertainty in $A(r)$ and $\xi$ comes from the finite size of the experimental puddle images. The error in $A(r)$ is estimated by dividing the image area into $N=4$ patches and analyzing each patch separately. Treating the $N$ patches as independent samples, we compute the variance of the set and estimate the variance of the whole sample as the variance of the $N$ independent samples, divided by $N-1$. The uncertainty in $A(r)$ is square root of the variance, and the uncertainty in $\xi$ then follows. We say that the experimentally determined $A(r)$ agrees with the calculated $A(r)$ for $d=1~{}{\rm nm}$ in Fig. 6a because theory curve lies within the experimental uncertainty of $A(r)$.
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* Tan et al. (2007) Y.-W. Tan, Y. Zhang, K. Bolotin, Y. Zhao, S. Adam, E. H. Hwang, S. Das Sarma, H. L. Stormer, and P. Kim, Phys. Rev. Lett. 99, 246803 (2007).
* Dean et al. (2010) C. R. Dean, A. F. Young, I. Meric, C. Lee, L. Wang, S. Sorgenfrei, K. Watanabe, T. Taniguchi, P. Kim, K. L. Shepard, et al., Nature Nano. 5, 722 (2010).
* Xue et al. (2011) J. Xue, J. Sanchez-Yamagishi, D. Bulmash, P. Jacquod, A. Deshpande, K. Watanabe, T. Taniguchi, P. Jarillo-Herrero, and B. J. Leroy, Nature Mat. 10, 282 (2011).
* Decker et al. (2011) R. Decker, Y. Wang, V. W. Brar, W. Regan, H.-Z. Tsai, Q. Wu, W. Gannett, A. Zettl, and M. F. Crommie, Nano Lett. 11, 2291 (2011).
|
arxiv-papers
| 2011-12-05T21:00:02 |
2024-09-04T02:49:24.989095
|
{
"license": "Public Domain",
"authors": "Shaffique Adam, Suyong Jung, Nikolai N. Klimov, Nikolai B. Zhitenev,\n Joseph A. Stroscio, M. D. Stiles",
"submitter": "Shaffique Adam",
"url": "https://arxiv.org/abs/1112.1070"
}
|
1112.1111
|
# Investigating the origin of the long range pseudo rapidity correlation in 2d
di-hadron measurements from STAR††thanks: Presented at Strange Quark Matter
conference, September, 2011, Krakow, Poland.
L. C. De Silva The Department of Physics, University of Houston, 617 SR
Building 1, Houston, Texas 77204-5005
###### Abstract
Angular di-hadron correlations reveal novel structures in central Au+Au
collisions at $\sqrt{S_{NN}}$ = 200 GeV. One of them, known as the ridge, is
elongated in pseudo rapidity and peaks on the near side ($\Delta\phi$
$\approx$ 0). Investigating the origin of the ridge structure helps to
understand the hot dense matter that is created in ultra relativistic heavy
ion collisions. Results showing the $\langle$$p_{T}$$\rangle$ dependence of
the ridge structure are presented. Evidence for possible jet and non-jet
contributions to the ridge structure will be discussed.
## 1 Introduction
Recent triggered di-hadron correlation studies by STAR report a ridge
structure in two dimensions ($\Delta\eta$,$\Delta\phi$)[1]. Based on their
studies, the near side was assumed to consist of two independent structures, a
jet-like structure and the ridge. A complementary approach, presented here, is
carried out by using all possible charged particle pairs. The approach does
not require a trigger particle but with appropriate kinematic cuts reproduces
qualitatively similar correlation structure to that of triggered analysis. In
this approach, two empirical 2d fit functional models are used to extract
ridge properties. Possible model dependent quantitative evidence for jet and
non-jet phenomena contribution to the near side structure can be extracted
from the fits.
## 2 Data and analysis
The data used in this analysis were collected during run 4, year 2004, from
the STAR detector at Relativistic Heavy Ion Collider (RHIC) Brookhaven
National Lab (BNL), Long Island, New York. 32M Au+Au collisions at
$\sqrt{S_{NN}}$ = 200 GeV were analyzed. Charged tracks reconstructed using
the Time Projection Chamber (TPC) with 2$\pi$ azimuthal coverage and $|\eta|$
$\leq$ 1 in pseudo rapidity were selected. An earlier centrality evolution
analysis[2] was based on particles in the full transverse momentum range above
0.15 GeV/c. For the study presented here the lower threshold of the transverse
momentum was raised for both particles. The selected event vertexes are within
$\pm$ 25cm of the detector center and primary tracks are selected to come from
within a distance of closet approach of 3cm to the event vertex.
The correlation function for this analysis is a construct of Pearson s
correlation coefficient. Mathematically it can be derived as[2]:
$\displaystyle\frac{\Delta\rho}{\sqrt{\rho_{ref}}}$
$\displaystyle=\frac{\rho_{sib}\\--\rho_{mix}}{\sqrt{\rho_{ref}}}$ (1)
$\rho_{sib}$ is the normalized charged particle pair density within a single
event, which carries both correlated and uncorrelated pairs. The uncorrelated
background is subtracted by mixed event pairs. The resulting uncorrelated
density is denoted as $\rho_{mix}$. The difference is the number of correlated
pairs which is normalized by the square root of mixed pair density to yield
the _correlated particle pairs per final state charged particle_ as the
equivalent correlation measure to the Pearson s correlation coefficient. In
order to avoid possible artifacts due to pseudo rapidity acceptance, the
correlation measure is carried out in sub bins of the z-vertex. The event
mixing artifacts were eliminated by mixing events of a particular centrality
within a multiplicity window of 50 charged tracks. To correct for pair loss a
pair cut is implemented on both mixed and same event pairs. The effects due to
pileup events were removed by one of the standard pileup eliminating
procedures in STAR[7]. Finally the conversion electron positron pair
background is reduced by using a 1.5$\sigma$ dE/dX cut on electrons in
momentum ranges, 0.2 <$p_{T}$ <0.45 GeV/c and 0.7 <$p_{T}$ <0.8 GeV/c.
Figure 1: Evolution of the raw normalized correlation structure for selected
$\langle$$p_{T}$$\rangle$ bins at 0-1% centrality bin. The lower momentum
threshold of both particles have been increased in generating the spectrum of
correlation plots.
## 3 Fit functions
The initial empirical fit function carries three mathematical model components
that lead to seven free fit parameters.
F = $c_{0}$ \+ $c_{1}$cos($\Delta\phi$) +
$c_{2}$exp(-0.5*(($\Delta\phi$/$c_{3}$)2 \+ ($\Delta\eta$/$c_{4}$)2) +
$c_{5}$exp(-0.5*(($\Delta\phi$/$c_{6}$)2 \+ ($\Delta\eta$/$c_{6}$)2))
The first component of the fit is an offset parameter followed by a
cos($\Delta\phi$) term to extract the away side structure as seen in Fig. 1.
The assumption is that this structure corresponds to correlations due to a
recoiling jet[3]. The two remaining terms are introduced based on
$\langle$$p_{T}$$\rangle$ evolution of data (see Fig. 1). The asymmetric 2d
Gaussian attempts to address the long range correlation properties excluding
the peak structure around (0,0) that appears as $\langle$$p_{T}$$\rangle$
increases. This structure is described via the symmetric 2d Gaussian component
of the fit.
Modifications to the model were introduced after indications of higher order
harmonics were first predicted [4, 5] and then observed (see Fig. 2).
F = $c_{0}$ \+ $c_{1}$cos($\Delta\phi$) + $c_{2}$cos(2$\Delta\phi$) +
$c_{3}$cos(3$\Delta\phi$) + $c_{4}$cos(4$\Delta\phi$) +
$c_{5}$cos(5$\Delta\phi$) + $c_{6}$exp(-0.5*(($\Delta\phi$/$c_{7}$)2 \+
($\Delta\eta$/$c_{8}$)2))
The resulting fit carries higher order harmonics up to 5th order and an
asymmetric 2d Gaussian. It is important to note that in such a Fourier
expansion, the first order term serves as a momentum conserving term
(correlation due to recoil jet). In essence, the new model help in further
constraining the knowledge gathered from the initial model study.
Figure 2: Correlations at high $p_{T}$ very central data show evidence for the
existence of higher order harmonics.
## 4 Results
The focus of comparing the initial model study to the study which includes fit
components for higher harmonics based on initial energy density fluctuations
is to determine whether any information can be extracted from the asymmetric
2d Gaussian fit to the away side structure.
### 4.1 Initial model study
Resulting asymmetric 2d Gaussian parameters are reported in Fig. 3. The long
correlation strength approaches zero at higher $\langle$$p_{T}$$\rangle$ as
indicated by the asymmetric 2d Gaussian amplitude and volume parameters. The
$\Delta\eta$ and $\Delta\phi$ width evolution clearly indicate the asymmetry
shown in data. Also the $\Delta\eta$ width suggests that within the STAR
acceptance the ridge is flat whereas the $\phi$ width shows a very smooth
narrowing. An important observation is that the $\langle$$p_{T}$$\rangle$ of
the ridge ($\langle$$p_{T}$$\rangle$ = 0.72 GeV/c) is greater than that of the
inclusive spectrum ($\langle$$p_{T}$$\rangle$ $\approx$ 0.42 GeV/c [8]). The
hardness of the ridge spectrum is an indication of possible Jet contributions
to the ridge.
(a) Amplitude
(b) $\Delta\eta$ width
(c) $\Delta\phi$ width
(d) Symmetric width
Figure 3: Asymmetric and Symmetric 2d Gaussian parameter evolution as a
function of $\langle$$p_{T}$$\rangle$.
On the other hand, the symmetric 2d Gaussian component reveals important
information about an unmodified jet emerging at higher momenta (Fig. 3). The
width is comparable to the jet width in PP collisions. In order to further
constrain the above conclusion, amplitude and yield comparisons need to be
made.
(a) $\frac{v_{3}}{v_{2}}$ ratio
(b) $v_{n}\hskip 2.84544ptscaling\hskip 2.84544ptat\hskip 2.84544pt0\\--10\%$
Figure 4: Comparison to predicted higher harmonic ratios[5] and scaling
relations[4]
### 4.2 Higher harmonics study
Recent advancements in heavy ion collision models, predict significant
contributions from higher order harmonics to the observed correlation spectra
due to initial energy density fluctuations[4, 5]. STAR data also support the
predicted observation in high $p_{T}$, very central AuAu 200 GeV data(Fig. 2).
The $\frac{v_{3}}{v_{2}}$ ratio predictions were tested and reasonable
agreement to theory is evident (see Fig. 4a). However it is to be noted at low
$p_{T}$, deviations from the trends are observed for $N_{part}$ <150\. One
should note, thought, that the data is extracted fitting 2d di-hadron
correlations, whereas, theory uses single particle spectra and event plane
calculations.
A predicted hydro scaling relation[6] using extracted higher order Fourier
coefficients is studied, and reasonable agreement to theory has been observed
(Fig. 4b).In summary, given the different approaches used by theory
predictions tested in this proceedings and experiment, the data trends seem
support the hydro based theory predictions to a reasonable extent.
The remaining structure on the near side, after subtracting the Fourier
components still shows a $\Delta\eta-\Delta\phi$ asymmetry (see Fig. 3, red
data markers). This $\langle$$p_{T}$$\rangle$ evolution indicates that this
$\Delta\eta$ elongation is visible up to 2.7 GeV/c and is strongest at 0.9
GeV/c. The structure thus suggests possible jet modification at low
$\langle$$p_{T}$$\rangle$ which evolves to an unmodified jet at high $p_{T}$.
The amplitude and volume of the asymmetric structure scales lineary with
Glauber scaling as a function of centrality, which hints at a jet origin. The
rise of amplitude at high $\langle$$p_{T}$$\rangle$ is expected due to an
increase in per charge particle pair yield in jets. Further studies will be
carried out using the comparison to p+p collisions.
## 5 Summary and discussion
Un-triggered 2d di-hadron correlation studies reproduce qualitatively similar
results to that of a triggered analysis. The observed near side correlation is
modeled via an empirical fit function which extracts short and long range
structure properties. The $\langle$$p_{T}$$\rangle$ evolution of the extracted
parameters suggests a possible correlation between jets and the observed long
range correlation. Further constraint to data has been imposed modifying the
empirical model to incorporate higher order Fourier model components
cos(n$\Delta\phi$), n = 1– 5, based on theoretical predictions and evidence in
the data. The extracted higher order component strengths show reasonable
agreement to predicted hydrodynamical trends. The remainder on the near side
still reveals an asymmetric 2d Gaussian suggestive of possible modified jet
phenomena. A Glauber linear scaling was successfully applied to the un-
triggered centrality evolution (see Fig. 5) which indicates a jet origin for
the observed near side Gaussian.
Figure 5: Comparison of near side asymmetric 2D Gaussian amplitude to Glauber
linear scaling as a function of centrality.
## References
* [1] B. I. Abelev et al. (STAR Collaboration), Phys. Rev. C 80, 064912 (2009)
* [2] M Daugherity et al. (STAR Collaboration), J. Phys. G 35, 104090 (2008)
* [3] J. Porter and T. Trainor, J. Phys.: Conf. Ser. 27, 98 (2005)
* [4] Xin Niang Wang et al., Phys. Rev. Lett. 106, 162301(2011)
* [5] B. Alver et al., Phys. Rev. C 81, 054905 (2010)
* [6] C. Gombeaud et al., Phys. Rev. C 81, 014901 (2010)
* [7] Nuclear Phys. Lab. Annual Report, University of Washington (2009) p. 58
* [8] J. Adams et al. STAR Collaboration
|
arxiv-papers
| 2011-12-05T22:13:47 |
2024-09-04T02:49:24.997087
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "L. C. De Silva (for the STAR collaboration)",
"submitter": "Chanaka De Silva Mr.",
"url": "https://arxiv.org/abs/1112.1111"
}
|
1112.1139
|
# Quantum Verification of Minimum Spanning Tree
Mark Heiligman Intelligence Advanced Research Projects Activity, Office of
the Director of National Intelligence, Washington, D.C.
mark.i.heiligmanugov.gov
(January 29, 2011 ; 1926 ; 2006 ; 1997 ; 1997 ; 1985 )
Previous studies has shown that for a weighted undirected graph having $n$
vertices and $m$ edges, a minimal weight spanning tree can be found with
$O^{*}\bigl{(}\sqrt{mn}\bigr{)}$ calls to the weight oracle. The present note
shows that a given spanning tree can be verified to be a minimal weight
spanning tree with only $O\bigl{(}n\bigr{)}$ calls to the weight oracle and
$O\bigl{(}n+\sqrt{m}\log n\bigr{)}$ total work.
quantum algorithms, graph theory, spanning tree
††support: Disclaimer. All statements of fact, opinion, or analysis expressed
in this paper are solely those of the author and do not necessarily reflect
the official positions or views of the Office of the Director of National
Intelligence (ODNI), the Intelligence Advanced Research Projects Activity
(IARPA), or any other government agency. Nothing in the content should be
construed as asserting or implying U.S. Government authentication of
information or ODNI endorsement of the author’s views.
## Introduction
### Problem Statement
The determination of a minimal weight spanning tree of a weighted undirected
graph is a central problem in computational graph theory and a number of well
known classical algorithms address the problem quite efficiently. This problem
has also shown up in the realm of quantum algorithms and the paper [DHHM]
provides nearly matching upper and lower bounds for the problem. (The term
“nearly matching” as used here means that the upper and lower bounds agree to
withing a power of the logarithm of the problem size.) The algorithm in [DHHM]
uses some of the constructs that occur in the classical minimal spanning tree
algorithms, along with a somewhat sophisticated version of the quantum minimum
algorithm (which itself is based on Grover’s algorithm).
A closely related problem deals with the verification of minimal spanning
tree. In this formulation of the problem, both a weighted graph and a spanning
tree of that graph are given as inputs, and the problem is to decide whether
the given spanning tree is of minimal weight (and if not to give a lower
weight spanning tree). Based on work of [Ko], a simple classical verification
algorithm was given in [Ki].
###### Problem
Given a graph $G=(V,E)$ consisting of $n={\left|{G}\right|}$ vertices and
$m={\left|{E}\right|}$ edges along with a weight function on the edges
$w:\,E\rightarrow{\mathbb{R}}^{+}$, and a spanning tree $T=(V,F)$ with
$F\subset E$ and ${\left|{F}\right|}=n-1$, verify that $T$ is a minimal weight
spanning tree.
The goal of this paper is to develop a quantum algorithm for the verification
problem. We build heavily on the graph theory methods given in [Ki] and
[KPRS]. Our quantum tool in this case is a fairly simple version of Grover’s
algorithm. Nevertheless we are able to show that verification is simpler than
finding the solution ab initio.
### Computational Models
There is a basic question of how the graphs $G$ and $T$ are presented, and
this can critically affect the efficiency of the algorithm. A graph can be
presented by an adjacency matrix or by a simple listing of its edges (and this
may be either a sorted or an unsorted list).
In the classical world, the problem statement is fairly simple. In the quantum
world, the graph is presented to the algorithm as an oracle, and the
complexity of the algorithm is measured in the number of oracle calls
necessary to solve the problem.
Oracles can be applied in the classical world, as well, but they are less
indicative of the computational complexity of the problem than in the quantum
world. In the classical world, the entire graph needs to be made available to
the algorithm, so in the adjacency matrix model there would be
$O\bigl{(}n^{2}\bigr{)}$ calls to the oracle specifying the graph, while in
the edge list model there would be $m$ calls to the oracle simply to get the
entire graph into the computer.
There also has to be an oracle that gives the weight of an edge, and in the
adjacency matrix model or the edge list model, there would be $m$ calls to the
weight oracle. It is useful to combine the graph oracle with the weight
oracle. In the adjacency matrix model, the graph is extended to a complete
graph and weight $+\infty$ is assigned to all non-graph edges, with the oracle
being given as a function $w:\,V\times
V\rightarrow{\mathbb{R}}^{+}\cup\\{\infty\\}$. In the edge list model the
oracle is a function, $e:\,[1,m]\rightarrow V\times V\times{\mathbb{R}}^{+}$
where the first two components give the endpoints of the $i^{\hbox{\sevenrm
th}}$ edge of the graph $G$ and the last component gives the weight of that
edge.
In this note, we consider both of these models, but from the quantum
perspective, the oracle has to be viewed as a reversible function that then
operates on quantum states. The two models to be considered here are:
(1) There is the weight oracle in the adjacency matrix model. For this model,
a call to the quantum weight oracle is
${\mid{a,b,x}\rangle}\rightarrow{\mid{a,b,x\oplus w(a,b)}\rangle}$ where
$a,b\in V$ are a pair of vertices. (Note that $x$ here is just some arbitrary
initial bit string.) For finding minimum weight spanning trees, this is bad if
the graph is moderately sparse. For checking the minimality of a spanning
tree, the input would consist of a simple listing of the edges and would be of
length $n-1$.
(2) There is the combined edge list and edge weight in the edge list model.
For this model, a call to the quantum oracle is
${\mid{i,x,y,z}\rangle}\rightarrow{\mid{i,x\oplus a,y\oplus b,z\oplus
w}\rangle}$ where $a,b\in V$ are a pair of vertices such that $(a,b)$ is the
$i^{\hbox{\sevenrm th}}$ edge of the graph $G$ and $w=w(a,b)$ is its weight.
(Note that $x$, $y$, and $z$ here is just some arbitrary initial bit strings.)
Note that the result for model (1) above will give an upper bound for model
(2), but both models will be considered in this note. Thus in the adjacency
matrix model, there will be a weight oracle given $w:\,V\times
V\rightarrow{\mathbb{R}}^{+}\cup\\{+\infty\\}$ and the spanning tree to be
checked for minimality will be (classically) input as a list of edges
$T=\bigl{\\{}(a_{1},b_{1}),(a_{2},b_{2}),\ldots,(a_{n-1},b_{n-1})\bigr{\\}}$
with $(a_{i},b_{i})\in V\times V$ for $i=1,\ldots,n-1$. We will also consider
the $e$ oracle in the edge list model. However, even there, the spanning tree
to be checked for minimality will still be classically input as a list of
edges, only now the spanning tree to be checked for minimality will be
(classically) input as a list of edge indices $T=\bigl{\\{}e_{i},e_{2}\,\ldots
e_{n-1}\bigr{\\}}$ with $e_{i}\in[1,m]$.
In both of the above formulations, the subtree to be tested for minimality by
the quantum algorithm is input classically. This gives a lower bound for the
complexity of the quantum algorithm of $O(n)$, since the algorithm has to at
least read in all the (classical) input. However, there are other possible
statements of the problem.
(3) Given an oracle for the weights of $G$ (which is by default, also an
oracle for querying whether a given pair of points of $V$ is an edge of $G$),
the input could be by an oracle for the putative minimal spanning tree. Thus,
in the adjacency matrix model, there is a function $mst:V\times
V\rightarrow\\{0,1\\}$ where $mst(a,b)=1$ if $(a,b)$ is an edge in $T$ and
$mst(a,b)=0$ if $(a,b)$ is not an edge in $T$, while in the edge list
model,there is a function $mst:V\times[1,m]\rightarrow\\{0,1\\}$ where
$mst(i)=1$ if $i$ is an edge index in $T$ and $mst(i)=0$ if $i$ is not an edge
in $T$. In either case, the problem then becomes to determine whether $mst$ is
a correct oracle. The complication is that there are now two oracles to count
calls to, and in principle there could be an operation curve of tradeoffs.
In fact, this is almost certainly the case, because on one extreme the minimal
spanning tree can be found simply by computing it with a quantum algorithm and
then checking the $mst$ oracle for mismatches with the minimal spanning tree
found. This reduces to the problem: Given a set $S$ and a subset $U\subset S$
and a (quantum) oracle $p:\,S\rightarrow\\{0,1\\}$, is it the case that
$p(s)=1$ if and only if $s\in U$? Counting oracle calls here would seem to be
a simple application of Grover’s algorithm.
## Minimal Weight Spanning Trees
### Checking a Spanning Tree
The key observation from [Ki] is the following. For a graph $G=(V,E)$ and any
spanning tree $T$ of $G$, there is a unique path between any two edges $u,v\in
V$. $T$ is a minimal weight spanning tree if and only if the weight of each
edge $(u,v)\in E-T$ is greater than or equal to the the heaviest edge in the
path in $T$ between $u$ and $v$. What is needed is an easy way to find the
weight of the heaviest edge in the path in $T$ between $u$ and $v$.
The idea for checking a putative spanning tree $T$ for minimality is to show
that for any other edge of $G$ not in $T$, in the cycle formed by including
this edge, the highest weight edge in the cycle is exactly this edge. There is
no way that this edge can be part of a minimum weight spanning tree.
This is to be checked for all edges of $G$, so by invoking Grover’s algorithm
in the quantum setting, the total work is $O\bigl{(}\sqrt{m}\bigr{)}$ times
the work of checking an edge. The problem is that for checking an edge
$(u,v)\in E-T$, the length of the path in $T$ between $u$ and $v$ could be
very large, perhaps even as big as $n-1$, so even using Grover’s algorithm to
find the maximal weight edge on this path is not adequately efficient.
### Boruvka Trees
What is needed is a a new data structure that allows the maximal weight edge
on any path in $T$ to be found efficiently. The basic idea for this comes from
one of the earliest papers in computational graph theory [B], that was the
forerunner of several modern spanning tree algorithms. The key properties of
the Boruvka tree built from a putative minimal spanning tree come from [Ki]
and will be summarized here without proof.
The idea of a Boruvka tree built from a spanning tree $T$ is that $B_{T}$ is a
tree whose leaves are the vertices $V$ and whose internal nodes are to be
viewed as subsets of $V$.
In general for any graph $G=(V,E)$ and any spanning tree $T$ of $G$, the
Boruvka graph is a rooted tree of depth at most
$\lceil\log_{2}({\left|{V}\right|})\rceil$. This Boruvka graph consists of
successively larger aggregations of elements of $V$. All nodes in a Boruvka
tree are subsets of $V$. The leaves are all the singleton sets $\\{v_{i}\\}$
as $v_{i}$ runs over all the elements of $V$, and eventually the root is
formed, which will be $V$ itself. Any intermediate node in a Boruvka tree is
the union of its children.
A Boruvka tree is built from the bottom up. At each stage or level, every node
computes its nearest neighbor (i.e. the node that it is closest to), and an
edge is formed for all such nodes. The nodes of the next level up are then the
connected components of the graph of the previous level. The weight of each
branch is just the weight of the edge of $G$ that was just added. The result
is a rooted tree with at most $2\,n$ nodes and $n$ leaves.
The key property of $B_{T}$ is that if $u,v\in V$ are a pair of vertices and
if $B_{T}(u,v)$ is the smallest subtree of $B_{T}$ that has both $u$ and $v$
as leaves, then the weight of the heaviest edge that connects $u$ and $v$ in
the original spanning tree $T$ is equal to the weight of the heaviest edge in
$B_{T}(u,v)$. The Boruvka tree $B_{T}$ is a full branching tree, which means
that it has a specified root, all its leaves are at the same level, and each
internal node has at least two children.
The height of $B_{T}$ is at most $\lceil\log_{2}n\rceil$. Therefore once
$B_{T}$ has been constructed, finding the heaviest edge in $B_{T}(u,v)$ costs
at most $O\bigl{(}\log n\bigr{)}$ operations. In fact, if $B_{T}$ has already
been built, then finding the heaviest edge in $B_{T}(u,v)$ requires no queries
of the edge weight oracle. Therefore, to check any edge in the original graph
requires only one oracle query, and total work at most $O\bigl{(}\log
n\bigr{)}$.
## The Quantum Algorithm
The Boruvka tree $B_{T}$ can be made with work $O\bigl{(}n\bigr{)}$, (not just
$O\bigl{(}n\log n\bigr{)}$ work), and can be done classically (see [Ki] and
[Ko]), the total number of oracle queries of the weight function being
$O\bigl{(}n\bigr{)}$, as well. This is what makes this algorithm so effective.
Once the Boruvka tree of the input spanning tree has been formed, it is
possible to check whether the input spanning tree is minimal.
To check any edge $(a,b)\in E$ from the original graph $G$, the maximal weight
of the edge in the path in $T$ that connects $a$ and $b$ is easily found.
Simply start with the leaves $a$ and $b$ and go up $B_{T}$ one level at a time
until they meet at a common internal node (which might be the root). Recording
the maximal weight found in the set of edges traversed in $B_{T}$ up the their
common internal node gives the maximal weight of the edge in the path that
connects $a$ and $b$. Since the height of $B_{T}$ is bounded by
$\lceil\log_{2}n\rceil$, the total work for this is $O\bigl{(}\log n\bigr{)}$,
and no oracle calls are required since $B_{T}$ is already built classically.
To check if $(a,b)$ is of lower weight than the maximal weight of the edge in
the path in $T$ that connects $a$ and $b$ only require one invocation of the
weight oracle. Of course, if this weight is less, then a lower weight spanning
tree than $T$ has been found by swapping out the maximum weight edge in the
path that connects $a$ and $b$ in $T$ with the edge $(a,b)$.
Using Grover’s algorithm over all vertex pairs $V\times V$ for the weight
oracle, therefore requires $O\bigl{(}n\bigr{)}$ oracle queries, and
$O\bigl{(}n\log n\bigr{)}$ total work.
If an edge weight oracle is given, then it is possible to run Grover’s
algorithm over the original edge set $E$, which only requires
$O\bigl{(}\sqrt{m}\bigr{)}$ oracle queries, and $O\bigl{(}\sqrt{m}\log
n\bigr{)}$ total work. Since $m<n^{2}$, it follows that the number of oracle
queries for the quantum part of the algorithm is less than the number of
classical oracle queries needed to construct $B_{T}$. Therefore the total
number of oracle queries is $O\bigl{(}n\bigr{)}$ and the total work is
$O\bigl{(}n+\sqrt{m}\log n\bigr{)}$.
In conclusion, it is interesting that the verification of a putatively correct
answer can be accomplished with considerably less work than that of finding
the answer.
## References
* B Otakar Borűvka, O jistém problému minimálním (About a certain minimal problem), Práca Moravské Pr̆írodoc̆edecké Spolec̆nosti 3, 37-58 . (in Czech, with German summary)
* DHHM Christoph Dürr, Mark Heiligman, Peter Høyer, and Mehdi Mhalla, Quantum query complexity of some graph problems, SIAM Journal on Computing 35(6), 1310–1328 , see also http://xxx.lanl.gov/abs/quant-ph/0401091.
* Ki Valerie King, A simpler minimum spanning tree verification algorithm., Algorithmica 18, 263-270 .
* KPRS V. King, C.K. Poon, V. Ramachandran, and S. Sinha, An optimal EREW PRAM algorithm for minimum spanning tree verification., Information Processing Letters 62, 153-159 .
* Ko Komlos, Linear verifcation for spanning trees, Combinatorica 5, 57-65 .
|
arxiv-papers
| 2011-12-05T15:26:51 |
2024-09-04T02:49:25.006298
|
{
"license": "Public Domain",
"authors": "Mark Heiligman",
"submitter": "Mark Heiligman",
"url": "https://arxiv.org/abs/1112.1139"
}
|
1112.1159
|
# Evolution of the single-mode squeezed vacuum state in amplitude dissipative
channel
Hong-Yi Fan1, Shuai Wang1 and Li-Yun Hu2∗ E-mail: hlyun2008@126.com.
1Department of Physics, Shanghai Jiao Tong University, Shanghai 200240,China
2College of Physics & Communication Electronics, Jiangxi Normal University,
Nanchang 330022, China
$\ast$Corresponding author.E-mail: hlyun2008@126.com.
###### Abstract
Using the way of deriving infinitive sum representation of density operator as
a solution to the master equation describing the amplitude dissipative channel
by virtue of the entangled state representation, we show manifestly how the
initial density operator of a single-mode squeezed vacuum state evolves into a
definite mixed state which turns out to be a squeezed chaotic state with
decreasing-squeezing. We investigate average photon number, photon statistics
distributions for this state.
## I Introduction
Squeezed states are such for which the noise in one of the chosen pair of
observables is reduced below the vacuum or ground-state noise level, at the
expense of increased noise in the other observable. The squeezing effect
indeed improves interferometric and spectroscopic measurements, so in the
context of interferometric detection and of gravitational waves the squeezed
state is very useful 01 ; 02 . In a very recently published paper, Agarwal r1
revealed that a vortex state of a two-mode system can be generated from a
squeezed vacuum by subtracting a photon, such a subtracting mechanism may
happen in a quantum channel with amplitude damping. Usually, in nature every
system is not isolated, dissipation or dephasing usually happens when a system
is immersed in a thermal environment, or a signal (a quantum state) passes
through a quantum channel which is described by a master equation 03 . For
example, when a pure state propagates in a medium, it inevitably interacts
with it and evolves into a mixed state 04 . Dissipation or dephasing will
deteriorate the degree of nonclassicality of photon fields, so physicists pay
much attention to it 05 ; 06 ; 07 .
In this present work we investigate how an initial single-mode squeezed vacuum
state evolves in an amplitude dissipative channel (ADC). When a system is
described by its interaction with a channel with a large number of degrees of
freedom, master equations are set up for a better understanding how quantum
decoherence is processed to affect unitary character in the dissipation or
gain of the system. In most cases people are interested in the evolution of
the variables associated with the system only. This requires us to obtain the
equations of motion for the system of interest only after tracing over the
reservoir variables. A quantitative measure of nonclassicality of quantum
fields is necessary for further investigating the system’s dynamical behavior.
For this channel, the associated loss mechanism in physical processes is
governed by the following master equation 03
$\frac{d\rho\left(t\right)}{dt}=\kappa\left(2a\rho
a^{\dagger}-a^{\dagger}a\rho-\rho a^{\dagger}a\right),$ (1)
where $\rho$ is the density operator of the system, and $\kappa$ is the rate
of decay. We have solved this problem with use of the thermo entangled state
representation 08 . Our questions are: What kind of mixed state does the
initial squeezed state turns into? How does the photon statistics
distributions varies in the ADC?
Thus solving master equations is one of the fundamental tasks in quantum
optics. Usually people use various quasi-probability representations, such as
P-representation, Q-representation, complex P-representation, and Wigner
functions, etc. for converting the master equations of density operators into
their corresponding c-number equations. Recently, a new approach 08 ; 09 ,
using the thermal entangled state representation 10 ; 11 to convert operator
master equations to their c-number equations is presented which can directly
lead to the corresponding Kraus operators (the infinitive representation of
evolved density operators) in many cases.
The work is arranged as follows. In Sec. 2 by virtue of the entangled state
representation we briefly review our way of deriving the infinitive sum
representation of density operator as a solution of the master equation. In
Sec. 3 we show that a pure squeezed vacuum state (with squeezing parameter
$\lambda)$ will evolves into a mixed state (output state), whose exact form is
derived, which turns out to be a squeezed chaotic state. We investigate
average photon number, photon statistics distributions for this state. The
probability of finding $n$ photons in this mixed state is obtained which turns
out to be a Legendre polynomial function relating to the squeezing parameter
$\lambda$ and the decaying rate $\kappa$. In Sec. 4 we discuss the photon
statistics distributions of the output state. In Sec. 5 and 6 we respectively
discuss the Wigner function and tomogram of the output state.
## II Brief review of deducing the infinitive sum representation of
$\rho\left(t\right)$
For solving the above master equation, in a recent review paper 12 we have
introduced a convenient approach in which the two-mode entangled state 10 ; 11
$|\eta\rangle=\exp(-\frac{1}{2}|\eta|^{2}+\eta
a^{{\dagger}}-\eta^{\ast}\tilde{a}^{{\dagger}}+a^{{\dagger}}\tilde{a}^{{\dagger}})|0\tilde{0}\rangle,$
(2)
is employed, where $\tilde{a}^{{\dagger}}$ is a fictitious mode independent of
the real mode $a^{\dagger},$ $[\tilde{a},a^{\dagger}]=0$. $|\eta=0\rangle$
possesses the properties
$\displaystyle a|\eta$
$\displaystyle=0\rangle=\tilde{a}^{{\dagger}}|\eta=0\rangle,$ $\displaystyle
a^{{\dagger}}|\eta$ $\displaystyle=0\rangle=\tilde{a}|\eta=0\rangle,$ (3)
$\displaystyle(a^{{\dagger}}a)^{n}|\eta$
$\displaystyle=0\rangle=(\tilde{a}^{{\dagger}}\tilde{a})^{n}|\eta=0\rangle.$
Acting the both sides of Eq.(1) on the state
$|\eta=0\rangle\equiv\left|I\right\rangle$, and denoting
$\left|\rho\right\rangle=\rho\left|I\right\rangle$, we have
$\displaystyle\frac{d}{dt}\left|\rho\right\rangle$
$\displaystyle=\kappa\left(2a\rho a^{\dagger}-a^{\dagger}a\rho-\rho
a^{\dagger}a\right)\left|I\right\rangle$
$\displaystyle=\kappa\left(2a\tilde{a}-a^{\dagger}a-\tilde{a}^{\dagger}\tilde{a}\right)\left|\rho\right\rangle,$
(4)
so its formal solution is
$\left|\rho\right\rangle=\exp\left[\kappa
t\left(2a\tilde{a}-a^{\dagger}a-\tilde{a}^{\dagger}\tilde{a}\right)\right]\left|\rho_{0}\right\rangle,$
(5)
where $\left|\rho_{0}\right\rangle\equiv\rho_{0}\left|I\right\rangle,$
$\rho_{0}$ is the initial density operator.
Noticing that the operators in Eq.(5) obey the following commutative relation,
$\left[a\tilde{a},a^{\dagger}a\right]=\left[a\tilde{a},\tilde{a}^{\dagger}\tilde{a}\right]=\tilde{a}a$
(6)
and
$\left[\frac{a^{\dagger}a+\tilde{a}^{\dagger}\tilde{a}}{2},a\tilde{a}\right]=-\tilde{a}a,$
(7)
as well as using the operator identity 13
$e^{\lambda\left(A+\sigma B\right)}=e^{\lambda
A}e^{\sigma\left(1-e^{-\lambda\tau}\right)B/\tau},$ (8)
(which is valid for $\left[A,B\right]=\tau B$), we have
$e^{-2\kappa
t\left(\frac{a^{\dagger}a+\tilde{a}^{\dagger}\tilde{a}}{2}-a\tilde{a}\right)}=e^{-\kappa
t\left(a^{\dagger}a+\tilde{a}^{\dagger}\tilde{a}\right)}e^{T^{\prime}a\tilde{a}},$
(9)
where $T^{\prime}=1-e^{-2\kappa t}.$
Then substituting Eq.(9) into Eq.(5) yields 12
$\displaystyle\left|\rho\right\rangle$ $\displaystyle=e^{-\kappa
t\left(a^{\dagger}a+\tilde{a}^{\dagger}\tilde{a}\right)}\sum_{n=0}^{\infty}\frac{T^{\prime
n}}{n!}a^{n}\tilde{a}^{n}\left|\rho_{0}\right\rangle$
$\displaystyle=e^{-\kappa ta^{\dagger}a}\sum_{n=0}^{\infty}\frac{T^{\prime
n}}{n!}a^{n}\rho_{0}a^{{\dagger}n}e^{-\kappa
t\tilde{a}^{\dagger}\tilde{a}}\left|I\right\rangle$
$\displaystyle=\sum_{n=0}^{\infty}\frac{T^{\prime n}}{n!}e^{-\kappa
ta^{\dagger}a}a^{n}\rho_{0}a^{{\dagger}n}e^{-\kappa
ta^{\dagger}a}\left|I\right\rangle,$ (10)
which leads to the infinitive operator-sum representation of$\ \rho$,
$\rho=\sum_{n=0}^{\infty}M_{n}\rho_{0}M_{n}^{\dagger},$ (11)
where
$M_{n}\equiv\sqrt{\frac{T^{\prime n}}{n!}}e^{-\kappa ta^{\dagger}a}a^{n}.$
(12)
We can prove
$\displaystyle\sum_{n}M_{n}^{\dagger}M_{n}$
$\displaystyle=\sum_{n}\frac{T^{\prime n}}{n!}a^{{\dagger}n}e^{-2\kappa
ta^{\dagger}a}a^{n}$ $\displaystyle=\sum_{n}\frac{T^{\prime n}}{n!}e^{2n\kappa
t}\colon a^{{\dagger}n}a^{n}\colon e^{-2\kappa ta^{\dagger}a}$
$\displaystyle=\left.:e^{T^{\prime}e^{2\kappa
t}a^{\dagger}a}:\right.e^{-2\kappa ta^{\dagger}a}$
$\displaystyle=\left.:e^{\left(e^{2\kappa
t}-1\right)a^{\dagger}a}:\right.e^{-2\kappa ta^{\dagger}a}=1,$ (13)
where $\colon\colon$ stands for the normal ordering. Thus $M_{n}$ is a kind of
Kraus operator, and $\rho$ in Eq.(11) is qualified to be a density operator,
i.e.,
$Tr\left[\rho\left(t\right)\right]=Tr\left[\sum_{n=0}^{\infty}M_{n}\rho_{0}M_{n}^{\dagger}\right]=Tr\rho_{0}.$
(14)
Therefore, for any given initial state $\rho_{0}$, the density operator
$\rho\left(t\right)$ can be directly calculated from Eq.(11). The entangled
state representation provides us with an elegant way of deriving the
infinitive sum representation of density operator as a solution of the master
equation.
## III Evolving of an initial single-mode squeezed vacuum state in ADC
It is seen from Eq.(11) that for any given initial state $\rho_{0}$, the
density operator $\rho\left(t\right)$ can be directly calculated. When
$\rho_{0}$ is a single-mode squeezed vacuum state,
$\rho_{0}=\text{sech}\lambda\exp\left(\frac{\tanh\lambda}{2}a^{{\dagger}2}\right)\left|0\right\rangle\left\langle
0\right|\exp\left(\frac{\tanh\lambda}{2}a^{2}\right),$ (15)
we see
$\displaystyle\rho\left(t\right)$ $\displaystyle=$
$\displaystyle\text{sech}\lambda\sum_{n=0}^{\infty}\frac{T^{\prime
n}}{n!}e^{-\kappa
ta^{\dagger}a}a^{n}\exp\left(\frac{\tanh\lambda}{2}a^{{\dagger}2}\right)\left|0\right\rangle$
(16) $\displaystyle\times\left\langle
0\right|\exp\left(\frac{\tanh\lambda}{2}a^{2}\right)a^{{\dagger}n}e^{-\kappa
ta^{\dagger}a}.$
Using the Baker-Hausdorff lemma 14 ,
$e^{\lambda\hat{A}}\hat{B}e^{-\lambda\hat{A}}=\hat{B}+\lambda\left[\hat{A},\hat{B}\right]+\frac{\lambda^{2}}{2!}\left[\hat{A},\left[\hat{A},\hat{B}\right]\right]+\cdots.$
(17)
we have
$\displaystyle
a^{n}\exp\left(\frac{\tanh\lambda}{2}a^{{\dagger}2}\right)\left|0\right\rangle$
$\displaystyle=$ $\displaystyle
e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}e^{-\frac{\tanh\lambda}{2}a^{{\dagger}2}}a^{n}e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}\left|0\right\rangle$
(18) $\displaystyle=$ $\displaystyle
e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}\left(a+a^{\dagger}\tanh\lambda\right)^{n}\left|0\right\rangle.$
Further employing the operator identity 15
$\left(\mu a+\nu
a^{\dagger}\right)^{m}=\left(-i\sqrt{\frac{\mu\nu}{2}}\right)^{m}\colon
H_{m}\left(i\sqrt{\frac{\mu}{2\nu}}a+i\sqrt{\frac{\nu}{2\mu}}a^{\dagger}\right)\colon,$
(19)
where $H_{m}(x)$ is the Hermite polynomial, we know
$\displaystyle\left(a+a^{\dagger}\tanh\lambda\right)^{n}$ (20)
$\displaystyle=$
$\displaystyle\left(-i\sqrt{\frac{\tanh\lambda}{2}}\right)^{n}\colon
H_{n}\left(i\sqrt{\frac{1}{2\tanh\lambda}}a+i\sqrt{\frac{\tanh\lambda}{2}}a^{\dagger}\right)\colon.$
From Eq.(18), it follows that
$\displaystyle
a^{n}e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}\left|0\right\rangle$
$\displaystyle=$
$\displaystyle\left(-i\sqrt{\frac{\tanh\lambda}{2}}\right)^{n}e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}$
(21) $\displaystyle\times
H_{n}\left(i\sqrt{\frac{\tanh\lambda}{2}}a^{\dagger}\right)\left|0\right\rangle.$
On the other hand, noting $e^{-\kappa ta^{\dagger}a}a^{\dagger}e^{\kappa
ta^{\dagger}a}=a^{\dagger}e^{-\kappa t},e^{\kappa ta^{\dagger}a}ae^{-\kappa
ta^{\dagger}a}=ae^{-\kappa t}$ and the normally ordered form of the vacuum
projector $\left|0\right\rangle\left\langle 0\right|=\colon
e^{-a^{\dagger}a}\colon,$ we have
$\displaystyle\rho\left(t\right)$
$\displaystyle=\text{sech}\lambda\sum_{n=0}^{\infty}\frac{T^{\prime
n}}{n!}e^{-\kappa
ta^{\dagger}a}a^{n}e^{\frac{\tanh\lambda}{2}a^{{\dagger}2}}\left|0\right\rangle$
$\displaystyle\times\left\langle
0\right|e^{\frac{\tanh\lambda}{2}a^{2}}a^{{\dagger}n}e^{-\kappa
ta^{\dagger}a}$
$\displaystyle=\text{sech}\lambda\sum_{n=0}^{\infty}\frac{\left(T^{\prime}\tanh\lambda\right)^{n}}{2^{n}n!}e^{\frac{e^{-2\kappa
t}a^{{\dagger}2}\tanh\lambda}{2}}$ $\displaystyle\times
H_{n}\left(i\sqrt{\frac{\tanh\lambda}{2}}a^{\dagger}e^{-\kappa
t}\right)\left|0\right\rangle\left\langle 0\right|$ $\displaystyle\times
H_{n}\left(-i\sqrt{\frac{\tanh\lambda}{2}}ae^{-\kappa
t}\right)e^{\frac{e^{-2\kappa t}a^{2}\tanh\lambda}{2}}$
$\displaystyle=\text{sech}\lambda\sum_{n=0}^{\infty}\frac{\left(T^{\prime}\tanh\lambda\right)^{n}}{2^{n}n!}\colon
e^{\frac{e^{-2\kappa
t}\left(a^{2}+a^{{\dagger}2}\right)\tanh\lambda}{2}-a^{\dagger}a}$
$\displaystyle\times
H_{n}\left(i\sqrt{\frac{\tanh\lambda}{2}}a^{\dagger}e^{-\kappa
t}\right)H_{n}\left(-i\sqrt{\frac{\tanh\lambda}{2}}ae^{-\kappa
t}\right)\colon$ (22)
then using the following identity 16
$\displaystyle\sum_{n=0}^{\infty}\frac{t^{n}}{2^{n}n!}H_{n}\left(x\right)H_{n}\left(y\right)$
(23) $\displaystyle=$
$\displaystyle\left(1-t^{2}\right)^{-1/2}\exp\left[\frac{t^{2}\left(x^{2}+y^{2}\right)-2txy}{t^{2}-1}\right],$
and $e^{\lambda a^{{\dagger}}a}=\colon
e^{\left(e^{\lambda}-1\right)a^{{\dagger}}a}\colon,$ we finally obtain the
expression of the output state
$\rho\left(t\right)=We^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}a^{2}},$
(24)
with $T^{\prime}=1-e^{-2\kappa t}$ and
$W\equiv\frac{\text{sech}\lambda}{\sqrt{1-T^{\prime
2}\tanh^{2}\lambda}},\text{\ss}\equiv\frac{e^{-2\kappa
t}\tanh\lambda}{1-T^{\prime 2}\tanh^{2}\lambda}.$ (25)
By comparing Eq.(15) with (23) one can see that after going through the
channel the initial squeezing parameter $\tanh\lambda$ in Eq.( 15) becomes to
ß$\equiv\frac{e^{-2\kappa t}\tanh\lambda}{1-T^{\prime 2}\tanh^{2}\lambda},$
and $\left|0\right\rangle\left\langle
0\right|\rightarrow\frac{1}{\sqrt{1-T^{\prime
2}\tanh^{2}\lambda}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)},$
a chaotic state (mixed state), due to $T^{\prime}>0,$ we can prove
$\frac{e^{-2\kappa t}}{1-T^{\prime 2}\tanh^{2}\lambda}<1,$ which means a
squeezing-decreasing process. When $\kappa t=0$, then $T^{\prime}=0$ and ß
$=\tanh\lambda$, Eq.(22) becomes the initial squeezed vacuum state as
expected.
It is important to check: if Tr$\rho(t)=1$. Using Eq.(22) and the completeness
of coherent state $\int\frac{d^{2}z}{\pi}\left|z\right\rangle\left\langle
z\right|=1$ as well as the following formula 17
$\int\frac{d^{2}z}{\pi}e^{\zeta\left|z\right|^{2}+\xi z+\eta
z^{\ast}+fz^{2}+gz^{\ast
2}}=\frac{1}{\sqrt{\zeta^{2}-4fg}}e^{\frac{-\zeta\xi\eta+f\eta^{2}+g\xi^{2}}{\zeta^{2}-4fg}},$
(26)
whose convergent condition is Re$\left(\zeta\pm f\pm g\right)<0$ and$\
\mathtt{Re}\left(\frac{\zeta^{2}-4fg}{\zeta\pm f\pm g}\right)<0$, we really
see
$\displaystyle\text{Tr}\rho\left(t\right)$ $\displaystyle=$ $\displaystyle
W\int\frac{d^{2}z}{\pi}\left\langle
z\right|e^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}a^{2}}\left|z\right\rangle$
(27) $\displaystyle=$
$\displaystyle\frac{W}{\sqrt{\left(\text{\ss}T^{\prime}\tanh\lambda-1\right)^{2}-\text{\ss}^{2}}}=1.$
so $\rho\left(t\right)$ is qualified to be a mixed state, thus we see an
initial pure squeezed vacuum state evolves into a squeezed chaotic state with
decreasing-squeezing after passing through an amplitude dissipative channel.
## IV Average photon number
Using the completeness relation of coherent state and the normally ordering
form of $\rho\left(t\right)$ in Eq. (22), and using
$e^{\frac{\text{\ss}}{2}a^{2}}a^{\dagger}e^{-\frac{\text{\ss}}{2}a^{2}}=a^{\dagger}+$ß$a$,
as well as
$e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}a^{\dagger}e^{-a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}$=$a^{\dagger}$ß$T^{\prime}\tanh\lambda,$
we have
$\displaystyle\mathtt{Tr}\left(\rho\left(t\right)a^{\dagger}a\right)$
$\displaystyle=W\int\frac{d^{2}z}{\pi}\left\langle
z\right|e^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}a^{2}}a^{\dagger}a\left|z\right\rangle$
$\displaystyle=W\int\frac{d^{2}z}{\pi}\left\langle
z\right|e^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}a^{2}}za^{\dagger}\left|z\right\rangle$
$\displaystyle=W\int\frac{d^{2}z}{\pi}z\left\langle
z\right|e^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}\left(a^{\dagger}+\text{\ss}a\right)e^{\frac{\text{\ss}}{2}a^{2}}\left|z\right\rangle$
$\displaystyle=W\text{\ss}\int\frac{d^{2}z}{\pi}ze^{\frac{\text{\ss}}{2}\left(z^{\ast
2}+z^{2}\right)}\left\langle
z\right|\left(a^{\dagger}T^{\prime}\tanh\lambda+z\right)e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}\left|z\right\rangle$
$\displaystyle=W\text{\ss}\int\frac{d^{2}z}{\pi}\left(|z|^{2}T^{\prime}\tanh\lambda+z^{2}\right)$
$\displaystyle\times\exp\left[\left(\text{\ss}T^{\prime}\tanh\lambda-1\right)|z|^{2}+\frac{\text{\ss}}{2}\left(z^{\ast
2}+z^{2}\right)\right].$ (28)
In order to perform the integration, we reform Eq.(28) as
$\displaystyle\mathtt{Tr}\left(\rho\left(t\right)a^{\dagger}a\right)$
$\displaystyle=$ $\displaystyle
W\text{\ss}\left\\{T^{\prime}\tanh\lambda\frac{\partial}{\partial
f}+\frac{2}{\text{ \ss}}\frac{\partial}{\partial s}\right\\}$ (29)
$\displaystyle\times\int\frac{d^{2}z}{\pi}\exp\left[\left(\text{\ss}T^{\prime}\tanh\lambda-1+f\right)|z|^{2}\right.$
$\displaystyle+\left.\frac{\text{\ss}}{2}\left(z^{\ast
2}+\left(1+s\right)z^{2}\right)\right]_{f=s=0}$ $\displaystyle=$
$\displaystyle\frac{1-\text{\ss}T^{\prime}\tanh\lambda}{\left(\text{\ss}T^{\prime}\tanh\lambda-1\right)^{2}-\text{\ss}^{2}}-1$
in the last step, we have used Eq.(27). Using Eq.(29), we present the time
evolution of the average photon number in Fig. 1, from which we find that the
average photon number of the single-mode squeezed vacuum state in the
amplitude damping channel reduces gradually to zero when decay time goes.
Figure 1: (Color online) The average $\bar{n}\left(\kappa t\right)$ as the
function of $\kappa t$ for different values of squeezing parameter $\lambda$
(from bottom to top $\lambda=0,0.1,0.3,0.5,1$.)
## V Photon statistics distribution
Next, we shall derive the photon statistics distributions of
$\rho\left(t\right)$. The photon number is given by
$p\left(n,t\right)=\left\langle
n\right|\rho\left(t\right)\left|n\right\rangle$. Noticing
$a^{{\dagger}m}\left|n\right\rangle=\sqrt{(m+n)!/n!}\left|m+n\right\rangle$
and using the un-normalized coherent state
$\left|\alpha\right\rangle=\exp[\alpha a^{{\dagger}}]\left|0\right\rangle$, 18
; 19 leading to
$\left|n\right\rangle=\frac{1}{\sqrt{n!}}\frac{\mathtt{d}^{n}}{\mathtt{d}\alpha^{n}}\left|\alpha\right\rangle\left|{}_{\alpha=0}\right.,$
$\left(\left\langle\beta\right.\left|\alpha\right\rangle=e^{\alpha\beta^{\ast}}\right)$,
as well as the normal ordering form of $\rho\left(t\right)$ in Eq. (22), the
probability of finding $n$ photons in the field is given by
$\displaystyle p\left(n,t\right)$ (30) $\displaystyle=$
$\displaystyle\left\langle n\right|\rho\left(t\right)\left|n\right\rangle$
$\displaystyle=$
$\displaystyle\frac{W}{n!}\frac{\mathtt{d}^{n}}{\mathtt{d}\beta^{\ast
n}}\frac{\mathtt{\
d}^{n}}{\mathtt{d}\alpha^{n}}\left.\left\langle\beta\right|e^{\frac{\text{\ss}}{2}\beta^{\ast
2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}\alpha^{2}}\left|\alpha\right\rangle\right|_{\alpha,\beta^{\ast}=0}$
$\displaystyle=$
$\displaystyle\frac{W}{n!}\frac{\mathtt{d}^{n}}{\mathtt{d}\beta^{\ast
n}}\frac{\mathtt{\
d}^{n}}{\mathtt{d}\alpha^{n}}\left.\exp\left[\beta^{\ast}\alpha\text{
\ss}T^{\prime}\tanh\lambda+\frac{\text{\ss}}{2}\beta^{\ast
2}+\frac{\text{\ss}}{2}\alpha^{2}\right]\right|_{\alpha,\beta^{\ast}=0}.$
Note that
$\left[e^{\frac{\text{\ss}}{2}a^{\dagger
2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}\alpha^{2}}\right]^{\dagger}=e^{\frac{\text{\ss}}{2}a^{\dagger
2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}\alpha^{2}}$
so
$\left\langle
n\right|\rho\left(t\right)\left|n\right\rangle^{\ast}=\left\langle
n\right|\rho\left(t\right)^{\dagger}\left|n\right\rangle=\left\langle
n\right|\rho\left(t\right)\left|n\right\rangle$
$\displaystyle\frac{\partial^{n+n}}{\partial t^{n}\partial t^{\prime
n}}\exp\left[2xtt^{\prime}-t^{2}-t^{\prime 2}\right]_{t=t^{\prime}=0}$ (31)
$\displaystyle=$ $\displaystyle
2^{n}n!\sum_{m=0}^{[n/2]}\frac{n!}{2^{2m}\left(m!\right)^{2}(n-2m)!}x^{n-2m},$
we derive the compact form for $\mathfrak{p}\left(n,t\right)$, i.e.,
$\displaystyle p\left(n,t\right)$ (32) $\displaystyle=$
$\displaystyle\frac{W}{n!}\left(-\frac{\text{\ss}}{2}\right)^{n}\frac{\mathtt{d}^{n}}{\mathtt{d}\beta^{\ast
n}}\frac{\mathtt{d}^{n}}{\mathtt{d}\alpha^{n}}\left.e^{-2T^{\prime}\tanh\lambda\beta^{\ast}\alpha-\beta^{\ast
2}-\alpha^{2}}\right|_{\alpha,\beta^{\ast}=0}$ $\displaystyle=$ $\displaystyle
W\left(\text{\ss}T^{\prime}\tanh\lambda\right)^{n}\sum_{m=0}^{[n/2]}\frac{n!\left(T^{\prime}\tanh\lambda\right)^{-2m}}{2^{2m}\left(m!\right)^{2}(n-2m)!}.$
Using the newly expression of Legendre polynomials found in Ref. 20
$x^{n}\sum_{m=0}^{[n/2]}\frac{n!}{2^{2m}\left(m!\right)^{2}(n-2m)!}\left(1-\frac{1}{x^{2}}\right)^{m}=P_{n}\left(x\right),$
(33)
we can formally recast Eq.(32) into the following compact form, i.e.,
$p\left(n,t\right)=W\left(e^{-\kappa
t}\sqrt{-\text{\ss}\tanh\lambda}\right)^{n}P_{n}\left(e^{\kappa
t}T^{\prime}\sqrt{-\text{\ss}\tanh\lambda}\right)$
note that since $\sqrt{-\text{\ss}\tanh\lambda}$ is pure imaginary, while
$p\left(n,t\right)$ is real, so we must still use the power-series expansion
on the right-hand side of Eq.(32) to depict figures of the variation of
$p\left(n,t\right)$. In particular, when $t=0$, Eq.(32 ) reduces to
$\displaystyle p\left(n,0\right)$ $\displaystyle=$
$\displaystyle\text{sech}\lambda\left(\tanh\lambda\right)^{n}\lim_{T^{\prime}\rightarrow
0}\sum_{m=0}^{[n/2]}\frac{n!\left(T^{\prime}\tanh\lambda\right)^{n-2m}}{2^{2m}\left(m!\right)^{2}(n-2m)!}$
(36) $\displaystyle=$
$\displaystyle\left\\{\begin{array}[]{cc}\frac{\left(2k\right)!}{2^{2k}k!k!}\text{sech}\lambda\tanh^{2k}\lambda,&n=2k\\\
0&n=2k+1\end{array}\right.,$
which just correspond to the number distributions of the squeezed vacuum state
21 ; 22 . From Eq.(36) it is not difficult to see that the photocount
distribution decreases as the squeezing parameter $\lambda$ increases. While
for $\kappa t\rightarrow\infty,$ we see that $p\left(n,\infty\right)=0.$ This
indicates that there is no photon when a system interacting with a amplitude
dissipative channel for enough long time, as expected. In Fig. 2, the photon
number distribution is shown for different $\kappa t$.
Figure 2: (Color online) Photon number distribution of the squeezed vacuum
state in amplitude damping channel for $\lambda=1$, and different $\kappa t$:
($a)$ $\kappa t=0$, ($b$) $\kappa t=0.5,$ ($c$) $\kappa t=1$ and ($d$) $\kappa
t=2$.
## VI Wigner functions
In this section, we shall use the normally ordering for of density operators
to calculate the analytical expression of Wigner function. For a single-mode
system, the WF is given by 23
$W\left(\alpha,\alpha^{\ast},t\right)=e^{2\left|\alpha\right|^{2}}\int\frac{d^{2}\beta}{\pi^{2}}\left\langle-\beta\right|\rho\left(t\right)\left|\beta\right\rangle
e^{-2\left(\beta\alpha^{\ast}-\beta^{\ast}\alpha\right)},$ (37)
where $\left|\beta\right\rangle$ is the coherent state 18 ; 19 . From Eq.(22)
it is easy to see that once the normal ordered form of $\rho\left(t\right)$ is
known, we can conveniently obtain the Wigner function of $\rho\left(t\right)$.
On substituting Eq.(24) into Eq.(37) we obtain the WF of the single-mode
squeezed state in the ADC,
$\displaystyle W\left(\alpha,\alpha^{\ast},t\right)$ (38) $\displaystyle=$
$\displaystyle
We^{2\left|\alpha\right|^{2}}\int\frac{d^{2}\beta}{\pi^{2}}\exp\left[-\left(1+\text{\ss}T^{\prime}\tanh\lambda\right)\left|\beta\right|^{2}\right.$
$\displaystyle\left.-2\left(\beta\alpha^{\ast}-\beta^{\ast}\alpha\right)+\frac{\text{\ss}}{2}\beta^{\ast
2}+\frac{\text{\ss}}{2}\beta^{2}\right]$ $\displaystyle=$
$\displaystyle\frac{W}{\pi\sqrt{\left(1+\text{\ss}T^{\prime}\tanh\lambda\right)^{2}-\text{\ss}^{2}}}\exp\left[2\left|\alpha\right|^{2}\right]$
$\displaystyle\times\exp\left[2\frac{-2\left(1+\text{\ss}T^{\prime}\tanh\lambda\right)\left|\alpha\right|^{2}+\text{\ss}\left(\alpha^{\ast
2}+\alpha^{2}\right)}{\left(1+\text{\ss}T^{\prime}\tanh\lambda\right)^{2}-\text{\ss}^{2}}\right]$
In particular, when $t=0$ and $t\rightarrow\infty$, Eq.(38) reduces to
$W\left(\alpha,\alpha^{\ast},0\right)=\frac{1}{\pi}\exp[-2\left|\alpha\right|^{2}\cosh
2\lambda+\left(\alpha^{\ast 2}+\alpha^{2}\right)\sinh 2\lambda]$, and
$W\left(\alpha,\alpha^{\ast},\infty\right)=\frac{1}{\pi}\exp\left[-2\left|\alpha\right|^{2}\right]$,
which are just the WF of the single-mode squeezed vacuum state and the vacuum
state, respectively. In Fig. 3, the WF of the single-mode squeezed vacuum
state in the amplitude damping channel is shown for different decay time
$\kappa t$.
Figure 3: (Color online) Wigner function of the squeezed vacuum state in
amplitude damping channel for $\lambda=1.0$, different $\kappa t$: ($a)$
$\kappa t=0.0$, ($b$) $\kappa t=0.5$, ($c$) $\kappa t=1$, and ($d$) $\kappa
t=2$.
## VII Tomogram
As we know, once the probability distributions
$P_{\theta}\left(\hat{x}_{\theta}\right)$ of the quadrature amplitude are
obtained, one can use the inverse Radon transformation familiar in tomographic
imaging to obtain the WF and density matrix 24 . Thus the Radon transform of
the WF is corresponding to the probability distributions
$P_{\theta}\left(\hat{x}_{\theta}\right)$. In this section we derive the
tomogram of $\rho\left(t\right)$.
For a single-mode system, the Radon transform of WF, denoted as $\mathcal{R}$
is defined by 25
$\displaystyle\mathcal{R}\left(q\right)_{f,g}$ $\displaystyle=$
$\displaystyle\int\delta\left(q-fq^{\prime}-gp^{\prime}\right)Tr\left[\Delta\left(\beta\right)\rho\left(t\right)\right]dq^{\prime}dp^{\prime}$
(39) $\displaystyle=$ $\displaystyle Tr\left[\left|q\right\rangle_{f,g\text{
}f,g}\left\langle q\right|\rho\left(t\right)\right]=_{f,g}\left\langle
q\right|\rho\left(t\right)\left|q\right\rangle_{f,g}$
where the operator $\left|q\right\rangle_{f,g\text{ }f,g}\left\langle
q\right|$ is just the Radon transform of single-mode Wigner operator
$\Delta\left(\beta\right)$, and
$\left|q\right\rangle_{f,g}=A\exp\left[\frac{\sqrt{2}qa^{{\dagger}}}{B}-\frac{B^{\ast}}{2B}a^{{\dagger}2}\right]\left|0\right\rangle,$
(40)
as well as $B=f-ig,$
$A=\left[\pi\left(f^{2}+g^{2}\right)\right]^{-1/4}\exp[-q^{2}/2\left(f^{2}+g^{2}\right)]$.
Thus the tomogram of a quantum state $\rho\left(t\right)$ is just the quantum
average of $\rho\left(t\right)$ in $\left|q\right\rangle_{f,g}$ representation
(a kind of intermediate coordinate-momentum representation) 26 .
Substituting Eqs.(24) and (40) into Eq.(39), and using the completeness
relation of coherent state, we see that the Radom transform of WF of
$\rho\left(t\right)$ is given by
$\displaystyle\mathcal{R}\left(q\right)_{f,g}$ (41) $\displaystyle=$
$\displaystyle W_{f,g}\left\langle
q\right|e^{\frac{\text{\ss}}{2}a^{{\dagger}2}}e^{a^{\dagger}a\ln\left(\text{\ss}T^{\prime}\tanh\lambda\right)}e^{\frac{\text{\ss}}{2}a^{2}}\left|q\right\rangle_{f,g}$
$\displaystyle=$
$\displaystyle\frac{WA^{2}}{\sqrt{E}}\exp\left\\{\frac{q^{2}\text{\ss}}{E\left|B\right|^{4}}\left(B^{2}+B^{\ast}{}^{2}\right)\right.$
$\displaystyle+\left.\frac{2q^{2}\text{\ss}}{E\left|B\right|^{2}}\left(T\tanh\lambda+\text{\ss}-\text{\ss}T^{2}\tanh^{2}\lambda\right)\right\\},$
where we have used the formula (26) and
$\left\langle\alpha\right|\left.\gamma\right\rangle=\exp[-\left|\alpha\right|^{2}/2-\left|\gamma\right|^{2}/2+\alpha^{\ast}\gamma]$,
as well as
$\displaystyle E$ $\displaystyle=$
$\displaystyle\left(1+\text{\ss}\frac{B}{B^{\ast}}\right)\left(1+\frac{B^{\ast}}{B}\text{\ss}-B^{\ast}\frac{\left(\text{\ss}T^{\prime}\tanh\lambda\right)^{2}}{B^{\ast}+\text{\ss}B}\right)$
(42) $\displaystyle=$
$\displaystyle\left|1+\frac{\text{\ss}B}{B^{\ast}}\right|^{2}-\left(\text{
\ss}T^{\prime}\tanh\lambda\right)^{2}.$
In particular, when $t=0,$ ($T=0$), then Eq.(41) reduces to
($\frac{B}{B^{\ast}}=e^{2i\phi}$)
$\displaystyle\mathcal{R}\left(q\right)_{f,g}$ $\displaystyle=$
$\displaystyle\frac{A^{2}\text{sech}\lambda}{\left|1+e^{2i\phi}\tanh\lambda\right|}$
(43)
$\displaystyle\times\exp\left\\{\frac{q^{2}\left(B^{2}+B^{\ast}{}^{2}+2\left|B\right|^{2}\tanh\lambda\right)\tanh\lambda}{\left|1+e^{2i\phi}\left|B\right|^{4}\tanh\lambda\right|^{2}}\right\\},$
which is a tomogram of single-mode squeezed vacuum state; while for $\kappa
t\rightarrow\infty,$($T=1$), then $\mathcal{R}\left(q\right)_{f,g}=A^{2},$
which is a Gaussian distribution corresponding to the vacuum state.
In summary, using the way of deriving infinitive sum representation of density
operator by virtue of the entangled state representation describing, we
conclude that in the amplitude dissipative channel the initial density
operator of a single-mode squeezed vacuum state evolves into a squeezed
chaotic state with decreasing-squeezing. We investigate average photon number,
photon statistics distributions, Wigner functions and tomogram for the output
state.
## Acknowledgments
This work is supported by the National Natural Science Foundation of China
(Grant No.11175113 and 11047133), Shandong Provincial Natural Science
Foundation in China (Gant No.ZR2010AQ024), and a grant from the Key Programs
Foundation of Ministry of Education of China (Grant No. 210115),as well as
Jiangxi Provincial Natural Science Foundation in China (No. 2010GQW0027).
## References
* (1) Caves C M 1981 Phys. Rev. D 23 1693
* (2) Loudon R 1981 Phys. Rev. Lett. 47 815
* (3) Agarwal G S 2011 New J. Phys. 13 073008
* (4) Gardner C W and Zoller P 2000 Quantum Noise (Berlin: Spinger)
* (5) Louisell W H 1973 Quantum Statistical Properties of Radiation (New York: Wiley)
* (6) Biswas A and Agarwal G S 2007 Phys. Rev. A 75 032104
* (7) Hu L Y and Fan H Y 2008 J. Opt. Soc. Am. B 25 1955
* (8) Hu L Y, Xu X X, Wang Z S and Xu X F 2010 Phys. Rev. A 82 043842
* (9) Fan H Y and Hu L Y 2008 Opt. Commun. 281 5571
* (10) Hu L Y and Fan H Y, Opt. Commun. 282, 4379 (2009)
* (11) Fan H Y and Fan Y 1998 Phys. Lett. A 246 242
* (12) Fan H Y and Fan Y 2001 Phys. Lett. A 282 269
* (13) Fan H Y and Hu L Y 2008 Mod. Phys. Lett. B 22 2435
* (14) Fan H Y 1997 Representation and Transformation Theory in Quantum Mechanics (Shanghai: Shanghai Scientific and Technical) (in Chinese)
* (15) Klauder J R and Skargerstam B S 1985 Coherent States (Singapore: World Scientific)
* (16) Xu X X, Yuan H C, Hu L Y and Fan H Y 2011 J. Phys. A: Math. Theor. 44 445306
* (17) Rainville E D 1960 Special Functions (New York: MacMillan Company)
* (18) Puri R R 2001 Mathematical Methods of Quantum Optics (Berlin/Heidelberg/New York: Springer-Verlag)
* (19) Glauber R J 1963 Phys. Rev. 130 2529
* (20) Glauber R J 1963 Phys. Rev. 131 2766
* (21) Fan H Y, Hu L Y and Xu X X 2009 Mod. Phys. Lett. A 24 1597
* (22) Kim M S, de Oliveira F A M and Knight P L 1989 Phys. Rev. A 40 2494
* (23) Paulina Marian 1992 Phys. Rev. A 45 2044
* (24) Fan H Y andZaidi H R 1987 Phys. Lett. A 124 303
* (25) Vogel K and Risken H 1989 Phys. Rev. A 40 2847
* (26) Fan H Y and Niu J B 2010 Opt. Commun. 283 3296
* (27) Fan H Y and Hu L Y 2009 Opt. Commun. 282 3734
|
arxiv-papers
| 2011-12-06T05:04:48 |
2024-09-04T02:49:25.012889
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hong-Yi Fan, Shuai Wang and Li-Yun Hu",
"submitter": "Liyun Hu",
"url": "https://arxiv.org/abs/1112.1159"
}
|
1112.1248
|
$\rm{D^{0}}$ cross section in pp collisions at $\sqrt{s}=7$ TeV, measured with
the ALICE experiment
Xianbao Yuan, for the ALICE Collaboration
Institute of Particle Physics, Central China Normal University
University and INFN, Padova, Italy
Contact:yuanxb@iopp.ccnu.edu.cn or xianbao.yuan@pd.infn.it
The measurement of the cross-section for charm production in pp collisions at
the LHC is not only a fundamental reference to investigate medium properties
in heavy-ion collisions, but also a key test of pQCD predictions in a new
energy domain.
The ALICE [1] experiment has measured the D meson production in pp collisions
at $\sqrt{s}=7$ TeV. We present the analysis procedure for $\rm D^{0}\to\rm
K^{-}\rm\pi^{+}$ and for the calculation of efficiency and acceptance
corrections. Finally, we show the preliminary results on $\rm D^{0}$ cross
section in pp collisions at $\sqrt{s}=7$ TeV, measured in the region $2<p_{\rm
t}<12$ GeV/$c$ at central rapidity $|y|<0.5$. These results are compared to
perturbative QCD predictions.
The analysis is based on an invariant mass analysis of opposite-charge pairs
of reconstructed tracks that can represent a $\rm D^{0}$ with a displaced
vertex (the mean proper decay length of the $\rm D^{0}$ is $c\tau\approx
123~{}\mu$m). The selection is based on topological cuts and particle
identification via specific energy deposit and time-of-flight measurement. The
cross section is calculated from the raw signal yield extracted with the
invariant mass analysis, $N^{\rm D^{0}~{}raw}(p_{\rm t})$, using the following
formula:
$\left.\frac{{\rm d}\sigma^{\rm D^{0}}}{{\rm d}p_{\rm
t}}\right|_{|y|<0.5}=\frac{1}{2}\frac{1}{2\,y_{\rm acc}\Delta p_{\rm
t}}\frac{\left.f_{prompt}\cdot N^{\rm D^{0}~{}raw}(p_{\rm
t})\right|_{|y|<y_{\rm acc}}}{\epsilon_{prompt}\cdot{\rm
BR}\cdot\mathit{L}_{int}}\cdot$ (1)
Here, $\epsilon_{prompt}$ means the efficiency of prompt mesons, which
accounts for selection cuts, for track and primary vertex reconstruction
efficiency, and for detector acceptance. The $f_{prompt}$ is the prompt
fraction of raw yield.
Figure 1 (Left) shows the invariant mass distribution for $p_{\rm t}>2$
GeV/$c$ after applying the cuts, which corresponds to $1.1\times 10^{8}$
minimum bias events collected by ALICE in 2010 at $\sqrt{s}=7$ TeV. Figure 1
(Right) shows the efficiencies for $\rm D^{0}\to\rm K^{-}\rm\pi^{+}$ with all
the decay particles in the acceptance $\left|\eta\right|<0.9$. The
efficiencies increase and flatten at about 0.1 at $p_{\rm t}>2$ GeV/$c$. The
efficiency without particle identification selection, shown for comparison, is
the same as with particle identification for $p_{\rm t}>2$ GeV/$c$, indicating
that this selection is essentially fully efficient for the signal. The
efficiencies for $\rm D^{0}$ meson from B meson decay, also shown for
comparison, are larger by a factor about 2, because this feed-down component
is more displaced from the primary vertex, due to the longer B life time. The
$10-15\%$ feed-down from B decays is subtracted based on pQCD prediction [2].
Several sources of systematic uncertainties were considered, namely those
affecting the signal extraction from the invariant mass spectra and all the
correction factors applied to obtain the $p_{\rm t}$-differential cross
sections. A summary of the estimated relative systematic errors is shown in
Fig 2 (Left).
The $p_{\rm t}$-differential cross section for prompt $\rm D^{0}$, obtained
from the yields extracted by fitting the invariant mass spectra and corrected
for efficiency and B feed-down, is shown in Fig 2 (Right). The error bars
represent the statistical errors, while the systematic errors are plotted as
rectangle areas around the data points. The measured $\rm D^{0}$ meson
production cross section is compared to two theoretical predictions, namely
FONLL [2] and GM-VFNS [3]. Our measurement of $\rm D^{0}$ at $\sqrt{s}=7$ TeV
is reproduced by both models within their theoretical uncertainties.
Figure 1: Left: $p_{\rm t}>$2 GeV/$c$ invariant mass distribution. Right:
efficiencies for $\rm D^{0}$ as a function of $p_{\rm t}$.
## References
* [1] B. Abelev, A. Abrahantes Quintana, $\mathit{et~{}al}$., [ALICE Coll.], __Arxiv: hep-ex/1111.1553v1.
* [2] M. Cacciari $\mathit{et~{}al}$., __JHEP 0407 (2004) 033; Private communication.
* [3] B.A. Kniehl $\mathit{et~{}al}$., __Phys. Rev. D77 (2008) 014011; Private communication.
Figure 2: Left: systematic errors summary plot. Right: $p_{\rm
t}$-differential cross section for prompt $\rm D^{0}$ in pp collisions at
$\sqrt{s}=7$ TeV compared with FONLL [2] and GM-VFNS [3] theoretical
predictions.
|
arxiv-papers
| 2011-12-06T12:17:04 |
2024-09-04T02:49:25.021371
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Xianbao Yuan",
"submitter": "Xianbao Yuan",
"url": "https://arxiv.org/abs/1112.1248"
}
|
1112.1387
|
# Systematic enumeration of configuration classes for entropic sampling of
Ising models
BRUNO JEFERSON LOURENÇO 111brunojl@fisica.ufmg.br
and
RONALD DICKMAN 222dickman@fisica.ufmg.br
_Departamento de Física, Instituto de Ciências Exatas, and National Institute
of Science and Technology for Complex Systems,
Universidade Federal de Minas Gerais
C.P. 702, 30123-970, Belo Horizonte, MG, Brazil
_
## Abstract
We describe a systematic method for complete enumeration of configuration
classes (CCs) of the spin-1/2 Ising model in the energy-magnetization plane.
This technique is applied to the antiferromagnetic Ising model in an external
magnetic field on the square lattice, which is simulated using the tomographic
entropic sampling algorithm. We estimate the number of configurations,
$\Omega(E,m,L)$, and related microcanonical averages, for all allowed energies
$E$ and magnetizations $m$ for $L=10$ to 30, with $\Delta L=2$. With prior
knowledge of the CCs, we can be sure that all allowed classes are sampled
during the simulation. Complete enumeration of CCs also enables us to use the
final estimate of $\Omega(E,m,L)$ to obtain good initial estimates,
$\Omega_{0}(E,m,L^{\prime})$, for successive system sizes ($L^{\prime}>L$)
through a two-dimensional interpolation. Using these results we calculate
canonical averages of the thermodynamic quantities of interest as continuous
functions of temperature $T$ and external field $h$. In addition, we determine
the critical line in the $h$-$T$ plane using finite-size scaling analysis, and
compare these results with several approximate theoretical expressions.
Keywords: Ising model; antiferromagnet; configuration classes; Monte Carlo
simulation; phase diagram.
## 1 Introduction
The most important task of entropic sampling algorithms [1]-[6] is to visit
the full configuration space (CS) to obtain good estimates of the number of
configurations, $\Omega$, as functions of the energy E and other quantities of
interest. In studies of the Ising model in an external field, for example, we
require $\Omega(E,m)$ with $m$ the magnetization; each allowed $(E,m)$ pair
defines a class of configurations (CC). Only if we know beforehand the
possible values of $(E,m)$ for a given system size, can we be sure that all
CCs are sampled during the simulation.
Figure 1 shows the CCs in the $n-m$ plane for the spin-1/2 Ising model with
nearest-neighbor (NN) interactions, on a square lattice of $L\times L$ sites
with periodic boundaries. [We use $n$ to denote the number of NN pairs of
spins with the same orientation; the interaction energy of the
antiferromagnetic (AF) Ising model is $E=-2(L^{2}-n)$.] Although the CCs tend
to fill in a triangular region, some gaps are evident near the lower apex and
along the upper edge. Knowing just which $(n,m)$ values are allowed for a
given lattice size is important if we are to implement entropic sampling with
confidence. In this paper we present a method for systematically enumerating
all CCs of Ising models in the $(n,m)$ plane.
Figure 1: Allowed configuration classes for system size $L=10$; $\eta\equiv
n/L^{2}$ and $\nu\equiv m/L^{2}$.
We study the spin-1/2 antiferromagnetic Ising model in an external magnetic
field, whose energy is given by
$\mathcal{H}=-J\sum_{<i,j>}\sigma_{i}\sigma_{j}-h\sum_{i=1}^{N}\sigma_{i}=-E-hm,$
(1)
where $\sigma_{i}=\pm 1$, $<i,j>$ indicates a sum over NN pairs of spins, $h$
is the external field, and $N$ is the number of spins; the model is defined on
a square lattice of $L\times L$ sites, with periodic boundary conditions.
Unlike the ferromagnetic Ising model ($J>0$), which exhibits a unique critical
point in the $h-T$ plane and has an exact solution [7], the AF model ($J<0$)
possesses a critical line, which is not completely understood.
Various approximate methods have been applied to determine the critical line
of the AF Ising model on the square lattice [8]-[18]; these results, however,
do not agree altogether. Binder and Landau [19] estimated the critical line
via Monte Carlo simulation, obtaining very good agreement with the approximate
closed-form expression of Müller-Hartmann and Zittartz [8], raising the
possibility that the latter expression was in fact exact. (It was later shown
that this is not so [12].) Hwang _et al._ [20] studied the AF Ising model on
the square and triangular lattices using a microcanonical transfer matrix
method and the Wang-Landau algorithm [3]. They performed an exact enumeration
of the number of configurations, $\Omega(E,m)$, and found CSs in the $E-m$
plane similar to that shown in Fig. 1.
Using the tomographic entropic sampling (TES) algorithm [6] we estimate
$\Omega(n,m,L)$, and associated microcanonical averages, for lattice sizes
$L=10$ to 30, with an increment $\Delta L=2$. We then calculate the canonical
averages of the thermodynamic quantities of interest. Using these data we map
out the critical line in the $h-T$ plane, and compare our results with several
theoretical expressions.
Prior determination of the set of allowed CCs is an important tool to verify
the quality of the sampling: we want to be sure that all CCs are visited
during the simulation. Since this algorithm uses an initial guess,
$\Omega_{0}(n,m,L)$, to begin the study, it is convenient to use the final
estimate $\Omega_{N}(n,m,L)$ (after the $N$-th iteration) to obtain the
initial guess $\Omega_{0}(n,m,L_{0})$ of the next system size to be studied
$(L_{0}>L)$. As we will show, good initial estimates, $\Omega_{0}(n,m,L)$, can
be obtained using a two-dimensional interpolation because we know a priori the
set of allowed CCs.
This paper is organized as follows. In Sec. 2 we define the basic CCs and the
respective allowed values of $(n,m)$ for the spin-1/2 Ising model on the
square lattice; the main goal of this section is to find all gaps in the
$(n,m)$ plane. In Sec. 3, this information is used in simulations of the AF
Ising model via the TES algorithm. There we describe the method used to
determine $\Omega_{0}(n,m,L_{2})$ via two-dimensional interpolation of the
final estimate, $\Omega_{N}(n,m,L_{1})$, of the previous system size studied.
Simulation results are reported in Sec. 4, for the order parameter and the
staggered susceptibility as functions of $h$ and $T$. Points along the
critical line in the $h-T$ plane are obtained using finite size scaling
analysis, and the results compared with several theoretical expressions. We
summarize our findings in Sec. 5.
## 2 Configuration Classes
As pointed out above, one of the main problems in entropic sampling methods is
the prior determination of the complete set of configuration classes for a
given system size. Let us denote by $N_{+}$ and $N_{-}$ the number of up and
down spins, respectively, on a square lattice with $N=N_{+}+N_{-}=L^{2}$
spins. The number of pairs of opposite spins, $u$, and $N_{+}$ are related to
$n$ and $m$, respectively, via
$n=2L^{2}-u$ (2)
and
$m=2N_{+}-L^{2}.$ (3)
Thus, once the possible values of $(N_{+},u)$ are determined, so are those of
$(n,m)$. Note that $u$, $n$ and $m$ can only take even values. The gaps in CS
fall near the maximum and minimum values of $n$ ($n_{max}$ and $n_{min}$,
respectively) for a given $m$. Therefore, we will identify the possible values
of $u$ near its maximum, $u_{max}$, and minimum, $u_{min}$, for a given
$N_{+}\in[0,\frac{L^{2}}{2}]$; note that the number of configurations is
symmetric under interchange of $N_{+}$ and $N_{-}$.
### 2.1 Determining $u_{min}$ and nearby classes
#### 2.1.1 Compact configurations
Compact configurations consist of a square or rectangular cluster with $N_{+}$
up spins. To begin, consider the case of a square cluster of size $l\times l$
$(2<l<L)$ with $N_{+}=l^{2}$, as is illustrated in Fig. 2. It is evident that
this configuration corresponds to the minimum value of $u$, $u^{(0)}=4l$. A
configuration with the same number of up spins, and $u=N_{+}+2$ is obtained by
transferring an up spin from one of the corners of the square to an edge, as
shown in Fig. 3. From this configuration, further rearrangements leading to
$u=u_{min}+4$, etc., are possible. When $N_{+}$ is not a square number, the
most compact configuration (i.e., with the smallest perimeter) is a rectangle,
or a square or rectangle with an incomplete layer of sites along one edge. For
$N_{+}=l(l-1)$ we have $u_{min}=4l-2$, while for $l(l-1)<N_{+}<l^{2}$,
$u_{min}=4l$. In all cases, moving a corner site to an edge, one constructs a
configuration with $u=u_{min}+2$, and further arrangements yield additional
increases in $u$.
Figure 2: Basic compact configuration. Up and down spins are represented by
“$\times$” and “$\bullet$”, respectively; wavy lines represent pairs of
opposite spins. Figure 3: Modified compact configuration. The new up and
down spins are represented by “$\otimes$” (previously “$\bullet$”) and
“$\odot$” (previously “$\times$”), respectively; double straight and double
wavy lines represent new pairs of identical and opposite NNs, respectively.
#### 2.1.2 Extended configurations
Thus there are no gaps in the large-$n$ region due to compact configurations.
Such compact configurations, however, are not necessarily the ones that
minimize $u$ for a given value of $N_{+}$. Consider for example the case
$N_{+}=kL$ with $k$ an integer. The cluster of up spins can be arranged to
wrap around the lattice, yielding $u_{min}=2L$, independent of $k$. We call
such a configuration extended. For $k$ greater than a certain value, on the
order of $L/4$, the configurations that minimize $u$ are extended rather than
compact. Suppose now that for a given $N_{+}$, $u_{min}$ is obtained with an
extended configuration, and that all compact configurations have $u$ strictly
greater than $u_{min}$. If $N_{+}$ lies between $kL$ and $(k+1)L$, then the
configuration that minimizes $u$ has at least two corners, and by the previous
construction, configurations with $u=u_{min}+2$ exist. But if $N_{+}=kL$, the
minimizing configuration has no corners and any modification yields a
configuration with $u\geq u_{min}+4$. This is how the gaps near the maximum
values of $n$ arise.
Summarizing, if $L$ and $N_{+}$ are such that $u_{min}$ is obtained with a
compact configuration, then there are configurations with $u_{min}+2$,
$u_{min}+4$, …, etc., and no gap exists. If, on the other hand, $u_{min}$ is
realized only for an extended configuration, and $N_{+}$ is an integer
multiple of $L$, then there are no configurations with $u=u_{min}+2$.
### 2.2 Determining $u_{max}$ and nearby classes
The largest possible value of $u$, $u_{max}$, occurs in a configuration such
that $N_{+}=L^{2}/2$ with spins arranged in a chess board (CB) configuration,
such that all up spins have down spins as NNs and vice versa. One readily
verifies that for $0\leq N+\leq L^{2}/2$, the maximum number of unlike NN
pairs is $u_{max}=4N_{+}$. Starting from the CB configuration, we can reduce
$u$ by exchanging an up and a down spin. If the exchanged spins are NNs, the
resulting configuration has $u=u_{max}-6$; otherwise one has $u=u_{max}-8$.
Thus, for $N_{+}=L^{2}/2$, there are no configurations with $u=u_{max}-2$ or
$u_{max}-4$. One readily verifies that configurations with $u=u_{max}-10$,
$u_{max}-12$, etc., can be obtained via further exchanges of spins.
For $N_{+}=L^{2}/2-1$, we have $u_{max}=2L^{2}-4$. Such configurations can be
constructed by flipping one up spin in the CB configuration. Starting from
this configuration, exchanging a pair of spins, one can reduce $u$ by 4, 6 or
8, but there is no rearrangement which reduces $u$ by just two. Thus for
$N_{+}=L^{2}/2-1$, there is no configuration having $u=u_{max}-2$;
configurations with $u=u_{max}-4$, $u_{max}-6$, etc. do exist.
Finally, we note that for $1<N+<L^{2}/2-1$, there are no gaps in the
neighborhood of $u_{max}$. To verify this, consider a configuration obtained
by flipping $k$ up spins in the CB configuration, so that $N_{+}=L^{2}/2-k$,
where $1<k<L^{2}/2$. By hypothesis there are now at least four more sites with
down spins than with up spins. Starting from a configuration with $u=u_{max}$,
in which each up spin is surrounded by down spins, we can create a single NN
pair of up spins, with all six neighbors down, and with all remaining up spins
completely surrounded by down spins. In this manner, $u$ is reduced by two.
Configurations with $u=u_{min}-4$, $u_{min}-6$, etc. can be obtained by
further exchanges of up and down spins.
Using the facts summarized above, it is straightforward to construct an
algorithm that determines which values of $u$ are possible, for a given $L$
and $N_{+}$, and thereby which values $(n,m)$ are accessible for a given
system size. In the simulations reported below, we have verified that our
entropic sampling scheme converges to visit all allowed classes.
## 3 Implementation
Using tomographic sampling, we study the antiferromagnetic Ising model in an
external field on the square lattice; we consider periodic boundary conditions
and NN interactions. The CCs of the systems are defined in the energy-
magnetization space $(n,m)$.
The TES method is applied in order to generate estimates of $\Omega(n,m,L)$.
For the smallest system size $(L=10)$ we begin with a guess,
$\Omega_{0}(n,m)$, obtained using a mean-field approximation. For subsequent
system sizes, however, we use a two-dimensional interpolation of
$\Omega_{N}(n,m,L)$ (the final estimate of $(n,m)$ for the smaller system
size, $L$) to construct $\Omega_{0}(n,m,L^{\prime})$, where $L^{\prime}>L$.
For most studies we use five iterations, each one with $N_{sim}=10$ initial
configurations, which are simulated for $N_{U}=10^{7}$ lattice updates or
Monte Carlo steps.
Let us denote by $\Gamma(n,m)$ and $\Gamma^{\prime}(n_{0}=n+\Delta
n,m_{0}=m+\Delta m)$ the CCs that contain configurations $\mathcal{C}$ and
$\mathcal{C^{\prime}}$, respectively. The simulation uses a single spin-flip
dynamics, so that the possible variations of $n$ and $m$ are $\Delta n=0,\pm
2,\pm 4$ and $\Delta m=\pm 2$. At iteration $j$, the acceptance probability
for the transition $(\mathcal{C}\to\mathcal{C^{\prime}})$ is
$p_{j}(\Gamma\to\Gamma^{\prime})=\textrm{min}\bigg{[}\frac{\Omega_{j-1}(\Gamma)}{\Omega_{j-1}(\Gamma^{\prime})},1\bigg{]}.$
(4)
These probabilities are stored in a table. For each configuration generated,
be it a new one (if it is accepted) or the same (if it is rejected), we update
the sums used to calculate the microcanonical and canonical averages of
$\phi$, $|\phi|$, $\phi^{2}$, and $\phi^{4}$, where
$\phi\equiv m_{A}-m_{B}$ (5)
is the order parameter (staggered magnetization); $m_{A,B}$ are the
magnetizations of the two sublattices. At the end of each iteration $j$ the
estimate of $\Omega_{j}(n,m)$ is refined according to
$\Omega_{j}(\Gamma)=\frac{H_{j}(\Gamma)}{\overline{H}_{j}(\Gamma)}\,\Omega_{j-1}(\Gamma),$
(6)
where $H_{j}(\Gamma)$ is the histogram containing the number of times the CC
$\Gamma$ is visited during the sampling, and $\overline{H}_{j}(\Gamma)$ is the
average of $H_{j}(\Gamma)$ over all accessible CCs; the acceptance probability
[Eq. (4)] is updated using the new estimate of $\Omega_{j}(n,m)$ and the
histogram, $H_{j}(n,m)$, is set to zero.
A single iteration of our method consists of ten independent simulations, each
involving $10^{7}$ lattice updates, and each beginning from a different
initial configuration. (By a lattice update we mean one attempted flip per
spin; the initial configurations include, high, low, and intermediate
interaction energies and both signs of the magnetization.)
### 3.1 Determining $\Omega_{0}(n,m,L)$: mean field approach
As noted above, for the smallest system size we use an estimate of
$\Omega_{0}(n,m,L)$ obtained via a mean-field approximation, specifically
$\frac{1}{L^{2}}\ln\Omega_{0}(n,m,L)=\frac{1}{L^{2}}\ln\Omega(m)-\frac{(n-\langle
n\rangle)^{2}}{2L^{2}\sigma^{2}}-\frac{\ln\sigma}{L^{2}}+const.\,,$ (7)
where $\sigma\equiv\sigma(m)=\sqrt{\textrm{var}(n)}$, and
$\frac{\ln\Omega(m)}{L^{2}}\simeq\ln
2-\frac{1}{2}[(1+m/L^{2})\ln(1+m/L^{2})+(1-m/L^{2})\ln(1-m/L^{2})].$ (8)
To obtain this expression, we first note that
$\Omega(m,L)=\sum_{n}\Omega(n,m,L)=\binom{L}{N_{+}}$, and use Stirling’s
approximation. The dependence on $n$ is then obtained by estimating $\langle
n\rangle$ and $\textrm{var}(n)$ for a given $m$ and $L$, using a random-mixing
approximation, and supposing that $n$ follows a Gaussian distribution.
In Fig. 4 we plot the estimate of $\ln\Omega_{0}(n,m,L=10)$ given by Eq. (7).
We present in Fig. 5 the final estimated value of $\ln\Omega_{5}(n,m,L=10)$
after the 5$th$ iteration of the simulation; this result is quite similar to
that presented by Hwang _et al._ [20] for the square lattice, which was
obtained using Wang-Landau sampling [3]. We note that the differences between
the initial estimate obtained via mean field approximation and the final
simulation result of $\ln\Omega(n,m,L=10)$ are more evident near the edges of
the CS, specially near to the maximum values of $n$. To have a better idea of
how close this initial estimate is to the final simulation result, we plot in
Fig. 6 the difference between that estimate and the final result of the
simulation for $L=10$. Analyzing this figure, it is again clear that the
differences are larger near the edges of the CS.
Figure 4: Estimate of $\ln\Omega_{0}(n,m,L=10)$ via mean-field approximation.
Figure 5: Final estimate of $\ln\Omega_{5}(n,m,L=10)$ after the 5$th$
iteration of the simulation. Figure 6: Difference between
$\ln\Omega_{0}(n,m,L=10)$, estimated via mean-field approximation, and the
final simulation result after the 5th iteration, $\ln\Omega_{5}(n,m,L=10)$.
### 3.2 Determining $\Omega_{0}(n,m,L)$: interpolation
We expect that the closer the estimate $\Omega_{0}(n,m)$ is to the (unknown)
exact value, $\Omega_{e}(n,m)$, the faster the simulation will converge, and
the more accurate our final estimate will be. Since the mean-field estimate
worsens as the system size grows, we only use it for the smallest system size
studied; initial estimates of subsequent system sizes are obtained by
interpolating the final estimate, $\Omega_{N}(n,m,L)$, of the previous system
size studied. Our procedure is based on the existence of the limiting
microcanonical entropy density as a function of the intensive parameters
$\eta$ and $\nu$ (the bond and magnetization densities, respectively),
$s(\eta,\nu)=\lim_{L\to\infty}\frac{1}{L^{2}}\ln\Omega(n_{L},m_{L},L),$ (9)
where $n_{L}\simeq\eta L^{2}$ and $m_{L}\simeq\nu L^{2}$. (Since $n$ and $m$
are restricted to even integers we have $n_{L}=\eta L^{2}+{\cal O}(1/L^{2})$
and similarly for $m_{L}$.) The idea is then to write
$\frac{1}{L^{\prime
2}}\ln\Omega_{0}(n^{\prime},m^{\prime},L^{\prime})=\frac{1}{L^{2}}\ln\Omega_{N}(\eta
L^{2},\nu L^{2},L),$ (10)
where $\eta=n^{\prime}/L^{\prime 2}$, $\nu=m^{\prime}/L^{\prime 2}$, and the
r.h.s. is evaluated by extending $\ln\Omega_{N}$ to noninteger $n$ and $m$ via
extrapolation and interpolation. Using this approach, we obtain better
estimates, as is shown in Fig. 7. (Note that the largest differences continue
to fall along the edges of the CS.)
Figure 7: Difference between $\ln\Omega_{0}(n,m,L=18)$, estimated by
interpolating the final result of $L=16$, and the final simulation result
after the 5th iteration, $\ln\Omega_{5}(n,m,L=18)$.
### 3.3 Extrapolation and interpolation
Suppose we wish to construct the initial estimate $\Omega_{0}(n,m,L_{2})$ on
the basis of the simulation results for a smaller system,
$\Omega_{N}(n,m,L_{1})$. We could do this via interpolation if every CC of the
larger system were surrounded by four points of the smaller one in the
$\eta$-$\mu$ plane. Fig. 8 shows, however, that along the edges of the CS, the
points corresponding to classes of the larger system are not surrounded by
four points of the smaller one. For those points, one could in principle use
extrapolation rather than interpolation. We found, however, that direct
extrapolation yields poor estimates for $\Omega$.
Figure 8: Region of the $(\eta,\nu)$ plane with every possible configuration
classes for system sizes $L_{1}=26$ and $L_{2}=28$.
We obtain better estimates by first extrapolating the points along the edges
of the CS for the smaller system, such that every point of the larger system
is surrounded by four points of the smaller. Figure 9 shows the CS with
accessible and extrapolated CCs for $L=10$. Following this extrapolation we
perform a linear two-dimensional interpolation as per Eq. (10).
Figure 9: Accessible and extrapolated configuration classes of a system of
size $L=10$.
## 4 Results
In this section we present results of the AF Ising model on the square lattice
in an external field; periodic boundary conditions are employed. We use TES to
simulate systems of sizes $L=10$ to 30, with $\Delta L=2$. To calculate the
uncertainties we perform five independent studies for each system size. We
plot in Fig. 10 the staggered magnetization (order parameter) per site,
$\phi$, as a function of $h$ at $T=0.2$. We can see that $\phi$ decreases
considerably between $h=3.85$ and $h=3.90$; this behavior suggests a critical
point, $h_{c}$, marking a phase transition from the AF to the paramagnetic
state. Figure 11 shows the staggered susceptibility per site,
$\chi_{\phi}\equiv\frac{1}{L^{2}}(\langle\phi^{2}\rangle-\langle\phi\rangle^{2}),$
(11)
as a function of $h$ at $T=0.02$. As $L$ grows the peaks tend to the critical
point, $h_{c}$. The specific heat per site, $c$, (not shown) has a similar
behavior.
Figure 10: Staggered magnetization per site as a function of $h$ at $T=0.2$,
for $L=10$ to 30. The absolute uncertainties are plotted in the inset. Figure
11: Staggered susceptibility per site as a function of $h$ at $T=0.02$, for
$L=10$ to 30. The highest peaks correspond to the largest system sizes. The
absolute uncertainties are plotted in the inset.
### 4.1 Phase diagram
Using finite size scaling analysis [21] we estimate the critical line,
$h_{c}(T)$, or, equivalently, $T_{c}(h)$, via the relations:
$h_{c}(Y_{max},T,L)=h_{c}(Y_{max},T)+a_{1}/L+a_{2}/L^{2},$ (12)
and
$T_{c}(Y_{max},h,L)=T_{c}(Y_{max},h)+b_{1}/L+b_{2}/L^{2},$ (13)
where $h_{c}(Y_{max},T,L)$ is the field at which $Y$ (the specific heat or the
staggered susceptibility) takes its maximum for a given temperature and system
size; $T_{c}(Y_{max},h,L)$ is defined in an analogous manner. The estimates
obtained using the maximum of $c$ and $\phi$ are averaged to yield $h_{c}(T)$
and $T_{c}(h)$; Figure 12 illustrates the procedure for estimating $h_{c}$,
for $T=0.02$. Our estimated points for the phase boundary are shown in Fig. 13
along with the approximate expression derived by Müller-Hartmann and Zittartz
[8]:
$\cosh\bigg{(}\frac{h}{T_{c}}\bigg{)}=\sinh^{2}\bigg{(}\frac{2J}{T_{c}}\bigg{)}.$
(14)
Our simulation results are in good agreement with their expression. For the
critical field, the greatest relative difference between theory and simulation
is about 0.9%, which occurs at $T_{c}=1.8$.
Figure 12: Critical field determination using finite size scaling analysis:
$\overline{h}_{c}(T=0.02)=3.98666(3)$. Figure 13: Phase diagram of the Ising
AF on the square lattice. Comparison between simulation and the theoretical
expression of Müller-Hartmann and Zittartz. Asterisks denote points obtained
varying $h$ with $T$ fixed; circles denote points obtained varying $T$, with
$h$ fixed. The error bars of our results are smaller than the symbols.
In Table 1 we compare our simulation estimates for the critical magnetic field
$h_{c}(T)$ with some theoretical approximations. For temperature $T\leq 1$ we
find good agreement with the estimates of Monroe [17] (whose analysis involves
a free parameter $\omega$), Wu and Wu [12], and Blöte and Wu [13], whereas
there are significant discrepancies in relation to the other approximations.
At higher temperatures, differences appear between simulation and the
predictions of Monroe, Wu and Wu, and Blöte and Wu. These may reflect a small
systematic error or an underestimate of simulation uncertainties. We intend to
examine this issue in greater detail in future work.
Table 1: Comparison between our simulation estimates of $h_{c}(T)$ with some theoretical approximations. $T$ | TES | Monroe | Monroe | MHZ | WW | BW | WK
---|---|---|---|---|---|---|---
| | $(\omega=0.92484)$ | $(\omega=0.93895)$ | | | |
0.1 | 3.93304(16) | 3.93307 | 3.93318 | 3.93069 | 3.93329 | 3.93330 | 3.93372
0.5 | 3.6648(8) | 3.66506 | 3.66561 | 3.65309 | 3.66611 | 3.66614 | 3.67589
1.0 | 3.2906(14) | 3.29303 | 3.29391 | 3.26843 | 3.29200 | 3.29261 | 3.31764
1.5 | 2.7258(14) | 2.73243 | 2.73396 | 2.70401 | 2.73094 | 2.73176 | 2.75099
2.0 | 1.696(2) | 1.71629 | 1.71872 | 1.69490 | 1.71492 | 1.71499 | 1.71512
TES: Tomographic entropic sampling [6]
MHZ: Müller-Hartmann and Zittartz [8]
WW: Wu and Wu [12]
BW: Blöte and Wu [13]
WK: Wang and Kim [16]
## 5 Conclusions
The complete enumeration of CCs is of fundamental importance to study the
antiferromagnetic Ising model using tomographic entropic sampling. The
determination of CCs for entropic sampling of Ising models also enables us to
obtain good initial estimates for the configuration numbers
$\Omega_{0}(n,m,L)$, using two-dimensional linear interpolation; the initial
estimate is reasonably close to the final estimate $\Omega_{N}$. Despite the
relatively small system sizes used in this study, we obtain a good estimate
for the critical line in the temperature-magnetic field plane. Further details
on critical behavior will be published elsewhere [22].
## Acknowledgments
We are grateful to CNPq and CAPES, Brazil, for financial support.
## References
* [1] B. A. Berg and T. Neuhaus, Phys. Rev. Lett. 68, 9 (1992).
* [2] J. Lee, Phys. Rev. Lett. 71, 211 (1993).
* [3] F. Wang and D.P. Landau, _Phys. Rev. Lett._ , 84, 10 (2001).
* [4] F. Wang and D.P. Landau, _Phys. Rev. E_ , 63, 056101 (2001).
* [5] M.E.J Newman and G.T. Barkema. _Monte Carlo Methods in Statistical Physics_ (Oxford University Press, New York, 2001), p. 161 - 169.
* [6] R. Dickman and A.G. Cunha-Netto, _Phis. Rev. E_ , 84, 026701 (2011).
* [7] L. Onsager, _Phys. Rev._ , 65, 117 (1944).
* [8] E. Müller-Hartmann and J. Zittartz, _Z. Phys._ , 73, 261 (1977).
* [9] L. Sneddon, _J. Phys. C_ , 12, 3051 (1979).
* [10] R.R. Santos, _J. Phys. C_ , 18, L1067 (1985).
* [11] M. Kaufman, _Phys. Rev. B_ , 36, 3697 (1987).
* [12] X.N. Wu and F.Y. Wu, _Phys. Lett. A_ , 144, 123 (1990).
* [13] H.W.J. Blöte and X-N. Wu, _J. Phys. A_ , 23, L627 (1990).
* [14] M. Badehdah, A. Benyoussef and M. Touzani, _J. Magn. Magn. Mater._ , 172, 254 (1997).
* [15] X-Z. Wang and J.S. Kim, _Phys. Rev. Lett._ , 78, 413 (1997).
* [16] X-Z. Wang and J.S. Kim, _Phys. Rev. E_ , 56, 2793 (1997).
* [17] J.L. Monroe, _Phys. Rev. E_ , 64, 016126 (2003).
* [18] S.J. Penney, V.K. Cumyn and D.D. Betts, _Physica A_ , 330, 507 (2003).
* [19] K. Binder and D.P. Landau, _Phys. Rev. B_ , 21, 5 (1980).
* [20] C-O. Hwang, S-Y. Kim, D. Kang and J.M Kim, _J. Stat. Mech: Theory Exp._ , 2007, L05001 (2007).
* [21] V. Privman. _Finite Size Scaling Analysis and Numerical Simulations of Statistical Systems_ (World Scientific, London, 1990).
* [22] B.J. Lourenço and R. Dickman, _in preparation_.
|
arxiv-papers
| 2011-12-06T19:51:20 |
2024-09-04T02:49:25.028746
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Bruno Jeferson Louren\\c{c}o and Ronald Dickman",
"submitter": "Bruno Louren\\c{c}o",
"url": "https://arxiv.org/abs/1112.1387"
}
|
1112.1409
|
# Separating the BL Lac and Cluster X-ray Emissions in Abell 689 with
_Chandra_
P. A. Giles, B. J. Maughan, M. Birkinshaw, D. M. Worrall, K. Lancaster
HH Wills Physics Laboratory, Tyndall Avenue, Bristol, BS8 1TL, UK E-mail:
P.Giles@bristol.ac.uk
(Accepted 2011 August 26. Received 2011 August 4; in original form 2011 March
10)
###### Abstract
We present the results of a _Chandra_ observation of the galaxy cluster Abell
689 (z=0.279). Abell 689 is one of the most luminous clusters detected in the
ROSAT All Sky Survey (RASS), but was flagged as possibly including significant
point source contamination. The small PSF of the _Chandra_ telescope allows us
to confirm this and separate the point source from the extended cluster X-ray
emission. For the cluster we determine a bolometric luminosity of Lbol =
(3.3$\pm$0.3)$\times$1044 erg s-1 and a temperature of
kT=5.1${}^{+2.2}_{-1.3}$ keV when including a physically motivated background
model. We compare our measured luminosity for A689 to that quoted in the Rosat
All Sky Survey (RASS) and find L0.1-2.4,keV = 2.8$\times$1044 erg s-1, a value
$\sim$10 times lower than the ROSAT measurement. Our analysis of the point
source shows evidence for significant pileup, with a pile-up fraction of
$\simeq$60%. SDSS spectra and HST images lead us to the conclusion that the
point source within Abell 689 is a BL Lac object. Using radio and optical
observations from the VLA and HST archives, we determine $\alpha_{\rm
ro}$=0.50, $\alpha_{\rm ox}$=0.77 and $\alpha_{\rm rx}$=0.58 for the BL Lac,
which would classify it as being of ‘High-energy peak BL Lac’ (HBL) type.
Spectra extracted of A689 show a hard X-ray excess at energies above 6 keV
that we interpret as inverse Compton emission from aged electrons that may
have been transported into the cluster from the BL Lac.
###### keywords:
galaxies: clusters: general – galaxies: clusters: individual: Abell 689 – BL
Lacertae objects: general
††pubyear: 2011
## 1 Introduction
Studies of clusters of galaxies, including measurements of their number
density and growth from the highest density perturbations in the early
Universe, offer insight into the underlying cosmology (Vikhlinin et al., 2003;
Allen et al., 2004). However, in order to use clusters as a cosmological probe
three essential tools are required (Del Popolo et al., 2010): (a) an efficient
method to find clusters over a wide redshift range, (b) an observational
method of determining the cluster mass, and (c) a method to compute the
selection function or the survey volume in which clusters are found. These
requirements are met by large surveys with well understood selection criteria.
Arguably the most effective method of building large, well defined cluster
samples has been via X-ray selection. The high X-ray luminosities of clusters
make it relatively easy to detect and study clusters out to high redshifts.
Many cluster samples have been constructed based upon large X-ray surveys such
as the Einstein Medium Sensitivity Survey (EMSS; Gioia et al., 1990) and the
ROSAT All Sky Survey (RASS; Voges, 1992). However due to the relatively poor
angular resolution of these X-ray observatories, observations of clusters were
susceptible to point source contamination. Indeed, within the ROSAT Brightest
Cluster Sample (BCS; Ebeling et al., 1998) and its low-flux extension (eBCS;
Ebeling et al., 2000), 9 out of 201 clusters and 8 out of 99 clusters
respectively were flagged as probably having a significant fraction of the
quoted flux from embedded point sources. Being able to resolve these point
sources is of crucial importance for the reliable estimation of cluster
properties, and indeed the nature of the point source contamination is of
independent interest.
The study of galaxy clusters has been transformed with the launch of powerful
X-ray telescopes such as _Chandra_ and XMM Newton, which have allowed the
study of the X-ray emitting intracluster medium (ICM) with unprecedented
detail and accuracy. With the launch of this new generation of X-ray
telescopes, we are able to uncover interesting features in the morphologies of
individual clusters. In particular, _Chandra’s_ high angular resolution
provides the means to examine individual cluster features with great detail.
Abell 689 (hereafter A689; Abell, 1958), at z=0.279 (Collins et al., 1995),
was detected in the RASS in an accumulated exposure time of 317s. It is
included in the BCS, with a measured X-ray luminosity of 3.0$\times$1045 erg
s-1 in the 0.1–2.4 keV band. This luminosity is the third highest in the BCS,
and thus A689 meets the selection criteria for various highly X-ray luminous
cluster samples (e.g. Dahle et al., 2002). However this cluster was noted as
having possible point source contamination, and for this reason has often been
rejected from many flux limited samples. In this study we present results of
_Chandra_ observations designed to separate any point sources and determine
uncontaminated cluster properties.
The outline of this paper is as follows. In $\S$ 2 we discuss the observation
and data analysis. Results of the X-ray cluster analysis is presented in $\S$
3. In $\S$ 4 we present our analysis of the X-ray point source through X-ray,
optical and radio observations. We interpret our results in $\S$ 5 and the
conclusions are presented in $\S$ 6. Throughout this paper we adopt a
cosmology with $\Omega_{M}$= 0.3, $\Omega_{\Lambda}$ = 0.7 and H0 = 70 km s-1
Mpc-1, so that 1′′ corresponds to 4.22 kpc at the redshift of A689. We define
spectral index, $\alpha$, in the sense S $\propto$ $\nu^{-\alpha}$.
## 2 Observations and data reduction
The _Chandra_ observation of A689 (ObsID 10415) was carried out January 01,
2009. A summary of the cluster’s properties is presented in Table 1. The
observation was taken in VFAINT mode, and the source was observed in an ACIS-I
configuration at the aim point of the I3 chip, with the ACIS S2 chip also
turned on.
Name | RA | Dec | z | NH,Gal | RASS LX,0.1-2.4keV
---|---|---|---|---|---
Abell 689 | 08h 37m 24s.70 | +14o 58′ 20′′.78 | 0.279 | 3.66$\times$1020 cm-2 | 30.4$\times$1044 erg s-1
Table 1: Basic properties. Columns: (1) = Source name; (2) = Right Ascension
at J2000 from Chandra; (3) = Declination at J2000 from Chandra; (4) =
Redshift; (5) = Galactic column density; (6) = Intrinsic X-ray luminosity in
the 0.1–2.4 keV band based upon ROSAT observations (Ebeling et al., 1998).
For the imaging analysis of the cluster we used the CIAO111See
http://cxc.harvard.edu/ciao/ 4.2 software package with CALDB222See
http://cxc.harvard.edu/caldb/ version 4.3.0 and followed standard reduction
methods. Since our observation was telemetered in VFAINT mode additional
background screening was applied by removing events with significantly
positive pixels at the border of the 5$\times$5 event island333See
http://cxc.harvard.edu/ciao/why/aciscleanvf.html. We inspected background
light curves of the observation following the recommendations given in
Markevitch et al. (2003), to search for possible background fluctuations. The
light curve was cleaned by 3$\sigma$ clipping and periods with count rates
$>$20$\%$ above the mean rate were rejected. curve with rejected periods
showed in red. The final level-2 event file had a total cleaned exposure time
of 13.862 ks.
As discussed in Sect. 4, there is a bright point source at the center of the
extended cluster emission which is affected by pileup. For the analysis of
this source we followed the same reduction method, with the exception that
VFAINT cleaning was not applied. Applying VFAINT cleaning leads to incorrect
rejection of piled-up events, introducing artifacts in the data.
Figure 1: Background light curve for the observation of A689 in the 0.3-12.0
keV band. The CCD on which the cluster falls (ACIS-I3) and all point sources
are excluded. The red bands show periods excluded by the Good Time Interval
(GTI) file.
## 3 X-RAY CLUSTER ANALYSIS
In this section we determine global cluster properties of Abell 689. Figure 2
shows a Gaussian smoothed image of the cleaned level 2 events file in the
0.7-2.0 keV band (the readout streak is removed using the CIAO tool
ACISREADCORR), with an inset image of the point source which lies at the
center of the cluster. The extent of the diffuse cluster emission was
determined by plotting an exposure-corrected radial surface brightness profile
(Fig 3), in the 0.7-2.0 keV band, of both the observation and blank-sky
background to determine where the cluster emission is lost against the
background. Figure 3 demonstrates that the diffuse cluster emission is
detectable to a radius r $\approx$ 570′′ ($\approx$ 2.41 Mpc). At large radii
(r$\geq$700′′) the curves rise due to vignetting corrections (larger at larger
radii) applied to all the counts, whereas in reality each curve contains a
component from particles that have not been focused by the telescope.
Figure 2: 0.7-2.0 keV image of A689, smoothed by a Gaussian ($\sigma$=1.5
pixels, where 1 pixel = 3.94′′), cleaned in VFAINT mode and with the readout
streak removed. Inset is a zoomed-in unbinned image of the central point
source within A689, cleaned in FAINT mode (see $\S$ 2). The inner black circle
(r = 208′′) is the region in which we extract the spectra for our analysis of
the cluster emission (see $\S$ 3.2), the outer black circle (r = 570′′)
represents the detected cluster radius (see $\S$ 3), and the region between
this and the black box was used for the local background (see $\S$ 3.2). The
white box displays the size of the inset, and the inner white circle (r =
26′′) shows the region excluded due to the central point source. Many point
sources are seen in the observation and are excluded from our analysis. An
extended source on the NE chip can be seen, that is unrelated to A689 and is
also excluded from our analysis. Figure 3: Exposure-corrected radial surface
brightness profiles (0.7-2.0 keV band) of the cluster (red) and the blank-sky
background (blue), with the green line representing the radius beyond which no
significant cluster emission is detected.
### 3.1 Background Subtraction
In order to take the background of the observation into account, appropriate
period E _Chandra_ blank-sky backgrounds were obtained, processed identically
to the cluster, and reprojected onto the sky to match the cluster observation.
We then followed a method similar to that outlined in Vikhlinin et al. (2006),
in order to improve the accuracy of the background by applying small
adjustments to the baseline model. Firstly we corrected for the rate of
charged particle events, which has a secular and short-term variation by as
much as 30%. We renormalise the background in the 9.5–12 keV band, where the
_Chandra_ effective area is nearly zero and the observed flux is due entirely
to the particle background events. The renormalisation factor was derived by
taking the ratio of the observed count rate in the source and background
observations respectively. The normalised spectrum of the blank-sky background
is shown in Figure 4, over-plotted on the local background for comparison. The
spectra agree well in the 9.5–12.0 keV band, and across the whole spectrum
with only slight differences. In addition to the particle background, the
blank-sky and source observations contain differing contributions from the
soft X-ray background, containing a mixture of the Galactic and geocoronal
backgrounds, significant at energies $\leq$1 keV. To take into account any
difference in this background component between the blank-sky and source
observations, the spectra were subtracted and residuals were modeled in the
0.4-1keV band using an APEC thermal plasma model (Smith et al., 2001). This
model was included in the spectral fitting for the cluster analysis. As can be
seen in Figure 4 this component is very weak in the case of A689.
Figure 4: Comparison of the local (black) and blank-sky (red) background
spectra, normalised to match in the 9.5-12.0 keV band.
### 3.2 X-ray Cluster Properties
The analysis of the diffuse X-ray emission allows us to determine the X-ray
environment surrounding the cluster central point source. Throughout this
process we excluded the central point source (Fig 2; r$\leq$26′′) and the
associated readout streaks to avoid contaminating the cluster emission.
To determine cluster properties we extract spectra out to a radius chosen so
that the cluster has the maximum possible signal-to-noise (SNR). The net
number of counts, corrected for background, is then 414 (with SNR = 15) and
the extraction annulus is 26′′$<$r$<$208′′ (see Fig 2), centered on the
cluster (at $\alpha$, $\delta$ = 08h 37m 24s.70, 14o 58′ 20′′.78). We fitted
the extracted spectrum in XSPEC with an absorbed thermal plasma model
(WABS$\times$APEC) and subtracted the background described in $\S$ 3.1. We
obtain a temperature of 13.6${}^{+13.2}_{-5.1}$ keV and a bolometric
luminosity of Lbol=(10.2$\pm$2.9)$\times$$10^{44}$ ergs s-1. The measured
temperature is far above what we would expect from the luminosity. Figure 5
shows the luminosity-temperature relation for a sample of 115 galaxy clusters
(Maughan et al., 2008), along with the luminosity and temperature derived for
A689 from our values above (pink square). As the Maughan et al. (2008) sample
of clusters covers a wide redshift range, the luminosities of the clusters
were corrected for the expected self-similar evolution, given by
LX$\times$E(z)-1, where:
$E(z)=\Omega_{\rm M0}(1+z)^{3}+(1-\Omega_{\rm
M0}-\Omega_{\Lambda})(1+z)^{2}+\Omega_{\Lambda}.$ (1)
The same correction was also applied to the A689 data for the plot. Our
luminosity was derived within the same annular region as the cluster
temperature and extrapolated both inward and outward in radius between
(0-1)r500 (where r500 represents the radius at which the density of the
cluster is 500 times the critical density of the Universe at that redshift)
using parameters from a $\beta$-profile fit to the surface brightness profile.
This takes into account our exclusion of the region near the central point
source within the cluster, and the extrapolation to r500 is in order to
compare with the Maughan et al. (2008) sample. Our r500 value was determined
from the temperature and using the relation between r500 and T given in
Vikhlinin et al. (2006).
The high temperature for A689 could be an indication that our observation
suffers from background flaring which would lead to an overestimate of the
cluster temperature. However, no evidence is found that the spectrum of the
blank-sky background differs from that of the local background (Figure 4), or
that the background of the observation suffers from periods of flaring (Figure
1). To investigate the sensitivity of the temperature to our choice of
background, we repeated the analysis using a local background region far from
the cluster emission (see Fig 2). We obtained a temperature of
10.0${}^{+13.8}_{-3.3}$ keV. This temperature is again anomalous given the
luminosity-temperature relation (Fig 5, green triangle. The luminosity is
derived using the same method as above, only this time using this new value
for the temperature to derive r500. The spectrum with a local background
subtracted is shown in Figure 6. We see an excess of photons in the $\sim$6-9
keV band, which might have an effect on the spectral fit. We perform the same
fit using a local background, however this time fitting in the 0.6-6.0 keV
band to ignore these excess high energy photons. We obtain a temperature of
9.4${}^{+8.7}_{-2.9}$ keV.
Figure 5 shows that the temperature is high for the luminosity, or the cluster
is X-ray under-luminous. Before considering a physical interpretation for the
apparently high cluster temperature, we investigated possible systematic
effects from the background subtraction. This was done by independently
modeling the background.
Figure 5: Luminosity-Temperature relation of a sample of 115 clusters of
Maughan et al. (2008)(blue open circles). The luminosities are measured within
[0 $<$ r $<$ 1]r500 and the temperatures within [0.15 $<$ r $<$ 1.0]r500, in
order to minimize the effect of cool cores on the derived cluster temperature.
Our derived temperatures for A689 are overplotted for comparison (pink square,
green triangle, red diamond) (see $\S$ 3.2). Figure 6: Spectrum of the cluster
with the local background subtracted fitted with an absorbed thermal plasma
model, with a reduced statistic of $\chi^{2}_{\rm\nu}$=1.45 ($\nu$=79).
We model the background based upon a physical representation of its
components. Our model consists of a thermal plasma APEC model, two power-law
components and five Gaussian components. The APEC model and one of the power-
law components are convolved with files describing the telescope and
instrument response (the Auxiliary Response File (ARF) and Redistribution
Matrix File (RMF)), and are intended to model soft X-ray thermal emission from
the Galaxy and unresolved X-ray background. The second power-law component is
not convolved with the ARF as it is used to describe the high energy particle
background, which does not vary with effective area. The five Gaussian
components are similarly not convolved, as these are used to model line
features in the spectrum caused by the fluorescence of material in the
telescope and focal plane caused by high energy particles.
As described below and summarized in Table 2, the parameters of this
background model were fit to the blank-sky background, local background or
high energy cluster regions or taken from the literature, in order to build
the most reliable model possible. We start by modeling blank-sky background in
the region we used to extract our cluster spectrum, with the model outlined
above. This allows us to place reasonable constraints on the line energy and
widths of each Gaussian component, and the slope of the unconvolved power-law.
We find line features at energies of 1.48, 1.75, 2.16, 7.48 and 8.29 keV. We
also model the convolved power-law component, fixing the slope at a value of
1.48 taken from Hickox & Markevitch (2006). The normalisation of this power-
law component will then be used throughout our modeling process, scaled by
area where necessary. A spectrum of the blank-sky background and a fit using
our model are shown in Figure 7.
Next we further constrained model parameters by fitting our high-energy
Gaussian and unconvolved power-law components in the 5.0-9.0 keV band within
the cluster region. At these energies, the emission is dominated by particle
background. We fix the line energies and widths to those found in the blank-
sky background and fit for the normalisations. This fit finds a slope of the
power-law consistent with that found in the blank-sky background. We therefore
fix the slope of the power-law at $\Gamma$=0.0061, as found for the blank-sky
background, and fit for the normalisation. We finally fit the low energy
Gaussians and APEC normalisation model parameters in a local background region
far from the cluster emission. The high energy Gaussians and unconvolved
power-law components are frozen at the values found in the blank-sky and
cluster regions, with the normalisations scaled by area. The convolved power-
law component is frozen at the values found in the blank-sky background and
the normalisation is scaled by area. The Gaussian features at 1.48, 1.75 and
2.16 keV are frozen at the energies and widths found in the blank-sky
background. The temperature of the APEC model is frozen at 0.177 keV (taken
from Hickox & Markevitch, 2006). We note that our APEC temperature is not well
constrained by our data, but this is a weak component. A spectrum of the local
background and the corresponding fit with the model are shown in Figure 8.
We now model the cluster with an absorbed thermal plasma (WABS$\times$APEC)
model, including the background model outlined above. The normalisations of
the background APEC component and the Gaussians at energies 1.48, 1.75 and
2.16 keV were scaled by the ratio of the areas from the local background
region to the cluster region. The normalisations of the fluorescent lines also
vary with detector position. To account for this effect in the low energy
($<$3 keV) lines (where we must fit the normalisations in the off-axis local
background region), we measure the relative change in the normalisations of
each line between the local background and source regions in the blank-sky
background data. The normalisation of each low energy Gaussian component in
the fit to the cluster data is scaled for the different detector of the
cluster region by the relative change in normailsiation determined above (in
addition to the geometrical scaling for size of the extraction region). The
unconvolved power-law and Gaussians at 7.48 and 8.29 keV are all frozen at the
values found in the 5.0-9.0 keV cluster region fit. The normalisation of the
convolved power-law is frozen at the value found in the blank-sky background.
All parameters of the model used to describe the background are frozen in the
corresponding cluster fit, we also freeze the redshift at 0.279 and the
abundance at 0.3. Our fit yields a temperature of 5.1${}^{+2.2}_{-1.3}$ keV
($\chi^{2}_{\rm\nu}$=1.15 ($\nu$=79)). We measure a bolometric luminosity of
Lbol = 1.7$\times$1044 erg s-1. The result is shown in Figure 5 (red diamond).
The spectrum with the corresponding fit to the cluster including the
background model is shown in Figure 9.
Component | Represents | Parameter | Value | Where measured
---|---|---|---|---
Convolved power-law | Unresolved X-ray bg | slope | 1.48 | Hickox & Markevitch (2006)
normalisation | 4.11$\times$10-6 | blank-sky bg
Unconvolved power-law | Particle bg | slope | 0.061 | blank-sky bg
normalisation | 0.015 | 5.0-9.0 keV cluster region
Gaussian 1 | Al K$\alpha$ fluorescence | energy | 1.48 keV | blank-sky bg
width | 0.022 keV | blank-sky bg
normalisation | 1.82$\times$10-4 | local bg
Gaussian 2 | Si K$\alpha$ fluorescence | energy | 1.75 keV | blank-sky bg
width | 0.95 keV | blank-sky bg
normalisation | 1.45$\times$10-2 | local bg
Gaussian 3 | Au M$\alpha\beta$ fluorescence | energy | 2.16 keV | blank-sky bg
width | 0.045 keV | blank-sky bg
normalisation | 2.24$\times$10-3 | local bg
Gaussian 4 | Ni K$\alpha$ fluorescence | energy | 7.48 keV | blank-sky bg
width | 0.022 keV | blank-sky bg
normalisation | 5.73$\times$10-3 | 5.0-9.0 keV cluster region
Gaussian 5 | Cu + Ni fluorescence | energy | 8.29 keV | blank-sky bg
width | 0.168 keV | blank-sky bg
normalisation | 4.56$\times$10-3 | 5.0-9.0 keV cluster region
APEC | Galactic foreground emission | kT | 0.177 keV | Hickox & Markevitch (2006)
abundance | 1.0 | solar abundance
redshift | 0 | Galactic
normalisation | 2.9$\times$10-5 | local bg
Table 2: Table of the individual model components used to represent the
background, with a brief interpretation of each component, individual
component parameters, parameter values and where each value is calculated. All
normalisations are measured in photons/keV/cm2/s at 1 keV, and scaled to the
cluster region. Blank-sky background parameters were derived in the same
region and the local background were measured in an area 2.45 times that of
the cluster and therefore the normalisations reduced by this factor. The low
energy Gaussians were also corrected due to the normalisation dependence with
position on the detector (see text). Figure 7: Spectrum of the blank-sky
background in the source region and corresponding the fit (see Sect. 3.2),
$\chi^{2}_{\rm\nu}$=1.17 ($\nu$=585). Figure 8: Spectrum of the local
background and corresponding fit (see Sect 3.2). $\chi^{2}_{\rm\nu}$=0.79
($\nu$=120). Figure 9: Spectrum of the cluster plus background fit with an
absorbed thermal plasma model including a background model (see Sect 3.2).
$\chi^{2}_{\rm\nu}$=1.15 ($\nu$=79).
## 4 The Central Point Source
The point source is displayed in Figure 10. The presence of strong readout
streaks indicate that the point source is likely to be affected by pile-up.
The readout streaks occur as X-rays from the source are received during the
ACIS parallel frame transfer, which provides 40$\mu$s exposure per frame in
each pixel along the streak. We detail our analysis of the point source and
estimate of the pileup fraction in the following section.
Figure 10: _Chandra_ image of A689 showing the regions used for extracting
spectra of the readout streak (inner rectangles) and the corresponding
background regions for the readout streak (outer rectangles).
### 4.1 X-ray Analysis of the Point Source
As a first test of the predicted pile-up, we compared the image of the point
source to the _Chandra_ Point Spread Function (PSF). We made use of CIAO tool
MKPSF to create an image of the on-axis PSF following the method outlined in
Donato et al. (2003). This consists of merging 7 different monochromatic PSFs
chosen and weighted on the basis of the source energy spectrum between 0.3 and
8 keV. This method can be summarized as follows:
1. 1.
We first extract the energy distributions of the photons from a circular
region centered on the peak brightness of the source with a radii of 2.5′′.
2. 2.
We choose seven discrete energy values at which to creating each PSF, with the
number of counts at each energy corresponding to that PSF’s ‘weight’.
3. 3.
Using MKPSF we create seven monochromatic PSFs at the position of the point
source on the detector and co-added them. Each PSF is weighted by its relative
normalisation (found in the previous step).
Figure 11: Surface brightness profiles of the point source (red squares) and
the composite PSF (blue circles) in the 0.3-8.0 keV band, normalised to agree
in the 2.46-4.92 arcsec radii region.
Once we obtained the composite PSF, we normalised it to the counts within an
annulus (inner and outer radius 2.46 and 4.92 arcsec respectively) in order to
avoid any effects of pile-up in the central region. We then compare surface
brightness profiles of the point source and PSF to look for evidence of pile-
up in the core of the point source image. Figure 11 shows the radial surface
brightness profiles of the point source (red) and the composite PSF (blue). We
find that the point source and PSF agree well in the wings of the PSF ($>$
2.46′′) and that there is an excess in the PSF surface brightness above that
of the point source at the peak of the source. This is consistent with the
flattening of the source profile relative to that of the PSF due to pileup in
the core. The PSF then gives an estimate of the non piled up count rate. Given
this count rate we use PIMMS444http://cxc.harvard.edu/toolkit/pimms.jsp to
estimate a pile-up fraction of 65%.
We also compute a second estimate of the core count rate using the ACIS
readout streak. By fitting a model to the spectrum extracted from the readout
streak we can compare this to a spectrum extracted in the core and fitted
using a pileup model. We follow the method outlined in Marshall et al. (2005)
in order to correct the exposure time of the readout streak spectrum. For an
observation of live time tlive, a section of the readout streak that is
$\theta_{s}$ arcsec long accumulates an exposure time of ts = 4
$\times$10-5tlive$\theta_{\rm s}$/(tf$\theta_{\rm x}$) s, where $\theta_{\rm
x}$ = 0.492′′ is the angular size of an ACIS pixel. For our observation, tlive
= 13.862 ks and the frame time parameter tf = 3.1 s, giving ts = 165s in a
streak segment that is 454′′ long. Figure 10 shows the regions used for
extracting spectra of the readout and an adjacent background region. This
choice of background region ensures the cluster emission is subtracted from
the readout streak spectrum. Using the SHERPA package (Refsdal et al., 2009)
we fitted an absorbed 1-D power-law model (WABS$\times$POWER-LAW) to the
extracted spectrum of the readout streak. We obtain fit parameters for the
photon index = 2.33${}^{+0.34}_{-0.30}$ and a normalisation of 0.0033$\pm
0.0005$ photons keV-1 cm-2 s-1 ($\chi^{2}_{\rm\nu}$=0.34 ($\nu$=63)). We then
extract a spectrum of the point source in a region of radius 5′′, and
subtracted the same background as for the readout streak. We once again fitted
an absorbed power-law model, including this time a pileup model (jdpileup). We
obtain fit parameters for the photon index = 2.22${}^{+0.05}_{-0.04}$,
consistent to that found from the readout streak, and a normalisation of
0.0015$\pm$0.0001 photons keV-1 cm-2 s-1 ($\chi^{2}_{\rm\nu}$=1.6
($\nu$=119)). The extracted spectrum and corresponding fit is shown in Figure
12. The pileup fraction is estimated to be 60%, which is consistent with that
found using the PSF count rate. The normalisation found in the fit can be
converted to an X-ray flux density for this source, f${}_{\rm
1~{}keV}$=0.99$\pm$0.07 $\mu$Jy.
Figure 12: Spectra of the point source fitted with an absorbed power-law,
including a pileup model. $\chi^{2}_{\rm\nu}$=1.6 ($\nu$=119).
### 4.2 Optical Observations
Abell 689 was observed with the Hubble Space Telescope (HST) with the F606W
filter (Ṽ-band) on January 20, 2008. Marking the position of the peak X-ray
emission of the point source on the HST image, we find that this corresponds
to an object resembling an active nucleus in a relatively bright galaxy (Fig
13). We searched the SDSS DR7 archive for information on the spectral
properties of this object. At the coordinates of the X-ray point source (SDSS
coordinates of $\alpha$,$\delta$ = 08h 37m 24.7, 14o 58′ 19′′.8) we find a
blue object with a corresponding relatively featureless spectrum (Fig 14).
From SDSS we quote an r-band magnitude of 17.18. The spectrum resembles that
of a BL Lac object, a type of AGN orientated such that the relativistic jet is
closely aligned to the line of sight. From the H and K lines in the spectrum
(dotted green lines in Fig 14), thought to be from the host galaxy, the
redshift is determined to be z=0.279, consistent with the redshift assigned to
the cluster (Collins et al., 1995). Using the HST observation we measure an
optical flux for the BL Lac of f5997Å=112mJy.
Figure 13: HST image of the point source with a cross marking the position of
the peak of the X-ray emission. Figure 14: Spectra of the BL Lac from SDSS
(SDSS J083724.71+145819.8). The vertical green lines represent the H and K
lines thought to be from the host galaxy and the green line at the bottom
represents the error spectrum.
### 4.3 Radio Observations
Archival radio observations of Abell 689 are available, allowing us to
determine the radio spectral index, $\alpha$. We obtained 8.46 GHz data taken
in March 1998 from the VLA archive which we mapped using AIPS. The source is
unresolved at 8.46 GHz, and we measure a flux density of 18.6$\pm$0.27 mJy. We
also obtained a 1.4 GHz radio image from the FIRST survey, from which we
measure a flux of 62.7$\pm$0.25 mJy. From these data we obtain a spectral
index of $\alpha_{\rm r}$$\sim$0.67$\pm$0.01. We note that a slight angular
extension in the FIRST survey suggests that the 8.46 GHz image may be missing
some flux density.
## 5 Discussion
### 5.1 The ICM properties of Abell 689
| r500 | TX | LX,bol | |
---|---|---|---|---|---
Background Subtraction | (arcsec) | (keV) | ($\times$1044 erg s-1) | reduced $\chi^{2}$ | degrees of freedom
Blank-sky | 1130 | 13.6${}^{+13.2}_{-5.1}$ | 10.1$\pm$2.9 | 1.19 | 74
Local | 696 | 10.0${}^{+13.8}_{-3.3}$ | 6.2$\pm$1.4 | 1.45 | 79
Physically motivated model† | 266 | 5.1${}^{+2.2}_{-1.3}$ | 3.3$\pm$0.3 | 1.15 | 79
Table 3: Table listing the derived spectral properties of A689 for each of the
background treatments we employ in our analysis. † indicates our favored
method for determining the cluster properties of A689.
We have derived the ICM properties of A689 using three methods of background
subtraction. Table 3 shows the spectral properties of the ICM for each of the
background treatments we employ with our favored method coming from a
physically motivated model of the background components. Through detailed
modeling of the local and blank-sky backgrounds for A689, and including this
background model in a spectral fit to the cluster, we have determined a
temperature and luminosity for A689 of T = 5.1${}^{+2.2}_{-1.3}$ keV and Lbol
= 3.3$\times$1044 erg s-1. Plotting these values on the luminosity-temperature
plot (Fig 5) and comparing to the large X-ray sample of Maughan et al. (2008),
we find that A689 is observed to be at the edge of the observed scatter in the
luminosity-temperature relation. This suggests that either the temperature of
the ICM has been enhanced or suppression of the luminosity has occurred.
It has been shown that systems that host a radio source are likely to have
higher temperatures at a given luminosity. Croston et al. (2005) showed that
this is the case for radio-loud galaxy groups, and Magliocchetti & Brüggen
(2007) showed that for clusters that host a radio source there is a departure
from the typical luminosity-temperature relation, particularly in the case of
low mass systems. Croston et al. (2005) also showed, through analysis of
Chandra and XMM-Newton observations, evidence for radio-source interaction
with the surrounding gas for many of the radio-loud groups. A similar process
could be occurring within A689, which has a confirmed radio source at the
center of the cluster. We note that more detailed X-ray and radio observations
would to be needed in order to test any interaction between the BL Lac and
ICM.
The other possible explanation for the offset of radio-loud systems from the
luminosity-temperature relation is the suppression of the luminosity, which
could be caused by displacement of large amounts of gas due to the interaction
of the radio source with the ICM. The interaction of the radio source with the
ICM will cause an increase in the entropy of the local ICM. This higher
entropy gas will be displaced so that it is in entropy equilibrium with the
surrounding gas. However, Magliocchetti & Brüggen (2007) showed that there is
a correlation between the radio luminosity and the heat input required to
produce the observed temperature increment in clusters hosting radio sources.
They note that this correlation favors an enhanced temperature scenario caused
by the radio galaxy induced heating.
### 5.2 Comparison with the BCS
A689 was noted in the BCS as having a significant fraction of its flux coming
from embedded point sources. We have confirmed that A689 contains a point
source at the center of the cluster and is that of a BL Lac object. We stated
that the measured X-ray luminosity for A689 as quoted in the BCS (L0.1-2.4keV
= 3$\times$1045 erg s-1, see $\S$ 1), is the third brightest in the BCS. From
our follow up observation with Chandra we calculate the luminosity and compare
with the BCS value. The luminosities in the BCS are calculated within a
standard radius of 1.43 Mpc, which at the redshift of A689 corresponds to a
radius of 338′′. We therefore employ the same method as in Sect 3.2 and
integrate under a beta model fitted to the derived surface brightness profile
and extrapolate inward and outward from 26-206′′ to 0-338′′. The unabsorbed
flux in the 0.1 - 2.4 keV band (observed frame) was f0.1-2.4,keV =
5.8$\times$10-13 erg s-1 cm-2. After k-correction the X-ray luminosity in the
0.1 - 2.4 keV band (rest frame) was L0.1-2.4,keV = 2.8$\times$1044 erg s-1.
Note that we assume an H0 of 50 for consistency with the BCS catalog. This
value is $\sim$10 times lower than that quoted in the BCS, and A689 is now
ranked 110th out of 201 in luminosity.
### 5.3 Classifying the BL Lac
BL Lac objects may be split into ‘High-energy peak BL Lacs’ (HBL) and ‘Low-
energy peak BL Lacs’ (LBL), for objects which emit most of their synchrotron
power at high (UV–soft-X-ray) or low (far-IR, near-IR) frequencies
respectively (Padovani & Giommi, 1995). HBL and LBL objects have radio-to-X-
ray spectral indies of $\alpha_{rx}$$\leq$0.75 and $\alpha_{rx}$$\geq$0.75,
respectively. We calculate a radio-to-X-ray spectral index for the BL Lac in
A689 of $\alpha_{rx}$=0.58$\pm$0.04. From this value we classify our BL Lac as
an HBL type. We also compared our BL Lac with those of Fossati et al. (1998),
who investigated the properties of large samples of BL lacs at radio to
$\gamma$-ray wavelengths. Our value of $\alpha_{rx}$=0.58 falls into the
region 0.35$\leq$$\alpha_{rx}$$\leq$0.7, dominated by X-ray selected BL Lacs
(XBL). This is consistent with the X-ray selection of this cluster. Using the
measured HST flux, we also calculate $\alpha_{ro}$=0.50 and
$\alpha_{ox}$=0.77. These values are not atypical for BL Lac objects (Figure 7
in Worrall et al., 1999).
Donato et al. (2003) tried to determine whether HBLs and LBLs were
characterized by different environments. They found that of 5 sources
exhibiting diffuse X-ray emission that 4 were HBLs and 1 was an LBL. The BL
Lac in A689 continues this trend, as it appears to be an HBL embedded within a
cluster environment.
### 5.4 Evidence for Inverse-Compton Emission
Our initial analysis of the extended emission in A689 yielded temperature
estimates that were significantly higher than expected based upon its X-ray
luminosity. We also found evidence for an hard excess of X-ray photons in the
6.0-9.0 keV band of the cluster spectrum. Before assigning a physical cause we
must be careful to eliminate systematic effects in the background subtraction,
as underestimating the particle background will give a hard excess. This is
unlikely to be the case here as a high temperature was obtained when either a
blank-sky and local background was used.
When using our background model instead of subtracting a background spectrum,
the cluster temperature was more consistent with its luminosity (Fig 5). This
appears to be due to the unconvolved power law fitting a hard excess in the
cluster region. In order to assess what impact the high-energy particle
component has on the cluster temperature, we varied the normalisation of the
unconvolved power-law component of our background model (sect. 3.2) within its
1$\sigma$ errors. The temperature ranged from 3.86 to 8.51 keV. Thus small
changes in the particle background have a significant influence on the cluster
temperature. Our power-law component was derived within the cluster region in
the 5.0-9.0 keV band, so our background model may be removing a physical hard
X-ray component associated with the cluster. We assess the possible properties
of such a component by re-deriving the normalisation of the unconvolved power-
law component in the local background region, scaling this value by the ratio
of the areas, and using this value for the unconvolved power-law normalisation
in the background model for our cluster fit. We find a power-law normalisation
of 0.0129 photons keV-1 cm-2 s-1 (as opposed to 0.0149 photons keV-1 cm-2 s-1
in the cluster region). Using this value and re-fitting (Fig 15), we find a
cluster temperature of 14.3${}^{+15.6}_{-5.4}$ keV. This suggests that the
hard component is spatially associated with the cluster.
Figure 15: Spectrum of the cluster including our physically motivated model,
with the normalisation of the unconvolved power-law component derived in the
local background region. $\chi^{2}_{\rm\nu}$=1.17 ($\nu$=79)
Bonamente et al. (2007) find a similar excess of X-ray emission in the cluster
A3112, which is known to have a central AGN. It is argued that this excess may
be due to emission of a non-thermal component. Relativistic electrons in the
intergalactic medium will cause CMB photons to scatter into the X-ray band
(inverse Compton scattering). The same process could occur in A689, with the
relativistic particles responsible for the inverse Compton scattering provided
by the jets of the BL Lac.
In order to test the assumption of inverse-Compton emission, we add a
convolved power-law component with $\alpha$=1.5 (appropriate for modeling IC
emission from aged electrons). We fit for the normalisation of this added
power-law component and freeze all other parameters of the background and
cluster model. We note that we use the normalisation of the unconvolved power-
law found in the local background region as found above (a value of 0.0129).
When fit, the $\chi$2 increases slightly but the fit is still acceptable at
the 95% confidence level. From the additional power-law component we measure a
1 keV flux density of $\sim$7 nJy. If the X-ray emission is from scattering of
the CMB by an aged population of electrons of power-law number index 4.0, we
can determine what the implied synchrotron emission would be at 1.4 GHz, given
plausible magnetic fields. Clusters typically have magnetic fields of a few
$\mu$G (Carilli & Taylor, 2002). The NVSS survey would have detected a flux
density of $\sim$10 mJy over the entire cluster. Adopting a 1.4 GHz flux
density limit of 10 mJy, we require a magnetic field of B $\leq$ 2 $\mu$G over
the cluster to avoid over-predicting radio emission. We also calculate a value
for the minimum-energy magnetic field, BE_min, that would give a 1.4 GHz flux
density of S${}_{\rm 1.4~{}GHz}$ = 10 mJy, which in this case is BE_min = 7.5
$\mu$G. For a flux density of S${}_{\rm 1.4~{}GHz}$ = 5 mJy, we would require
a magnetic field of 1.6 $\mu$G and BE_min = 6.5 $\mu$G. Note that B $\propto$
S${}_{\rm 1.4~{}GHz}$1/1+α and BE_min $\propto$ S${}_{\rm 1.4~{}GHz}$1/3+α. We
conclude that the excess X-ray emission can be attributed to inverse-Compton
scattering without over-predicting the radio emission if the magnetic field
strength is in a range typical of clusters and within a factor of a few of the
minimum energy value.
## 6 Conclusions
We have used a 14 ks _Chandra_ observation of the galaxy cluster A689 in order
to determine the nature of the cluster’s point source contamination and to
analyze the cluster properties excluding the central point source. Our main
conclusions are as follows.
* 1\.
Using background subtraction of both local and blank-sky backgrounds, we
obtain temperatures which are high relative to the luminosity-temperature
relation.
* 2\.
We construct a physically motivated model for the background and include this
model in a fit to the cluster spectrum. If the particle background in the
cluster is allowed to exceed that in the local and blank-sky backgrounds we
obtain a temperature of 5.1${}^{+2.2}_{-1.3}$ keV. However, there is no reason
for there to be a higher particle rate in the specific region of the CCD in
which the cluster lies. A hard excess needed to bring the temperature to a
reasonable value must have a different origin.
* 3\.
We confirm the presence of a point source within A689 as suspected in the BCS.
When excluding the point source and using our derive background model we find
a luminosity of L0.1-2.4,keV = 2.8$\times$1044 erg s-1, a value $\sim$10 times
lower than quoted in the BCS.
* 4\.
From the X-ray analysis of the point source we find a “flat-topped” point
source with a pileup fraction of $\simeq$ 60%.
* 5\.
Optical observations of the cluster from SDSS and HST lead us to conclude that
the point source is a BL Lac type AGN.
* 6\.
We classify the BL Lac as an ‘High-energy peak BL Lac’ with
$\alpha$rx=0.58$\pm$0.04.
* 7\.
We interpret the hard X-ray excess needed to bring the cluster temperature to
a reasonable value as inverse-Compton emission from aged electrons that may
have been transported into the cluster from the BL Lac.
We have shown here not only the importance of resolving and excluding point
sources in cluster observations, but also the effect these point sources can
have when determining the ICM properties of galaxy clusters. The detailed
analysis we have performed here may not however be suitable for all clusters
as it is unclear whether this analysis can be performed at higher redshifts.
Separating the point source and cluster emissions becomes increasingly
difficult at high redshifts, however Chandra has proved capable at resolving
point source emission in clusters to z$\sim$2 (e.g. Andreon et al., 2009).
Resolving point source and cluster emission out to high redshift becomes more
important with redshift due to the evolution of the number density of point
sources within clusters (Galametz et al., 2009). The area used to model the
background components associated with the high energy cluster regions will
decrease spatially with redshift, and separating the between thermal and
inverse Compton emissions at higher redshift will become increasingly
difficult (Fabian et al., 2003), due to the increase in the energy density of
the CMB with redshift.
## Acknowledgments
We thank J.Price for useful discussions regarding the SDSS and Hubble data. We
thank Ewan O’Sullivan and Dominique Sluse for useful discussions on the nature
of the point source. We thank the anonymous referee for valuable comments and
suggestions. PG also acknowledges support from the UK Science and Technology
Facilities Council.
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|
arxiv-papers
| 2011-12-06T21:00:04 |
2024-09-04T02:49:25.036485
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "P. A. Giles, B. J. Maughan, M. Birkinshaw, D. M. Worrall, K. Lancaster",
"submitter": "Paul Giles",
"url": "https://arxiv.org/abs/1112.1409"
}
|
1112.1472
|
# Decaying Holographic Dark Energy and Emergence of Friedmann Universe
Titus K Mathew
Department of Physics,
Cochin University of Science and Technology,
Kochi-22, India.
E-mail: tituskmathew@gmail.com, titus@cusat.ac.in
###### Abstract
A universe started in almost de Sitter phase, with time varying holographic
dark energy corresponding to a time varying cosmological term is considered.
The time varying cosmological dark energy and the created matter are
consistent with the Einstein’s equation. The general conservation law for the
decaying dark energy and the created matter is stated. By assuming that the
initial matter were created is in relativistic form, we have analyzed the
possibility of evolving the universe from de Sitter phase to Friedmann
universe.
Keywords: dark energy, Friedmann Universe, cosmology
## 1 Introduction
Recent Astrophysical data have shown that the present universe is undergoing
an accelerated expansion [1]. This shows that the present universe is
dominated by some kind of very smooth form of energy with negative pressure
has been called dark energy, which accounts for about 75$\%$ of the total
energy density of the present universe. Various models have been proposed to
explain this phenomenon, for example there are models based on the dynamics of
scalar or multi-scalar filed, called quintessence models [2]. Another dark
energy candidate is the cosmological constant, which was initially introduced
by Einstein. In the cosmological constant model, the dark energy density,
$\rho_{\Lambda}$, remains constant throughout the entire history of the
universe, while the matter density decreases during the expansion. The
equation of the state for cosmological constant as dark energy is
$w=p/\rho_{\Lambda}=-1$. While in Phantom models [7], it is possible to have
an equation of state with $w<-1.$
An alternative approach to the dark energy problem arises from the holographic
principle. According to the principle of holography the number of degrees of
freedom in a bounded system should be finite and has relations with area of
its boundary. By applying the principle to cosmology, one can obtain the upper
bound of the entropy contained in the universe. For a system with size $L$ and
UV cut-off $\Lambda$, without decaying into a black hole, it is required that
the total energy in a region of size $L$ should not exceed the mass of a black
hole of the same size, thus $L^{3}\rho_{\Lambda}\leq LM_{P}^{2}$ . The largest
L allowed is the one saturating this inequality, thus
$\rho_{\Lambda}=3c^{2}M_{P}^{2}L^{-2}$ (1)
where $c$ is numerical constant having value close to one, we will take it as
one in our analysis. and $M_{P}$ is the reduced Planck Mass $M_{P}^{-2}=8\pi
G.$ When we take the whole universe into account,the vacuum energy related to
this holographic principle can be viewed as dark energy. $L$ can be taken as
the large scale of the universe, for example Hubble horizon, future event
horizon or particle horizon which were discussed by many [5, 6, 7, 8, 9].
In this paper we assume a decaying cosmological term. We also assume that the
universe is started in de-Sitter phase. While in the de-Sitter phase the
universe is completely dominated with the cosmological term. As the universe
expands the dynamical cosmological term decaying in to matter and the universe
will subsequently enter the Friedman phase. As it expands further, the
universe enter a matter dominated phase with decelerated expansion. In section
two we have shown that the decaying cosmological dark energy and created
matter are consistent with the Einstein’s equation. In section 3, we have
obtained the Friedmann equations for the decaying dark energy, and analyzed
the possibility of the evolution of the universe in to the Friedmann phase. We
have also obtained the time evolution of the decaying dark energy and its
equation of state. In section 4, we have presented a comprehensive discussion
of our analysis.
## 2 Dynamical dark energy and horizon
In the presence of cosmological constant the Einstein’s field equation is
$G^{\mu\nu}=R^{\mu\nu}-\frac{1}{2}Rg^{\mu\nu}={8\pi G\over
c^{4}}T^{\mu\nu}_{total}$ (2)
where $G^{\mu\nu}$ is the Einstein tensor, $R^{\mu\nu}$ is Ricci tensor, $R$
is the Ricci scalar (except in this equation, we will refer $R$ as the scale
factor of the expanding universe)and $T^{\mu\nu}_{total}$ is the total energy
momentum tensor comprising matter and cosmological term, and is
$T^{\mu\nu}_{total}=T^{\mu\nu}+\rho_{\Lambda}g^{\mu\nu}$ (3)
in which $T^{\mu\nu}$ is the energy momentum tensor due to matter in perfect
fluid form and $\rho_{\Lambda}$ is the density due to cosmological term, given
as,
$\rho_{\Lambda}={c^{4}\Lambda\over 8\pi G}$ (4)
with $\Lambda$ as the so called “cosmological constant”.
Einstein’s equation satisfies the covarient conservation condition,
$\nabla_{\mu}G^{\mu\nu}=0$ (5)
In the conventional case this implies that, $\nabla_{\mu}T^{\mu\nu}=0.$ As
such this condition doesn’t give any time conserved charge. If the matter is
being created form an independent source, say form the cosmological term, the
conservation law will then take the general form
$\nabla_{\mu}\left(T^{\mu\nu}+\rho_{\Lambda}g^{\mu\nu}\right)=0$ (6)
This conservation law implies that the energy and momentum of matter alone is
not conserved, bur energy and momentum of matter and cosmological term or dark
energy are together be conserved. This general conservation law allows the
exchange of energy and momentum between matter and dark energy and acting as a
controlling condition for this exchange. The existing theories predicts a very
large value for the cosmological term [3, 4] in the early stage of the
universe, but the present observations points towards a very low value for the
cosmological term for the late universe. In this light it is inevitable to
consider that, there must be a transference of energy form the dark energy or
cosmological term sector to the matter sector.
Let us assume that, the term $\Lambda$ correspondingly $\rho_{\Lambda}$ is a
function of time, since a space dependent $\Lambda$ will lead to an
anisotropic universe. The covarient conservation law will then give the
equations,
$\nabla_{\mu}T^{\mu i}=0$ (7)
and
$\nabla_{\mu}T^{\mu 0}=-{c^{3}\over 8\pi G}{d\Lambda\over dt}$ (8)
where $i=1,2,3$ for the spatial part and $i=0$ for the time part.
In reference [10] authors have considered the energy transference between
decaying cosmological term and matter. It is important to realise that the
covarient conservation law given above is drastically different from that
appearing in in some quintessence model [11, 12, 13], where energy-momentum
tensor of the scalar field that replaces the cosmological term is itself
covariently conserved, but no matter creation. In the present paper we have
considered that the cosmological term decaying into matter which is consistent
with the Friedmann model of the universe.
The energy density $\rho_{\Lambda}$ corresponds to the time varying
cosmological term is taken as the holographic dark energy as defined in
equation (1). A simple holographic dark energy model is by taking $L=H^{-1}$,
where $H$ is the Hubble’s constant is considered by Hsu et al [5] and they
have shown that the Friedmann model with $\rho_{\Lambda}=3c^{2}M^{2}_{p}H^{2}$
makes the dark energy behave like ordinary matter rather than a negative
pressure fluid, and prohibits accelerating expansion of the universe. We have
adopt an equation for holographic dark energy energy, where the future event
horizon $(R_{h})$ is used instead of the Hubble horizon as the IR cut-off L,
which was shown to lead an accelerating universe by Li [14]. Thus the time
varying cosmological energy density is
$\rho_{\Lambda}=3c^{2}M^{2}_{p}R_{h}^{-2}$ (9)
where $c$ is a constant have values $O(1)$ and the event horizon $R_{h}(t)$, a
function of cosmological time, is given by
$R_{h}(t)=R(t)\int^{\infty}_{t}{dR(t^{{}^{\prime}})\over
H(t^{{}^{\prime}})R(t^{{}^{\prime}})^{2}}$ (10)
where $R(t)$ is the expansion factor and $H(t)$ is the Hubble constant.
## 3 Cosmic evolution of dark energy and Friedmann Universe
Let us consider an empty universe in de Sitter phase, with very large [15]
decaying cosmological term, corresponds to dark energy density as given by
equation (LABEL:eq:rho12). If the matter and energy are created from the
decaying dark energy term are homogeneous and isotropic, then the geometry of
the universe can be that of Friedmann-Robertson-Walker form,
$ds^{2}=c^{2}dt^{2}-R^{2}\left[{dr^{2}\over
1-kr^{2}}+r^{2}\left(d\theta^{2}+\sin^{2}\theta\,d\phi^{2}\right)\right]$ (11)
where $k$ is the curvature parameter $R$ is the scale factor of expansion, $t$
is the cosmological time and $(r,\theta,\phi)$ are the co-moving coordinates.
By taking the energy momentum tensor of matter as
$T^{\mu}_{\nu}=(\rho+p)u^{\mu}u_{\nu}-p\delta^{\mu}_{\nu}$ (12)
where $\rho$ is the sum of energy densities of the created components of due
to the decay of cosmological term and $p$ is the pressure of the matter
components. Under these conditions, the covarient conservation law (8) leads
to (here we consider only one component of matter)
${d\rho_{m}\over dt}+3H\left(\rho_{m}+p_{m}\right)=-{d\rho_{\Lambda}\over dt}$
(13)
where $H={dR\over dt}/R$ the Hubble parameter, $\rho_{m}$ is the density of
the created matter and $p_{m}$ is its pressure. This equation obtained form
the general conservation law is found to followed from the combinations the
standard Friedmann equations,
$\left({dR\over dt}\right)^{2}={8\pi G\over
3c^{2}}\left(\rho_{m}+\rho_{\Lambda}\right)R^{2}-kR^{2}$ (14)
and
${d^{2}R\over dt^{2}}={8\pi G\over
3c^{2}}\left(\rho_{\Lambda}-\frac{1}{2}\left(\rho_{m}+3p_{m}\right)\right)R$
(15)
provided the cosmological term $\rho_{\Lambda}$ is time dependent. If one
assumes ordinary pressureless matter as
$\rho_{m}=\rho_{m0}R^{-3}$ (16)
where $\rho_{m0}$ is the present density of matter. then equation (13) will
lead to the result that, the cosmological term will be independent of time. On
the other hand this shows that the time dependent cosmological term does not
decay in to pressureless matter.
Let us assume that the cosmological term can possibly decay into some form of
matter with equation of state $p_{m}=\omega_{m}\rho_{m},$ where the parameter
$\omega_{m}$ is assumed to be in the range $0\leq\omega_{m}\leq 1,$ the exact
value is depends on particular matter component. In this paper we are
considering only one component of matter. The covarient conservation law (8)
can now be written for the possibility of cosmological term decaying into
matter as
${d\rho_{m}\over
dt}+3H\left(1+\omega_{m}\right)\rho_{m}=-{d\rho_{\Lambda}\over dt}$ (17)
Since density behaviour of ordinary pressureless matter does not work for a
varying cosmological dark energy, we will assume the form for $\rho_{m}$ which
is slightly different from its canonical form, as [16, 17]
$\rho_{m}=\rho_{m0}R^{-3+\delta}$ (18)
where $\rho_{m0}$ is the present value of $\rho_{m}$ and $\delta$ is a
parameter which is effectively depends upon the state of the universe. From
equation (18) and (LABEL:eq:rho12), equation (17) become,
$\rho_{\Lambda}\left({HR_{h}-1\over
HR_{h}}\right)=\left({3\omega_{m}+\delta\over 2}\right)\rho_{m}$ (19)
where we have assumed a vanishing integration constant and also with the
condition that lim t$\rightarrow\infty$, $R(t)\rightarrow\infty$ and equation
for time rate of horizon can be cast into the differential form
${dR_{h}\over dt}=HR_{h}-1.$ (20)
The equation (19) suggest that depending on the parameters $\delta$ and
$\omega_{m}$ the energy densities $\rho_{m}$ and $\rho_{\Lambda}$ may
eventually be of the same order, as suggested by the present observations
[18]. This relation rather suggest the relation between dark energy and
matter, when $\delta=3$, the case corresponds to a constant matter density at
which there exist a equilibrium between matter creation and universe expansion
In general event horizon is not existing for Friedmann universe. But for de
Sitter universe there exists event horizon, which satisfies the relation
$R_{h}\sim H^{-1}.$ Consequently for de Sitter universe, $HR_{h}\sim 1$ which
implies that in de Sitter phase the energy density is almost completely
dominated by the cosmological term or dark energy. For Friedmann universe we
will hence take $HR_{h}$ as very large In general we will take the value of
$HR_{h}$ is equal to one or large.
### 3.1 Friedmann Universe
Friedmann universe is a homogeneous and istropic universe, satisfying the
conditions (14) and (15). With the relation between decaying cosmological term
and created matter (19), the second Firedmann equation become,
${d^{2}R\over dt^{2}}={(\left(1+3\omega_{m}\right)\beta^{2}\over
2}\left[{3\omega_{m}+\delta\over 1+3\omega_{m}}\left({HR_{h}\over
HR_{h}-1}\right)-1\right]R^{\delta-2}$ (21)
where $\beta^{2}={8\pi G\rho_{mo}/3c^{2}}$ . For de Sitter phase the
acceleration is very large, for which $HR_{h}\sim 1.$ As it enters the
Friedmann phase by the decay of cosmological or the dark energy, the
acceleration can be negative or positive, depends on the value of the term in
the parenthesis of the right hand side of the above equation. The condition
for acceleration is,
${3\omega_{m}+\delta\over 3\omega_{m}+1}>1-{1\over HR_{h}}$ (22)
The factor $\displaystyle{1-{1\over HR_{h}}}$ is in the range $\displaystyle
0\leq{1-{1\over HR_{h}}}\leq 1.$ The extreme limits are corresponds to de
Sitter phase and matter dominated universes respectively. For the transition
period from de Sitter phase to Friedmann phase, $HR_{h}$ is near to one, and
assuming that the created matter have the equation state
$\omega_{m}=\frac{1}{3}$, where the created matter is in relativistic form
then
$\delta>2\alpha-1$ (23)
where $\displaystyle\alpha=1-{1\over HR_{h}}$ having value less than one. This
implies that during the period of decay of the cosmological dark energy term
the density of created matter is diluted slowly as the universe expand, than
the decreasing of the density of non-relativistic matter in the Friedmann
universe. This shows that even in the friedmann universe it is possible to
have an initial accelerating phase, where the cosmological dark energy is
start its decay into matter and is still dominating over matter. As universe
proceeds, the created matter will subsequently dominate and hence the universe
will come to a matter dominated phase, at which the universe is expanding with
deceleration.
For decelerating expansion, where matter is dominating over the cosmological
term, the condition in the limit where $HR_{h}$ is very large is,
$\delta<1$ (24)
This condition is true irrespective of whether the created matter is
relativistic or non-relativistic. However as the universe enter the
decelerating phase, the matter will become non-relativistic, satisfying the
extreme condition that $\delta\rightarrow 0$ so that $\rho_{m}\sim R^{-3}$
### 3.2 Flat Universe
For flat universe, where the curvature parameter $k=0,$ the Friedmann equation
(LABEL:eq:frw1) become
$\left({dR\over dt}\right)^{2}=\beta^{2}\left({3\omega_{m}+\delta\over
2}{HR_{h}\over HR_{h}-1}+1\right)R^{\delta-1}$ (25)
On integration and avoiding the integration constant, the solution would be,
$R\sim t^{2\over 3-\delta}$ (26)
By considering the relation for matter creating out of decaying dark energy,
i.e. $\rho_{m}=\rho_{m0}R^{-3+\delta},$ then the relation between dark energy
density and matter will have beheaviour
$\rho_{m}\sim\rho_{\Lambda}\sim t^{-2}$ (27)
This is the time dependence of $\rho_{\Lambda}$ for any value of $\omega_{m}$
and $\delta.$ This time dependence shows that $\rho_{\Lambda}$ diverge at the
initial time, which implies the existence of initial singularity.
### 3.3 An equation of state for the decaying dark energy
An equation of state for the time decaying cosmological term can be written as
[19]
$\omega_{\Lambda}^{eff}=-1-{1\over 3}\left({dln\rho_{\Lambda}\over d\ln
R}\right)$ (28)
With the equation for the relation between decaying dark energy and creating
matter [fang], the equation of state become
$\omega_{\Lambda}^{eff}=-{\delta\over 3}-{1\over 3\left(HR_{h}-1\right)^{2}}$
(29)
This shows that for large values of $HR_{h}$ the equation of sate become
$\omega_{\Lambda}=-{\delta\over 3}$. From the above analysis, it is seen that
for an accelerating universe $\omega_{\Lambda}$ is less than $-{1\over 3},$
but for a decelerating Friedmann universe it is greater than $-{1\over 3},$
which is similar to the latest analysis by many.
## 4 Discussion
In the presence of a time varying cosmological term, assumed to be of
holographic dark energy form, it is possible that the universe may starts with
the de Sitter phase, exhibiting horizon and where the energy density is
completely dominated by the dark energy. The decaying holographic dark energy
cause the primordial inflation. If the horizon $R_{h}$ is assumed to be equal
to the plank length at very early stage, then $\Lambda$ would have a value of
the order
$\Lambda\sim 10^{66}\,cm^{-2}$ (30)
the corresponding dark energy density would be $\rho_{\Lambda}\sim
10^{112}\,ergcm^{-3}.$ This enormous dark energy decay into matter as the
universe is evolved to the Friedmann phase, and the dark energy reached the
present value
$\rho_{\Lambda}^{0}=10^{-8}\,erg\,cm^{-3}$ (31)
The evolution of the de Sitter phase in to the Friedmann universe is in such a
way that the total energy density comprising the dark energy, created matter
and the gravitational field together be conserved. In this paper we have
considered the decay of the dark energy into matter. During the initial phase
of decay, the universe might be in the accelerating phase, where the parameter
$\delta$ characterizing the equation $\rho_{m}\sim R^{-3+\delta},$ is greater
than one. This implies that the dilution in the density of created matter in
slower compared to the non-relativistic matter. As the universe proceeds with
expansion, the matter created will come to dominate, and the universe
eventually go over to matter dominated phase, with the condition $\delta<1.$
In the extreme limit this condition may go the limit $\delta\rightarrow 0$,
which emphasises that, the created matted will eventually become non-
relativistic, with behaviour, $\rho_{m}\sim R^{-3}.$ In section 3.3, the
equation of the state of the decaying holographic dark energy, shows that, in
the early phase of universe too, $\omega_{\Lambda}<-{1\over 3}$ as it’s
equation of state in late accelerating universe [20]. To explain why
$\rho_{\Lambda}$ and $\rho_{m}$ are of the same order today, it is essential
to have a specific time evolution for dark energy. We argued that a dynamical
dark energy, endowed with an appropriate time evolution can contain the
possibility of the development of a Friedmann universe from a de Sitter
universe. As the de Sitter phase evolved in to the Friedmann universe, the
value of the dark energy is decreased gradually to a low value, which
eventually lead to matter dominating phase with decelerating expansion. But on
the other hand the recent observations indicating that the present universe is
in a accelerating expansion. This fact indicating the possibility that at
present the dark energy is increasing at the expense of decaying matter. This
time increasing dark energy, in other words, implies that the universe may
evolving in to stage where the whole energy density is coming to dominate
completely by the dark energy. The ultimate clarity regarding these, of course
be given by the proper quantum effects, which is still an open question.
## References
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D 72 (2005) 103503 [SPIRES]
* [2] Ratra B and Peebles P J E, Cosmological Consequences of a Rolling Homogeneous Scalar field Phys. Rev. D 37 (1998) 3406 [SPIRES]
Zlatev I, Wang L and Steinhardt P J, Quintessence, Cosmic Coincidence and
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* [3] Weinberg S, Rev. Mod. Phys. The Cosmological Constant Problem, 61 (1989) 1
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* [6] M. Li, A model of holographic dark energy, Phys. Lett. B 603 (2004) 1 [hep-th/0403127] [SPIRES];
* [7] Q.-G. Huang and M. Li, Anthropic principle favors the holographic dark energy, JCAP 03 (2005) 001 [hep-th/0410095] [SPIRES];
Q. -G. Huang and M. Li, The holographic dark energy in a non-flat universe,
JCAP 08 (2004) 013 [astro-ph/0404229] [SPIRES];
B. Chen, M. Li and Y. Wang, Inflation with Holographic Dark Energy, Nucl. Phys.
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J.-f. Zhang, X. Zhang and H.-y. Liu, Holographic dark energy in a cyclic
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S. Nojiri and S.D. Odintsov, Unifying phantom inflation with late-time
acceleration: Scalar phantom-non-phantom transition model and generalized
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[SPIRES].
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* [9] R. Horvat, Holography and variable cosmological constant, Phys. Rev. D 70 (2004) 087301 [astro-ph/0404204] [SPIRES].
* [10] R.Aldrovandi, J.P.Beltran Almeida and J.G.Pereira, Time-Varying Cosmological Term: Emergence and Fate of a FRW Universe, Grav. Cosmol. 11 (2005) 277, arXiV:gr-qc/0312017v3 12 Apr 2005.
* [11] B.Ratra and P.J.E.Peebles, Cosmological Consequences of a Rolling Homogeneous Scalar field Phys. Rev. D 37 (1998) 3406 [SPIRES]
* [12] J.Friman, C.T.Hill and R.Watkins, Late Time Cosmological Phase Transitions,1. Particle Physics Models and Cosmic Evolution Phys. Rev D 46 (1992) 1226
* [13] V.Sahni and A.A.Starobisky, The Case for a Positive Cosmological Lambda-term Int. J. Mod. Phys. D9 (2000) 373.
* [14] Huang Q-G and Li M, Holographic Dark Energy in Non-Flat Universe 2004 J. Cosmol. Astropart. Phys. JCAP 0408 (2004) 013 [SPIRES]
* [15] B.Mashhoon and P.S.Wesson, Gauge Dependent Cosmological Constant Class. Quant. Grav. 21, (2004) 3611
* [16] B.Guberina, R.Horvat and Hrvoje Nikolic, Non-Structured Holographic Dark Energy, JCAP 0701 (2007) 012
* [17] B.Guberina, R.Horvat and Hrvoje Nikolic, Dynamical Dark Energy with a constant vacuum energy density, Phys. Lett.B 636: (2006) 80-85
* [18] S M Carrol, why is the universe Accelerating, in - Measuring and modelling the universe. ed. by W.L.Freeman, Cambridge University Press, Cambridge, 2003 [astro-ph/0310342]
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|
arxiv-papers
| 2011-12-07T05:18:22 |
2024-09-04T02:49:25.047143
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Titus K. Mathew",
"submitter": "Titus K Mathew",
"url": "https://arxiv.org/abs/1112.1472"
}
|
1112.1482
|
# Omnidirectionally Bending to the Normal in $\epsilon$-near-Zero
Metamaterials
Simin Feng simin.feng@navy.mil Michelson Lab, Physics Division, Naval Air
Warfare Center, China Lake, California 93555
###### Abstract
Contrary to conventional wisdom that light bends away from the normal at the
interface when it passes from high to low refractive index media, here we
demonstrate an exotic phenomenon that the direction of electromagnetic power
bends towards the normal when light is incident from arbitrary high refractive
index medium to $\epsilon$-near-zero metamaterial. Moreover, the direction of
the transmitted beam is close to the normal for all angles of incidence. In
other words, the electromagnetic power coming from different directions in air
or arbitrary high refractive index medium can be redirected to the direction
almost parallel to the normal upon entering the $\epsilon$-near-zero
metamaterial. This phenomenon is counterintuitive to the behavior described by
conventional Snell’s law and resulted from the interplay between
$\epsilon$-near-zero and material loss. This property has potential
applications in communications to increase acceptance angle and energy
delivery without using optical lenses and mechanical gimbals.
###### pacs:
42.25.Bs, 42.79.Wc, 78.67.Pt
Bending of light towards the normal when it passes from low to high refractive
index media is one of the fundamental phenomena in optics. As a manifestation
of this phenomenon, directive emission into air by a source inside the
material with vanishingly small permittivity, known as $\epsilon$-near-zero
(ENZ) metamaterials, has been demonstrated Enoch . With other intriguing
properties, such as ultrathin waveguides Silveir ; Liu ; Edwards ; Adams ;
Alu1 , diffraction-suppressed propagation and self-collimation Feng ; Mocella1
; Mocella2 ; Polles , the ENZ materials have gained prominence as useful
components to tailor antenna radiations Alu2 ; Halterman ; Saenz ; Mart .
Previous studies on ENZ-directive emission have been focused on the radiation
from low ($\epsilon\approx 0$) to high (air) refractive index media Enoch ;
Ziolkow ; Yuan ; Jin ; Lovat , where the directive transmission can be
intuitively understood from Snell’s law that dictates the light bending
towards the normal as it passes from low to high refractive index media. From
the reciprocal theorem, for the radiation from high to low refractive index
materials, the transmitted beam should spread out into grazing angles as the
result of bending away from the normal. Contrary to this conventional
behavior, in this paper we will demonstrate an exotic phenomenon that
electromagnetic (EM) power can bend towards the normal when light passes from
arbitrary high ($\epsilon_{1}\gg 1$) to low ($\epsilon_{2}\approx 0$)
refractive index media as shown in Fig. 1a. Furthermore, the direction of the
transmitted beam is close to the normal for all angles of incidence. More
interestingly, this counterintuitive to conventional Snell’s law behavior is
induced by material loss. The interplay between ENZ and loss leads to unusual
wave interaction. This phenomenon can be used to project EM power coming from
different directions to one direction to the receivers as shown in Fig. 1b,
where the incoming waves all bend to the normal pointing to the receptors upon
entering the ENZ medium. A plasmonic thin film is superimposed on the ENZ
material to enhance the transmission through structural resonances. For all
the incoming directions including grazing angles, the transmitted powers
impinge normally to the receptors or photocells embedded in the ENZ medium,
effectively increasing the acceptance angle and energy transfer.
Figure 1: (Color online) (a) A plane wave is incident from arbitrary high
permittivity ($\epsilon_{1}\gg 1$) medium to ENZ ($\epsilon_{2}\approx 0$)
metamaterial. (b) Incoming waves in air from different directions all bend to
the normal upon entering the ENZ medium. A nanoplasmonic thin film is
superimposed on the ENZ material to enhance the transmission. Receptors or
photocells are embedded in the ENZ metamaterial.
Our derivation is based on anisotropic media. The results can be applied to
isotropic materials. Assuming a harmonic time dependence $\exp(-i\omega t)$
for the EM field, from Maxwell’s equations, we have
$\displaystyle\begin{split}\nabla\times\bigl{(}\bar{\bar{\mu}}_{n}^{-1}\cdot\nabla\times{\bm{E}}\bigr{)}&=\,k_{0}^{2}\bigl{(}\bar{\bar{\epsilon}}_{n}\cdot{\bm{E}}\bigr{)}\,,\\\
\nabla\times\bigl{(}\bar{\bar{\epsilon}}_{n}^{-1}\cdot\nabla\times{\bm{H}}\bigr{)}&=\,k_{0}^{2}\bigl{(}\bar{\bar{\mu}}_{n}\cdot{\bm{H}}\bigr{)}\,,\end{split}$
(1)
where $k_{0}=\omega/c$; and the $\bar{\bar{\epsilon}}_{n}$ and
$\bar{\bar{\mu}}_{n}$ are, respectively, the permittivity and permeability
tensors for each uniform region ($n=1,2,\cdots$), which in the principal
coordinates can be described by
$\displaystyle\bar{\bar{\epsilon}}_{n}$ $\displaystyle=$
$\displaystyle\epsilon_{nx}\hat{\bm{x}}\hat{\bm{x}}+\epsilon_{ny}\hat{\bm{y}}\hat{\bm{y}}+\epsilon_{nz}\hat{\bm{z}}\hat{\bm{z}}\,,$
(2) $\displaystyle\bar{\bar{\mu}}_{n}$ $\displaystyle=$
$\displaystyle\mu_{nx}\hat{\bm{x}}\hat{\bm{x}}+\mu_{ny}\hat{\bm{y}}\hat{\bm{y}}+\mu_{nz}\hat{\bm{z}}\hat{\bm{z}}\,.$
(3)
Consider transverse magnetic (TM) modes, corresponding to non-zero field
components $H_{y}$, $E_{x}$, and $E_{z}$. The magnetic field $H_{y}$ satisfies
the following wave equation:
$\frac{1}{\epsilon_{z}}\frac{\partial^{2}H_{y}}{\partial
x^{2}}+\frac{1}{\epsilon_{x}}\frac{\partial^{2}H_{y}}{\partial
z^{2}}+k_{0}^{2}\mu_{y}H_{y}=0\,,$ (4)
which permits solutions of the form $\psi(z)\exp(i\beta x)$. Here the
transverse wave number $\beta$ is determined by the incident wave, and is
conserved across the interface of different regions,
$\beta^{2}=k_{0}^{2}\epsilon_{nz}\mu_{ny}-\alpha_{n}^{2}\frac{\epsilon_{nz}}{\epsilon_{nx}}\,,\hskip
14.45377pt(n=1,2,\cdots)\,,$ (5)
where $\alpha_{n}$ is the wave number in the $z$ direction. The functional
form of $\psi(z)$ is either a simple exponential $\exp(i\alpha_{n}z)$ for the
semi-infinite regions or a superposition of $\cos(\alpha_{n}z)$ and
$\sin(\alpha_{n}z)$ terms for the bounded regions along the $z$ direction. The
other two components $E_{x}$ and $E_{z}$ can be solved from $H_{y}$ using
Maxwell’s equations. By matching boundary conditions at the interfaces, i.e.,
the continuity of $H_{y}$ and $E_{x}$, the electromagnetic field can be
derived in each region; and then the Poynting vector $\bm{S}$ can be computed
from $\bm{S}=\Re(\bm{E}\times\bm{H}^{*})$. In anisotropic materials, the
direction of the Poynting vector is different from that of the phase front of
the field. Here, only the direction of the Poynting vector is considered since
it is associated with the energy transport. The angle ($\theta_{S}$) of the
Poynting vector is measured from the Poynting vector to the surface normal,
and is given by $\theta_{S}=\tan^{-1}(S_{x}/S_{z})$. In Fig. 1a, the input
medium is isotropic material with permittivity $\epsilon_{1}$; the output
medium is ENZ material ($\epsilon_{2}\approx 0$). In the following, both
anisotropic ($\epsilon_{2x}\neq\epsilon_{2z}$) and isotropic
($\epsilon_{2x}=\epsilon_{2z}$) ENZ materials will be considered. Figure 2
illustrates the effect of loss of the ENZ-materials on the transmission angle
(TA), which is plotted against angle of incidence (AOI) with and without loss
for the different permittivity ($\epsilon_{1}$) of the input medium. In the
top panels, when the loss is zero $\bigl{(}\Im(\epsilon_{2z})=0\bigr{)}$, the
transmission angle is the grazing angle $90^{\circ}$ except for the normal
incidence. This behavior is complied with conventional Snell’s law. In the
bottom panels, with a moderate loss $\Im(\epsilon_{2z})=0.6$, the transmission
angle switches to near zero (normal direction) for all angles of incidence,
leading to collimated transmission in the normal direction. This switching
phenomenon persists even for the much higher permittivity ($\epsilon_{1}=100$)
of the input medium (middle and right panels).
Figure 2: (Color online) Transmission angle of the Poynting vector versus AOI.
Left and middle panels: anisotropic ENZ material with $\epsilon_{2x}=1$ and
$\epsilon_{2z}=0.001+i\epsilon_{2z}^{i}$. Right panels: isotropic ENZ material
with $\epsilon_{2x}=\epsilon_{2z}=0.001+i\epsilon_{2z}^{i}$. Top panels:
$\epsilon_{2z}^{i}=0$. Bottom panels: $\epsilon_{2z}^{i}=0.6$. Left panels:
$\epsilon_{1}=1$. Middle and right panels: $\epsilon_{1}=100$. A good
agreement between the numerical (blue-solid) and analytical (green-circles)
results. The material loss switches the TA from the grazing angle $90^{\circ}$
(top panels) to the near-zero angle (bottom panels) for all the AOI.
To understand this loss-induced switching behavior, let’s analyze the
transmission angle ($\theta_{S}$), which is given by
$\tan(\theta_{S})=\frac{S_{x}}{S_{z}}=\frac{\Re\left(\dfrac{\bar{\beta}}{\epsilon_{2z}}\right)}{\Re\left[\sqrt{\dfrac{\mu_{2y}}{\epsilon_{2x}}-\dfrac{\bigl{(}\bar{\beta}\bigr{)}^{2}}{\epsilon_{2x}\epsilon_{2z}}}\right]}\,,$
(6)
where $\bar{\beta}\equiv\beta/k_{0}$, and $\bar{\beta}$ (real) is determined
by the incidence angle. The transmission angle of the Poynting vector depends
only on the input and output media. In the case of $\epsilon_{2x}\rightarrow
0$ and $\epsilon_{2z}$ finite, Eq. (6) indicates $\theta_{S}\rightarrow
0^{\circ}$ (normal direction). For the case of $\epsilon_{2z}\rightarrow 0$
and $\epsilon_{2x}$ finite and the case of isotropic ENZ material with
$\epsilon_{2x}=\epsilon_{2z}\rightarrow 0$, the analysis is more involved. The
numerator of Eq. (6) can be written as
$\Re\left(\frac{\bar{\beta}}{\epsilon_{2z}}\right)=\frac{\bar{\beta}\,\epsilon_{2z}^{r}}{|\epsilon_{2z}|^{2}}\,,$
(7)
where $\epsilon_{2z}^{r}\equiv\Re(\epsilon_{2z})$. Assuming $\mu_{2y}$ is
real, the denominator of Eq. (6) becomes
$\Re\left[\sqrt{\dfrac{\mu_{2y}}{\epsilon_{2x}}-\dfrac{\bigl{(}\bar{\beta}\bigr{)}^{2}}{\epsilon_{2x}\epsilon_{2z}}}\right]=\frac{a\,\bar{\beta}}{|\epsilon_{2x}\epsilon_{2z}|}\,,$
(8)
where
$a^{2}=\frac{1}{2}\left(A\epsilon_{2x}^{r}+B|\epsilon_{2x}|-\epsilon_{2x}^{r}\epsilon_{2z}^{r}+\epsilon_{2x}^{i}\epsilon_{2z}^{i}\right)\,,$
(9)
where $\epsilon_{2z}^{i}\equiv\Im(\epsilon_{2z})$,
$\epsilon_{2x}^{r}\equiv\Re(\epsilon_{2x})$,
$\epsilon_{2x}^{i}\equiv\Im(\epsilon_{2x})$, and
$A\equiv\frac{|\epsilon_{2z}|^{2}\mu_{2y}}{\bigl{(}\bar{\beta}\bigr{)}^{2}}\,,\hskip
8.67204ptB=\sqrt{|\epsilon_{2z}|^{2}-2A\epsilon_{2z}^{r}+A^{2}}\,.$ (10)
Thus, the transmission angle ($\theta_{S}$) becomes
$\tan(\theta_{S})=\frac{|\epsilon_{2x}|\epsilon_{2z}^{r}}{a\,|\epsilon_{2z}|}\,.$
(11)
The loss-induced angle switching observed in Fig. 2 can be explained from Eq.
(11). For the anisotropic material with $\epsilon_{2x}\neq\epsilon_{2z}$ and
finite $\epsilon_{2x}$, if $\epsilon_{2z}^{i}=0$, when
$\epsilon_{2z}^{r}\rightarrow 0$,
$\epsilon_{2z}^{r}/|\epsilon_{2z}|\rightarrow 1$ and $a\rightarrow 0$, thus
$\theta_{S}\rightarrow 90^{\circ}$. If $\epsilon_{2z}^{i}\neq 0$, when
$\epsilon_{2z}^{r}\rightarrow 0$,
$\epsilon_{2z}^{r}/|\epsilon_{2z}|\rightarrow 0$ and $a$ is finite, thus
$\theta_{S}\rightarrow 0^{\circ}$. On the other hand, if $\epsilon_{2z}$ is
finite, when $\epsilon_{2x}\rightarrow 0$, $a\rightarrow\sqrt{\epsilon_{2x}}$,
thus $\theta_{S}\rightarrow 0^{\circ}$. For the isotropic case, let
$\epsilon_{2x}=\epsilon_{2z}\equiv\epsilon_{2}^{r}+i\epsilon_{2}^{i}$. If
$\epsilon_{2}^{i}=0$, when $\epsilon_{2}^{r}\rightarrow 0$,
$\epsilon_{2z}^{r}/|\epsilon_{2z}|\rightarrow 1$ and
$a\rightarrow(\epsilon_{2}^{r})^{3/2}$, thus $\theta_{S}\rightarrow
90^{\circ}$. If $\epsilon_{2}^{i}\neq 0$, when $\epsilon_{2}^{r}\rightarrow
0$, $\epsilon_{2z}^{r}/|\epsilon_{2z}|\rightarrow 0$ and $a$ is finite,
therefore $\theta_{S}\rightarrow 0^{\circ}$. To validate Eq. (11), in Fig. 2
the TAs calculated from Eq. (11) (green-circles) are compared to those
computed numerically (blue-solid), showing a perfect agreement.
To validate the loss-induced switching behavior is a robust feature, in Fig. 3
the transmission angle versus AOI is plotted for the different real parts of
$\epsilon_{2z}$ and $\epsilon_{2x}$ and the material loss. In essence, the
transmission angle decreases with increasing the loss $\Im(\epsilon_{2z})$ and
decreasing the $\Re(\epsilon_{2z})$. When $\Re(\epsilon_{2z})\rightarrow 0$,
the angular width of the transmission can be estimated from
$\Delta\theta_{S}\approx\left\\{\begin{matrix}\dfrac{\sqrt{2}\,\,|\epsilon_{x}|\,\epsilon_{z}^{r}}{|\epsilon_{z}|^{3/2}\sqrt{|\epsilon_{x}|+\epsilon_{x}^{i}+\eta\,\epsilon_{x}^{r}}}\,,&\mbox{if
}\eta\leq 1\\\ \\\
\dfrac{\sqrt{2}\,\,|\epsilon_{x}|\,\epsilon_{z}^{r}}{|\epsilon_{z}|^{3/2}\sqrt{\epsilon_{x}^{i}+\eta\bigl{(}|\epsilon_{x}|+\epsilon_{x}^{r}\bigr{)}}}\,,&\mbox{if
}\eta\geq 1\end{matrix}\right.\,,$ (12)
where $\eta\equiv\dfrac{|\epsilon_{z}|\mu_{y}}{\epsilon_{1}\mu_{1}}$, and the
subscript $2$ in the $\epsilon_{x}$, $\epsilon_{z}$, and $\mu_{y}$ was omitted
in above equation.
Figure 3: (Color online) Transmission angle versus AOI when the
$\Re(\epsilon_{2z})=0.001$ (blue-solid) and $\Re(\epsilon_{2z})=0.01$ (green-
dashed). Top panels: $\Im(\epsilon_{2z})=1.5$. Bottom panels:
$\Im(\epsilon_{2z})=3.2$. Left and middle panels: anisotropic ENZ material
with $\epsilon_{2x}=1.0$ (left panels) and $\epsilon_{2x}=2.0$ (middle
panels). Right panels: isotropic ENZ material. The permittivity of the input
medium $\epsilon_{1}=36$. Figure 4: (Color online) Transmission angle versus
material loss $\Im(\epsilon_{2z})$ computed for AOI $=0.1^{\circ}$ (blue-
solid) and AOI $=89^{\circ}$ (green-dashed). Top panels:
$\Re(\epsilon_{2z})=0.001$. Bottom panels: $\Re(\epsilon_{2z})=0.01$. Left and
middle panels: anisotropic ENZ material with $\epsilon_{2x}=1.0$ (left panels)
and $\epsilon_{2x}=2.0$ (middle panels). Right panels: isotropic ENZ material.
The permittivity of the input medium $\epsilon_{1}=36$. TA quickly converges
to zero in all the scenarios.
Figure 4 demonstrates how rapidly the transmission angle converges to zero as
the loss $\Im(\epsilon_{2z})$ increases for the different values of
$\Re(\epsilon_{2z})$ and $\epsilon_{2x}$. The blue-solid curves represent the
transmission angles calculated for the near-zero angle of incidence, while the
green-dashed curves for the grazing angle of incidence. The difference between
the green-dashed and blue-solid curves corresponds to the angular width of the
transmission. The angular width in the isotropic ENZ medium (right panels) is
usually smaller than that in the anisotropic medium (left and middle panels).
This is implicated in Eq. (12) as well. It is well-known that loss is
inextricable to metamaterial. Many fascinating effects diminish as the result
of high loss Dimmock ; Nieto . However, for the effect demonstrated here, the
material loss plays a positive role, resulting in the omnidirectional bending
of light towards the normal upon entering the ENZ medium.
This phenomenon may have many applications, such as directive antennas.
Instead of radiation applications, we will explore this phenomenon from a
receiving perspective, i.e., redirect the incoming EM power from different
directions to the direction of the receivers to enhance the acquisition power,
as illustrated in Fig. 1b. To increase the coupling, a matching coating can be
deposited on the surface of the ENZ medium such that the effective impedance
of the overall structure is matched to the free-space impedance. For
simplicity, here a dielectric-metal-dielectric thin film is superimposed on
the ENZ material. This sandwich structure possesses coupled surface plasmon
modes due to closely spaced two dielectric-metal interfaces. By exciting the
plasmonic resonances of the structure, the transmission can be enhanced. The
resonant frequency of the transmission can be tuned by changing the thickness
of the layers. In our simulation, the materials of the dielectric and metallic
layers are, respectively, amorphous polycarbonate (APC) and silver (Ag). The
refractive index of the APC is given by Roberts
$n_{p}=1.5567+\frac{8.0797\times 10^{-3}}{\lambda^{2}}+\frac{3.5971\times
10^{-4}}{\lambda^{4}}\,,$ (13)
where $\lambda$ is the wavelength in $\mu m$. The loss of the APC is very
small in the wavelength range of the simulation, and thus is neglected Roberts
. The absorption of Ag is included through the complex permittivity given from
Palik Palik .
Figure 5: (Color online) Top panels: Transmittance (blue-solid) and
reflectance (green-dashed) of the APC-Ag-APC thin film versus AOI when the
medium after the film is the anisotropic ENZ material with $\epsilon_{2x}=1$
(left panels) or the isotropic ENZ material (right panels).
$\epsilon_{2z}=0.001+0.6i$ for both cases. Bottom panels: transmission angle
(corresponding to the top panels) versus AOI. The thickness of the APC
$d=100\,$nm (left panels) and $d=80\,$nm (right panels). The thickness of Ag
is 10 nm for both cases. Figure 6: (Color online) Transmittance of the APC-Ag-
APC thin film versus AOI and wavelength when the medium after the film is the
anisotropic (left panel) or the isotropic (right panel) ENZ material. Color-
bars represent the magnitude of the transmittance. Simulation parameters are
the same as those in Fig. 5.
Shown in Fig. 5 are the transmission and reflection (top panels) of a plane
wave incident from air to the APC-Ag-APC structure, along with the
corresponding transmission angle (bottom panels). At the resonance, the
thickness of the APC $d=100\,$nm with the resonant wavelength
$\lambda=0.95\,\mu$m for the anisotropic ENZ medium (left panels); and the
$d=80\,$nm with the $\lambda=0.64\,\mu$m for the isotropic ENZ medium (right
panels). About $90\%$ transmittance are achieved for a wide range of the
incidence angle up to $70^{\circ}$ (see top panels) with nearly-collimated
transmission in the normal direction (see bottom panels). Transmittance of the
APC-Ag-APC thin film as a function of AOI and wavelength is presented in Fig.
6 when the medium at the back of the film is the anisotropic (left panel) or
the isotropic (right panel) ENZ material. In both cases, wide-angle $90\%$
transmittance are observed. It is worth mentioning that the loss of the ENZ
medium does not affect the transmittance of the APC-Ag-APC structure since the
transmitted power is computed right after the thin film, i.e., before
traveling through the ENZ medium. If the receptors are embedded very close to
the back of the film, the propagation loss in the ENZ medium can be minimized.
However, the loss of the ENZ material plays an important role on controlling
the direction of the transmission, no matter where the receptors are located.
In conclusions, we have demonstrated omnidirectionally transmitting the
electromagnetic power to one direction in the ENZ materials. This phenomenon
is realized based on two mechanisms. One is the loss-assistant bending of the
EM power to the normal for all angles of incidence. The other is the enhanced
transmission through structural resonances. This phenomenon may have
applications in communications, directive antennas, as well as detectors and
sensors to increase acceptance angle and redirect electromagnetic power
without using optical lenses and mechanical gimbals. The concept of employing
metamaterial loss to control the direction of the transmission brings a
positive perspective for material loss and may open up a new avenue for
metamaterial designs and applications.
The author gratefully acknowledge the sponsorship of ONR’s N-STAR and NAVAIR’s
ILIR programs.
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|
arxiv-papers
| 2011-12-07T06:18:04 |
2024-09-04T02:49:25.054327
|
{
"license": "Public Domain",
"authors": "Simin Feng",
"submitter": "Simin Feng",
"url": "https://arxiv.org/abs/1112.1482"
}
|
1112.1600
|
EUROPEAN ORGANIZATION FOR NUCLEAR RESEARCH (CERN)
CERN-PH-EP-2011-186 LHCb-PAPER-2011-025
Search for the rare decays $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
$B^{0}\rightarrow\mu^{+}\mu^{-}$
The LHCb Collaboration111Authors are listed on the following pages.
Abstract
A search for the decays $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
$B^{0}\rightarrow\mu^{+}\mu^{-}$ is performed with 0.37 fb-1 of $pp$
collisions at $\sqrt{s}$ = 7 TeV collected by the LHCb experiment in 2011. The
upper limits on the branching fractions are ${\cal
B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ $<1.6\times 10^{-8}$ and ${\cal
B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ $<3.6\times 10^{-9}$ at 95 % confidence
level. A combination of these results with the LHCb limits obtained with the
2010 dataset leads to ${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$
$<1.4\times 10^{-8}$ and ${\cal B}(B^{0}\rightarrow\mu^{+}\mu^{-})$
$<3.2\times 10^{-9}$ at 95 % confidence level.
Keywords: LHC, $b$-hadron, FCNC, rare decays, leptonic decays.
R. Aaij23, C. Abellan Beteta35,n, B. Adeva36, M. Adinolfi42, C. Adrover6, A.
Affolder48, Z. Ajaltouni5, J. Albrecht37, F. Alessio37, M. Alexander47, G.
Alkhazov29, P. Alvarez Cartelle36, A.A. Alves Jr22, S. Amato2, Y. Amhis38, J.
Anderson39, R.B. Appleby50, O. Aquines Gutierrez10, F. Archilli18,37, L.
Arrabito53, A. Artamonov 34, M. Artuso52,37, E. Aslanides6, G. Auriemma22,m,
S. Bachmann11, J.J. Back44, D.S. Bailey50, V. Balagura30,37, W. Baldini16,
R.J. Barlow50, C. Barschel37, S. Barsuk7, W. Barter43, A. Bates47, C. Bauer10,
Th. Bauer23, A. Bay38, I. Bediaga1, S. Belogurov30, K. Belous34, I.
Belyaev30,37, E. Ben-Haim8, M. Benayoun8, G. Bencivenni18, S. Benson46, J.
Benton42, R. Bernet39, M.-O. Bettler17, M. van Beuzekom23, A. Bien11, S.
Bifani12, T. Bird50, A. Bizzeti17,h, P.M. Bjørnstad50, T. Blake37, F. Blanc38,
C. Blanks49, J. Blouw11, S. Blusk52, A. Bobrov33, V. Bocci22, A. Bondar33, N.
Bondar29, W. Bonivento15, S. Borghi47,50, A. Borgia52, T.J.V. Bowcock48, C.
Bozzi16, T. Brambach9, J. van den Brand24, J. Bressieux38, D. Brett50, M.
Britsch10, T. Britton52, N.H. Brook42, H. Brown48, A. Büchler-Germann39, I.
Burducea28, A. Bursche39, J. Buytaert37, S. Cadeddu15, O. Callot7, M.
Calvi20,j, M. Calvo Gomez35,n, A. Camboni35, P. Campana18,37, A. Carbone14, G.
Carboni21,k, R. Cardinale19,i,37, A. Cardini15, L. Carson49, K. Carvalho
Akiba2, G. Casse48, M. Cattaneo37, Ch. Cauet9, M. Charles51, Ph.
Charpentier37, N. Chiapolini39, K. Ciba37, X. Cid Vidal36, G. Ciezarek49,
P.E.L. Clarke46,37, M. Clemencic37, H.V. Cliff43, J. Closier37, C. Coca28, V.
Coco23, J. Cogan6, P. Collins37, A. Comerma-Montells35, F. Constantin28, G.
Conti38, A. Contu51, A. Cook42, M. Coombes42, G. Corti37, G.A. Cowan38, R.
Currie46, B. D’Almagne7, C. D’Ambrosio37, P. David8, P.N.Y. David23, I. De
Bonis4, S. De Capua21,k, M. De Cian39, F. De Lorenzi12, J.M. De Miranda1, L.
De Paula2, P. De Simone18, D. Decamp4, M. Deckenhoff9, H. Degaudenzi38,37, M.
Deissenroth11, L. Del Buono8, C. Deplano15, D. Derkach14,37, O. Deschamps5, F.
Dettori24, J. Dickens43, H. Dijkstra37, P. Diniz Batista1, F. Domingo
Bonal35,n, S. Donleavy48, F. Dordei11, P. Dornan49, A. Dosil Suárez36, D.
Dossett44, A. Dovbnya40, F. Dupertuis38, R. Dzhelyadin34, A. Dziurda25, S.
Easo45, U. Egede49, V. Egorychev30, S. Eidelman33, D. van Eijk23, F. Eisele11,
S. Eisenhardt46, R. Ekelhof9, L. Eklund47, Ch. Elsasser39, D. Elsby55, D.
Esperante Pereira36, L. Estève43, A. Falabella16,14,e, E. Fanchini20,j, C.
Färber11, G. Fardell46, C. Farinelli23, S. Farry12, V. Fave38, V. Fernandez
Albor36, M. Ferro-Luzzi37, S. Filippov32, C. Fitzpatrick46, M. Fontana10, F.
Fontanelli19,i, R. Forty37, M. Frank37, C. Frei37, M. Frosini17,f,37, S.
Furcas20, A. Gallas Torreira36, D. Galli14,c, M. Gandelman2, P. Gandini51, Y.
Gao3, J-C. Garnier37, J. Garofoli52, J. Garra Tico43, L. Garrido35, D.
Gascon35, C. Gaspar37, N. Gauvin38, M. Gersabeck37, T. Gershon44,37, Ph.
Ghez4, V. Gibson43, V.V. Gligorov37, C. Göbel54, D. Golubkov30, A.
Golutvin49,30,37, A. Gomes2, H. Gordon51, M. Grabalosa Gándara35, R. Graciani
Diaz35, L.A. Granado Cardoso37, E. Graugés35, G. Graziani17, A. Grecu28, E.
Greening51, S. Gregson43, B. Gui52, E. Gushchin32, Yu. Guz34, T. Gys37, G.
Haefeli38, C. Haen37, S.C. Haines43, T. Hampson42, S. Hansmann-Menzemer11, R.
Harji49, N. Harnew51, J. Harrison50, P.F. Harrison44, J. He7, V. Heijne23, K.
Hennessy48, P. Henrard5, J.A. Hernando Morata36, E. van Herwijnen37, E.
Hicks48, K. Holubyev11, P. Hopchev4, W. Hulsbergen23, P. Hunt51, T. Huse48,
R.S. Huston12, D. Hutchcroft48, D. Hynds47, V. Iakovenko41, P. Ilten12, J.
Imong42, R. Jacobsson37, A. Jaeger11, M. Jahjah Hussein5, E. Jans23, F.
Jansen23, P. Jaton38, B. Jean-Marie7, F. Jing3, M. John51, D. Johnson51, C.R.
Jones43, B. Jost37, M. Kaballo9, S. Kandybei40, M. Karacson37, T.M. Karbach9,
J. Keaveney12, I.R. Kenyon55, U. Kerzel37, T. Ketel24, A. Keune38, B. Khanji6,
Y.M. Kim46, M. Knecht38, P. Koppenburg23, A. Kozlinskiy23, L. Kravchuk32, K.
Kreplin11, M. Kreps44, G. Krocker11, P. Krokovny11, F. Kruse9, K.
Kruzelecki37, M. Kucharczyk20,25,37,j, T. Kvaratskheliya30,37, V.N. La Thi38,
D. Lacarrere37, G. Lafferty50, A. Lai15, D. Lambert46, R.W. Lambert24, E.
Lanciotti37, G. Lanfranchi18, C. Langenbruch11, T. Latham44, C. Lazzeroni55,
R. Le Gac6, J. van Leerdam23, J.-P. Lees4, R. Lefèvre5, A. Leflat31,37, J.
Lefrançois7, O. Leroy6, T. Lesiak25, L. Li3, L. Li Gioi5, M. Lieng9, M.
Liles48, R. Lindner37, C. Linn11, B. Liu3, G. Liu37, J.H. Lopes2, E. Lopez
Asamar35, N. Lopez-March38, H. Lu38,3, J. Luisier38, A. Mac Raighne47, F.
Machefert7, I.V. Machikhiliyan4,30, F. Maciuc10, O. Maev29,37, J. Magnin1, S.
Malde51, R.M.D. Mamunur37, G. Manca15,d, G. Mancinelli6, N. Mangiafave43, U.
Marconi14, R. Märki38, J. Marks11, G. Martellotti22, A. Martens8, L. Martin51,
A. Martín Sánchez7, D. Martinez Santos37, A. Massafferri1, Z. Mathe12, C.
Matteuzzi20, M. Matveev29, E. Maurice6, B. Maynard52, A. Mazurov16,32,37, G.
McGregor50, R. McNulty12, C. Mclean14, M. Meissner11, M. Merk23, J. Merkel9,
R. Messi21,k, S. Miglioranzi37, D.A. Milanes13,37, M.-N. Minard4, J. Molina
Rodriguez54, S. Monteil5, D. Moran12, P. Morawski25, R. Mountain52, I. Mous23,
F. Muheim46, K. Müller39, R. Muresan28,38, B. Muryn26, B. Muster38, M. Musy35,
J. Mylroie-Smith48, P. Naik42, T. Nakada38, R. Nandakumar45, I. Nasteva1, M.
Nedos9, M. Needham46, N. Neufeld37, C. Nguyen-Mau38,o, M. Nicol7, V. Niess5,
N. Nikitin31, A. Nomerotski51, A. Novoselov34, A. Oblakowska-Mucha26, V.
Obraztsov34, S. Oggero23, S. Ogilvy47, O. Okhrimenko41, R. Oldeman15,d, M.
Orlandea28, J.M. Otalora Goicochea2, P. Owen49, K. Pal52, J. Palacios39, A.
Palano13,b, M. Palutan18, J. Panman37, A. Papanestis45, M. Pappagallo47, C.
Parkes47,37, C.J. Parkinson49, G. Passaleva17, G.D. Patel48, M. Patel49, S.K.
Paterson49, G.N. Patrick45, C. Patrignani19,i, C. Pavel-Nicorescu28, A. Pazos
Alvarez36, A. Pellegrino23, G. Penso22,l, M. Pepe Altarelli37, S.
Perazzini14,c, D.L. Perego20,j, E. Perez Trigo36, A. Pérez-Calero Yzquierdo35,
P. Perret5, M. Perrin-Terrin6, G. Pessina20, A. Petrella16,37, A.
Petrolini19,i, A. Phan52, E. Picatoste Olloqui35, B. Pie Valls35, B.
Pietrzyk4, T. Pilař44, D. Pinci22, R. Plackett47, S. Playfer46, M. Plo
Casasus36, G. Polok25, A. Poluektov44,33, E. Polycarpo2, D. Popov10, B.
Popovici28, C. Potterat35, A. Powell51, T. du Pree23, J. Prisciandaro38, V.
Pugatch41, A. Puig Navarro35, W. Qian52, J.H. Rademacker42, B.
Rakotomiaramanana38, M.S. Rangel2, I. Raniuk40, G. Raven24, S. Redford51, M.M.
Reid44, A.C. dos Reis1, S. Ricciardi45, K. Rinnert48, D.A. Roa Romero5, P.
Robbe7, E. Rodrigues47,50, F. Rodrigues2, P. Rodriguez Perez36, G.J. Rogers43,
S. Roiser37, V. Romanovsky34, M. Rosello35,n, J. Rouvinet38, T. Ruf37, H.
Ruiz35, G. Sabatino21,k, J.J. Saborido Silva36, N. Sagidova29, P. Sail47, B.
Saitta15,d, C. Salzmann39, M. Sannino19,i, R. Santacesaria22, C. Santamarina
Rios36, R. Santinelli37, E. Santovetti21,k, M. Sapunov6, A. Sarti18,l, C.
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Zhang3, A. Zhelezov11, L. Zhong3, E. Zverev31, A. Zvyagin37.
1Centro Brasileiro de Pesquisas Físicas (CBPF), Rio de Janeiro, Brazil
2Universidade Federal do Rio de Janeiro (UFRJ), Rio de Janeiro, Brazil
3Center for High Energy Physics, Tsinghua University, Beijing, China
4LAPP, Université de Savoie, CNRS/IN2P3, Annecy-Le-Vieux, France
5Clermont Université, Université Blaise Pascal, CNRS/IN2P3, LPC, Clermont-
Ferrand, France
6CPPM, Aix-Marseille Université, CNRS/IN2P3, Marseille, France
7LAL, Université Paris-Sud, CNRS/IN2P3, Orsay, France
8LPNHE, Université Pierre et Marie Curie, Université Paris Diderot,
CNRS/IN2P3, Paris, France
9Fakultät Physik, Technische Universität Dortmund, Dortmund, Germany
10Max-Planck-Institut für Kernphysik (MPIK), Heidelberg, Germany
11Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg, Heidelberg,
Germany
12School of Physics, University College Dublin, Dublin, Ireland
13Sezione INFN di Bari, Bari, Italy
14Sezione INFN di Bologna, Bologna, Italy
15Sezione INFN di Cagliari, Cagliari, Italy
16Sezione INFN di Ferrara, Ferrara, Italy
17Sezione INFN di Firenze, Firenze, Italy
18Laboratori Nazionali dell’INFN di Frascati, Frascati, Italy
19Sezione INFN di Genova, Genova, Italy
20Sezione INFN di Milano Bicocca, Milano, Italy
21Sezione INFN di Roma Tor Vergata, Roma, Italy
22Sezione INFN di Roma La Sapienza, Roma, Italy
23Nikhef National Institute for Subatomic Physics, Amsterdam, The Netherlands
24Nikhef National Institute for Subatomic Physics and Vrije Universiteit,
Amsterdam, The Netherlands
25Henryk Niewodniczanski Institute of Nuclear Physics Polish Academy of
Sciences, Kraców, Poland
26AGH University of Science and Technology, Kraców, Poland
27Soltan Institute for Nuclear Studies, Warsaw, Poland
28Horia Hulubei National Institute of Physics and Nuclear Engineering,
Bucharest-Magurele, Romania
29Petersburg Nuclear Physics Institute (PNPI), Gatchina, Russia
30Institute of Theoretical and Experimental Physics (ITEP), Moscow, Russia
31Institute of Nuclear Physics, Moscow State University (SINP MSU), Moscow,
Russia
32Institute for Nuclear Research of the Russian Academy of Sciences (INR RAN),
Moscow, Russia
33Budker Institute of Nuclear Physics (SB RAS) and Novosibirsk State
University, Novosibirsk, Russia
34Institute for High Energy Physics (IHEP), Protvino, Russia
35Universitat de Barcelona, Barcelona, Spain
36Universidad de Santiago de Compostela, Santiago de Compostela, Spain
37European Organization for Nuclear Research (CERN), Geneva, Switzerland
38Ecole Polytechnique Fédérale de Lausanne (EPFL), Lausanne, Switzerland
39Physik-Institut, Universität Zürich, Zürich, Switzerland
40NSC Kharkiv Institute of Physics and Technology (NSC KIPT), Kharkiv, Ukraine
41Institute for Nuclear Research of the National Academy of Sciences (KINR),
Kyiv, Ukraine
42H.H. Wills Physics Laboratory, University of Bristol, Bristol, United
Kingdom
43Cavendish Laboratory, University of Cambridge, Cambridge, United Kingdom
44Department of Physics, University of Warwick, Coventry, United Kingdom
45STFC Rutherford Appleton Laboratory, Didcot, United Kingdom
46School of Physics and Astronomy, University of Edinburgh, Edinburgh, United
Kingdom
47School of Physics and Astronomy, University of Glasgow, Glasgow, United
Kingdom
48Oliver Lodge Laboratory, University of Liverpool, Liverpool, United Kingdom
49Imperial College London, London, United Kingdom
50School of Physics and Astronomy, University of Manchester, Manchester,
United Kingdom
51Department of Physics, University of Oxford, Oxford, United Kingdom
52Syracuse University, Syracuse, NY, United States
53CC-IN2P3, CNRS/IN2P3, Lyon-Villeurbanne, France, associated member
54Pontifícia Universidade Católica do Rio de Janeiro (PUC-Rio), Rio de
Janeiro, Brazil, associated to 2
55University of Birmingham, Birmingham, United Kingdom
aP.N. Lebedev Physical Institute, Russian Academy of Science (LPI RAS),
Moscow, Russia
bUniversità di Bari, Bari, Italy
cUniversità di Bologna, Bologna, Italy
dUniversità di Cagliari, Cagliari, Italy
eUniversità di Ferrara, Ferrara, Italy
fUniversità di Firenze, Firenze, Italy
gUniversità di Urbino, Urbino, Italy
hUniversità di Modena e Reggio Emilia, Modena, Italy
iUniversità di Genova, Genova, Italy
jUniversità di Milano Bicocca, Milano, Italy
kUniversità di Roma Tor Vergata, Roma, Italy
lUniversità di Roma La Sapienza, Roma, Italy
mUniversità della Basilicata, Potenza, Italy
nLIFAELS, La Salle, Universitat Ramon Llull, Barcelona, Spain
oHanoi University of Science, Hanoi, Viet Nam
## 1 Introduction
Measurements of low-energy processes can provide indirect constraints on
particles that are too heavy to be produced directly. This is particularly
true for Flavour Changing Neutral Current (FCNC) processes which are highly
suppressed in the Standard Model (SM) and can only occur through higher-order
diagrams. The SM predictions for the branching fractions of the FCNC
decays222Inclusion of charged conjugated processes is implied throughout.
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and $B^{0}\rightarrow\mu^{+}\mu^{-}$ are
${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ = $(3.2\pm 0.2)\times 10^{-9}$
and ${\cal B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ = $(0.10\pm 0.01)\times
10^{-9}$ [1, *Buras:2010wr]. However, contributions from new processes or new
heavy particles can significantly enhance these values. For example, within
Minimal Supersymmetric extensions of the SM (MSSM), in the large $\tan\beta$
regime, ${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ is found to be
approximately proportional to $\tan^{6}\beta$ [3, *Hamzaoui:1998nu,
*Babu:1999hn, *Hall:1993gn, *Huang:1998vb], where $\tan\beta$ is the ratio of
the vacuum expectation values of the two neutral $C\\!P$-even Higgs fields.
The branching fractions could therefore be enhanced by orders of magnitude for
large values of $\tan\beta$.
The best published limits from the Tevatron are ${\cal
B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ $<~{}5.1~{}\times~{}10^{-8}$ at 95%
confidence level (CL) by the D0 collaboration using 6.1 fb-1 of data [8], and
${\cal B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ $<6.0\times 10^{-9}$ at 95% CL by
the CDF collaboration using 6.9 fb-1 of data [9]. In the same dataset the CDF
collaboration observes an excess of $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$
candidates compatible with ${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ =
$(1.8^{+1.1}_{-0.9})\times 10^{-8}$ and with an upper limit of ${\cal
B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ $<4.0\times 10^{-8}$ at 95% CL. The
CMS collaboration has recently published ${\cal
B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ $<1.9\times 10^{-8}$ at 95% CL and
${\cal B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ $<4.6\times 10^{-9}$ at 95% CL
using 1.14 fb-1 of data [10]. The LHCb collaboration has published the limits
[11] ${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ $<5.4\times 10^{-8}$ and
${\cal B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ $<1.5\times 10^{-8}$ at 95$\%$ CL
based on about 37 pb-1 of integrated luminosity collected in the 2010 run.
This Letter presents an analysis of the data recorded by LHCb in the first
half of 2011 which correspond to an integrated luminosity of $\sim$ 0.37 fb-1.
The results of this analysis are then combined with those published from the
2010 dataset.
## 2 The LHCb detector
The LHCb detector [12] is a single-arm forward spectrometer designed to study
production and decays of hadrons containing $b$ or $c$ quarks. The detector
consists of a vertex locator (VELO) providing precise locations of primary
$pp$ interaction vertices and detached vertices of long lived hadrons.
The momenta of charged particles are determined using information from the
VELO together with the rest of the tracking system, composed of a large area
silicon tracker located before a warm dipole magnet with a bending power of
$\sim$ 4 Tm, and a combination of silicon strip detectors and straw drift
chambers located after the magnet. Two Ring Imaging Cherenkov (RICH) detectors
are used for charged hadron identification in the momentum range 2–100
${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$. Photon, electron and hadron candidates
are identified by electromagnetic and hadronic calorimeters. A muon system of
alternating layers of iron and drift chambers provides muon identification.
The two calorimeters and the muon system provide the energy and momentum
information to implement a first level (L0) hardware trigger. An additional
trigger level (HLT) is software based, and its algorithms are tuned to the
experimental operating condition.
Events with a muon final states are triggered using two L0 trigger decisions:
the single-muon decision, which requires one muon candidate with a transverse
momentum $p_{\rm T}$ larger than 1.5 ${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$,
and the di-muon decision, which requires two muon candidates with transverse
momenta $p_{{\rm T},1}$ and $p_{{\rm T},2}$ satisfying the relation
$\sqrt{p_{{\rm T},1}\cdot p_{{\rm T},2}}>1.3$
${\mathrm{\,Ge\kern-1.00006ptV\\!/}c}$.
The single muon trigger decision in the second trigger level (HLT) includes a
cut on the impact parameter (${\rm IP}$) with respect to the primary vertex,
which allows for a lower $p_{\rm T}$ requirement ($p_{\rm T}>1.0$ GeV/$c$,
$\rm IP>0.1$ mm). The di-muon trigger decision requires muon pairs of opposite
charge with $p_{\rm T}>500$ MeV/c, forming a common vertex and with an
invariant mass $m_{\mu\mu}>4.7$ GeV/$c^{2}$. A second trigger decision,
primarily to select ${J\mskip-3.0mu/\mskip-2.0mu\psi\mskip 2.0mu}$ events,
requires $2.97<m_{\mu\mu}<3.21$ GeV/$c^{2}$. The remaining region of the di-
muon invariant mass range is also covered by trigger decisions that in
addition require the di-muon secondary vertex to be well separated from the
primary vertex.
Events with purely hadronic final states are triggered by the L0 trigger if
there is a calorimeter cluster with transverse energy $E_{\rm T}>3.6$ GeV.
Other HLT trigger decisions select generic displaced vertices, providing high
efficiency for purely hadronic decays.
## 3 Analysis strategy
Assuming the branching fractions predicted by the SM, and using the $b\bar{b}$
cross-section measured by LHCb in the pseudorapidity interval $2<\eta<6$ and
integrated over all transverse momenta of $\sigma_{b\overline{b}}=75\pm
14\,\mu$b [13], approximately $3.9$ $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
0.4 $B^{0}\rightarrow\mu^{+}\mu^{-}$ events are expected to be triggered,
reconstructed and selected in the analysed sample embedded in a large
background.
The general structure of the analysis is based upon the one described in Ref.
[11]. First a very efficient selection removes the biggest amount of
background while keeping most of the signal within the LHCb acceptance. The
number of observed events is compared to the number of expected signal and
background events in bins of two independent variables, the invariant mass and
the output of a multi-variate discriminant. The discriminant is a Boosted
Decision Tree (BDT) constructed using the TMVA package [14]. It supersedes the
Geometrical Likelihood (GL) used in the previous analysis [11] as it has been
found more performant in discriminating between signal and background events
in simulated samples. No data were used in the choice of the multivariate
discriminant in order not to bias the result.
The combination of variables entering the BDT discriminant is optimized using
simulated events. The probability for a signal or background event to have a
given value of the BDT output is obtained from data using
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ candidates (where
$h^{(^{\prime})}$ can be a pion or a kaon) as signal and sideband
$B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}$ candidates as background.
The invariant mass line shape of the signals is described by a Crystal Ball
function [15] with parameters extracted from data control samples. The central
values of the masses are obtained from $B^{0}\rightarrow K^{+}\pi^{-}$ and
$B^{0}_{s}\rightarrow K^{+}K^{-}$ samples. The $B^{0}_{s}$ and $B^{0}$ mass
resolutions are estimated by interpolating those obtained with di-muon
resonances ($J/\psi,\psi(2S)$ and $\Upsilon(1S,2S,3S)$) and cross-checked with
a fit to the invariant mass distributions of both inclusive
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ decays and exclusive
$B^{0}\rightarrow K^{+}\pi^{-}$ decays. The central values of the masses and
the mass resolution are used to define the signal regions.
The number of expected signal events, for a given branching fraction
hypothesis, is obtained by normalizing to channels of known branching
fractions: $B^{+}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}K^{+}$, $B^{0}_{s}\\!\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}\phi$ and $B^{0}\rightarrow K^{+}\pi^{-}$. These channels are selected
in a way as similar as possible to the signals in order to minimize the
systematic uncertainty related to the different phase space accessible to each
final state.
The BDT output and invariant mass distributions for combinatorial background
events in the signal regions are obtained using fits of the mass distribution
of events in the mass sidebands in bins of the BDT output.
The two-dimensional space formed by the invariant mass and the BDT output is
binned. For each bin we count the number of candidates observed in the data,
and compute the expected number of signal events and the expected number of
background events. The binning is unchanged with respect to the 2010 analysis
[11]. The compatibility of the observed distribution of events in all bins
with the distribution expected for a given branching fraction hypothesis is
computed using the $\textrm{CL}_{\textrm{s}}$ method [16], which allows a
given hypothesis to be excluded at a given confidence level.
## 4 Selection
The $B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}$ selections require two muon
candidates of opposite charge. Tracks are required to be of good quality and
to be displaced with respect to any primary vertex. The secondary vertex is
required to be well fitted ($\chi^{2}/{\rm nDoF}<9$) and must be separated
from the primary vertex in the forward direction by a distance of flight
significance ($L/\sigma(L)$) greater than 15. When more than one primary
vertex is reconstructed, the one that gives the minimum impact parameter
significance for the candidate is chosen. The reconstructed candidate has to
point to this primary vertex (${\rm IP}/\sigma({\rm IP})<5$).
Improvements have been made to the selection developed for 2010 data [11]. The
RICH is used to identify kaons in the $B^{0}_{s}\rightarrow J/\psi\phi$
normalization channel and the Kullback-Leibler (KL) distance [17,
*KLdistance2] is used to suppress duplicated tracks created by the
reconstruction. This procedure compares the parameters and correlation
matrices of the reconstructed tracks and where two are found to be similar, in
this case with a symmetrized KL divergence less than 5000, only the one with
the higher track fit quality is considered.
The inclusive $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ sample is the
main control sample for the determination from data of the probability
distribution function (PDF) of the BDT output. This sample is selected in
exactly the same way as the $B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}$ signals
apart from the muon identification requirement. The same selection is also
applied to the $B^{0}\rightarrow K^{+}\pi^{-}$ normalization channel.
The muon identification efficiency is uniform within $\sim 1\%$ in the
considered phase space therefore no correction is added to the BDT PDF
extracted from the $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ sample. The
remaining phase space dependence of the muon identification efficiency is
instead taken into account in the computation of the normalization factor when
the $B^{0}\rightarrow K^{+}\pi^{-}$ channel is considered.
The $J/\psi\rightarrow\mu\mu$ decay in the $B^{+}\rightarrow J/\psi K^{+}$ and
$B^{0}_{s}\rightarrow J/\psi\phi$ normalization channels is selected in a very
similar way to the $B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}$ channels, apart from
the pointing requirement. $K^{\pm}$ candidates are required to be identified
by the RICH detector and to pass track quality and impact parameter cuts.
To avoid pathological events, all tracks from selected candidates are required
to have a momentum less than 1 TeV/$c$. Only $B$ candidates with decay times
less than $5\,\tau_{B_{(s)}^{0}}$, where $\tau_{B^{0}_{(s)}}$ is the $B$
lifetime [19], are accepted for further analysis. Di-muon candidates coming
from elastic di-photon production are removed by requiring a minimum
transverse momentum of the $B$ candidate of
500${\mathrm{\,Me\kern-1.00006ptV\\!/}c}$.
## 5 Determination of the mass and BDT distributions
The variables entering the BDT discriminant are the six variables used as
input to the ${\rm GL}$ in the 2010 analysis plus three new variables. The six
variables used in the 2010 analysis are the $B$ lifetime, impact parameter,
transverse momentum, the minimum impact parameter significance (${\rm
IP}/\sigma({\rm IP})$) of the muons, the distance of closest approach between
the two muons and the isolation of the two muons with respect to any other
track in the event. The three new variables are:
1. 1.
the minimum $p_{\rm T}$ of the two muons;
2. 2.
the cosine of the angle between the muon momentum in the $B$ rest frame and
the vector perpendicular to the $B$ momentum and the beam axis:
$\cos P=\frac{p_{y,\mu 1}\,p_{x,B}-p_{x,\mu 1}\,p_{y,B}}{p_{{\rm
T},B}\,(m_{\mu\mu}/2)}$ (1)
where $\mu_{1}$ labels one of the muons and $m_{\mu\mu}$ is the reconstructed
$B$ candidate mass333As the $B$ is a (pseudo)-scalar particle, this variable
is uniformely distributed for signal candidates while is peaked at zero for
$b\bar{b}\rightarrow\mu^{+}\mu^{-}X$ background candidates. In fact, muons
from semi-leptonic decays are mostly emitted in the direction of the $b$’s
and, therefore, lie in a plane formed by the $B$ momentum and the beam axis.;
3. 3.
the $B$ isolation [20]
$I_{B}=\frac{p_{\rm T}(B)}{p_{\rm T}(B)+\sum_{i}p_{{\rm T},i}},$ (2)
where $p_{\rm T}(B)$ is the $B$ transverse momentum with respect to the beam
line and the sum is over all the tracks, excluding the muon candidates, that
satisfy $\sqrt{\delta\eta^{2}+\delta\phi^{2}}~{}<~{}1.0$, where $\delta\eta$
and $\delta\phi$ denote respectively the difference in pseudorapidity and
azimuthal angle between the track and the $B$ candidate.
The BDT output is found to be independent of the invariant mass for both
signal and background and is defined such that the signal is uniformly
distributed between zero and one and the background peaks at zero. The BDT
range is then divided in four bins of equal width. The BDT is trained using
simulated samples ($B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}$ for signals and
$b\bar{b}\rightarrow\mu^{+}\mu^{-}X$ for background where $X$ is any other set
of particles) and the PDF obtained from data as explained below.
### 5.1 Combinatorial background PDFs
The BDT and invariant mass shapes for the combinatorial background inside the
signal regions are determined from data by interpolating the number of
expected events using the invariant mass sidebands for each BDT bin. The
boundaries of the signal regions are defined as $m_{B^{0}}\pm 60$
${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ and $m_{B^{0}_{s}}\pm 60$
${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ and the mass sidebands as
$[m_{B^{0}}-600{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}},m_{B^{0}}-60{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}]$
and
$[m_{B^{0}_{s}}+60{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}},m_{B^{0}_{s}}+600{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}]$.
Figure 1 shows the invariant mass distribution for events that lie in each BDT
output bin. In each case the fit model used to estimate the expected number of
combinatorial background events in the signal regions is superimposed.
Aside from combinatorial background, the low-mass sideband is potentially
polluted by two other contributions: cascading $b\rightarrow
c\mu\nu\rightarrow\mu\mu X$ decays below
4900${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ and peaking background from
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ candidates with the two hadrons
misidentified as muons above 5000${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$.
To avoid these contaminations, the number of expected combinatorial background
events is obtained by fitting a single exponential function to the events in
the reduced low-mass sideband [4900, 5000]
${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ and in the full high-mass sideband.
As a cross-check, two other models, a single exponential function and the sum
of two exponential functions, have been used to fit the events in different
ranges of sidebands providing consistent background estimates inside the
signal regions.
Figure 1: Distribution of the $\mu^{+}\mu^{-}$ invariant mass for events in
each BDT output bin. The curve shows the model used to fit the sidebands and
extract the expected number of combinatorial background events in the
$B^{0}_{s}$ and $B^{0}$ signal regions, delimited by the vertical dotted
orange and dashed green lines respectively. Only events in the region in which
the line is solid have been considered in the fit.
### 5.2 Peaking background PDFs
The peaking backgrounds due to $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$
events in which both hadrons are misidentified as muons have been evaluated
from data and simulated events to be $N_{B^{0}_{s}}=1.0\pm 0.4$ events and
$N_{B^{0}}=5.0\pm 0.9$ events within the two mass windows and in the whole BDT
output range. The mass line shape of the peaking background is obtained from a
simulated sample of doubly-misidentified $B^{0}_{(s)}\rightarrow
h^{+}h^{{}^{\prime}-}$ events and normalized to the number of events expected
in the two search windows from data, $N_{B^{0}_{s}}$ and $N_{B^{0}}$. The BDT
PDF of the peaking background is assumed to be the same as for the signal.
### 5.3 Signal PDFs
The BDT PDF for signal events is determined using an inclusive
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ sample. Only events which are
triggered independently on the signal candidates have been considered (TIS
events).
The number of $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ signal events in
each BDT output bin is determined by fitting the $hh^{\prime}$ invariant mass
distribution under the $\mu\mu$ mass hypothesis [21]. Figure 2 shows the fit
to the mass distribution of the full sample and for the three highest BDT
output bins for $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ TIS events. The
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ exclusive decays, the
combinatorial background and the physical background components are drawn
under the fit to the data; the physical background is due to the partial
reconstruction of three-body $B$ meson decays.
In order to cross-check this result, two other fits have been performed on the
same dataset. The signal line shape is parametrized either by a single or a
double Crystal Ball function [15], the combinatorial background by an
exponential function and the physical background by an ARGUS function [22]. In
addition, exclusive
$B^{0}_{(s)}\rightarrow\pi^{-}K^{+},\pi^{-}\pi^{+},K^{-}K^{+}$ channels,
selected using the $K-\pi$ separation capability of the RICH system, are used
to cross-check the calibration of the BDT output both using the
$\pi^{-}K^{+},\pi^{-}\pi^{+},K^{-}K^{+}$ inclusive yields without separating
$B$ and $B^{0}_{s}$ and using the $B^{0}\rightarrow K^{+}\pi^{-}$ exclusive
channel alone. The maximum spread in the fractional yield obtained among the
different models has been used as a systematic uncertainty in the signal BDT
PDF. The BDT PDFs for signals and combinatorial background are shown in Fig.
3.
Figure 2: Invariant mass distributions of $B^{0}_{(s)}\rightarrow
h^{+}h^{{}^{\prime}-}$ candidates in the $\mu^{+}\mu^{-}$ mass hypothesis for
the whole sample (top left) and for the samples in the three highest bins of
the BDT output (top right, bottom left, bottom right). The
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ exclusive decays, the
combinatorial background and the physical background components are drawn
under the fit to the data (solid blue line). Figure 3: BDT probability
distribution functions of signal events (solid squares) and combinatorial
background (open circles): the PDF for the signal is obtained from the
inclusive sample of TIS $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ events,
the PDF for the combinatorial background is obtained from the events in the
mass sidebands.
The invariant mass shape for the signal is parametrized as a Crystal Ball
function. The mean value is determined using the $B^{0}\rightarrow
K^{+}\pi^{-}$ and $B^{0}_{s}\rightarrow K^{+}K^{-}$ exclusive channels and the
transition point of the radiative tail is obtained from simulated events [11].
The central values are
$\displaystyle m_{B^{0}_{s}}$ $\displaystyle=$ $\displaystyle 5358.0\pm
1.0{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}},$ $\displaystyle m_{B^{0}}$
$\displaystyle=$ $\displaystyle 5272.0\pm
1.0{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}.$
The measured values of $m_{B^{0}}$ and $m_{B^{0}_{s}}$ are $7-8$
${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ below the PDG values [19] due to
the fact that the momentum scale is uncalibrated in the dataset used in this
analysis. The mass resolutions are extracted from data with a linear
interpolation between the measured resolution of charmonium and bottomonium
resonances decaying into two muons: ${J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}$, $\psi{(2S)}$, $\mathchar 28935\relax{(1S)}$, $\mathchar
28935\relax{(2S)}$ and $\mathchar 28935\relax{(3S)}$. The mass line shapes for
quarkonium resonances are shown in Fig. 4. Each resonance is fitted with two
Crystal Ball functions with common mean value and common resolution but
different parameterization of the tails. The background is fitted with an
exponential function.
The results of the interpolation at the $m_{B^{0}_{s}}$ and $m_{B^{0}}$ masses
are
$\displaystyle\sigma(m_{B^{0}_{s}})$ $\displaystyle=$ $\displaystyle 24.6\pm
0.2_{\rm(stat)}\pm 1.0_{\rm(syst)}{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}},$
$\displaystyle\sigma(m_{B^{0}})$ $\displaystyle=$ $\displaystyle 24.3\pm
0.2_{\rm(stat)}\pm 1.0_{\rm(syst)}{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}.$
This result has been checked using both the fits to the
$B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ inclusive decay line shape and
the $B^{0}\rightarrow K^{+}\pi^{-}$ exclusive decay. The results are in
agreement within the uncertainties.
Figure 4: Di-muon invariant mass spectrum in the ranges (2.9 – 3.9)
${\mathrm{\,Ge\kern-1.00006ptV\\!/}c^{2}}$ (left) and (9–11)
${\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ (right).
## 6 Normalization
To estimate the signal branching fraction, the number of observed signal
events is normalized to the number of events of a channel with a well known
branching fraction. Three complementary normalization channels are used:
$B^{+}\rightarrow J/\psi(\mu^{+}\mu^{-})K^{+}$,
$B^{0}_{s}\rightarrow{J\mskip-3.0mu/\mskip-2.0mu\psi\mskip
2.0mu}(\mu^{+}\mu^{-})\phi(K^{+}K^{-})$ and $B^{0}\rightarrow K^{+}\pi^{-}$.
The first two channels have similar trigger and muon identification
efficiencies to the signal but different number of particles in the final
state. The third channel has a similar topology but is selected by different
trigger lines.
The numbers of $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
$B^{0}\rightarrow\mu^{+}\mu^{-}$ candidates are translated into a branching
fractions (${\cal B}$) using the equation
${\cal B}={\cal B}_{\rm
norm}\times\frac{\rm\epsilon_{norm}^{REC}\epsilon_{norm}^{SEL|REC}\epsilon_{norm}^{TRIG|SEL}}{\rm\epsilon_{sig}^{REC}\epsilon_{sig}^{SEL|REC}\epsilon_{sig}^{TRIG|SEL}}\times\frac{f_{\rm
norm}}{f_{d(s)}}\times\frac{N_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}}}{N_{\rm
norm}}=\alpha^{\rm norm}_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}}\times
N_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}},$ (3)
where $f_{d(s)}$ and $f_{\rm norm}$ are the probabilities that a $b$ quark
fragments into a $B^{0}_{(s)}$ and into the $b$ hadron involved for the chosen
normalization mode. LHCb has measured $f_{s}/f_{d}=0.267^{+0.021}_{-0.020}$
[23]. ${\cal B}_{\rm norm}$ is the branching fraction and $N_{\rm norm}$ is
the number of selected events of the normalization channel. The efficiency is
the product of three factors: $\epsilon^{\rm REC}$ is the reconstruction
efficiency of all the final state particles of the decay including the
geometric acceptance of the detector; $\epsilon^{\rm SEL|REC}$ is the
selection efficiency for reconstructed events; $\epsilon^{\rm TRIG|SEL}$ is
the trigger efficiency for reconstructed and selected events. The subscript
($\rm sig,norm$) indicates whether the efficiency refers to the signal or the
normalization channel. Finally, $\alpha^{\rm
norm}_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}}$ is the normalization factor (or
single event sensitivity) and $N_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}}$ the
number of observed signal events.
For each normalization channel $N_{\rm norm}$ is obtained from a fit to the
invariant mass distribution. The invariant mass distributions for
reconstructed $B^{+}\rightarrow J/\psi K^{+}$ and $B^{0}_{s}\rightarrow
J/\psi\phi$ candidates are shown in Fig. 5, while the $B^{0}\rightarrow
K^{+}\pi^{-}$ yield is obtained from the full $B^{0}_{(s)}\rightarrow
h^{+}h^{{}^{\prime}-}$ fit as shown in the top left of Fig. 2.
Figure 5: Invariant mass distributions of the $B^{+}\rightarrow J/\psi K^{+}$
(left) and $B^{0}_{s}\rightarrow J/\psi\phi$ (right) candidates used in the
normalization procedure.
The numbers used to calculate the normalization factors are summarized in
Table 1. A weighted average of the three normalization channels, assuming the
tracking and trigger uncertainties to be correlated between the two $J/\psi$
normalization channels and the uncertainty on $f_{d}/f_{s}$ to be correlated
between the $B^{+}\rightarrow J/\psi K^{+}$ and $B^{0}\rightarrow
K^{+}\pi^{-}$, gives
$\displaystyle\alpha^{\rm norm}_{B^{0}_{s}\rightarrow\mu^{+}\mu^{-}}=(8.38\pm
0.74)\times 10^{-10}\,,$ $\displaystyle\alpha^{\rm
norm}_{B^{0}\rightarrow\mu^{+}\mu^{-}}=(2.20\pm 0.11)\times 10^{-10}\,.$
These normalization factors are used to determine the limits.
Table 1: Summary of the quantities and their uncertainties required to calculate the normalization factors ($\alpha^{\rm norm}_{B^{0}_{(s)}\rightarrow\mu^{+}\mu^{-}}$) for the three normalization channels considered. The branching fractions are taken from Refs. [19, 24]. The trigger efficiency and the number of $B^{0}\rightarrow K^{+}\pi^{-}$ candidates correspond to TIS events. | $\cal B$ | $\frac{\rm\epsilon_{norm}^{REC}\epsilon_{norm}^{SEL|REC}}{\rm\epsilon_{sig}^{REC}\epsilon_{sig}^{SEL|REC}}$ | $\frac{\rm\epsilon_{norm}^{TRIG|SEL}\rule{0.0pt}{7.71552pt}}{\rm\epsilon_{sig}^{TRIG|SEL}\rule[-4.33998pt]{0.0pt}{0.0pt}}$ | $N_{\rm norm}$ | $\alpha^{\rm norm}_{B^{0}\rightarrow\mu^{+}\mu^{-}}$ | $\alpha^{\rm norm}_{B^{0}_{s}\rightarrow\mu^{+}\mu^{-}}$
---|---|---|---|---|---|---
| $(\times 10^{-5})$ | | | | | | | $(\times 10^{-10})$ | $(\times 10^{-9})$
$B^{+}\rightarrow J/\psi K^{+}$ | $6.01\,\pm$ | $\,0.21$ | $0.48\,\pm$ | $\,0.014$ | $0.95\,\pm$ | $\,0.01$ | $124\,518\,\pm$ | $\,2\,025$ | $2.23\,\pm$ | $\,0.11$ | $0.83\,\pm$ | $\,0.08$
$B^{0}_{s}\rightarrow J/\psi\phi$ | $3.4\,\pm$ | $\,0.9$ | $0.24\,\pm$ | $\,0.014$ | $0.95\,\pm$ | $\,0.01$ | $6\,940\,\pm$ | $\,93$ | $2.96\,\pm$ | $\,0.84$ | $1.11\,\pm$ | $\,0.30$
$B^{0}\rightarrow K^{+}\pi^{-}$ | $1.94\,\pm$ | $\,0.06$ | $0.86\,\pm$ | $\,0.02$ | $0.049\,\pm$ | $\,0.004$ | $4\,146\,\pm$ | $\,608$ | $1.98\,\pm$ | $\,0.34$ | $0.74\,\pm$ | $\,0.14$
## 7 Results
The results for $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
$B^{0}\rightarrow\mu^{+}\mu^{-}$ are summarized in Table 2 and Table 3
respectively and in each of the bins the expected number of combinatorial
background, peaking background, signal events, with the SM prediction assumed,
is shown together with the observations on the data. The uncertainties in the
signal and background PDFs and normalization factors are used to compute the
uncertainties on the background and signal predictions.
The two dimensional (mass, BDT) distribution of selected events can be seen in
Fig. 6. The distribution of the invariant mass in the four BDT bins is shown
in Fig. 7 for $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and in Fig. 8 for
$B^{0}\rightarrow\mu^{+}\mu^{-}$ selected candidates.
Figure 6: Distribution of selected di-muon events in the invariant mass–BDT
plane. The orange short-dashed (green long-dashed) lines indicate the $\pm
60{\mathrm{\,Me\kern-1.00006ptV\\!/}c^{2}}$ search window around the mean
$B^{0}_{s}$ ($B^{0}$) mass. Figure 7: Distribution of selected di-muon events
in the $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ mass window for the four BDT
output bins. The black dots are data, the light grey histogram shows the
contribution of the combinatorial background, the black filled histogram shows
the contribution of the $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$
background and the dark grey filled histogram the contribution of
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ signal events according to the SM rate.
The hatched area depicts the uncertainty on the sum of the expected
contributions. Figure 8: Distribution of selected di-muon events in the
$B^{0}\rightarrow\mu^{+}\mu^{-}$ mass window for the four BDT output bins. The
black dots are data, the light grey histogram shows the contribution of the
combinatorial background, the black filled histogram shows the contribution of
the $B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$ background and the dark
grey filled histogram shows the cross-feed of
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ events in the $B^{0}$ mass window
assuming the the SM rate. The hatched area depicts the uncertainty on the sum
of the expected contributions.
The compatibility of the distribution of events inside the search window in
the invariant mass–BDT plane with a given branching fraction hypothesis is
evaluated using the $\textrm{CL}_{\textrm{s}}$ method [16]. This method
provides three estimators: $\textrm{CL}_{\textrm{s+b}}$, a measure of the
compatibility of the observed distribution with the signal and background
hypotheses, $\textrm{CL}_{\textrm{b}}$, a measure of the compatibility with
the background-only hypothesis and $\textrm{CL}_{\textrm{s}}$, a measure of
the compatibility of the observed distribution with the signal and background
hypotheses normalized to the background-only hypothesis.
The expected $\textrm{CL}_{\textrm{s}}$ values are shown in Fig. 9 for
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and for $B^{0}\rightarrow\mu^{+}\mu^{-}$
as dashed black lines under the hypothesis that background and SM events are
observed. The shaded areas cover the region of $\pm 1\sigma$ of compatible
observations. The observed values of $\textrm{CL}_{\textrm{s}}$ as a function
of the assumed branching ratio is shown as dotted blue lines on both plots.
The expected limits and the measured limits for
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and $B^{0}\rightarrow\mu^{+}\mu^{-}$ at
90 % and 95 % CL are shown in Table 4 and Table 5, respectively. For the
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ decay, the expected limits are computed
allowing the presence of $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ events according
to the SM branching fraction. For the $B^{0}\rightarrow\mu^{+}\mu^{-}$ decay
the expected limit is computed in the background-only hypothesis and also
allowing the presence of $B^{0}\rightarrow\mu^{+}\mu^{-}$ events with the SM
rate: the two results are identical. In the determination of the limits, the
cross-feed of $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$
($B^{0}\rightarrow\mu^{+}\mu^{-}$) events in the $B^{0}$ ($B^{0}_{s}$) mass
window has been taken into account assuming the SM rates.
The observed $\textrm{CL}_{\textrm{b}}$ values are shown in the same tables.
The comparison of the observed distribution of events with the expected
background distribution results in a p-value $(1-\textrm{CL}_{\textrm{b}})$ of
5 % for the $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and 32 % for the
$B^{0}\rightarrow\mu^{+}\mu^{-}$ decay. For the
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ decay, the probability that the observed
events are compatible with the sum of expected background events and signal
events according to the SM rate is measured by
$1-$$\textrm{CL}_{\textrm{s+b}}$ and it is 33%.
The result obtained in 2011 with 0.37 fb-1 has been combined with the
published result based on $\sim 37$ pb-1[11]. The expected and observed limits
for 90 % and 95 % CL for the combined results are shown in Table 4 for the
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ decay and in Table 5 for the
$B^{0}\\!\rightarrow\mu^{+}\mu^{-}$ decay.
Figure 9: $\textrm{CL}_{\textrm{s}}$ as a function of the assumed $\cal B$.
Expected (observed) values are shown by dashed black (dotted blue) lines. The
expected $\textrm{CL}_{\textrm{s}}$ values have been computed assuming a
signal yield corresponding to the SM branching fractions. The green (grey)
shaded areas cover the region of $\pm 1\sigma$ of compatible observations. The
measured upper limits at 90% and 95% CL are also shown. Left:
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$, right: $B^{0}\rightarrow\mu^{+}\mu^{-}$.
Table 2: Expected combinatorial background events, expected peaking
($B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$) background events, expected
signal events assuming the SM branching fraction prediction, and observed
events in the $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ search window.
| | | BDT
---|---|---|---
| | | 0\. – 0.25 | 0.25 – 0.5 | 0.5 – 0.75 | 0.75 – 1.
Invariant mass [$\,{\rm MeV}/c^{2}$ ] | 5298 – 5318 | Expected comb. bkg | $575.5^{+6.5}_{-6.0}$ | $6.96^{+0.63}_{-0.57}$ | $1.19^{+0.39}_{-0.35}$ | $0.111^{+0.083}_{-0.066}$
Expected peak. bkg | $0.126^{+0.037}_{-0.030}$ | $0.124^{+0.037}_{-0.030}$ | $0.124^{+0.037}_{-0.030}$ | $0.127^{+0.038}_{-0.031}$
Expected signal | $0.059^{+0.023}_{-0.022}$ | $0.0329^{+0.0128}_{-0.0095}$ | $0.0415^{+0.0120}_{-0.0085}$ | $0.0411^{+0.0135}_{-0.0099}$
Observed | $533$ | $10$ | $1$ | $0$
5318 – 5338 | Expected comb. bkg | $566.8^{+6.3}_{-5.8}$ | $6.90^{+0.61}_{-0.55}$ | $1.16^{+0.38}_{-0.34}$ | $0.109^{+0.079}_{-0.063}$
Expected peak. bkg | $0.052^{+0.023}_{-0.018}$ | $0.054^{+0.026}_{-0.019}$ | $0.052^{+0.024}_{-0.018}$ | $0.051^{+0.023}_{-0.018}$
Expected signal | $0.205^{+0.073}_{-0.074}$ | $0.114^{+0.040}_{-0.031}$ | $0.142^{+0.036}_{-0.025}$ | $0.142^{+0.042}_{-0.031}$
Observed | $525$ | $9$ | $0$ | $1$
5338 – 5358 | Expected comb. bkg | $558.2^{+6.1}_{-5.6}$ | $6.84^{+0.59}_{-0.54}$ | $1.14^{+0.37}_{-0.33}$ | $0.106^{+0.075}_{-0.060}$
Expected peak. bkg | $0.024^{+0.028}_{-0.012}$ | $0.025^{+0.026}_{-0.012}$ | $0.024^{+0.027}_{-0.012}$ | $0.025^{+0.025}_{-0.012}$
Expected signal | $0.38^{+0.14}_{-0.14}$ | $0.213^{+0.075}_{-0.058}$ | $0.267^{+0.065}_{-0.047}$ | $0.265^{+0.077}_{-0.058}$
Observed | $561$ | $6$ | $2$ | $1$
5358 – 5378 | Expected comb. bkg | $549.8^{+6.0}_{-5.4}$ | $6.77^{+0.57}_{-0.52}$ | $1.11^{+0.36}_{-0.32}$ | $0.103^{+0.073}_{-0.057}$
Expected peak. bkg | $0.0145^{+0.0220}_{-0.0091}$ | $0.0151^{+0.0230}_{-0.0091}$ | $0.0153^{+0.0232}_{-0.0098}$ | $0.015^{+0.023}_{-0.010}$
Expected signal | $0.38^{+0.14}_{-0.14}$ | $0.213^{+0.075}_{-0.057}$ | $0.267^{+0.065}_{-0.047}$ | $0.265^{+0.077}_{-0.057}$
Observed | $515$ | $7$ | $0$ | $0$
5378 – 5398 | Expected comb. bkg | $541.5^{+5.8}_{-5.3}$ | $6.71^{+0.55}_{-0.51}$ | $1.09^{+0.34}_{-0.31}$ | $0.101^{+0.070}_{-0.054}$
Expected peak. bkg | $0.0115^{+0.0175}_{-0.0086}$ | $0.0116^{+0.0177}_{-0.0090}$ | $0.0118^{+0.0179}_{-0.0090}$ | $0.0118^{+0.0179}_{-0.0088}$
Expected signal | $0.204^{+0.073}_{-0.074}$ | $0.114^{+0.040}_{-0.031}$ | $0.142^{+0.036}_{-0.026}$ | $0.141^{+0.042}_{-0.031}$
Observed | $547$ | $10$ | $1$ | $1$
5398 – 5418 | Expected comb. bkg | $533.4^{+5.7}_{-5.2}$ | $6.65^{+0.53}_{-0.49}$ | $1.07^{+0.34}_{-0.30}$ | $0.098^{+0.068}_{-0.051}$
Expected peak. bkg | $0.0089^{+0.0136}_{-0.0065}$ | $0.0088^{+0.0133}_{-0.0066}$ | $0.0091^{+0.0138}_{-0.0070}$ | $0.0090^{+0.0137}_{-0.0065}$
Expected signal | $0.058^{+0.024}_{-0.021}$ | $0.0323^{+0.0128}_{-0.0093}$ | $0.0407^{+0.0120}_{-0.0087}$ | $0.0402^{+0.0137}_{-0.0097}$
Observed | $501$ | $4$ | $1$ | $0$
Table 3: Expected combinatorial background events, expected peaking
($B^{0}_{(s)}\rightarrow h^{+}h^{{}^{\prime}-}$) background events, expected
$B^{0}\\!\rightarrow\mu^{+}\mu^{-}$ signal events assuming the SM branching
fraction, expected cross-feed events from $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$
assuming the SM branching fraction and observed events in the
$B^{0}\rightarrow\mu^{+}\mu^{-}$ search window.
| | | BDT
---|---|---|---
| | | 0\. – 0.25 | 0.25 – 0.5 | 0.5 – 0.75 | 0.75 – 1.
Invariant mass [$\,{\rm MeV}/c^{2}$ ] | 5212 – 5232 | Expected comb. bkg | $614.2^{+7.5}_{-7.0}$ | $7.23^{+0.77}_{-0.68}$ | $1.31^{+0.46}_{-0.40}$ | $0.123^{+0.107}_{-0.072}$
Expected peak. bkg | $0.203^{+0.038}_{-0.034}$ | $0.206^{+0.038}_{-0.034}$ | $0.203^{+0.037}_{-0.034}$ | $0.205^{+0.038}_{-0.034}$
Cross-feed | $0.0056^{+0.0021}_{-0.0020}$ | $0.00312^{+0.00119}_{-0.00087}$ | $0.00391^{+0.00107}_{-0.00078}$ | $0.00387^{+0.00122}_{-0.00092}$
Expected signal | $0.0070^{+0.0027}_{-0.0026}$ | $0.0039^{+0.0015}_{-0.0011}$ | $0.0049^{+0.0014}_{-0.0010}$ | $0.0048^{+0.0016}_{-0.0012}$
| Observed | $554$ | $6$ | $0$ | $2$
5232 – 5252 | Expected comb. bkg | $605.0^{+7.2}_{-6.8}$ | $7.17^{+0.74}_{-0.65}$ | $1.29^{+0.44}_{-0.39}$ | $0.121^{+0.102}_{-0.072}$
Expected peak. bkg | $0.281^{+0.056}_{-0.049}$ | $0.279^{+0.056}_{-0.049}$ | $0.280^{+0.056}_{-0.049}$ | $0.280^{+0.058}_{-0.050}$
Cross-feed | $0.0071^{+0.0027}_{-0.0026}$ | $0.0039^{+0.0015}_{-0.0011}$ | $0.00496^{+0.00134}_{-0.00099}$ | $0.0049^{+0.0016}_{-0.0012}$
Expected signal | $0.0241^{+0.0086}_{-0.0087}$ | $0.0135^{+0.0048}_{-0.0037}$ | $0.0169^{+0.0042}_{-0.0031}$ | $0.0167^{+0.0050}_{-0.0037}$
| Observed | $556$ | $4$ | $2$ | $1$
5252 – 5272 | Expected comb. bkg | $595.9^{+7.0}_{-6.5}$ | $7.10^{+0.71}_{-0.63}$ | $1.26^{+0.42}_{-0.37}$ | $0.119^{+0.097}_{-0.072}$
Expected peak. bkg | $0.323^{+0.075}_{-0.061}$ | $0.326^{+0.074}_{-0.061}$ | $0.324^{+0.072}_{-0.060}$ | $0.325^{+0.075}_{-0.062}$
Cross-feed | $0.0097^{+0.0036}_{-0.0035}$ | $0.0054^{+0.0021}_{-0.0015}$ | $0.0068^{+0.0018}_{-0.0013}$ | $0.0067^{+0.0021}_{-0.0016}$
Expected signal | $0.045^{+0.016}_{-0.016}$ | $0.0252^{+0.0088}_{-0.0067}$ | $0.0317^{+0.0077}_{-0.0057}$ | $0.0313^{+0.0093}_{-0.0068}$
| Observed | $588$ | $11$ | $1$ | $0$
5272 – 5292 | Expected comb. bkg | $586.9^{+6.7}_{-6.3}$ | $7.04^{+0.68}_{-0.60}$ | $1.23^{+0.41}_{-0.36}$ | $0.117^{+0.092}_{-0.071}$
Expected peak. bkg | $0.252^{+0.058}_{-0.047}$ | $0.252^{+0.056}_{-0.046}$ | $0.253^{+0.059}_{-0.048}$ | $0.250^{+0.056}_{-0.046}$
Cross-feed | $0.0154^{+0.0058}_{-0.0055}$ | $0.0086^{+0.0033}_{-0.0024}$ | $0.0108^{+0.0029}_{-0.0021}$ | $0.0106^{+0.0033}_{-0.0025}$
Expected signal | $0.045^{+0.016}_{-0.016}$ | $0.0251^{+0.0089}_{-0.0067}$ | $0.0317^{+0.0077}_{-0.0057}$ | $0.0313^{+0.0092}_{-0.0069}$
| Observed | $616$ | $5$ | $2$ | $1$
5292 – 5312 | Expected comb. bkg | $578.1^{+6.5}_{-6.1}$ | $6.98^{+0.66}_{-0.58}$ | $1.20^{+0.39}_{-0.35}$ | $0.114^{+0.087}_{-0.067}$
Expected peak. bkg | $0.124^{+0.023}_{-0.021}$ | $0.124^{+0.023}_{-0.021}$ | $0.123^{+0.023}_{-0.021}$ | $0.124^{+0.023}_{-0.021}$
Cross-feed | $0.038^{+0.015}_{-0.014}$ | $0.0214^{+0.0086}_{-0.0061}$ | $0.0270^{+0.0080}_{-0.0056}$ | $0.0266^{+0.0089}_{-0.0064}$
Expected signal | $0.0241^{+0.0086}_{-0.0087}$ | $0.0134^{+0.0048}_{-0.0036}$ | $0.0169^{+0.0042}_{-0.0030}$ | $0.0167^{+0.0050}_{-0.0037}$
| Observed | $549$ | $7$ | $0$ | $0$
5312 – 5332 | Expected comb. bkg | $569.3^{+6.3}_{-5.9}$ | $6.92^{+0.63}_{-0.57}$ | $1.18^{+0.38}_{-0.34}$ | $0.111^{+0.083}_{-0.064}$
Expected peak. bkg | $0.047^{+0.023}_{-0.012}$ | $0.047^{+0.022}_{-0.012}$ | $0.047^{+0.021}_{-0.012}$ | $0.047^{+0.021}_{-0.012}$
Cross-feed | $0.149^{+0.055}_{-0.054}$ | $0.083^{+0.031}_{-0.022}$ | $0.104^{+0.027}_{-0.019}$ | $0.103^{+0.031}_{-0.023}$
Expected signal | $0.0068^{+0.0028}_{-0.0026}$ | $0.0038^{+0.0015}_{-0.0011}$ | $0.0048^{+0.0014}_{-0.0010}$ | $0.0048^{+0.0016}_{-0.0012}$
| Observed | $509$ | $10$ | $1$ | $1$
Table 4: Expected and observed limits on the $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ branching fraction for the 2011 data and for the combination of 2010 and 2011 data. The expected limits are computed allowing the presence of $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ events according to the SM branching fraction. | | at 90% CL | at 95% CL | $\textrm{CL}_{\textrm{b}}$
---|---|---|---|---
2011 | expected limit | $1.1\times 10^{-8}$ | $1.4\times 10^{-8}$ |
| observed limit | $1.3\times 10^{-8}$ | $1.6\times 10^{-8}$ | 0.95
2010+2011 | expected limit | $1.0\times 10^{-8}$ | $1.3\times 10^{-8}$ |
| observed limit | $1.2\times 10^{-8}$ | $1.4\times 10^{-8}$ | 0.93
Table 5: Expected and observed limits on the $B^{0}\rightarrow\mu^{+}\mu^{-}$ branching fraction for 2011 data and for the combination of 2010 and 2011 data. The expected limits are computed in the background only hypothesis. | | at 90% CL | at 95% CL | $\textrm{CL}_{\textrm{b}}$
---|---|---|---|---
2011 | expected limit | $2.5\times 10^{-9}$ | $3.2\times 10^{-9}$ |
| observed limit | $3.0\times 10^{-9}$ | $3.6\times 10^{-9}$ | 0.68
2010+2011 | expected limit | $2.4\times 10^{-9}$ | $3.0\times 10^{-9}$ |
| observed limit | $2.6\times 10^{-9}$ | $3.2\times 10^{-9}$ | 0.61
## 8 Conclusions
With 0.37$\mbox{\,fb}^{-1}$ of integrated luminosity, a search for the rare
decays $B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ and
$B^{0}\rightarrow\mu^{+}\mu^{-}$ has been performed and sensitivities better
than the existing limits have been obtained. The observed events in the
$B^{0}_{s}$ and in the $B^{0}$ mass windows are compatible with the background
expectations at 5% and 32% confidence level, respectively. For the
$B^{0}_{s}\rightarrow\mu^{+}\mu^{-}$ decay, the probability that the observed
events are compatible with the sum of expected background events and signal
events according to the SM rate is 33%. The upper limits for the branching
fractions are evaluated to be
$\displaystyle{\cal B}(B^{0}_{s}\\!\rightarrow\mu^{+}\mu^{-})$
$\displaystyle<$ $\displaystyle 1.3\,(1.6)\times 10^{-8}~{}{\rm
at}~{}90\,\%\,(95\,\%)~{}{\rm CL},$ $\displaystyle{\cal
B}(B^{0}\\!\rightarrow\mu^{+}\mu^{-})$ $\displaystyle<$ $\displaystyle
3.0\,(3.6)\times 10^{-9}~{}{\rm at}~{}90\,\%\,(95\,\%)~{}{\rm CL}.$
The ${\cal B}(B^{0}_{s}\rightarrow\mu^{+}\mu^{-})$ and ${\cal
B}(B^{0}\rightarrow\mu^{+}\mu^{-})$ upper limits have been combined with those
published previously by LHCb [11] and the results are
$\displaystyle{\cal B}(B^{0}_{s}\\!\rightarrow\mu^{+}\mu^{-})(2010+2011)$
$\displaystyle<$ $\displaystyle 1.2\,(1.4)\times 10^{-8}~{}{\rm
at}~{}90\,\%\,(95\,\%)~{}{\rm CL},$ $\displaystyle{\cal
B}(B^{0}\\!\rightarrow\mu^{+}\mu^{-})(2010+2011)$ $\displaystyle<$
$\displaystyle 2.6\,(3.2)\times 10^{-9}~{}{\rm at}~{}90\,\%\,(95\,\%)~{}{\rm
CL}.$
The above 90% (95%) CL upper limits are still about 3.8 (4.4) times the SM
branching fractions for the $B^{0}_{s}$ and 26 (32) times for the $B^{0}$.
These results represent the best upper limits to date.
## Acknowledgements
We express our gratitude to our colleagues in the CERN accelerator departments
for the excellent performance of the LHC. We thank the technical and
administrative staff at CERN and at the LHCb institutes, and acknowledge
support from the National Agencies: CAPES, CNPq, FAPERJ and FINEP (Brazil);
CERN; NSFC (China); CNRS/IN2P3 (France); BMBF, DFG, HGF and MPG (Germany); SFI
(Ireland); INFN (Italy); FOM and NWO (The Netherlands); SCSR (Poland); ANCS
(Romania); MinES of Russia and Rosatom (Russia); MICINN, XuntaGal and GENCAT
(Spain); SNSF and SER (Switzerland); NAS Ukraine (Ukraine); STFC (United
Kingdom); NSF (USA). We also acknowledge the support received from the ERC
under FP7 and the Region Auvergne.
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|
arxiv-papers
| 2011-12-07T15:40:33 |
2024-09-04T02:49:25.067414
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "LHCb Collaboration: R. Aaij, C. Abellan Beteta, B. Adeva, M. Adinolfi,\n C. Adrover, A. Affolder, Z. Ajaltouni, J. Albrecht, F. Alessio, M. Alexander,\n G. Alkhazov, P. Alvarez Cartelle, A. A. Alves Jr, S. Amato, Y. Amhis, J.\n Anderson, R. B. Appleby, O. Aquines Gutierrez, F. Archilli, L. Arrabito, A.\n Artamonov, M. Artuso, E. Aslanides, G. Auriemma, S. Bachmann, J. J. Back, D.\n S. Bailey, V. Balagura, W. Baldini, R. J. Barlow, C. Barschel, S. Barsuk, W.\n Barter, A. Bates, C. Bauer, Th. Bauer, A. Bay, I. Bediaga, S. Belogurov, K.\n Belous, I. Belyaev, E. Ben-Haim, M. Benayoun, G. Bencivenni, S. Benson, J.\n Benton, R. Bernet, M.-O. Bettler, M. van Beuzekom, A. Bien, S. Bifani, T.\n Bird, A. Bizzeti, P. M. Bj{\\o}rnstad, T. Blake, F. Blanc, C. Blanks, J.\n Blouw, S. Blusk, A. Bobrov, V. Bocci, A. Bondar, N. Bondar, W. Bonivento, S.\n Borghi, A. Borgia, T. J. V. Bowcock, C. Bozzi, T. Brambach, J. van den Brand,\n J. Bressieux, D. Brett, M. Britsch, T. Britton, N. H. Brook, H. Brown, A.\n B\\\"uchler-Germann, I. Burducea, A. Bursche, J. Buytaert, S. Cadeddu, O.\n Callot, M. Calvi, M. Calvo Gomez, A. Camboni, P. Campana, A. Carbone, G.\n Carboni, R. Cardinale, A. Cardini, L. Carson, K. Carvalho Akiba, G. Casse, M.\n Cattaneo, Ch. Cauet, M. Charles, Ph. Charpentier, N. Chiapolini, K. Ciba, X.\n Cid Vidal, G. Ciezarek, P. E. L. Clarke, M. Clemencic, H. V. Cliff, J.\n Closier, C. Coca, V. Coco, J. Cogan, P. Collins, A. Comerma-Montells, F.\n Constantin, G. Conti, A. Contu, A. Cook, M. Coombes, G. Corti, G. A. Cowan,\n R. Currie, B. D'Almagne, C. D'Ambrosio, P. David, P. N. Y. David, I. De\n Bonis, S. De Capua, M. De Cian, F. De Lorenzi, J. M. De Miranda, L. De Paula,\n P. De Simone, D. Decamp, M. Deckenhoff, H. Degaudenzi, M. Deissenroth, L. Del\n Buono, C. Deplano, D. Derkach, O. Deschamps, F. Dettori, J. Dickens, H.\n Dijkstra, P. Diniz Batista, F. Domingo Bonal, S. Donleavy, F. Dordei, P.\n Dornan, A. Dosil Su\\'arez, D. Dossett, A. Dovbnya, F. Dupertuis, R.\n Dzhelyadin, A. Dziurda, S. Easo, U. Egede, V. Egorychev, S. Eidelman, D. van\n Eijk, F. Eisele, S. Eisenhardt, R. Ekelhof, L. Eklund, Ch. Elsasser, D.\n Elsby, D. Esperante Pereira, L. Est\\'eve, A. Falabella, E. Fanchini, C.\n F\\\"arber, G. Fardell, C. Farinelli, S. Farry, V. Fave, V. Fernandez Albor, M.\n Ferro-Luzzi, S. Filippov, C. Fitzpatrick, M. Fontana, F. Fontanelli, R.\n Forty, M. Frank, C. Frei, M. Frosini, S. Furcas, A. Gallas Torreira, D.\n Galli, M. Gandelman, P. Gandini, Y. Gao, J-C. Garnier, J. Garofoli, J. Garra\n Tico, L. Garrido, D. Gascon, C. Gaspar, N. Gauvin, M. Gersabeck, T. Gershon,\n Ph. Ghez, V. Gibson, V. V. Gligorov, C. G\\\"obel, D. Golubkov, A. Golutvin, A.\n Gomes, H. Gordon, M. Grabalosa G\\'andara, R. Graciani Diaz, L. A. Granado\n Cardoso, E. Graug\\'es, G. Graziani, A. Grecu, E. Greening, S. Gregson, B.\n Gui, E. Gushchin, Yu. Guz, T. Gys, G. Haefeli, C. Haen, S. C. Haines, T.\n Hampson, S. Hansmann-Menzemer, R. Harji, N. Harnew, J. Harrison, P. F.\n Harrison, J. He, V. Heijne, K. Hennessy, P. Henrard, J. A. Hernando Morata,\n E. van Herwijnen, E. Hicks, K. Holubyev, P. Hopchev, W. Hulsbergen, P. Hunt,\n T. Huse, R. S. Huston, D. Hutchcroft, D. Hynds, V. Iakovenko, P. Ilten, J.\n Imong, R. Jacobsson, A. Jaeger, M. Jahjah Hussein, E. Jans, F. Jansen, P.\n Jaton, B. Jean-Marie, F. Jing, M. John, D. Johnson, C. R. Jones, B. Jost, M.\n Kaballo, S. Kandybei, M. Karacson, T. M. Karbach, J. Keaveney, I. R. Kenyon,\n U. Kerzel, T. Ketel, A. Keune, B. Khanji, Y. M. Kim, M. Knecht, P.\n Koppenburg, A. Kozlinskiy, L. Kravchuk, K. Kreplin, M. Kreps, G. Krocker, P.\n Krokovny, F. Kruse, K. Kruzelecki, M. Kucharczyk, T. Kvaratskheliya, V. N. La\n Thi, D. Lacarrere, G. Lafferty, A. Lai, D. Lambert, R. W. Lambert, E.\n Lanciotti, G. Lanfranchi, C. Langenbruch, T. Latham, C. Lazzeroni, R. Le Gac,\n J. van Leerdam, J.-P. Lees, R. Lef\\'evre, A. Leflat, J. Lefran\\c{c}ois, O.\n Leroy, T. Lesiak, L. Li, L. Li Gioi, M. Lieng, M. Liles, R. Lindner, C. Linn,\n B. Liu, G. Liu, J. H. Lopes, E. Lopez Asamar, N. Lopez-March, H. Lu, J.\n Luisier, A. Mac Raighne, F. Machefert, I. V. Machikhiliyan, F. Maciuc, O.\n Maev, J. Magnin, S. Malde, R. M. D. Mamunur, G. Manca, G. Mancinelli, N.\n Mangiafave, U. Marconi, R. M\\\"arki, J. Marks, G. Martellotti, A. Martens, L.\n Martin, A. Mart\\'in S\\'anchez, D. Martinez Santos, A. Massafferri, Z. Mathe,\n C. Matteuzzi, M. Matveev, E. Maurice, B. Maynard, A. Mazurov, G. McGregor, R.\n McNulty, C. Mclean, M. Meissner, M. Merk, J. Merkel, R. Messi, S.\n Miglioranzi, D. A. Milanes, M.-N. Minard, J. Molina Rodriguez, S. Monteil, D.\n Moran, P. Morawski, R. Mountain, I. Mous, F. Muheim, K. M\\\"uller, R. Muresan,\n B. Muryn, B. Muster, M. Musy, J. Mylroie-Smith, P. Naik, T. Nakada, R.\n Nandakumar, I. Nasteva, M. Nedos, M. Needham, N. Neufeld, C. Nguyen-Mau, M.\n Nicol, V. Niess, N. Nikitin, A. Nomerotski, A. Novoselov, A.\n Oblakowska-Mucha, V. Obraztsov, S. Oggero, S. Ogilvy, O. Okhrimenko, R.\n Oldeman, M. Orlandea, J. M. Otalora Goicochea, P. Owen, K. Pal, J. Palacios,\n A. Palano, M. Palutan, J. Panman, A. Papanestis, M. Pappagallo, C. Parkes, C.\n J. Parkinson, G. Passaleva, G. D. Patel, M. Patel, S. K. Paterson, G. N.\n Patrick, C. Patrignani, C. Pavel-Nicorescu, A. Pazos Alvarez, A. Pellegrino,\n G. Penso, M. Pepe Altarelli, S. Perazzini, D. L. Perego, E. Perez Trigo, A.\n P\\'erez-Calero Yzquierdo, P. Perret, M. Perrin-Terrin, G. Pessina, A.\n Petrella, A. Petrolini, A. Phan, E. Picatoste Olloqui, B. Pie Valls, B.\n Pietrzyk, T. Pila\\v{r}, D. Pinci, R. Plackett, S. Playfer, M. Plo Casasus, G.\n Polok, A. Poluektov, E. Polycarpo, D. Popov, B. Popovici, C. Potterat, A.\n Powell, T. du Pree, J. Prisciandaro, V. Pugatch, A. Puig Navarro, W. Qian, J.\n H. Rademacker, B. Rakotomiaramanana, M. S. Rangel, I. Raniuk, G. Raven, S.\n Redford, M. M. Reid, A. C. dos Reis, S. Ricciardi, K. Rinnert, D. A. Roa\n Romero, P. Robbe, E. Rodrigues, F. Rodrigues, P. Rodriguez Perez, G. J.\n Rogers, S. Roiser, V. Romanovsky, M. Rosello, J. Rouvinet, T. Ruf, H. Ruiz,\n G. Sabatino, J. J. Saborido Silva, N. Sagidova, P. Sail, B. Saitta, C.\n Salzmann, M. Sannino, R. Santacesaria, C. Santamarina Rios, R. Santinelli, E.\n Santovetti, M. Sapunov, A. Sarti, C. Satriano, A. Satta, M. Savrie, D.\n Savrina, P. Schaack, M. Schiller, S. Schleich, M. Schlupp, M. Schmelling, B.\n Schmidt, O. Schneider, A. Schopper, M.-H. Schune, R. Schwemmer, B. Sciascia,\n A. Sciubba, M. Seco, A. Semennikov, K. Senderowska, I. Sepp, N. Serra, J.\n Serrano, P. Seyfert, B. Shao, M. Shapkin, I. Shapoval, P. Shatalov, Y.\n Shcheglov, T. Shears, L. Shekhtman, O. Shevchenko, V. Shevchenko, A. Shires,\n R. Silva Coutinho, T. Skwarnicki, A. C. Smith, N. A. Smith, E. Smith, K.\n Sobczak, F. J. P. Soler, A. Solomin, F. Soomro, B. Souza De Paula, B. Spaan,\n A. Sparkes, P. Spradlin, F. Stagni, S. Stahl, O. Steinkamp, S. Stoica, S.\n Stone, B. Storaci, M. Straticiuc, U. Straumann, V. K. Subbiah, S. Swientek,\n M. Szczekowski, P. Szczypka, T. Szumlak, S. T'Jampens, E. Teodorescu, F.\n Teubert, C. Thomas, E. Thomas, J. van Tilburg, V. Tisserand, M. Tobin, S.\n Topp-Joergensen, N. Torr, E. Tournefier, M. T. Tran, A. Tsaregorodtsev, N.\n Tuning, M. Ubeda Garcia, A. Ukleja, P. Urquijo, U. Uwer, V. Vagnoni, G.\n Valenti, R. Vazquez Gomez, P. Vazquez Regueiro, S. Vecchi, J. J. Velthuis, M.\n Veltri, B. Viaud, I. Videau, X. Vilasis-Cardona, J. Visniakov, A. Vollhardt,\n D. Volyanskyy, D. Voong, A. Vorobyev, H. Voss, S. Wandernoth, J. Wang, D. R.\n Ward, N. K. Watson, A. D. Webber, D. Websdale, M. Whitehead, D. Wiedner, L.\n Wiggers, G. Wilkinson, M. P. Williams, M. Williams, F. F. Wilson, J. Wishahi,\n M. Witek, W. Witzeling, S. A. Wotton, K. Wyllie, Y. Xie, F. Xing, Z. Xing, Z.\n Yang, R. Young, O. Yushchenko, M. Zavertyaev, F. Zhang, L. Zhang, W. C.\n Zhang, Y. Zhang, A. Zhelezov, L. Zhong, E. Zverev, A. Zvyagin",
"submitter": "Gaia Lanfranchi",
"url": "https://arxiv.org/abs/1112.1600"
}
|
1112.1818
|
# Bubble tree of a class of conformal mappings and applications to Willmore
functional
Jingyi Chen & Yuxiang Li
###### Abstract.
We develop a bubble tree construction and prove compactness results for
$W^{2,2}$ branched conformal immersions of closed Riemann surfaces, with
varying conformal structures whose limit may degenerate, in a compact
Riemannian manifold with uniformly bounded areas and Willmore energies. The
compactness property is applied to construct Willmore type surfaces in compact
Riemannian manifolds. This includes (a) existence of a Willmore 2-sphere in
${\mathbb{S}}^{n}$ with at least 2 nonremovable singular points (b) existence
of minimizers of the Willmore functional with prescribed area in a compact
manifold $N$ provided (i) the area is small when genus is 0 and (ii) the area
is close to that of the area minimizing surface of Schoen-Yau and Sacks-
Uhlenbeck in the homotopy class of an incompressible map from a surface of
positive genus to $N$ and $\pi_{2}(N)$ is trivial (c) existence of smooth
minimizers of the Willmore functional if a Douglas type condition is
satisfied.
The first author is partially supported by NSERC. The second author was
partially supported by NSERC during his visit to UBC in the spring of 2011
where most part of this work was done
## 1\. Introduction
Let $\Sigma$ be a smooth Riemann surface and $\
f:\Sigma\rightarrow\mathbb{R}^{n}\ $ be a smooth immersion. The Willmore
functional of $f$ is defined by
$W(f)=\frac{1}{4}\int_{\Sigma}|H_{f}|^{2}d\mu_{f}$
where $H_{f}=\Delta_{g_{f}}f$ denotes the mean curvature vector of $f$, and
$\Delta_{g_{f}}$ is the Laplace operator in the induced metric $g_{f}$ and
$d\mu_{f}$ the induced area element on $\Sigma$.
For a sequence of immersions $f_{k}$ of a compact surfaces $\Sigma$ in
${\mathbb{R}}^{n}$ with uniformly bounded areas $\mu(f_{k})$ and Willmore
functionals $W(f_{k})$, a subsequence of $\Sigma_{k}=f_{k}(\Sigma)$ converges,
as Radon measures, to a two dimensional integral varifold, by Allard’s
integral compactness theorem. The second fundamental forms $A_{f_{k}}$ are
uniformly bounded in the $L^{2}$-norm as
$\int_{\Sigma}|A_{f_{k}}|^{2}d\mu_{f_{k}}=4W(f_{k})-4\pi\chi(\Sigma)$
from the Gauss equation and the Gauss-Bonnet formula. In general
$\|f_{k}\|_{W^{2,2}}$ are not uniformly bounded: we can find diffeomorphisms
$\phi_{k}$ from $\Sigma$ to $\Sigma$ such that $f_{k}=f\circ\phi_{k}$ diverge
in $C^{0}$, while a uniform bound on $\|f_{k}\|_{W^{2,2}}$ would imply
sequential convergence in $C^{0}$ (in fact $C^{\alpha},0<\alpha<1$) norm by
the Rellich-Kondrachov embedding theorem.
A recent advance in understanding the limit process is given in [13], where
each $f_{k}$ is a conformal immersion from a Riemann surface $(\Sigma,h_{k})$
into $\mathbb{R}^{n}$ and $h_{k}$ is the smooth metric of constant curvature:
(1.1) $\begin{array}[]{l}h_{k}\mbox{ has Gauss curvature }\pm 1,\mbox{ or
}(\Sigma,h_{k})=\mathbb{C}/\\{1,a+bi\\}\mbox{ with }\\\
-\frac{1}{2}<a\leq\frac{1}{2},\,\,\,\,b\geq 0,a^{2}+b^{2}\geq 1\mbox{ and
}a\geq 0\mbox{ whenever }a^{2}+b^{2}=1.\end{array}$
There are two reasons to use conformal immersions. One is that the conformal
diffeomorphism group of $(\Sigma,h_{k})$ is rather small comparing with the
group of diffeomorphisms. Secondly, if we set $g_{f_{k}}=e^{2u_{k}}g_{euc}$ in
an isothermal coordinate system, then we can estimate $\|u_{k}\|_{L^{\infty}}$
from the compensated compactness property of $K_{f_{k}}e^{2u_{k}}$. Thus it is
possible to get an upper bound of $\|f_{k}\|_{W^{2,2}}$ via the equation
$\Delta_{h_{k}}f_{k}=H_{f_{k}}$. When the conformal structures determined by
$f_{k}$ do not go to the boundary of the moduli space, convergence of $f_{k}$
is treated in [13]: if the conformal classes induced by $f_{k}$ converge in
the moduli space, then there exist Mobius transformations $\sigma_{k}$, such
that $\sigma_{k}\circ f_{k}$ converge locally in weak $W^{2,2}$ sense on
$\Sigma$ minus finitely many concentration points. The weak limit $f_{0}$ is a
$W^{2,2}$ branched conformal immersion.
The $W^{2,2}$ conformal immersions and $W^{2,2}$ branched conformal immersions
are as follows:
###### Definition 1.
Let $(\Sigma,h)$ be a connected Riemann surface with $h$ satisfies (1.1). A
map $f\in W^{2,2}(\Sigma,h,\mathbb{R}^{n})$ is called a conformal immersion of
$(\Sigma,h)$, if
$df\otimes
df=e^{2u}h\,\,\,\,\hbox{with}\,\,\,\,\|u\|_{L^{\infty}(\Sigma)}<+\infty.$
We denote the set of all such immersions by
$W^{2,2}_{conf}(\Sigma,h,\mathbb{R}^{n})$. It can be shown that for $f\in
W^{2,2}_{conf}(\Sigma,\mathbb{R}^{n})$ the corresponding $u$ is continuous
(see Appendix). When $f\in W^{2,2}_{loc}(\Sigma,h,\mathbb{R}^{n})$ with
$df\otimes df=e^{2u}h$ and $u\in L^{\infty}_{loc}(\Sigma)$, we say $f\in
W^{2,2}_{conf,loc}(\Sigma,h,\mathbb{R}^{n})$.
###### Definition 2.
We say $f$ is a $W^{2,2}$ branched conformal immersion of $(\Sigma,h)$ with
possible branch points $x_{1},\dots,x_{m}$, if $f\in
W^{2,2}_{conf,loc}(\Sigma\backslash\\{x_{1},\dots,x_{m}\\},h,\mathbb{R}^{n})$
and
$\int_{\Sigma\backslash\\{x_{1},\dots,x_{m}\\}}(1+|A_{f}|^{2})d\mu_{f}<+\infty.$
The set of $W^{2,2}$ branched conformal immersions is denoted by
$W^{2,2}_{b,c}(\Sigma,h,\mathbb{R}^{n})$. We say
$f\in\widetilde{W}^{2,2}(\Sigma,\mathbb{R}^{n})$, if we can find a smooth
metric $h$ satisfying (1.1) over $\Sigma$, such that $f\in
W^{2,2}_{b,c}(\Sigma,h,\mathbb{R}^{n})$.
The first part of the paper is a study of a sequence of $W^{2,2}$ branched
conformal immersions and the main goal is to establish compactness in
Hausdorff distance for such immersions with uniformly bounded areas and
Willmore functionals (cf. Theorem 1).
Our compactness result holds not only when $h_{k}$ converges, but also when
the conformal classes $c_{k}$ of $h_{k}$ diverge in the moduli space
$\mathcal{M}_{g}$. Bubbles develop near points where the Willmore energy
concentrates, and if $c_{k}$ go to a point in the boundary
$\overline{\mathcal{M}}_{g}\backslash\mathcal{M}_{g}$ additional complication
arises as the topology of the limit may be different from that of $\Sigma$ and
stratified surfaces are used as possible limits. The main idea to deal with
degenerating conformal structures in the limit process is as follows. First,
pull the immersions $f_{k}$ to the immersions from components of $\Sigma_{0}$,
by composing $f_{k}$ with diffeomorphisms from the regular parts of the limit
$\Sigma_{0}$ of $(\Sigma,h_{k})$ in $\overline{\mathcal{M}_{p}}$. Then we
study the limit of $f_{k}$ and construct bubble trees at the energy
concentration points and collars, and investigate behavior between bubbles. In
particular, we will prove that there is no loss of measure in the limit and
there is no neck between the bubbles. Then the limit $f_{0}$ of $f_{k}$ is a
union of conformal maps from some components
$\overline{\Sigma_{0}^{1}},\dots,\overline{\Sigma_{0}^{m}}$ of $\Sigma_{0}$
(we delete those components whose images are points) and finitely many
2-spheres $S_{1},\dots,S_{l}$ into $\mathbb{R}^{n}$. “no neck” means that we
can glue $\overline{\Sigma_{0}^{i}}$’s and $S_{j}$’s to a stratified surface
$\Sigma_{\infty}$ (see definition below), and $f_{0}$ is a continuous map from
$\Sigma_{\infty}$ into $\mathbb{R}^{n}$. Then we will apply a result of Hélein
[10] and a removable singularity theorem in [13] to show that for a sequence
of branched conformal immersions with uniformly bounded measures and Willmore
functionals, the limit we get in section 2 is in fact a branched conformal
immersion of a stratified surface.
We point out that the “no loss of measure” and “no neck” phenomenon must occur
whenever the following two equations hold:
(1.2) $-\Delta f_{k}=\frac{1}{2}|\nabla
f_{k}|^{2}H_{k},\,\,\,\,\mbox{with}\,\,\,\,\sup_{k}\int|\nabla
f_{k}|^{2}(1+|H_{k}|^{2})<\infty,$ (1.3) $\partial f_{k}\otimes\partial
f_{k}=0\,\,\,(\hbox{weakly conformal})$
where $\Delta,\nabla,\partial=\partial/\partial z$ are the operators in
$h_{k}$ and its conformal structure $c_{k}$. In section 2, we study the blowup
behavior of a sequence satisfies (1.2) and (1.3).
The equation (1.2) looks similar to the equation of harmonic maps
$-\Delta u=A(u)(du,du).$
In fact, the arguments in section 2 are originated from the “energy identity”
and “no neck” arguments of harmonic maps [2, 5, 16, 18, 19, 23, 24, 25, 28].
When conformal structures go to the boundary of ${\mathcal{M}}_{g}$, non-
trivial necks exist for harmonic map ( [2, 3, 23, 34]); in our case, however,
there is no non-trivial neck due to conformality although (1.2) is much weaker
than the harmonic map equation.
###### Definition 3.
Let $(\Sigma,d)$ be a connected compact metric space. We say $\Sigma$ is a
stratified surface with singular set $P$ if $P\subset\Sigma$ is finite set
such that
1\. $(\Sigma\backslash P,d)$ is a smooth Riemann surface without boundary
(possibly disconnected) and $d$ is a smooth metric $h=d|_{\Sigma\backslash
P}$, and
2\. For each $p\in P$, there is $\delta$, such that $B_{\delta}(p)\cap
P=\\{p\\}$ and
$B_{\delta}(p)\backslash\\{p\\}=\bigcup\limits_{i=1}^{m(p)}\Omega_{i}$, where
$1<m(p)<+\infty$, and each $\Omega_{i}$ is topologically a disk with its
center deleted. Moreover, on each $\Omega_{i}$, $h$ can be extended to be a
smooth metric on the disk.
The genus of $\Sigma$ is defined by
$g(\Sigma)=\frac{2-\chi(\Sigma)+\sum\limits_{p\in P}(m(p)-1)}{2}.$
When $g(\Sigma)=0$, $\Sigma$ is called a stratified sphere. A stratified
surface with singular set $P=\emptyset$ is a smooth Riemann surface.
Figure 1. Stratified torus
For a stratified surface $\Sigma$ with singular set $P$, we can write
$\Sigma\backslash P=\bigcup_{i}\Sigma^{i}$ where $\Sigma^{i}$’s are the
disjoint connected components of $\Sigma$, and each $\Sigma^{i}$ is a
punctured Riemann surface when there are more than one components. The
topological closure of $\Sigma^{i}$ is denoted by $\Sigma_{i}$, so as a point-
set $\Sigma=\bigcup_{i}\Sigma_{i}$. By 2 in the above definition, each
component $\Sigma^{i}$ can be extended to a closed Riemann surface
$\overline{\Sigma^{i}}$ by adding finitely many points. To illustrate the
difference of these notations, take, for example, the stratified torus on the
left in Figure 1: $P$ contains two points, $\Sigma^{1}$ is the “torus” with
two points deleted and $\Sigma^{2}$ is a 2-sphere with one point removed,
$\Sigma_{1}$ is the “torus” and $\Sigma_{2}$ is the 2-sphere, while
$\overline{\Sigma^{1}}$ is a Riemann sphere (adding 3 points at the punctures)
and $\overline{\Sigma^{2}}$ is also a Riemann sphere (adding 1 point at the
puncture).
When $\Sigma$ is a stratified surface we define $f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n})$ if $f$ is a $W^{2,2}$ non-trivial
branched conformal immersion on each $\overline{\Sigma^{i}}$.
We now state the main result in the first part of the paper:
###### Theorem 1.
Suppose that $\\{f_{k}\\}$ is a sequence of $W^{2,2}$ branched conformal
immersions of $(\Sigma,h_{k})$ in $\mathbb{R}^{n}$, where $h_{k}$ satisfies
(1.1). If $f_{k}(\Sigma)\cap B_{R_{0}}\neq\emptyset$ for some fixed $R_{0}$
and
$\sup_{k}\left(\mu(f_{k})+W(f_{k})\right)<+\infty$
then either $f_{k}$ converges to a point, or there is a stratified surface
$\Sigma_{\infty}$ with $g(\Sigma_{\infty})\leq g(\Sigma)$, a map $f_{0}\in
W^{2,2}_{b,c}(\Sigma_{\infty},\mathbb{R}^{n})$, such that a subsequence of
$f_{k}(\Sigma)$ converges to $f_{0}(\Sigma_{\infty})$ in Hausdorff distance
with
$\mu(f_{0})=\lim_{k\rightarrow+\infty}\mu(f_{k})\,\,\,\,and\,\,\,\,W(f_{0})\leq\lim_{k\rightarrow+\infty}W(f_{k}).$
Moreover, if $y_{1},\dots,y_{m}\in f_{k}(\Sigma)$ for all $k$, then
$y_{1},\dots,y_{m}\in f_{0}(\Sigma_{\infty}).$
In fact, we will prove that $f_{k}$ converges to $f_{0}$ in the sense of
bubble tree. For each $k$, we can find a domain $U_{k}$ of $\Sigma$ and a
domain $V_{k}$ of $\Sigma_{\infty}$, such that
1) $V_{k}\subset V_{k+1}$, and $P=\Sigma_{\infty}\backslash\bigcup_{k}V_{k}$
is a finite set which contains all singular points of $\Sigma_{\infty}$.
Moreover, $\Sigma_{\infty}\backslash V_{k}$ is a union of topological disks
and $H^{1}_{1}(\Sigma_{\infty}\backslash V_{k})\rightarrow 0$, where
$H^{1}_{1}$ is the Hausdorff measure:
$H_{1}^{1}(S)=\inf\left\\{\sum_{i=1}^{\infty}\mbox{\rm
diam}(U_{i}):S\subset\bigcup_{i=1}^{\infty}U_{i},\,\,\,\,\mbox{\rm
diam}(U_{i})<1\right\\}.$
2) $\Sigma\backslash U_{k}$ is a smooth surface with boundary, possibly
disconnected, $H^{1}_{1}(f_{k}(\Sigma\backslash U_{k}))\rightarrow 0$.
Moreover, $f_{k}(\Sigma\backslash U_{k})$ converge to $P$ in Hausdorff
distance.
3) There is a sequence of diffeomorphisms $\phi_{k}:V_{k}\rightarrow U_{k}$,
such that for any $\Omega\subset\subset\Sigma_{\infty}\backslash P$,
$f_{k}\circ\phi_{k}$ converge in $W^{2,2}(\Omega,\mathbb{R}^{n})$ weakly.
In Theorem 1, the singular points of $f_{0}(\Sigma_{\infty})$ arise in three
ways: (a) the limit point to which a sequence of closed geodesics that are not
null-homotopic in $f_{k}(\Sigma)$ pinches, (b) a bubble point of $f_{k}$, so
belonging to a 2-sphere (the bubble), (c) a point where both (a) and (b)
happen.
In the second part of the paper, we apply Theorem 1 to obtain several
existence results of Willmore surfaces in compact Riemannian manifolds. Here
we note that Theorem 1 is applicable for surfaces immersed in a compact
Riemannian manifold $N$. To see this, for $\Sigma$ immersed in $N$ which is
isometrically embedded in $\mathbb{R}^{n}$, direct calculation shows that the
Willmore functional of $\Sigma$ in $\mathbb{R}^{n}$ is dominated by its
Willmore functional in $N$ together with the area $\mu(\Sigma)$, see Lemma
4.1.
We first consider 2-spheres immersed in the round unit sphere
${\mathbb{S}}^{n},n\geq 3$. Fix at least two distinct points
$y_{1},\dots,y_{m},m\geq 2$ on ${\mathbb{S}}^{n}$. Define
$\beta^{n}_{0}(y_{1},\dots,y_{m})=\inf\left\\{W_{n}(f):f\in
W^{2,2}_{conf}(S^{2},{\mathbb{S}}^{n}),y_{1},\dots,y_{m}\in
f({S}^{2})\right\\}$
where
$W_{n}(f)=\int_{S^{2}}\left(1+\frac{1}{4}\left|H_{f}\right|^{2}\right)d\mu_{f}$
and $H_{f}$ is the mean curvature vector of $f(S^{2})$ in ${\mathbb{S}}^{n}$.
We show
###### Theorem 2.
If $\beta^{n}_{0}(y_{1},\dots,y_{m})<8\pi$, then there is a $W^{2,2}$
conformal immersion of $S^{2}$ in ${\mathbb{S}}^{n}$ without self-
intersections realizing $\beta^{n}_{0}(y_{1},\dots,y_{m})$. For any
$\epsilon>0$, there exists a Willmore sphere in ${\mathbb{S}}^{n}$ with
$W_{n}(f)<4\pi+\epsilon$, which has at least 2 nonremovable singularities.
By results in [15], [27], a singular point of a Willmore surface with density
$\theta^{2}<2$ in $\mathbb{R}^{n}$ can be removed if its residue is 0. Kuwert
and Schätzle also point out that the removability can not be true generally,
for example, 0 is the true singular point of an inverted half catenoid ([15],
P. 337). The second statement in Theorem 2 provides examples of embedded
Willmore surface which has a nonremovable singular point with density
$\theta^{2}=1$, and it is an application of the first statement with five
points prescribed in ${\mathbb{S}}^{n}$.
We then consider minimizers of the Willmore functional subject to area
constraint. A fundamental existence result for incompressible minimal surfaces
due to Schoen-Yau [30] and Sacks-Uhlenbeck [29] asserts: If $\varphi$ induces
an injection from the fundamental groups to $\Sigma$ and $N$, then there is a
branched minimal immersion $f:\Sigma\rightarrow N$ so that $f$ induces the
same map between fundamental groups as $\varphi$ and $f$ has least area among
all such maps. We denote the area of the minimizer by $a_{\varphi}$.
###### Theorem 3.
Let $N$ be a compact Riemannian manifold and $\Sigma$ be a closed surface of
genus $g$. Then
(1) For $\beta_{0}(N,a)=\inf\\{W(f):\mu(f)=a>0,f\in
W^{2,2}_{conf}(S^{2},N)\\}$, $\lim_{a\to 0}\beta_{0}(N,a)=4\pi$, and there is
an embedding realizing $\beta_{0}(N,a)$ for all sufficiently small $a$.
(2) Suppose $\varphi:\Sigma\to N$ induces an injection
$\varphi_{\\#}:\pi_{1}(\Sigma)\to\pi_{1}(N)$ and $\pi_{2}(N)=0$. Let
$\beta_{g}(N,a,\varphi)=\inf\\{W(f):f\in\widetilde{W}^{2,2}(\Sigma,N),\mu(f)=a,f\hbox{
is homotopic to}\,\,\,\varphi\\}$. Then there is $\delta>0$, such that for any
$a\in[a_{\varphi},a_{\varphi}+\delta)$, there is a branched conformal
immersion $f$ of $(\Sigma,h)$ attaining $\beta_{g}(N,a,\varphi)$. Moreover,
when $\dim N=3$, $f$ is an immersion for small $\delta$.
For $\beta_{0}(N,a)$, Lamm and Metzger showed in [17] that if it is attained
by a surface with positive mean curvature in the sufficiently small geodesic
ball around a point $p$, then the scalar curvature of $N$ must have a critical
point at $p$.
When $N$ has negative sectional curvature, the area of an immersed surface is
dominated by the Willmore functional. We now describe a sufficient condition
of Douglas type for existence. Let $S(g)$ be the set of connected stratified
Riemann surfaces $\Sigma=\bigcup_{i}\Sigma_{i}$ satisfying (a) genus of
$\Sigma_{i}<g$ if $g>0$ and (b) $i>1$ if $g=0$. Note that a surface in $S(g)$
has genus at most $g$ and smooth surfaces of genus $g$ are not in $S(g)$.
Isometrically embed $N$ into ${\mathbb{R}}^{n}$. Define
$\alpha^{*}(g)=\inf\\{W(f):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),\Sigma\in{S}(g)\\}$
$\alpha(g)=\inf\\{W(f):f\in
W^{2,2}_{b,c}(\Sigma,{\mathbb{R}}^{n}),\hbox{$\Sigma$ is a smooth surface of
genus $g$}\\}.$
Similarly, for $0<a<\infty$, define
$\displaystyle\gamma^{*}(g,a)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset N,\Sigma\in
S(g),\mu(f(\Sigma))\leq a\\}$ $\displaystyle\gamma(g,a)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset
N,\Sigma\in{\mathcal{M}}_{g},\mu(f(\Sigma))\leq a\\}.$
###### Theorem 4.
Let $N$ be a compact Riemannian manifold. If $0<\alpha(g)<\alpha^{*}(g)$ and
$N$ has negative sectional curvature, then there is a $W^{2,2}$ branched
conformal immersion $f$ from a smooth closed Riemann surface of genus $g$ with
$W(f)=\alpha(g)$. If $0<\gamma(g,a)<\gamma^{*}(g,a)$ then there is a $W^{2,2}$
branched conformal immersion $f$ from a smooth closed Riemann surface of genus
$g$ with $W(f)=\gamma(g)$.
## 2\. blowup analysis - energy identity and absence of neck
Let $(\Sigma,h)$ be a smooth Riemann surface which may not be compact, where
$h$ is the metric compatible with the complex structure of $\Sigma$. For given
$p>1$ and $R>0$, let $\mathcal{F}^{p}(\Sigma,h,R)$ be the set of mappings
$f:\Sigma\rightarrow\mathbb{R}^{n}$ which satisfy
1. (1)
$f\in W^{2,p}_{loc}(\Sigma,h)$;
2. (2)
$f(\Sigma)$ is contained in the closed ball centered at the origin with radius
$R$ in ${\mathbb{R}}^{n}$;
3. (3)
$\Delta_{h}f=F(f)$ with $|F(f)|\leq\beta\,|\nabla_{h}f|^{2}$ a.e. on $\Sigma$,
where $\beta$ is a nonnegative measurable function on $\Sigma$ with
$\int_{\Sigma}\beta^{2}|\nabla_{h}f|^{2}d\mu_{h}<+\infty.$
When $f\in\mathcal{F}^{p}(\Sigma,h,R)$, we introduce a notation by
$H(f)=\left\\{\begin{array}[]{ll}2\frac{\Delta_{h}f}{|\nabla_{h}f|^{2}},&\hbox{if
$|\nabla_{h}f|\neq 0$}\\\ 0,&\hbox{if $|\nabla_{h}f|=0$}.\end{array}\right.$
The Willmore functional of $f$ is defined to be
$W(f)=\frac{1}{4}\int_{\Sigma}|H(f)|^{2}|\nabla_{h}f|^{2}d\mu_{h}.$
That $W(f)<\infty$ for $f\in\mathcal{F}^{p}(\Sigma,h,R)$ follows from (3) as
$\Delta_{h}f=F(f)=\frac{1}{2}H(f)|\nabla_{h}f|^{2}.$
We denote by $\mathcal{F}^{p}_{conf}(\Sigma,h,R)$ the set of
$f\in\mathcal{F}^{p}(\Sigma,h,R)$ and $f$ is weakly conformal a.e., i.e.
$\partial f\otimes\partial f=0$ almost everywhere on $\Sigma$, where $\partial
f=\frac{\partial f}{\partial z}dz$ in a local complex coordinate system on
$\Sigma$.
Note that when $f$ is a smooth conformal immersion
$H(f),\frac{1}{2}|\nabla_{h}f|^{2}d\mu_{h},W(f)$ are the mean curvature
vector, the area element and the Willmore functional of $f(\Sigma)$,
respectively.
By the Kondrachov embedding theorem, functions in
${\mathcal{F}}^{p}(\Sigma,h,R)$ are also in $W^{1,2}$. The right hand side of
the equation $\Delta_{h}f=F(f)$ is not necessarily in $L^{2}$ under the
assumption (3).
We point out that $H(f)$, $\mathcal{F}^{p},\mathcal{F}^{p}_{conf}$ are
conformal invariant, in the sense that if $h^{\prime}=e^{2u}h$ for some smooth
function $u$ on $\Sigma$, we always have
$H_{h}(f)=H_{h^{\prime}}(f),\,\,\,\,\mathcal{F}^{p}(\Sigma,h,R)=\mathcal{F}^{p}(\Sigma,h^{\prime},R),\,\,\,\,\mathcal{F}^{p}_{conf}(\Sigma,h,R)=\mathcal{F}^{p}_{conf}(\Sigma,h^{\prime},R).$
Thus we may select preferred metrics $h$, e.g. the ones with constant
curvature.
In this section, we will study regularity, compactness and the blowup behavior
of a sequence $\\{f_{k}\\}\subset\mathcal{F}^{p}$.
### 2.1. $\epsilon$-regularity, removable singularity and weak limit
In this subsection, we will show that some well-known results for harmonic
maps still hold for mappings in $\mathcal{F}^{p}$.
Let $D$ be the unit 2-disk centered at 0. For simplicity, write
$\mathcal{F}^{p}(D,dx^{2}+dy^{2},R)$ as $\mathcal{F}^{p}(D,R)$.
###### Proposition 2.1.
($\epsilon$-regularity) There is an $\epsilon_{0}$ such that for any
$f\in\mathcal{F}^{p}(D,R)$, $1<p<2$, if $W(f)<\epsilon_{0}^{2}$, then
$\|\nabla f\|_{W^{1,p}(D_{\frac{1}{2}})}\leq C\,\|\nabla f\|_{L^{2}(D)}.$
###### Proof.
Set $\bar{f}=\frac{1}{|D|}\int_{D}fd\sigma$ and let $\eta$ be a cut-off
function which is 1 in $D_{1/2}$, 0 in $D\backslash D_{3/4}$ and
$0\leq\eta\leq 1$. Then for the equation
$\Delta\left(\eta(f-\bar{f})\right)=(f-\bar{f})\Delta\eta+2\nabla\eta\nabla
f+\frac{1}{2}\eta H(f)|\nabla f|^{2}:=\phi$
we have
$\displaystyle|\phi|$ $\displaystyle\leq$ $\displaystyle
C_{1}\left(|f-\bar{f}|+|\nabla
f|\right)+\frac{1}{2}\eta\left|H(f)\right||\nabla f|^{2}$ $\displaystyle\leq$
$\displaystyle C_{1}\left(|f-\bar{f}|+|\nabla
f|\right)+C_{2}\left|H(f)\right||\nabla
f|\left(|\nabla\left(\eta(f-\bar{f})\right)|+|f-\bar{f}|\right)$
since
$\begin{array}[]{lll}\frac{1}{2}\eta\left|H(f)\right||\nabla
f|^{2}&=&\frac{1}{2}\eta\left|H(f)\right|\nabla(f-\bar{f})\nabla
f\\\\[8.61108pt]
&=&\frac{1}{2}\left|H(f)\right|\nabla\left(\eta(f-\bar{f})\right)\nabla
f-\frac{1}{2}\left|H(f)\right|(f-\bar{f})\nabla\eta\nabla f\\\\[8.61108pt]
&\leq&C_{2}\left|H(f)\right||\nabla
f|\left(|\nabla\left(\eta(f-\bar{f})\right)|+|f-\bar{f}|\right).\end{array}$
By the $L^{p}$ estimates for elliptic equations,
$\displaystyle\left\|\eta(f-\bar{f})\right\|_{W^{2,p}(D)}$ $\displaystyle\leq$
$\displaystyle C_{3}\left(\left\|f-\bar{f}\right\|_{L^{p}(D)}+\left\|\nabla
f\right\|_{L^{p}(D)}\right.$ $\displaystyle\left.+\left\|H(f)|\nabla
f|\left(|\nabla\left(\eta(f-\bar{f})\right)|+|f-\bar{f}|\right)\right\|_{L^{p}(D)}\right).$
For $1<p<2$, the Hölder inequality and the Sobolev inequality imply
$\begin{array}[]{lll}&&\left\|H(f)|\nabla
f|\left(\left|\nabla\left(\eta(f-\bar{f})\right)\right|+\left|f-\bar{f}\right|\right)\right\|_{L^{p}(D)}\\\
&&\leq\|H(f)\nabla
f\|_{L^{2}(D)}\left(\|\nabla\left(\eta(f-\bar{f})\right)\|_{L^{\frac{2p}{2-p}}(D)}+\|f-\bar{f}\|_{L^{\frac{2p}{2-p}}(D)}\right)\\\\[8.61108pt]
&&\leq\epsilon_{0}\,C_{4}\,\|\eta(f-\bar{f})\|_{W^{2,p}(D)}+\epsilon_{0}\,C_{5}\,\|f-\bar{f}\|_{W^{1,p}(D)}\end{array}$
since $W(f)<\epsilon_{0}$. Applying the Poincaré inequality and noting
$1<p<2$, we get
$\|f-\bar{f}\|_{L^{p}(D)}+\|\nabla
f\|_{L^{p}(D)}+\epsilon_{0}\,C_{5}\,\|f-\bar{f}\|_{W^{1,p}(D)}\leq
C_{6}\,\|\nabla f\|_{L^{2}(D)}.$
Choose $\epsilon_{0}$ so that $C_{3}C_{4}\,\epsilon_{0}<1/2$, then we get
$\|\eta(f-\bar{f})\|_{W^{2,p}(D)}<C_{7}\,\|\nabla f\|_{L^{2}(D)}$
which completes the proof. $\hfill\Box$
###### Proposition 2.2.
(Gap constant) Let $\Sigma$ be a closed surface. There is a constant
${\epsilon}_{1}$ which depends on $\Sigma$ and $R$, such that for any
$f\in\mathcal{F}^{p}(\Sigma,h,R)$ where $1<p<2$, if $W(f)<{\epsilon}_{1}^{2}$,
then $f$ is constant.
###### Proof.
Let $\bar{f}=\frac{1}{|\Sigma|}\int_{\Sigma}f$. It follows from the equation
$\Delta_{h}(f-\bar{f})=\frac{1}{2}H(f)|\nabla f|^{2}$
that
$\begin{array}[]{lll}\displaystyle\int_{\Sigma}|\nabla(f-\bar{f})|^{2}&\leq&\displaystyle\frac{1}{2}\int_{\Sigma}|f-\bar{f}|\,|{H(f)}|\,|\nabla
f|^{2}\\\\[8.61108pt] &\leq&\displaystyle\left(\int_{\Sigma}{H(f)}^{2}|\nabla
f|^{2}\right)^{\frac{1}{2}}\left(\int_{\Sigma}|f-\bar{f}|^{\frac{2p}{2p-2}}\right)^{\frac{2p-2}{2p}}\left(\int_{\Sigma}|\nabla
f|^{\frac{2p}{2-p}}\right)^{\frac{2-p}{2p}}\\\\[8.61108pt] &\leq&\displaystyle
C_{1}W(f)^{\frac{1}{2}}\|\nabla f\|_{L^{2}(\Sigma)}\|\nabla
f\|_{L^{\frac{2p}{2-p}}(\Sigma)}\end{array}$
where we used the Sobolev inequality, the Poincaré inequality and $1<p<2$.
Then we get
$\|\nabla f\|_{L^{2}(\Sigma)}\leq C_{1}W(f)^{\frac{1}{2}}\|\nabla
f\|_{L^{\frac{2p}{2-p}}(\Sigma)}.$
Using the Poincaré inequality and $1<p<2$ again, we have
$\|f-\bar{f}\|_{L^{p}(\Sigma)}\leq C_{2}\|\nabla f\|_{L^{2}(\Sigma)}\leq
C_{1}C_{2}\,W(f)^{\frac{1}{2}}\|\nabla f\|_{L^{\frac{2p}{2-p}}(\Sigma)}.$
Since
$\left\|\frac{1}{2}H(f)|\nabla
f|^{2}\right\|_{L^{p}(\Sigma)}\leq\left(\int_{\Sigma}\frac{1}{4}{H(f)}^{2}|\nabla
f|^{2}\right)^{\frac{1}{2}}\left(\int_{\Sigma}|\nabla
f|^{\frac{2p}{2-p}}\right)^{\frac{2-p}{2p}}=W(f)^{\frac{1}{2}}\|\nabla
f\|_{L^{\frac{2p}{2-p}}(\Sigma)},$
it follows from the $L^{p}$ estimates for elliptic equations that
$\displaystyle\left\|f-\bar{f}\right\|_{W^{2,p}(\Sigma)}$ $\displaystyle\leq$
$\displaystyle C_{3}\,\left(\left\|\frac{1}{2}H(f)|\nabla
f|^{2}\right\|_{L^{p}(\Sigma)}+\left\|f-\bar{f}\right\|_{L^{p}(\Sigma)}\right)$
$\displaystyle\leq$ $\displaystyle
C_{3}(1+C_{1}C_{2})W(f)^{\frac{1}{2}}\left\|\nabla
f\right\|_{L^{\frac{2p}{2-p}}(\Sigma)}$ $\displaystyle\leq$ $\displaystyle
C_{4}W(f)^{\frac{1}{2}}\left\|f-\bar{f}\right\|_{W^{2,p}(\Sigma)}$
where the Sobolev inequality was used in the last step. By choosing
$\epsilon_{1}<1/C_{4}$ we immediately have $f=\overline{f}$. $\hfill\Box$
We now derive a key estimate for later applications. Set
$E(f,Q(t))=\int_{Q(t)}|\nabla
f|^{2},\,\,\hbox{where}\,\,Q(t)=S^{1}\times[-t,t]$
and denote $\mathcal{F}^{p}(Q(t),dt^{2}+d\theta^{2},R)$ by
$\mathcal{F}^{p}(Q(t),R)$. We will prove the following energy decay estimate:
###### Proposition 2.3.
(Decay estimate) Let $f\in\mathcal{F}^{p}_{conf}(Q(T),R)$ with $T\gg 1,1<p<2$.
Then there is a constant $\epsilon_{2}<\epsilon_{0}$, where $\epsilon_{0}$ is
the constant in Proposition 2.1, such that if
$\sup_{t\in[-T,T-1]}W(f,S^{1}\times[t,t+1])<\epsilon^{2}\leq\epsilon_{2}^{2}$
then
$\int_{Q(t)}|\nabla f|^{2}<{C}E(f,Q(T))e^{-(1-C\epsilon)(T-t)}$
for some positive constant $C$ independent of $T$ and $f$.
###### Proof.
Define
$f^{*}(t)=\frac{1}{2\pi}\int_{0}^{2\pi}f(t,\theta)d\theta.$
We have
$\displaystyle\int_{Q(t)}\left|\frac{\partial f^{*}}{\partial t}\right|^{2}$
$\displaystyle=$
$\displaystyle\int_{-t}^{t}\int_{0}^{2\pi}\left(\frac{1}{2\pi}\int_{0}^{2\pi}\frac{\partial
f}{\partial t}d\theta\right)^{2}d\theta dt$ $\displaystyle\leq$
$\displaystyle\displaystyle\frac{1}{2\pi}\int_{-t}^{t}\left(\int_{0}^{2\pi}\left|\frac{\partial
f}{\partial t}\right|^{2}d\theta\int_{0}^{2\pi}d\theta\right)dt$
$\displaystyle=$
$\displaystyle\int_{-t}^{t}\int_{0}^{2\pi}\left|\frac{\partial f}{\partial
t}\right|^{2}dtd\theta$ $\displaystyle=$
$\displaystyle\int_{Q(t)}\left|\frac{\partial f}{\partial
t}\right|^{2}dtd\theta.$
Then
(2.2) $\begin{array}[]{lll}\displaystyle\int_{Q(t)}\nabla(f-f^{*})\nabla
f&=&\displaystyle\int_{Q(t)}|\nabla
f|^{2}-\displaystyle\int_{Q(t)}\frac{\partial f}{\partial t}\frac{\partial
f^{*}}{\partial t}\\\\[8.61108pt] &\geq&\displaystyle\int_{Q(t)}|\nabla
f|^{2}-\frac{1}{2}\left(\displaystyle\int_{Q(t)}\left|\frac{\partial
f}{\partial t}\right|^{2}+\displaystyle\int_{Q(t)}\left|\frac{\partial
f^{*}}{\partial t}\right|^{2}\right)\\\\[8.61108pt]
&\geq&\displaystyle\int_{Q(t)}|\nabla f|^{2}-\int_{Q(t)}\left|\frac{\partial
f}{\partial t}\right|^{2}\\\\[8.61108pt]
&=&\displaystyle\frac{1}{2}\int_{Q(t)}|\nabla f|^{2}\end{array}$
where in the last step we used the fact that $|f_{t}|^{2}=|f_{\theta}|^{2}$
a.e. as $f$ is conformal a.e. On the other hand,
(2.3) $\begin{array}[]{lll}\displaystyle\int_{Q(t)}\nabla(f-f^{*})\nabla
f&=&-\displaystyle\int_{Q(t)}(f-f^{*})\Delta f-\displaystyle\int_{\partial
Q(t)}\frac{\partial f}{\partial t}(f-f^{*})\\\\[8.61108pt]
&\leq&\displaystyle\int_{Q(t)}|f-f^{*}||\nabla
f|^{2}\frac{|{H(f)}|}{2}+\left|\displaystyle\int_{\partial Q(t)}\frac{\partial
f}{\partial t}(f-f^{*})\right|.\end{array}$
Let $m\in[t,t+1)$ be an integer. Then for each $i=-m,-m+1,\dots,m-1$, by (2.1)
and the hypothesis in the proposition
$\sup_{t\in[-T,T-1]}\frac{1}{4}\int_{S^{1}\times[t,t+1]}|H(f)|^{2}<\epsilon^{2}\leq\epsilon_{0}^{2}$
it follows from Proposition 2.1 that
(2.4) $\|f-f^{*}\|_{L^{\infty}(S^{1}\times[i,i+1])}\leq C\,\|\nabla
f\|_{L^{2}(S^{1}\times[i-1,i+2])}.$
In fact, to see (2.4), denote the average of $f$ over $S\times[i-1,i+2]$ by
$\overline{f}$ and observe that from Proposition 2.1
$\|\nabla(f-\overline{f})\|_{W^{1,p}(S^{1}\times[i,i+1])}=\|\nabla
f\|_{W^{1,p}(S^{1}\times[i,i+1])}\leq C\|\nabla
f\|_{L^{2}(S^{1}\times[i-1,i+2])}$
and from the Poincaré inequality
$\|f-\overline{f}\|_{L^{2}(S^{1}\times[i,i+1])}\leq\|f-\overline{f}\|_{L^{2}(S^{1}\times[i-1,i+2])}\leq
C\|\nabla f\|_{L^{2}(S^{1}\times[i-1,i+2])}$
hence
$\|f-\overline{f}\|_{W^{2,p}(S^{1}\times[i,i+1])}\leq C\|\nabla
f\|_{L^{2}(S^{1}\times[i-1,i+2])}.$
The Sobolev embedding theorem then implies
$\|f-\overline{f}\|_{C^{0}(S^{1}\times[i,i+1])}\leq
C\|f-\overline{f}\|_{W^{2,p}(S^{1}\times[i,i+1])}\leq C\|\nabla
f\|_{L^{2}(S^{1}\times[i-1,i+2])}.$
Therefore for any $t\in[i,i+1]$
$\left|f^{*}(t)-\overline{f}\right|=\left|\frac{1}{2\pi}\int^{2\pi}_{0}\left(f(t,\theta)-\overline{f}\right)d\theta\right|\leq
C\|\nabla f\|_{L^{2}(S^{1}\times[i-1,i+2])}.$
It follows
$\|f-f^{*}\|_{C^{0}(S^{1}\times[i,i+1])}\leq\|f-\bar{f}\|_{C^{0}(S^{1}\times[i,i+1])}+\|f^{*}-\bar{f}\|_{C^{0}(S^{1}\times[i,i+1])}\leq
C\|\nabla f\|_{L^{2}(S^{1}\times[i-1,i+2])}.$
Clearly, by the mean value theorem, $f-f^{*}$ equals 0 somewhere in
$S^{1}\times[i,i+1]$, thus
$\|f-f^{*}\|_{L^{\infty}(S^{1}\times[i,i+1])}=\|f-f^{*}\|_{C^{0}(S^{1}\times[i,i+1])}\leq
C\|\nabla f\|_{L^{2}(S^{1}\times[i-1,i+2])}$
which shows (2.4) holds.
Then
$\begin{array}[]{l}\displaystyle\int_{S^{1}\times[i,i+1]}|f-f^{*}||\nabla
f|^{2}\frac{|{H(f)}|}{2}\\\\[8.61108pt]
\begin{array}[]{lll}&\leq&\displaystyle\left\|f-f^{*}\right\|_{L^{\infty}(S^{1}\times[i,i+1])}\times\left(W(f_{k},S^{1}\times[i,i+1])\int_{S^{1}\times[i,i+1]}|\nabla
f|^{2}\right)^{\frac{1}{2}}\\\\[8.61108pt] &\leq&\displaystyle
C\epsilon\,\left(\int_{S^{1}\times[i-1,i+2]}|\nabla
f|^{2}\int_{S^{1}\times[i,i+1]}|\nabla
f|^{2}\right)^{\frac{1}{2}}\\\\[8.61108pt] &\leq&\displaystyle
C\epsilon\int_{S^{1}\times[i-1,i+2]}|\nabla f|^{2}.\end{array}\end{array}$
Then
(2.5) $\begin{array}[]{lll}\displaystyle{\int}_{Q(t)}|f-f^{*}||\nabla
f|^{2}{H(f)}&\leq&\displaystyle\sum\limits_{i=-m}^{m-1}\int_{S^{1}\times[i,i+1]}|f-f^{*}||\nabla
f|^{2}{H(f)}\\\\[8.61108pt] &\leq&\displaystyle
C\epsilon\sum\limits_{i=-m}^{m-1}\int_{S^{1}\times[i-1,i+2]}|\nabla
f|^{2}\\\\[8.61108pt] &\leq&\displaystyle 3C\epsilon\int_{Q(t)}|\nabla
f|^{2}\\\\[8.61108pt]
&&\displaystyle+C\epsilon\left(\int_{S^{1}\times[-m-1,-m]}|\nabla
f|^{2}+\int_{S^{1}\times[m,m+1]}|\nabla f|^{2}\right)\\\\[8.61108pt]
&\leq&3C\epsilon\displaystyle\int_{Q(t+2)}|\nabla f|^{2}.\end{array}$
From (2.2), (2.3), (2.5), we have
(2.6) $\frac{1}{2}\int_{Q(t)}|\nabla
f|^{2}\leq\epsilon^{\prime}\int_{Q(t+2)}|\nabla f|^{2}+\left|\int_{\partial
Q(t)}(f-f^{*})\frac{\partial f}{\partial t}\right|$
where $\epsilon^{\prime}=3C\epsilon/2$. Moreover,
(2.7)
$\begin{array}[]{lll}\displaystyle\left|\int_{S^{1}\times\\{t\\}}\frac{\partial
f}{\partial
t}(f-f^{*})\right|&\leq&\displaystyle\left(\int_{0}^{2\pi}(f(\theta,t)-f^{*}(t))^{2}d\theta\right)^{\frac{1}{2}}\left(\int_{0}^{2\pi}\left|\frac{\partial
f}{\partial t}(\theta,t)\right|^{2}d\theta\right)^{\frac{1}{2}}\\\\[8.61108pt]
&\leq&\displaystyle\left(\int_{0}^{2\pi}\left|\frac{\partial
f}{\partial\theta}(\theta,t)\right|^{2}d\theta\right)^{\frac{1}{2}}\left(\int_{0}^{2\pi}\left|\frac{\partial
f}{\partial t}(\theta,t)\right|^{2}d\theta\right)^{\frac{1}{2}}\\\\[8.61108pt]
&=&\displaystyle{\frac{1}{2}}\displaystyle{\int}_{S^{1}\times\\{t\\}}|\nabla
f|^{2}d\theta\end{array}$
here we used the Poincaré inequality on $S^{1}$ and the fact that
$|\frac{\partial f}{\partial t}|^{2}=|\frac{\partial f}{\partial\theta}|^{2}$
a.e. Let
$\varphi(t)=\frac{1}{2}\int_{Q(t)}|\nabla f|^{2}.$
By (2.6) and (2.7), we have
$\varphi(t)\leq\varphi^{\prime}(t)+\epsilon^{\prime}\varphi(t+2).$
Then
$-(e^{-t}\varphi(t))^{\prime}\leq\epsilon^{\prime}\varphi(t+2)e^{-t},$
and integrating the inequality from $t$ to $T-2$ leads to
(2.8) $\begin{array}[]{lll}\displaystyle e^{-t}\varphi(t)&\leq&\displaystyle
e^{-T+2}\varphi(T-2)+\epsilon^{\prime}\int_{t}^{T-2}\varphi(s+2)e^{-s}ds\\\\[8.61108pt]
&=&\displaystyle
e^{-T+2}\varphi(T-2)+\epsilon^{\prime}\int_{t+2}^{T}\varphi(s)e^{-s+2}ds\\\\[8.61108pt]
&=&\displaystyle
e^{-T+2}\varphi(T-2)+\epsilon^{\prime}e^{2}\int_{t+2}^{T-2}\varphi(s)e^{-s}ds+\epsilon^{\prime}e^{2}\int_{T-2}^{T}\varphi(s)e^{-s}ds\\\\[8.61108pt]
&\leq&\displaystyle
e^{-T+2}\varphi(T)+\epsilon^{\prime}e^{2}\int_{t}^{T-2}\varphi(s)e^{-s}ds+\epsilon^{\prime}e^{2}\varphi(T)\left(e^{-T+2}-e^{-T}\right)\end{array}$
as $\varphi(t)$ is increasing in $t$. Let
$F(t)=\int_{t}^{T-2}\varphi(s)e^{-s}ds$ and
$\epsilon_{2}=\epsilon^{\prime}e^{2}<1$. Now (2.8) leads to
$-F^{\prime}(t)\leq 2\,\varphi(T)e^{-T+2}+\epsilon_{2}F(t)$
or equivalently
$\left(e^{\epsilon_{2}t}F(t)\right)^{\prime}+2\varphi(T)e^{-T+2}e^{\epsilon_{2}t}\geq
0.$
Integrating over $[t,T-2]$ and noting $F(T-2)=0$, we have
(2.9)
$F(t)\leq\frac{2\varphi(T)}{\epsilon_{2}}e^{2-T}\left(e^{\epsilon_{2}(T-2)}-e^{\epsilon_{2}t}\right)e^{-\epsilon_{2}t}.$
Substitute (2.9) into (2.8):
$\displaystyle\varphi(t)$ $\displaystyle\leq$ $\displaystyle
e^{2-T+t}\varphi(T)+2\varphi(T)e^{\epsilon_{2}(T-t-2)}e^{2-T+t}+\epsilon_{2}\varphi(T)e^{t-T+2}$
$\displaystyle\leq$ $\displaystyle C\varphi(T)e^{(1-\epsilon_{2})(T-t)}$
$\displaystyle=$ $\displaystyle C\varphi(T)e^{(1-C\epsilon)(T-t)}$
for some positive constant $C$ independent of $T$ and $f$. $\hfill\Box$
###### Proposition 2.4.
(Removability of point singularity) Let
$f\in\mathcal{F}^{p}_{conf}(D\backslash\\{0\\},R)$, where $1<p<2$. If
$\int_{D}|\nabla f|^{2}<+\infty$, then
$f\in\mathcal{F}^{p^{\prime}}_{conf}(D,R)$ for any
$p^{\prime}\in(1,\frac{4}{3})\cap(1,p]$.
###### Proof.
We may assume that $W(f)<\epsilon^{2}<\epsilon_{2}^{2}$, otherwise, we can
replace $f$ with $f(\lambda x)$ for some $\lambda<1$. Let
$\phi:\mathbb{R}^{1}\times S^{1}\rightarrow\mathbb{R}^{2}$ be the conformal
mapping given by $r=e^{-t},\theta=\theta$. Then $f^{\prime}=f(\phi)$ is a map
from $[0,+\infty)\times S^{1}$ into $\mathbb{R}^{n}$. By translating
$S^{1}\times[t-1,t+1]\subset S^{1}\times[0,2t]$ to $S^{1}\times[-1,1]\subset
S^{1}\times[-t,t]$, from Proposition 2.3 we conclude
$\int_{S^{1}\times[t-1,t+1]}|\nabla f^{\prime}|^{2}\leq C_{1}e^{-\delta
t},\,\,\,\,\hbox{where}\,\,\,\,\delta=1-C\epsilon.$
Then for any $r_{k}=e^{-k}$, we have $t_{k}=k$ and
(2.10) $\int_{D_{r_{k-1}}\backslash D_{r_{k+1}}}|\nabla
f|^{2}<C_{1}r_{k}^{-\delta}.$
Set $f_{k}(x)=f(r_{k}x)$. Applying Proposition 2.1 and (2.10), we get
$\|\nabla f_{k}\|_{W^{1,p}(D_{1}\backslash D_{e^{-1}})}\leq C_{2}\,\|\nabla
f_{k}\|_{L^{2}(D_{e}\backslash D_{e^{-2}})}\leq
C_{3}\,r_{k}^{\frac{\delta}{2}}.$
By the Sobolev inequality, we have
$\left(\int_{D_{1}\backslash D_{e^{-1}}}|\nabla
f_{k}|^{q}\right)^{\frac{1}{q}}\leq C_{4}\,\|\nabla
f_{k}\|_{W^{1,p}(D_{1}\backslash D_{e^{-1}})}\leq
C_{5}\,r_{k}^{-\frac{\delta}{2}},\,\,\,\,where\,\,\,\,q\leq\frac{2p}{2-p}.$
Then
$\int_{D_{1}\backslash D_{e^{-1}}}|\nabla f_{k}|^{q}\leq
C_{6}\,e^{-qk\frac{\delta}{2}}.$
Since
$r_{k}^{2-q}\int_{D_{1}\backslash D_{e^{-1}}}|\nabla
f_{k}|^{q}=\int_{D_{r_{k}}\backslash D_{r_{k+1}}}|\nabla f|^{q},$
we have
$\int_{D_{r_{k}}\backslash D_{r_{k+1}}}|\nabla f|^{q}\leq
C_{6}\,e^{-qk\frac{\delta}{2}+(q-2)k}=C_{6}\,e^{k(-2+q(1-\frac{\delta}{2}))}.$
When $q<4$, we can choose $\epsilon$ suitably such that
$q(1-\frac{\delta}{2})<2$, which yields
$\int_{D}|\nabla f|^{q}\leq
C_{6}\sum_{k}2^{-qk\frac{\delta}{2}+(q-2)k}<C_{7}<\infty.$
For any $p^{\prime}\in(1,\frac{4}{3})$, set
$q=\frac{2p^{\prime}}{2-p^{\prime}}$, so $q\in(2,4)$. We have
$\int_{D}{H(f)}^{p^{\prime}}|\nabla
f|^{2p^{\prime}}\leq\left(\int_{D}{H(f)}^{2}|\nabla
f|^{2}\right)^{\frac{p^{\prime}}{2}}\left(\int_{D}|\nabla
f|^{q}\right)^{\frac{p^{\prime}}{q}}<C_{8}.$
Therefore, $F(f)\in L^{p^{\prime}}(D)$ with $p^{\prime}>1$ and then there
exists $v$ which solves the equation
$-\Delta v=F(f),\,\,\,\,v|_{\partial D}=0,$
and $v\in W^{2,p^{\prime}}(D)$. Obviously, $f-v$ is a harmonic function on
$D\backslash\\{0\\}$ with
$\|\nabla(f-v)\|_{L^{2}(D)}+\|f-v\|_{L^{2}(D)}<+\infty.$
Then $f-v$ is smooth on $D$. Now $f\in{\mathcal{F}}^{p^{\prime}}_{conf}(D,R)$
is evident for $p^{\prime}\leq p$ and $1<p^{\prime}<\frac{4}{3}$. $\hfill\Box$
We now consider weak compactness property of a bounded sequence in
${\mathcal{F}}^{p}(D,R)$. Let $\\{f_{k}\\}\subset\mathcal{F}^{p}(D,R)$. The
blowup set of $\\{f_{k}\\}$ is defined to be
$\mathcal{C}(\\{f_{k}\\})=\left\\{z\in D:\lim_{r\rightarrow
0}\varliminf_{k\rightarrow+\infty}W(f_{k},D_{r}(z))>\epsilon_{2}^{2}\right\\}.$
Then for any $z\in D\backslash\mathcal{C}(\\{f_{k}\\})$, we can find $r$ and a
subsequence of $\\{f_{k}\\}$ which is still denoted by $\\{f_{k}\\}$ for
simplicity, such that
$\lim_{k\rightarrow+\infty}W(f_{k},D_{r}(z))<\epsilon_{0}^{2}.$
Then we get from Proposition 2.1 that
$\|f_{k}\|_{W^{2,p}(D_{r/2}(z))}<C\|\nabla f_{k}\|_{L^{2}(D_{r})}$. Thus we
may assume $f_{k}$ converges weakly in
$W^{2,p}(D\backslash\mathcal{C}(\\{f_{k}\\}))$.
###### Corollary 2.5.
Let $\\{f_{k}\\}\subset\mathcal{F}^{p}_{conf}(D,R)$ with
$\sup_{k}\\{E(f_{k},D)+W(f_{k},D)\\}<\Lambda<\infty$
and $f_{0}$ be the weak limit of $f_{k}$ in
$W^{2,p}_{loc}(D\backslash\mathcal{C}(\\{f_{k}\\}))$. If
$p\in(1,\frac{4}{3})$, then $f_{0}\in\mathcal{F}^{p}_{conf}(D,R)$ and
(2.11) $W(f_{0},D)\leq\varliminf_{k\rightarrow+\infty}W(f_{k},D).$
###### Proof.
Set $\Delta f_{k}=F_{k},k\in{\mathbb{N}}$. For any $\Omega\subset\subset
D\backslash\mathcal{C}(f_{k})$, we have
$\|f_{k}\|_{W^{2,p}(\Omega)}<C(\Omega)$. Then by the Hölder inequality and the
Sobolev inequality
$\|F_{k}\|_{L^{p}(\Omega)}\leq\left\|\frac{1}{2}H(f_{k})\nabla
f_{k}\right\|_{L^{2}(\Omega)}\left\|\nabla
f_{k}\right\|_{L^{\frac{2p}{2-p}}(\Omega)}\leq C\Lambda^{\frac{1}{2}}\|\nabla
f_{k}\|_{W^{1,p}(\Omega)}<C^{\prime}(\Omega,\Lambda).$
We may assume, by selecting subsequences if necessary, that
$F_{k}\rightharpoonup
F_{0}\,\,\,\,\hbox{in}\,\,\,\,L^{p}(\Omega)\,\,\,\,\hbox{and}\,\,\,\,|H(f_{k})||\nabla
f_{k}|\rightharpoonup\alpha\,\,\,\,\hbox{in}\,\,\,\,L^{2}.$
Since we may also assume $|\nabla f_{k}|\rightarrow|\nabla f_{0}|$ in
$L^{2}(\Omega)$ because $f_{k}\rightharpoonup f_{0}$ in $W^{2,p}(\Omega)$, we
have
$|H(f_{k})||\nabla f_{k}|^{2}\rightharpoonup\alpha|\nabla f_{0}|$
in the sense of measures in $\Omega$. Define
$\beta_{0}=\left\\{\begin{array}[]{ll}\frac{\alpha}{|\nabla f_{0}|}&\hbox{
when}\,|\nabla f_{0}|\neq 0\\\ 0&\,\,\hbox{otherwise.}\end{array}\right.$
Clearly, $\beta_{0}|\nabla f_{0}|^{2}=\alpha|\nabla f_{0}|$. Let
$F_{k}^{+}=\max\\{F_{k},0\\}$ and $F_{k}^{-}=-\min\\{F_{k},0\\}$. Then
$F_{k}=F_{k}^{+}-F_{k}^{-}$ and $|F_{k}|=F_{k}^{+}+F_{k}^{-}$. We may assume
that
$F_{k}^{+}\rightharpoonup
F_{0}^{1}\,\,\,\,\hbox{and}\,\,\,\,F_{k}^{-}\rightharpoonup
F_{0}^{2}\,\,\,\,\hbox{in}\,\,\,\,L^{p}(\Omega).$
Obviously $F_{0}=F_{0}^{1}-F_{0}^{2}$. Then for any nonnegative function
$\varphi\in C_{0}^{\infty}(\Omega)$,
$\int_{\Omega}\varphi|F_{0}|\leq\int_{\Omega}\varphi(F_{0}^{1}+F_{0}^{2})=\lim_{k\rightarrow+\infty}\int_{\Omega}\varphi|F_{k}|\leq\lim_{k\rightarrow+\infty}\int_{\Omega}\frac{1}{2}\varphi|H(f_{k})||\nabla
f_{k}|^{2}=\int_{\Omega}\frac{1}{2}\varphi\beta_{0}|\nabla f_{0}|^{2}.$
Hence we conclude
$|F_{0}|\leq\frac{1}{2}\beta_{0}|\nabla f_{0}|^{2},\,\,\,\,\hbox{a.e.}\,z\in
D.$
Then, we have
$\int_{\Omega}\beta_{0}^{2}|\nabla
f_{0}|^{2}\leq\int_{\Omega}\alpha^{2}\leq\varliminf_{k\rightarrow+\infty}\int_{\Omega}|H_{k}(f_{k})|^{2}|\nabla
f_{k}|^{2}.$
Moreover, as $f_{k}$ converge in $L^{2}(\Omega)$, it follows from $\partial
f_{k}\otimes\partial f_{k}=0$ a.e. in $D$ that $\partial f_{0}\otimes\partial
f_{0}=0$ a.e. in $D$ as $\Omega$ is arbitrary. Since
$\sup_{k}\\{E(f_{k})+W(f_{k})\\}<\infty$, there are at most finitely many
points in ${\mathcal{C}}(f_{k})$. Then we conclude that
$f_{0}\in\mathcal{F}^{p}_{conf}(D,R)$ if $p\in(1,\frac{4}{3})$ by removing the
point singularity across $\mathcal{C}(f_{k})$ ensured by Proposition 2.4.
Furthermore, we have $H(f_{0})\leq\beta_{0}$ whenever $|\nabla f_{0}|\neq 0$,
hence we get (2.11). $\hfill\Box$
### 2.2. A criterion for absence of bubbles along cylinders
Let $f_{k}\in\mathcal{F}^{p}_{conf}(Q(T_{k}),R)$, with
$\sup_{k}\\{E(f_{k})+W(f_{k})\\}<\Lambda<\infty.$
Given a sequence $t_{k}\in(-T_{k},T_{k})$ with
(2.12)
$T_{k}-t_{k}\rightarrow+\infty\,\,\,\hbox{and}\,\,\,t_{k}-(-T_{k})\rightarrow+\infty,$
we say the limit $f_{0}$ of a subsequence of $f_{k}(\theta,t+t_{k})$, as in
Corollary 2.5, is nontrivial if $E(f_{0})>0$. When $f_{0}$ is nontrivial, it
is a bubble of $\\{f_{k}\\}$.
###### Proposition 2.6.
Let $f_{k}\in{\mathcal{F}}^{p}_{conf}(S^{1}\times(-T_{k},T_{k}),R)$ with
$\sup_{k}\\{E(f_{k})+W(f_{k})\\}=\Lambda<\infty.$
Let $\epsilon_{2}$ be the constant in Proposition 2.3. If
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[-T_{k}+T,T_{k}-T]}W(f_{k},S^{1}\times[t,t+1])<\epsilon_{2}^{2},$
then we have the following
1. (1)
$\\{f_{k}\\}$ has no bubble;
2. (2)
there is no energy loss, i.e.
(2.13)
$\lim_{T\rightarrow+\infty}\lim_{k\rightarrow+\infty}\int_{S^{1}\times[-T_{k}+T,T_{k}-T]}|\nabla
f_{k}|^{2}=0,$
3. (3)
there is no neck, i.e.
(2.14)
$\lim_{t\rightarrow+\infty}\lim_{k\rightarrow+\infty}f_{k}(\theta,-T_{k}+t)=\lim_{t\rightarrow+\infty}\lim_{k\rightarrow+\infty}f_{k}(\theta,T_{k}-t).$
###### Proof.
We may assume that $f_{k}(\theta,-T_{k}+t)$ and $f_{k}(\theta,T_{k}-t)$
converge to $f_{0}^{+}(\theta,t)$ and $f_{0}^{-}(\theta,t)$ weakly in
$W^{2,p}_{loc}(S^{1}\times[0,+\infty))$, respectively. Then $f_{0}^{+}(\phi)$
and $f_{0}^{-}(\phi)\in\mathcal{F}^{p}(D\backslash\\{0\\},R)$ with
$E(f^{\pm}_{0}(\phi))+W(f^{\pm}_{0}(\phi))\leq\Lambda$
where $\phi$ is the conformal diffeomorphism between $D\backslash\\{0\\}$ with
$S^{1}\times(0,+\infty)$. By removability of point singularities asserted in
Proposition 2.4, they are in $\mathcal{F}^{p^{\prime}}(D,R)$ for some
$p^{\prime}>1$. It then follows from the compact embedding
$W^{2,p^{\prime}}\subset L^{2}$:
$\lim_{T\to\infty}\int_{S^{1}\times[T,T+1]}\left(|\nabla
f_{0}^{+}|^{2}+|\nabla f_{0}^{-}|^{2}\right)=0.$
Define $f_{k}^{*}(t)=\frac{1}{2\pi}\int_{0}^{2\pi}f_{k}(t,\theta)d\theta$. It
is easy to check that
$\lim_{t\rightarrow+\infty}\lim_{k\rightarrow+\infty}\left|\int_{\partial
Q(T_{k}-t)}(f_{k}-f_{k}^{*})\frac{\partial f_{k}}{\partial t}\right|=0.$
In fact, this can be seen as follows:
$\sup_{S^{1}\times\\{T_{k}-T\\}}|f_{k}-f_{k}^{*}|\leq
C\,\underset{S^{1}\times\\{T_{k}-T\\}}{\hbox{osc}}f_{k}$
which will converge to $C\hbox{osc}_{S^{1}\times\\{-T\\}}f_{0}^{+}$ as
$k\to\infty$. By removability of singularity,
$\lim_{T\to\infty}\hbox{osc}_{S^{1}\times\\{-T\\}}f_{0}^{+}=0.$
By the Sobolev trace embedding,
$\int_{S^{1}\times\\{T_{k}-T\\}}|\nabla f_{k}|\leq C\|\nabla
f_{k}\|_{W^{1,p}(S^{1}\times[T_{k}-T-1,T_{k}-T+1])}$
By $\epsilon$-regularity,
$\left\|\nabla f_{k}\right\|_{W^{1,p}(S^{1}\times[T_{k}-T-1,T_{k}-T+1])}\leq
C\|\nabla f_{k}\|_{L^{2}(S^{1}\times[T_{k}-T-2,T_{k}-T+2])}<C.$
Then (2.13) follows from (2.6).
Let $m_{k}$ be the integer in $[T_{k}-T,T_{k}-T+1)$. For $0\leq i\leq
m_{k}-2$, applying Proposition 2.3 on $S^{1}\times[i-m_{k},m_{k}]$ (by
shifting the center circle to $S^{1}\times\\{i\\}$, and the same below), we
have
$\int_{S^{1}\times[i-2,i+2]}|\nabla
f_{k}|^{2}<CE(f_{k},Q(T_{k}-T))e^{-\delta(m_{k}-i)},\,\,\,\,\delta=1-C\epsilon_{2}.$
Then from (2.4)
$\underset{S^{1}\times[i-1,i+1]}{\hbox{osc}}f_{k}\leq
C\sqrt{E(f_{k},Q(T_{k}-T))}e^{-\frac{\delta}{2}(m_{k}-i)}.$
When $-m_{k}+2\leq i\leq 0$, applying Proposition 2.3 on
$S^{1}\times[-m_{k},m_{k}+i]$, we get
$\int_{S^{1}\times[i-2,i+2]}|\nabla
f_{k}|^{2}<CE(f_{k},Q(T_{k}-T))e^{-\frac{\delta}{2}(m_{k}-|i|)},$
then we obtain
$\underset{S^{1}\times[i-1,i+1]}{\hbox{osc}}f_{k}\leq
C\sqrt{E(f_{k},Q(T_{k}-T))}e^{-\frac{\delta}{2}(m_{k}-|i|)}.$
Hence,
$\underset{Q(T_{k}-T)}{\hbox{osc}}f_{k}\leq
2C\sqrt{E(f_{k},Q(T_{k}-T))}\sum_{i=1}^{m_{k}}e^{-\frac{\delta}{2}(m_{k}-i)}\leq
C^{\prime}\sqrt{E(f_{k},Q(T_{k}-T))}.$
Then (2.14) can be deduced from (2.13). $\hfill\Box$
### 2.3. Bubble trees for a sequence of maps from the disk $D$
Let $f_{k}\in\mathcal{F}^{p}_{conf}(D,R)$ with
$\sup_{k}\left\\{E(f_{k},D)+W(f_{k},D)\right\\}=\Lambda<\infty.$
We assume $0$ is the only blowup point of $\\{f_{k}\\}$, i.e. the only point
such that
$\lim_{r\rightarrow
0}\varliminf_{k\rightarrow+\infty}W(f_{k},D_{r}(z))\geq\epsilon_{2}^{2}.$
We assume that $f_{k}$ converges to $f_{\infty}$ weakly in
$W^{2,p}_{loc}(D\backslash\\{0\\})$. The construction of the bubble tree at
$0$ will be divided into the following steps:
Step 1. Construct the first level of the bubble tree.
There exists a sequence of points $z_{k}\in D$ and a sequence of radii
$r_{k}\to 0$ such that
(2.15) $W(f_{k},D_{r_{k}}(z_{k}))=\frac{\epsilon_{2}^{2}}{2}$
and $W(f_{k},D_{r}(z))<{\epsilon_{2}^{2}/2}$ for any $r<r_{k}$ and
$D_{r}(z)\subset D$. It is easy to check that $z_{k}\rightarrow 0$ as $0$ is
the only blowup point of $\\{f_{k}\\}$.
We set $f_{k}^{\prime}(z)=f_{k}(z_{k}+r_{k}z)$. Since
$\mathcal{C}(\\{f_{k}^{\prime}\\})=\emptyset$, $f_{k}^{\prime}(z)$ converge
weakly in $W^{2,p}_{loc}(\mathbb{C})$. We denote the limit by $f^{F}$. Note
that it may be a trivial mapping.
Let $(r,\theta)$ be the polar coordinates centered at $z_{k}$, and set
$T_{k}=-\ln r_{k}$. Let
$\phi_{k}:S^{1}\times[0,T_{k}]\rightarrow\mathbb{R}^{2}$ be the conformal
mapping given by $\phi_{k}(t,\theta)=(e^{-t},\theta).$ Then
$\phi_{k}^{*}(dx^{1}\otimes dx^{1}+dx^{2}\otimes
dx^{2})=\frac{1}{r^{2}}(dt^{2}+d\theta^{2}).$
Thus $f_{k}\circ\phi_{k}\in\mathcal{F}^{p}_{conf}(S^{1}\times[0,T_{k}],R)$. We
will also denote $f_{k}\circ\phi_{k}$ by $f_{k}$ for simplicity of notations.
###### Lemma 2.7.
There exists a subsequence of $\\{f_{k}\\}$ and
$0=d_{k}^{0}<d_{k}^{1}<\cdots<d_{k}^{l}=T_{k}$ with
$l<{\Lambda/\epsilon_{2}^{2}}+1$, such that
(2.16) $\lim_{k\rightarrow+\infty}d_{k}^{j}-d_{k}^{j-1}=\infty,$ (2.17)
$W(f_{k},S^{1}\times[d_{k}^{j},d_{k}^{j}+1])\geq\epsilon_{2}^{2},\,\,\,\,j\neq
0,l$
and
(2.18)
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[d_{k}^{j-1}+T,d_{k}^{j}-T]}W(f_{k},S^{1}\times[t,t+1])\leq\epsilon_{2}^{2},\,\,\,\,j=1,...,l.$
###### Proof.
Suppose
$(m-1)\epsilon_{2}^{2}<W(f_{k},S^{1}\times[0,T_{k}])\leq\epsilon_{2}^{2}\,m,$
where $m$ is a positive integer. We prove the lemma by induction on $m$.
When $m=1$, the lemma is obvious by taking $d^{0}_{k}=0,d^{1}_{k}=T_{k}$ and
(2.17) is vacuous. Assuming the lemma is true for $m-1$, we will prove it also
true for $m$. First of all, if
(2.19)
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[T,T_{k}-T]}W(f_{k},S^{1}\times[t,t+1])\leq\epsilon_{2}^{2},$
then the lemma follows since $[d^{j-1}_{k}+T,d^{j}_{k}-T]\subset[T,T_{k}-T]$.
If (2.19) does not hold, we can find $t_{k}$ such that
$t_{k}\rightarrow+\infty,\,\,\,\,T_{k}-t_{k}\rightarrow+\infty,$
and
$W(f_{k},S^{1}\times[t_{k},t_{k}+1])\geq\epsilon_{2}^{2}.$
Then
$W(f_{k},S^{1}\times[0,t_{k}])\leq\epsilon_{2}^{2}\,(m-1)\,\,\,\,\hbox{and}\,\,\,\,W(f_{k},S^{1}\times[t_{k}+1,T_{k}])\leq\epsilon_{2}^{2}\,(m-1).$
Using the induction hypothesis on $[0,t_{k}]$ and $[t_{k}+1,T_{k}]$, we can
find
$0=\bar{d}_{k}^{0}<\bar{d}_{k}^{1}<\cdots<\bar{d}_{k}^{\bar{l}}=t_{k},\,\,\,\,\hbox{and}\,\,\,\,t_{k}+1=\widehat{d}_{k}^{0}<\widehat{d}_{k}^{1}<\cdots<\widehat{d}_{k}^{\hat{l}}=T_{k},$
such that
$\bar{d}_{k}^{i}-\bar{d}_{k}^{i-1}\rightarrow+\infty,\,\,\,\,\,\,\,\,\widehat{d}_{k}^{i}-\widehat{d}_{k}^{i-1}\rightarrow+\infty,$
$W(f_{k},S^{1}\times[\bar{d}_{k}^{j},\bar{d}_{k}^{j}+1])\geq\epsilon_{2}^{2},\,\,\,\,W(f_{k},S^{1}\times[\widehat{d}_{k}^{j},\widehat{d}_{k}^{j}+1])\geq\epsilon_{2}^{2},$
and
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[\bar{d}_{k}^{j-1}+T,\bar{d}_{k}^{j}-T]}W(f_{k},S^{1}\times[t,t+1])\leq\epsilon_{2}^{2},$
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[\widehat{d}_{k}^{j-1}+T,\widehat{d}_{k}^{j}-T]}W(f_{k},S^{1}\times[t,t+1])\leq\epsilon_{2}^{2}.$
Put
$d_{k}^{i}=\left\\{\begin{array}[]{ll}\bar{d}_{k}^{i}&i\leq\bar{l},\\\
\widehat{d}_{k}^{i-\bar{l}}&i>\bar{l}.\end{array}\right.$
The induction is complete. $\hfill\Box$
We now start to construct the bubble tree at the first level. In Lemma 2.7, if
$l=1$, in view of Proposition 2.6, we do not do anything as there is no bubble
developing in $S^{1}\times[0,T_{k}]$ when $k\to\infty$. If $l>1$, we set
$f_{k}^{i}(\theta,t)=f_{k}(\theta,d_{k}^{i}+t)$. We may assume
$\\{f_{k}^{i}\\}$ converges weakly in $W^{2,p}$ to a bubble $f_{\infty}^{i}$
in any compact set outside the blowup points of $\\{f_{k}^{i}\\}$. By
Proposition 2.6, there is no other bubble of $f_{k}$ between $f_{\infty}^{i}$
and $f_{\infty}^{i+1}$ and $f_{\infty}^{i}\cup f_{\infty}^{i+1}$ is connected.
Clearly, $\\{f_{k}^{0}\\}$ and $\\{f_{k}^{l}\\}$ have no blowup points.
Moreover $f_{\infty}^{0}$ is just a part of $f_{\infty}$ and $f_{\infty}^{l}$
is just a part of $f^{F}$. Removing the point singularity by Proposition 2.4,
$f_{\infty}^{1}$, $\cdots$, $f_{\infty}^{l-1}$ and $f^{F}$ can be considered
as conformal mappings from $S^{2}$ into $\mathbb{R}^{n}$.
Figure 2. Bubble tree: First level (dots denote concentration points)
For a stratified sphere, we can define a dual graph as following: 1) Associate
one vertex for each component of the stratified sphere; 2) Vertices are
connected by edges if the corresponding components meet at a point.
Let $S_{1}$ be the stratified sphere with $l$ components whose dual graph is a
tree, i.e. no loops. We define $F^{1}$ to the continuous map from $S_{1}$ into
$\mathbb{R}^{n}$, such that $F^{1}$ is $f_{\infty}^{i}$ on the $i$-th
component when $i<l$ and $f^{F}$ on the $l$-th component. We call $F^{1}$ the
first level of bubble tree of $\\{f_{k}\\}$.
We define $E(F^{1})$ and $W(F^{1})$ by
$E(F^{1})=\sum_{i=1}^{l-1}\int_{S^{1}\times\mathbb{R}}|\nabla
f_{\infty}^{i}|^{2}+\int_{S^{2}}|\nabla
f^{F}|,\,\,\,\,W(F^{1})=\sum_{i=1}^{l-1}W(f_{\infty}^{i})+W(f^{F}).$
Then
$\lim_{\delta\rightarrow 0}\lim_{k\rightarrow+\infty}\int_{D_{\delta}}|\nabla
f_{k}|^{2}=E(F^{1})+\sum_{i}\sum_{p\in\mathcal{C}(\\{f_{k}^{i}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}\int_{B_{r}(p)}|\nabla f_{k}^{i}|^{2}$
and
$\lim_{\delta\rightarrow 0}\lim_{k\rightarrow+\infty}W(f_{k},D_{\delta})\geq
W(F^{1})+\sum_{i}\sum_{p\in\mathcal{C}(\\{f_{k}^{i}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}W(f_{k}^{i},B_{r}(p)).$
Step 2. We consider the convergence of $\\{f_{k}^{i}\\}$ near its blow up
points.
For each $p\in\mathcal{C}(\\{f_{k}^{i}\\})$, we find a small $r$ such that
$B_{r}(p)\subset S^{1}\times\mathbb{R}$ contains only one blowup point. Then
for each $p$, using the arguments in Step 1, we will get the first level of
bubble tree of $\\{f_{k}^{i}\\}$, which is a map $F_{p}$ from a stratified
sphere $S_{p}$ into $\mathbb{R}^{n}$. Each $S_{p}$ is attached to $S_{1}$ at
$p$. Taking union over $p\in\mathcal{C}(\\{f_{k}^{i}\\})$ gives us a
continuous map $F^{2}$ from $S_{2}$, which is a union of stratified spheres,
into $\mathbb{R}^{n}$. We call $F^{2}$ the second level of the bubble tree of
$\\{f_{k}\\}$.
Figure 3. Bubble tree: Second level
Step 3. In the same way, we can build the third and higher levels of the
bubble tree.
Since each step will take away at least $\epsilon_{2}^{2}$ from the Willmore
functional, the construction will stop after finite many steps. In the end we
get a stratified surface $S$ which is the union of all levels and a mapping
$F$ from $S$ into $\mathbb{R}^{n}$. We shrink all the components of $S$ on
which $F$ is trivial into points, i.e. deleting the ghost bubbles, then we get
a new stratified surface $S^{\prime}$ and a continuous map $F^{\prime}$ from
$S^{\prime}$ into $\mathbb{R}^{n}$, such that $F^{\prime}$ is nontrivial on
each component of $S^{\prime}$. We call $F^{\prime}$ is the bubble tree of
$\\{f_{k}\\}$ at $0$. Moreover, we have
$\lim_{\delta\rightarrow 0}\lim_{k\rightarrow+\infty}\int_{D_{\delta}}|\nabla
f_{k}|^{2}=E(F^{\prime}),$
and
$W(F^{\prime})\leq\lim_{\delta\rightarrow
0}\lim_{k\rightarrow+\infty}W(f_{k},D_{\delta}).$
### 2.4. Bubble trees for a sequence of maps from cylinders $Q(T_{k})$ with
$T_{k}\rightarrow+\infty$
In this subsection, we show that in the situation that blowup occurs along a
long cylinder we can divide the long cylinder into smaller ones and apply the
results in the subsection 2.3 on each.
Let $f_{k}\in\mathcal{F}^{p}_{conf}(Q(T_{k}),R)$ with
$T_{k}\rightarrow+\infty$. In light of Proposition 2.6, we only need to
consider the case that the following happens:
(2.20)
$\lim_{T\rightarrow+\infty}\varliminf_{k\rightarrow+\infty}\sup_{t\in[-T_{k}+T,T_{k}-T]}W(f_{k},S^{1}\times[t,t+1])\geq\epsilon_{2}^{2}$
since otherwise there will be no bubbles, no necks and no energy loss. When
(2.20) holds, there exist $t_{k}>0$ such that $T_{k}-t_{k}\to+\infty$ as
$k\to\infty$ and
$W(f_{k},S^{1}\times[t_{k},t_{k}+1])\geq\epsilon^{2}_{2}.$
Case I. If $\\{t_{k}\\}$ contains a bounded subsequence, we choose this
subsequence and apply the bubble tree construction in subsection 2.3 at the
blowup points.
Case II. If $\\{t_{k}\\}$ does not contain any bounded subsequences, then by
Lemma 2.7, we can find (by translations)
$-T_{k}=d_{k}^{0}<d_{k}^{1}<\cdots<d_{k}^{l}=T_{k}$
which satisfy (2.16), (2.17) and (2.18). Recall that $l$ is independent of
$k$. We may assume $f_{k}^{i}(t,\theta)=f_{k}(d_{k}^{i}+t,\theta)$ converge
weakly to $f_{\infty}^{i}$ in $W^{2,p}$ outside the blowup points
${\mathcal{C}}(\\{f^{i}_{k}\\})$ of $\\{f_{k}^{i}\\}$. Let
$\Sigma_{\infty}^{1}$ be the stratified surface with $l-1$ components whose
dual graph is a tree. Then we get a continuous map $F^{1}$ from
$\Sigma^{1}_{\infty}$ into $\mathbb{R}^{n}$, and $F^{1}$ is $f_{\infty}^{i}$
on the $i$-th component. Moreover, we have
$\lim_{T\rightarrow+\infty}\lim_{k\rightarrow+\infty}\int\limits_{S^{1}\times[-T_{k}+T,T_{k}-T]}|\nabla
f_{k}|^{2}=E(F^{1})+\sum_{i=1}^{l-1}\sum_{p\in\mathcal{C}(\\{f_{k}^{i}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}\int_{B_{r}(p)}|\nabla f_{k}^{i}|^{2}$
and
$\lim_{T\rightarrow+\infty}\lim_{k\rightarrow+\infty}W(f_{k},Q(T_{k}-T))\geq
W(F^{1})+\sum_{i=1}^{l-1}\sum_{p\in\mathcal{C}(\\{f_{k}^{i}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}W(f_{k}^{i},B_{r}(p)).$
The first level of the bubble tree of $\\{f_{k}\\}$ is $F^{1}$ together with
the bubble tree in Case I. Then we repeat this process to construct the second
level of the bubble tree at $\bigcup^{l}_{i=1}{\mathcal{C}}(\\{f^{i}_{k}\\})$,
and similarly the third level and so on. The construction stops in finite
steps.
### 2.5. Convergence in Hausdorff distance
The main aim of this subsection is to prove the following:
###### Theorem 2.8.
Assume that $(\Sigma,h_{k})$ are a sequence of close Riemann surface of genus
$g$, where $h_{k}$ satisfies (1.1). Suppose that
$f_{k}\in\mathcal{F}^{p}_{conf}(\Sigma,h_{k},R)$ with $p\in(1,\frac{4}{3})$
and
(2.21) $\sup_{p}\\{E(f_{k})+W(f_{k})\\}<\Lambda<\infty.$
Then either $f_{k}$ converges to a point, or there is a stratified surface
$\Sigma_{\infty}$ with $g(\Sigma_{\infty})\leq g$, an
$f_{0}\in\mathcal{F}^{p}_{conf}(\Sigma_{\infty},R)$, such that a subsequence
of $f_{k}(\Sigma_{k})$ converges to $f_{0}(\Sigma_{\infty})$ in the Hausdorff
distance with
$E(f_{0})=\lim_{k\rightarrow+\infty}E(f_{k}),\,\,\,\,and\,\,\,\,W(f_{0})\leq\lim_{k\rightarrow+\infty}W(f_{k}).$
Remark. Here $f_{0}\in\mathcal{F}^{p}_{conf}(\Sigma_{\infty},R)$ means that
$f_{0}\in C^{0}(\Sigma_{\infty},\mathbb{R}^{n})$, and for any component
$\Sigma^{i}_{\infty}$ of $\Sigma_{\infty}$, $F$ is nontrivial on
$\Sigma^{i}_{\infty}$ and
$F|_{\Sigma^{i}_{\infty}}\in\mathcal{F}^{p}(\overline{\Sigma^{i}_{\infty}},h_{i},R)$.
Proof of Theorem 2.8: The proof will be divided into three cases according to
the genus of $\Sigma$.
Spherical case. When $\Sigma$ is a sphere, as there is only one conformal
structure on a 2-sphere, we may let $h_{k}\equiv h$. Let
$\mathcal{C}(\\{f_{k}\\})=\\{p_{1},\dots,p_{m}\\}.$ We can choose $\delta$,
such that $B_{\delta}(p_{i})\cap B_{\delta}(p_{j})=\emptyset$. Using
isothermal coordinates, each $B_{\delta}(p_{i})$ with metric $h$ is conformal
to a Euclidean disk, the results can be deduced from subsection 2.3 directly.
Toric case. Suppose that $(\Sigma,h)$ is induced by lattice $\\{1,a+bi\\}$ in
$\mathbb{C}$, where $-\frac{1}{2}<a\leq\frac{1}{2}$, $b>0$, $a^{2}+b^{2}\geq
1$, and $a\geq 0$ whenever $a^{2}+b^{2}=1$. Then the conformal map $f$ from
$(\Sigma,h)$ into $\mathbb{R}^{n}$ can be lifted to a conformal map
$\widetilde{f}$ from $\mathbb{C}$ into $\mathbb{R}^{n}$ which satisfies
$\widetilde{f}(z+a+bi)=\widetilde{f}(z).$
Let
$\begin{array}[]{lll}\Pi:&\mathbb{C}\rightarrow S^{1}\times\mathbb{R}\\\
&a+bi\rightarrow(\pi a,\pi b),\end{array}$
be the conformal covering map. Then $(\Sigma,h)$ is conformal to
$(S^{1}\times\mathbb{R})/G$, where $G\cong\mathbb{Z}$ is the transformation
group of $S^{1}\times\mathbb{R}$ generalized by
$(\theta,t)\rightarrow(\theta+2\pi a,t+2\pi b).$
Then $f$ can be lifted to a map $f^{\prime}:S^{1}\times\mathbb{R}^{n}$, which
satisfies $f^{\prime}(\Pi)=\widetilde{f}$.
Now we assume $(\Sigma_{k},h_{k})=S^{1}\times\mathbb{R}/G_{k}$, where $G_{k}$
is generalized by
$(\theta,t)\rightarrow(\theta+\theta_{k},t+b_{k}),\,\,\,\,where\,\,\,\,b_{k}\geq\sqrt{\pi^{2}-\theta_{k}^{2}},\,\,\,\,and\,\,\,\,\theta_{k}\in[-\frac{\pi}{2},\frac{\pi}{2}].$
In the moduli space $\mathcal{M}_{1}$ of genus 1 surfaces, $(\Sigma,h_{k})$
diverges if and only if $b_{k}\rightarrow+\infty$.
For $f_{k}\in\mathcal{F}^{p}_{conf}(\Sigma,h_{k},R)$ with (2.21), we lift each
$f_{k}$ to a mapping
$f_{k}^{\prime}:S^{1}\times\mathbb{R}\rightarrow\mathbb{R}^{n}$ which
satisfies
$f_{k}^{\prime}(\theta,t)=f_{k}^{\prime}(\theta+\theta_{k},t+a_{k}).$
After translations, we may assume that
$f_{k}^{\prime}(\theta,t+\frac{a_{k}}{2})$ and
$f_{k}^{\prime}(\theta,t-\frac{a_{k}}{2})$ have no blowup points as
$k\to\infty$. Then $f_{k}^{\prime}$ satisfies the conditions in subsection 2.4
for $T_{k}=a_{k}/2$. Since
$f_{k}^{\prime}(\theta,-T_{k}+t)=f_{k}^{\prime}(\theta+\theta_{k},T_{k}+t)$,
the weak limit of $f_{k}^{\prime}(\theta,-T_{k}+t)$ in
$W^{2,p}_{loc}(S^{1}\times[0,+\infty))$ and the weak limit of
$f_{k}^{\prime}(\theta,T_{k}+t)$ in $W^{2,p}_{loc}(S^{1}\times(-\infty,0])$
are just the two parts of a conformal map from $S^{1}\times\mathbb{R}$ into
$\mathbb{R}^{n}$. So the Hausdorff limit of $f_{k}(\Sigma)$ is the image of a
continuous map $F$ from a stratified surface $S$ of genus 1 into
$\mathbb{R}^{n}$ with
$E(F)=\lim_{k\rightarrow+\infty}E(f_{k}),\,\,\,\,W(F)\leq\lim_{k\rightarrow+\infty}W(f_{k}).$
Hyperbolic case. For the hyperbolic case, we first briefly review the
compactness of moduli space.
Let $\Sigma_{0}$ be a stable surface in $\overline{\mathcal{M}}_{g}$ with
nodal points $\mathcal{N}=\\{a_{1},\dots,a_{m^{\prime}}\\}$. Geometrically,
$\Sigma_{0}$ is obtained by pinching $m^{\prime}$ non null homotopy curves in
a surface with genus $g>1$ to points $a_{1},\dots,a_{m^{\prime}}$, thus
$\Sigma_{0}\backslash\mathcal{N}$ can be divided to finite components
$\Sigma_{0}^{1},\dots,\Sigma_{0}^{s}$. For each $\Sigma_{0}^{i}$, we can
extend $\Sigma_{0}^{i}$ to a smooth closed Riemann surface
$\overline{\Sigma_{0}^{i}}$ by adding a point at each puncture. Moreover, the
complex structure of $\Sigma_{0}^{i}$ can be extended smoothly to a complex
structure of $\overline{\Sigma_{0}^{i}}$.
We say $h_{0}$ determines a hyperbolic structure on $\Sigma_{0}$ if $h_{0}$ is
a smooth complete metric on $\Sigma_{0}\backslash\mathcal{N}$ with finite
volume and Gauss curvature $-1$. We define a neighborhood around each nodal
point $a_{j}$ in $\Sigma_{0}$ by
$\Sigma_{0}(a_{j},\delta)=\left\\{p\in\Sigma_{0}:\,\,\,\,\hbox{injrad}_{\Sigma_{0}\backslash\mathcal{N}}^{h_{0}}(p)<\delta,\,\,\,\,\forall
p\in\Sigma_{0}(a_{j},\delta)\backslash\\{a_{j}\\}\right\\}\bigcup\,\\{a_{j}\\}.$
Let $h_{0}^{i}$ be the metric on $\overline{\Sigma_{0}^{i}}$ which has Gauss
curvature $\pm 1$ or curvature 0, and is conformal to $h_{0}$ on
$\Sigma_{0}^{i}$.
Now, we let $\Sigma_{k}$ be a sequence of compact Riemann surfaces of fixed
genus $g$ with hyperbolic structures $h_{k}$, such that
$\Sigma_{k}\rightarrow\Sigma_{0}$ in the moduli space
$\overline{\mathcal{M}_{g}}$. By Proposition 5.1 in [11], there exists a
maximal collection
$\Gamma_{k}=\\{\gamma_{k}^{1},\ldots,\gamma_{k}^{m^{\prime}}\\}$ of pairwise
disjoint, simple closed geodesics in $\Sigma_{k}$ with
$\ell^{j}_{k}=L(\gamma_{k}^{j})\to 0$, such that after passing to a
subsequence the following holds:
* (1)
There are maps $\varphi_{k}\in C^{0}(\Sigma_{k},\Sigma_{0})$, such that
$\varphi_{k}:\Sigma_{k}\backslash\Gamma_{k}\to\Sigma_{0}\backslash\mathcal{N}$
is diffeomorphic and $\varphi_{k}(\gamma_{k}^{j})=a_{j}$ for
$j=1,\ldots,m^{\prime}$.
* (2)
For the inverse diffeomorphisms
$\psi_{k}:\Sigma_{0}\backslash\mathcal{N}\to\Sigma_{k}\backslash\Gamma_{k}$,
we have $\psi_{k}^{\ast}(h_{k})\to h_{0}$ in
$C^{\infty}_{loc}(\Sigma_{0}\backslash\mathcal{N})$, where $h_{0}$ determine a
hyperbolic structure over $\Sigma_{0}\backslash\mathcal{N}$.
* (3)
Let $c_{k}$ be the complex structure over $\Sigma_{k}$, and $c_{0}$ be the
complex structure over $\Sigma_{0}\backslash\mathcal{N}$. Then
$\psi_{k}^{*}(c_{k})\rightarrow
c_{0}\,\,\,\,in\,\,\,\,C^{\infty}_{loc}(\Sigma_{0}\backslash\mathcal{N}).$
Moreover, we have the following Collar Lemma [9, 12, 20, 26]:
###### Lemma 2.9.
For each $\gamma_{k}^{j}$ as above, there is a collar $U_{k}^{j}$ containing
$\gamma_{k}^{j}$, which is isometric to the cylinder
$Q_{k}^{j}=Q(\frac{\pi^{2}}{l_{k}^{j}})$ with metric
(2.22)
$h_{k}^{j}=\left(\frac{1}{2\pi\sin(\frac{l_{k}^{j}}{2\pi}t+\theta_{k})}\right)^{2}(dt^{2}+d\theta^{2}),$
where $\theta_{k}=\arctan(\sinh(\frac{l_{k}^{j}}{2}))+\frac{\pi}{2}$.
Moreover, for any $(\theta,t)\in
S^{1}\times(-\frac{\pi^{2}}{l_{k}^{j}},\frac{\pi^{2}}{l_{k}^{j}})$, we have
(2.23)
$\sinh(\it{injrad}_{\Sigma_{k}}(\theta,t))\sin(\frac{l_{k}^{j}t}{2\pi}+\theta_{k})=\sinh\frac{l_{k}^{j}}{2}.$
Let $\phi_{k}^{j}$ be the isometry between $Q_{k}^{j}$ and $U_{k}^{j}$. Then
$\varphi_{k}\circ\phi_{k}^{j}(\theta,\frac{\pi^{2}}{l_{k}^{j}}+t)\bigcup\varphi_{k}\circ\phi_{k}^{j}(\theta,-\frac{\pi^{2}}{l_{k}^{j}}+t)$
converges in $C^{\infty}_{loc}(S^{1}\times(-\infty,0)\cup
S^{1}\times(0,\infty))$ to an isometry from $S^{1}\times(-\infty,0)\cup
S^{1}\times(0,+\infty)$ to $\Sigma_{0}(a_{j},1)\backslash\\{a_{j}\\}$.
The Collar Lemma can be found in [12], [9] and [11].
We also need the following local existence and compactness of conformal
diffeomorphisms.
###### Theorem 2.10.
[4] Let $h_{k},h_{0}$ be smooth Riemannian metrics on a surface $M$, such that
$h_{k}\to h_{0}$ in $C^{s,\alpha}(M)$, where $s\in\mathbb{N}$,
$\alpha\in(0,1)$. Then for each point $z\in M$ there exist neighborhoods
$U_{k},U_{0}$ of $z$ and smooth conformal diffeomorphisms $\vartheta_{k}:D\to
U_{k},\vartheta_{0}:D\rightarrow U$, such that $\vartheta_{k}\to\vartheta_{0}$
in $C^{s+1,\alpha}(\overline{D},M)$.
Proof of Theorem 2.8 (continued): For a sequence
$f_{k}\in\mathcal{F}^{p}_{conf}(\Sigma,h_{k},R)$ satisfying the energy bound
(2.21), let
$\widetilde{f}_{k}=f_{k}\circ\psi_{k}$
which is a mapping from $\Sigma_{0}\backslash\mathcal{N}$ to $\mathbb{R}^{n}$.
It is easy to check that
$\widetilde{f}_{k}\in\mathcal{F}^{p}_{conf}(\Sigma_{0}\backslash\mathcal{N},\psi_{k}^{\ast}(h_{k}),R)$.
First, we show $\widetilde{f}_{k}$ converge in
$W^{2,p}_{loc}(\Sigma_{0}\backslash(\mathcal{N}\cup\mathcal{C}(\\{f_{k}\\})))$.
Given a point
$z\in\Sigma_{0}\backslash(\mathcal{N}\cup\mathcal{C}(\\{\widetilde{f}_{k}\\}))$,
we choose $U_{k},U,\vartheta_{k},\vartheta$ as in Theorem 2.10 and $U_{k}$,
$U\subset\Sigma_{0}\backslash(\mathcal{N}\cup\mathcal{C}(\\{\widetilde{f}_{k}\\}))$.
Let
$\widehat{f}_{k}=\widetilde{f}_{k}\circ\varphi_{k}$
and note that $\widehat{f}_{k}\in\mathcal{F}^{p}_{conf}(D,R)$. We can assume
that $\widehat{f}_{k}$ converge to $\widehat{f}_{\infty}$ in
$W^{2,p}_{loc}(D_{3/4})$ with
$\partial\widehat{f}_{\infty}\otimes\partial\widehat{f}_{\infty}=0$. Let
$V=\vartheta(D_{1/2})$. Since $\vartheta_{k}$ converge to $\vartheta$,
$\vartheta_{k}^{-1}(V)\subset D_{3/4}$ for sufficiently large $k$,
$\widetilde{f}_{k}=\widehat{f}_{k}(\vartheta_{k}^{-1})$ converge to
$\widetilde{f}_{\infty}=\widehat{f}_{\infty}(\vartheta_{0}^{-1})$ weakly in
$W^{2,p}(V,h_{0})$. Then
$\widetilde{f}_{\infty}\in\mathcal{F}^{p}_{conf}(V,h_{0},R)$. Moreover, for
any nonnegative function $\varphi$ with support in $V$, from Fatou’s lemma
(2.24)
$\lim_{k\rightarrow+\infty}\int_{V}\varphi|H(\widetilde{f}_{k})|^{2}|\nabla\widetilde{f}_{k}|^{2}=\lim_{k\rightarrow+\infty}\int_{D}\varphi(\vartheta_{k})|H(\widehat{f}_{k})|^{2}|\nabla\widehat{f}_{k}|^{2}\geq\int_{D}\varphi(\vartheta_{0})|H(\widehat{f}_{0})|^{2}|\nabla\widehat{f}_{0}|^{2}.$
We may thus assume $\widetilde{f}_{k}$ converge weakly to
$\widetilde{f}_{\infty}$ in
$W^{2,p}_{loc}(\Sigma_{0}\backslash(\mathcal{N}\cup\,\mathcal{C}(\\{f_{k}\\})))$.
Then $\widetilde{f}_{\infty}|_{\Sigma_{0}^{i}}\in
W^{2,p}_{loc}(\Sigma_{0}^{i},h_{0}^{i})$. So for $p\in(1,{4/3})$,
$\widetilde{f}_{\infty}|_{\Sigma_{0}^{i}}$ can be extended to a map in
$\mathcal{F}^{p}_{conf}(\overline{\Sigma_{0}^{i}},h_{0}^{i},R)$. Further,
$\lim_{k\rightarrow+\infty}E(f_{k})=E(\widetilde{f}_{\infty})+\sum_{z\in\mathcal{C}(\\{\widetilde{f}_{k}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}E(\widetilde{f}_{k},B_{r}(z,h_{0}))+\sum_{j}\lim_{\delta\rightarrow
0}\lim_{k\rightarrow+\infty}E(\widetilde{f}_{k},\Sigma_{0}(a_{j},\delta))$
and from (2.24)
$\lim_{k\rightarrow+\infty}W(\widetilde{f}_{k})\geq
W(\widetilde{f}_{\infty})+\sum_{z\in\mathcal{C}(\\{\widetilde{f}_{k}\\})}\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}W(\widetilde{f}_{k},B_{r}(z,h_{0}))+\sum_{j}\lim_{\delta\rightarrow
0}\lim_{k\rightarrow+\infty}W(\widetilde{f}_{k},\Sigma_{0}(a_{j},\delta)).$
Next, we construct bubble trees at a point
$z\in\mathcal{C}(\\{f_{k}\\})\backslash\mathcal{N}$. We have a bubble tree $F$
of $\widehat{f}_{k}$ at $z$. We define it to be a bubble tree of
$\widetilde{f}_{k}$ at $z$. By the arguments in subsection 2.3, we have
$\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}E(\widetilde{f}_{k},B_{r}(z,h_{0}))=\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}E(\widehat{f}_{k},D_{r})=E(F),$
and
$\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}W(\widetilde{f}_{k},B_{r}(z,h_{0}))=\lim_{r\rightarrow
0}\lim_{k\rightarrow+\infty}W(\widehat{f}_{k},D_{r})\geq W(F).$
Lastly, we consider the convergence of $f_{k}$ at the collars. Set
$\check{f}_{k}^{j}=f_{k}\circ\phi_{k}^{j}$
and $T_{k}^{j}={\pi^{2}/l_{k}^{j}}-T$. We may choose $T$ to be sufficiently
large such that the two sequences $\check{f}_{k}^{j}(T_{k}^{j}-t,\theta)$ and
$\check{f}_{k}^{j}(-T_{k}^{j}+t,\theta)$ have no blowup points in $[0,T]$ and
$[-T,0]$ respectively (otherwise, return to the previous case as in Case I,
subsection 2.4). Then $\check{f}_{k}^{j}$ satisfies the conditions in
subsection 2.4. We get a bubble tree $F^{j}$. So the convergence of
$\check{f}_{k}^{j}$ is clear. Since
$\check{f}_{k}^{j}=f_{k}\circ\phi_{k}^{j}=f_{k}\circ\psi_{k}\circ(\varphi_{k}\circ\phi_{k}^{j})=\widetilde{f}_{k}(\varphi_{k}\circ\phi_{k}^{j}),$
we have
$\displaystyle\check{f}_{k}^{j}(T_{k}^{j}-t,\theta)$ $\displaystyle=$
$\displaystyle\widetilde{f}_{k}(\varphi_{k}\circ\phi_{k}^{j}(T_{k}^{j}-t,\theta)),$
$\displaystyle\check{f}_{k}^{j}(t-T_{k}^{j},\theta)$ $\displaystyle=$
$\displaystyle\widetilde{f}_{k}(\varphi_{k}\circ\phi_{k}^{j}(t-T_{k}^{j},\theta)).$
By the convergence statement of the Collar Lemma,
$\varphi_{k}\circ\phi^{j}_{k}(t+T-\pi^{2}/l^{j}_{k},\theta)$ and
$\varphi_{k}\circ\phi_{k}^{j}({\pi^{2}/l_{k}^{j}}-T-t,\theta)$ converge in
$C^{\infty}_{loc}(S^{1}\times(0,\infty))$ to an isometry from
$S^{1}\times(-\infty,0)\cup S^{1}\times(0,+\infty)$ to
$\Sigma_{0}(a_{j},1)\setminus\\{a_{j}\\}$. We conclude that the image of the
limit of $\check{f}_{k}^{j}(T_{k}^{j}-t,\theta)$ and that of
$\check{f}_{k}^{j}(-T_{k}^{j}+t,\theta)$ are both contained in the image of
$\widetilde{f}_{\infty}$. Moreover,
$\lim_{\delta\rightarrow
0}\lim_{k\rightarrow+\infty}E(\widetilde{f}_{k},\Sigma_{0}(a_{j},\delta))=\lim_{T\rightarrow+\infty}\lim_{k\rightarrow+\infty}E(\check{f},Q(T_{k}-T))=E(F^{j}),$
and
$\lim_{\delta\rightarrow
0}\lim_{k\rightarrow+\infty}W(\widetilde{f}_{k},\Sigma_{0}(a_{j},\delta))=\lim_{T\rightarrow+\infty}\lim_{k\rightarrow+\infty}W(\check{f},Q(T_{k}-T))=E(F^{j}).$
Thus, we complete the proof. $\hfill\Box$
###### Remark 2.11.
When $\Sigma_{0}\in\mathcal{M}_{p}$, i.e. $\mathcal{N}=\emptyset$, $\psi_{k}$
is just a smooth diffeomorphism sequence from $\Sigma$ to $\Sigma$. In this
case, $g(\Sigma_{\infty})=g(\Sigma)$, and
$\Sigma_{\infty}=\Sigma_{\infty}^{\prime}\cup S_{1}\cup S_{2}\cdots\cup
S_{m},$
where $\Sigma_{\infty}^{\prime}$ is a smooth Riemann surface of genus $p$, and
each $S_{i}$ is a sphere.
Figure 4. $\Sigma_{0}$ (limit of $(\Sigma,h_{k}))$ and $\Sigma_{\infty}$ (some
components of $\Sigma_{0}$ do not appear in $\Sigma_{\infty}$)
We now generalize Theorem 2.8 to surfaces with marked points. Let us briefly
review the compactification of the moduli space of surfaces with marked
points. Let $\overline{\mathcal{M}}_{g,m}$ be the moduli space of closed
Riemann surfaces of genus $g$ with $m$ marked points. Let
$(\Sigma_{0},x_{0,1},\dots,x_{0,m})\in\partial\overline{\mathcal{M}}_{g,m}$
with nodal points $\mathcal{N}=\\{a_{1},\dots,a_{m^{\prime}}\\}$.
Geometrically, $\Sigma_{0}$ is obtained by pinching some homotopically
nontrivial closed curves which do not pass any of $x_{0,1},\dots,x_{0,m}$ into
the points in ${\mathcal{N}}$, and $\Sigma\backslash\mathcal{N}$ can be
divided to connected components $\Sigma_{0}^{1}$, $\cdots$, $\Sigma_{0}^{s}$.
For each $\Sigma_{0}^{i}$, we can extend $\Sigma_{0}^{i}$ to a smooth closed
Riemann surface $\overline{\Sigma_{0}^{i}}$ by adding a point at each
puncture. Moreover, the complex structure of $\Sigma_{0}^{i}$ can be extended
smoothly to a complex structure of $\overline{\Sigma^{i}_{0}}$.
We say $h$ is a hyperbolic structure on
$(\Sigma,x_{1},\dots,x_{m})\in\mathcal{M}_{g,m}$ if $h$ is a smooth complete
metric on $\Sigma\backslash\\{x_{1},\dots,x_{m}\\}$ with curvature $-1$ and
finite volume. We say $h_{0}$ is a hyperbolic structure on
$(\Sigma_{0},x_{0,1},\dots,x_{0,m})\in\overline{\mathcal{M}}_{g,m}\backslash\mathcal{M}_{g,m}$
if $h_{0}$ is a smooth complete metric on
$\Sigma\backslash\\{a_{1},\dots,a_{m^{\prime}},x_{0,1},...,x_{0,m}\\}$
with curvature $-1$ and finite volume.
For a surface $\Sigma$ with hyperbolic structure $h$ and with marked points
$x_{1},\dots,x_{m}$, we define
$\Sigma^{*}=\Sigma\backslash\\{x_{1},\dots,x_{m}\\}$, and $h^{*}$ to be the
hyperbolic structure on $(\Sigma,x_{1},\dots,x_{m})$ which is conformal to $h$
on $\Sigma^{*}$.
Let $\\{(\Sigma_{k},x_{k,1},\dots,x_{k,m}\\}$ be a sequence of marked surfaces
in $\mathcal{M}_{g,m}$ with hyperbolic structures $h_{k}$ and
$(\Sigma_{k},x_{k,1},\dots,x_{k,m})\rightarrow(\Sigma_{0},x_{0,1},\dots,x_{0,m})\,\,\,\,in\,\,\,\,\overline{\mathcal{M}}_{g,m}.$
By Proposition 5.1 in [11] again, there exists a maximal collection
$\Gamma_{k}=\\{\gamma_{k}^{1},\ldots,\gamma_{k}^{m^{\prime}}\\}$ of pairwise
disjoint, simple closed geodesics in $\Sigma_{k}$ with
$\ell^{j}_{k}=L(\gamma_{k}^{j})\to 0$ as $k\to\infty$, such that, after
passing to a subsequence if necessary, the following holds:
* (1)
There are maps $\varphi_{k}\in C^{0}(\Sigma_{k},\Sigma_{0})$, such that
$\varphi_{k}:\Sigma_{k}\backslash\Gamma_{k}\to\Sigma_{0}\backslash\mathcal{N}$
is diffeomorphic and $\varphi_{k}(\gamma_{k}^{i})=a_{i}$ for
$i=1,\dots,m^{\prime}$, and $\varphi_{k}(x_{k,j})=x_{0,j}$ for $j=1,\dots,m$.
* (2)
For the inverse diffeomorphisms
$\psi_{k}:\Sigma_{0}\backslash\mathcal{N}\to\Sigma_{k}\backslash\Gamma_{k}$,
we have $\psi_{k}^{\ast}(h_{k})\to h$ in
$C^{\infty}_{loc}(\Sigma_{0}^{*}\backslash\mathcal{N})$.
* (3)
Let $c_{k}$ be the complex structure over $\Sigma_{k}$, and $c_{0}$ be the
complex structure over $\Sigma_{0}\backslash\mathcal{N}$. Then
$\psi_{*}(c_{k})\rightarrow
c_{0}\,\,\,\,in\,\,\,\,C^{\infty}_{loc}(\Sigma_{0}\backslash\mathcal{N}).$
Moreover, the Collar Lemma also holds for the moduli space with marked points.
###### Theorem 2.12.
In addition to the assumptions in Theorem 2.8, we assume for $m\geq 2$
$y_{1},\dots,y_{m}\in f_{k}(\Sigma).$
Then there is a stratified surface $\Sigma_{\infty}$ with
$g(\Sigma_{\infty})\leq g$, and an
$f_{0}\in\mathcal{F}^{p}_{conf}(\Sigma_{\infty},R)$ with
$y_{1},\dots,y_{m}\in f_{0}(\Sigma_{\infty}),$
such that a subsequence of $f_{k}(\Sigma_{k})$ converges to
$f_{0}(\Sigma_{\infty})$ in Hausdorff distance with
$E(f_{0})=\lim_{k\rightarrow+\infty}E(f_{k})\,\,\,\,and\,\,\,\,W(f_{0})\leq\lim_{k\rightarrow+\infty}W(f_{k}).$
###### Proof.
Let $\widetilde{f}_{k}=f_{k}\circ\psi_{k}$. In view of Theorem 2.8, we only
need to consider convergence of $\\{\widetilde{f}_{k}\\}$ near
$x_{0,j},j=1,\dots,m$.
Choose a complex coordinate $\\{U,(x,y)\\}$ on $\Sigma_{0}$ compatible with
$c_{0}$, with $x_{0,j}=(0,0)$. Let $c_{k}^{\prime}=\psi_{k}^{\ast}(c_{k})$. We
set
$e_{1}=\frac{\partial}{\partial x},\,\,\,\,e_{2}=c_{k}^{\prime}(e_{1}),$
and $h_{k}^{\prime}$ to be the metric on $U$ defined by
$h_{k}^{\prime}(e_{1},e_{1})=h_{k}^{\prime}(e_{2},e_{2})=1,\,\,\,\,h_{k}^{\prime}(e_{1},e_{2})=0.$
Then $h_{k}^{\prime}$ is compatible with $c_{k}^{\prime}$, and converges
smoothly to a metric which is compatible with $c_{0}$ in $U$. Then we consider
the weak convergence of $\\{\widetilde{f}_{k}\\}$ in
$U\backslash\mathcal{C}(\\{\widetilde{f}_{k}\\})$, using the arguments in
subsection 2.3.
It remains to check that each marked point $y_{i}$ is on the image of
$f_{\infty}$ or one of the bubbles. If $x_{0,j}$ is not a blow up point of
$\\{\widetilde{f}_{k}\\}$, it is obvious that $y_{j}\in f_{\infty}(U)$. Now
assume $x_{0,j}$ is the only blow-up point in $D$. We take $U_{k}$, $U_{0}$,
$\vartheta_{k},\vartheta_{0},\widehat{f}_{k},V$ as in the proof of Theorem 2.8
for $z=x_{0,j}$. We will prove it by induction on the number of the levels of
the bubble tree. We take $z_{k}$, $r_{k}$, $\phi_{k}$ and $d_{k}$ as in
subsection 2.3 for $\widehat{f}_{k}$. If ${|z_{k}|/r_{k}}<L$ for some fixed
$L$, then we may assume $-{z_{k}/r_{k}}\rightarrow z_{\infty}$, by selecting a
subsequence if necessary. Recalling that $\widehat{f}_{k}(0)\equiv y_{j}$, we
get $y_{j}=\widehat{f}^{F}(z_{\infty})$. Let $(r,\theta)$ be the polar
coordinates centered at $z_{k}$, $T_{k}=-\ln r_{k}$ and
$\phi_{k}:[0,T_{k}]\times S^{1}\rightarrow\mathbb{R}^{2}$ be the conformal
mapping given by $\phi_{k}(t,\theta)=(e^{-t},\theta)$. We set
$\phi_{k}^{-1}(0)=(t_{k},\theta_{k})$. Then
${|z_{k}|/r_{k}}\rightarrow+\infty$ means that $t_{k}\in[0,T_{k}]$ and
$T_{k}-t_{k}\rightarrow+\infty$. Thus we may assume
$t_{k}\in[d_{k}^{i},d_{k}^{i+1}]$ for some $i$, where $d_{k}^{i}$ are defined
in Lemma 2.7. Then, if $t_{k}-d_{k}^{i}\rightarrow+\infty$ and
$d_{k}^{i+1}-t_{k}\rightarrow+\infty$, we have
$y_{j}=f_{\infty}^{i}(+\infty)=f_{\infty}^{i+1}(-\infty).$ If at least one of
$t_{k}-d_{k}^{i}$ and $d_{k}^{i+1}-t_{k}$ is bounded above for all large $k$,
then we repeat the above argument at the second level of the bubble tree, and
proceed in this way for the finitely many levels of the bubble tree if
necessary, and we conclude that $y^{j}$ is on one of bubbles of
$\widetilde{f}_{k}^{i}$ or $\widetilde{f}_{k}^{i+1}$.
Finally, as $m\geq 2$ and all $y_{i}\in f_{0}(\Sigma_{\infty})$, $f_{k}$
cannot converge to a single point. $\hfill\Box$
## 3\. Branched conformal immersions and proof of Theorem 1
For a branched conformal immersion, we have the following result:
###### Theorem 3.1.
[13] Suppose that $f\in W^{2,2}_{conf,loc}(D\backslash\\{0\\},\mathbb{R}^{n})$
satisfies
$\int_{D}|A_{f}|^{2}\,d\mu_{g}<\infty\quad\mbox{ and }\quad\mu_{g}(D)<\infty,$
where $g_{ij}=e^{2u}\delta_{ij}$ is the induced metric. Then $f\in
W^{2,2}(D,\mathbb{R}^{n})$ and we have
$\displaystyle u(z)$ $\displaystyle=$
$\displaystyle\lambda\log|z|+\omega(z)\quad\mbox{ where }\lambda\geq
0,\quad\lambda\in\mathbb{Z},\quad\omega\in C^{0}\cap W^{1,2}(D),$
$\displaystyle-\Delta u$ $\displaystyle=$
$\displaystyle-2\lambda\pi\delta_{0}+K_{g}e^{2u}\quad\mbox{ in }D.$
The density of $f(D_{\sigma})$ as varifolds at $f(0)$ is given by $\lambda+1$
for any small $\sigma>0$.
The classical Gauss-Bonnet formula is generalized in [6] for smooth branched
surface. Following arguments in [6], we provide a version for $W^{2,2}$
branched conformal immersions.
###### Lemma 3.2.
Let $(\Sigma,g)$ be a smooth closed compact Riemann surface. Then for any
$f\in W_{b,c}^{2,2}(\Sigma,g,\mathbb{R}^{n})$, there holds
(3.1) $\int_{\Sigma}K_{f}d\mu_{f}=2\pi\chi(\Sigma)+2\pi b,$
where $b$ is the number of branch points counted with multiplicities and at
each branch point $p$ the branching order is $\lambda=\theta^{2}(p)-1$.
###### Proof.
Without loss of generality, we only prove the case that $f$ has only one
branch point $p$. Let $g_{f}=e^{2u}g$ be the metric induced by $f$ and $K_{f}$
be its Gauss curvature. In [13], we proved that the equation
$-\Delta_{g}u=K_{f}e^{2u}-K_{g}$ holds weakly in $\Sigma\backslash\\{p\\}$:
for any smooth $\varphi$ with support in $\Sigma\backslash\\{p\\}$, it holds
$\int_{\Sigma}\nabla_{g}u\nabla_{g}\varphi\,d\mu_{g}=\int_{\Sigma}\varphi
K_{f}e^{2u}d\mu_{g}-\int_{\Sigma}\varphi K_{g}du_{g}.$
Take a complex coordinate chart $\\{U;z\\}$ around $p$ with $p=0$. For any
small $\epsilon>0$, we choose a function
$\varphi_{\epsilon}(z)=\varphi_{\epsilon}(|z|)$ between 0 and 1 with
$|\varphi_{\epsilon}^{\prime}|<{C/\epsilon}$ and equals 1 outside
$D_{\epsilon}$ and $0$ in $D_{\epsilon/2}$. Then we have
$\int_{D_{\epsilon}}\frac{\partial u}{\partial
r}\varphi_{\epsilon}^{\prime}dx=\int_{\Sigma}\varphi_{\epsilon}K_{f}e^{2u}d\mu_{g}-\int_{\Sigma}\varphi_{\epsilon}K_{g}du_{g}.$
By Theorem 3.1, $u=\lambda\log|z|+\omega$, we get
$\int_{D_{\epsilon}}\frac{\partial u}{\partial
r}\varphi_{\epsilon}^{\prime}dx=\int_{D_{\epsilon}}\frac{\partial\omega}{\partial
r}\varphi_{\epsilon}^{\prime}dx+2\pi\lambda\left(\varphi_{\epsilon}(\epsilon)-\varphi_{\epsilon}(0)\right)=\int_{D_{\epsilon}}\frac{\partial\omega}{\partial
r}\varphi_{\epsilon}^{\prime}dx+2\pi\lambda.$
Since
$\int_{D_{\epsilon}}\left|\frac{\partial\omega}{\partial
r}\varphi_{\epsilon}^{\prime}\right|\leq C\left(\int_{D_{\epsilon}\backslash
D_{\epsilon/2}}\left|\frac{\partial\omega}{\partial
r}\right|^{2}\right)^{1/2}\left(\int_{D_{\epsilon}\backslash
D_{\epsilon/2}}\frac{1}{r^{2}}\right)^{1/2}\leq
C\|\nabla\omega\|_{L^{2}(D_{\epsilon})}\rightarrow
0,\,\,\,\,\hbox{as}\,\,\,\,\epsilon\rightarrow 0,$
we conclude, by applying the classical Gauss-Bonnet theorem on $(\Sigma,g)$,
that
$\int_{\Sigma}K_{f}d\mu_{f}=\lim_{\epsilon\rightarrow
0}\int_{\Sigma}\varphi_{\epsilon}K_{f}d\mu_{f}=\lim_{\epsilon\rightarrow
0}\int_{\Sigma}\varphi_{\epsilon}K_{g}d\mu_{g}+\lim_{\epsilon\rightarrow
0}\int_{\Sigma}\frac{\partial u}{\partial
r}\varphi_{\epsilon}^{\prime}d\mu_{g}=2\pi\chi(\Sigma)+2\lambda\pi$
and complete the proof. $\hfill\Box$
###### Remark 3.3.
Since $\int_{\Sigma}K_{f}d\mu_{f}\leq W(f)$, it follows from Lemma 3.2
$b\leq\frac{1}{2\pi}W(f)-\chi(\Sigma).$
Moreover,
$\int_{\Sigma}|A_{f}|^{2}d\mu_{f}=4W(f)-2\int K_{f}\leq
4W(f)-2\pi\chi(\Sigma).$
Then $\sup_{k}W(f_{k})<+\infty$ implies that $\sup_{k}b_{k}<+\infty$ and
$\sup_{k}\int_{\Sigma}|A_{f_{k}}|^{2}d\mu_{f_{k}}<+\infty$.
To study the convergence conformal immersions, we recall an important result
of Hélein.
###### Theorem 3.4.
[10] Let $f_{k}\in W^{2,2}_{conf}(D,\mathbb{R}^{n})$ be a sequence of
conformal immersions with induced metrics
$(g_{k})_{ij}=e^{2u_{k}}\delta_{ij}$, and assume
$\int_{D}|A_{f_{k}}|^{2}\,d\mu_{g_{k}}\leq\gamma<\gamma_{n}=\begin{cases}8\pi&\mbox{
for }n=3,\\\ 4\pi&\mbox{ for }n\geq 4.\end{cases}$
Assume also that $\mu_{g_{k}}(D)\leq C$ and $f_{k}(0)=0$. Then $f_{k}$ is
bounded in $W^{2,2}_{loc}(D,\mathbb{R}^{n})$, and there is a subsequence such
that one of the following two alternatives holds:
* (a)
$u_{k}$ is bounded in $L^{\infty}_{loc}(D)$ and $f_{k}$ converges weakly in
$W^{2,2}_{loc}(D,\mathbb{R}^{n})$ to a conformal immersion $f\in
W^{2,2}_{conf,loc}(D,\mathbb{R}^{n})$.
* (b)
$u_{k}\to-\infty$ and $f_{k}\to 0$ locally uniformly on $D$.
Hélein first proved the above result for $\gamma_{n}={8\pi/3}$ [10, Theorem
5.1.1]. In [13] $\gamma_{n}$ is shown to be optimal.
Before proving Theorem 1, we recall a monotonicity formula (for more details,
see [15, 31]). Let $\mu$ be a 2-dimensional integral varifold with square
integrable weak mean curvature $H_{\mu}\in L^{2}(\mu)$. Then we have
$g_{x_{0}}(\varrho)\leq
g_{x_{0}}(\sigma),\,\,\,\,\hbox{when}\,\,\,\,\varrho<\sigma,$
where
$g_{x_{0}}(r)=\frac{\mu(B_{r}(x_{0}))}{\pi
r^{2}}+\frac{1}{4\pi}W(\mu,B_{r}(x_{0}))+\frac{1}{2\pi
r^{2}}\int_{B_{r}(x_{0})}\langle x-x_{0},H\rangle\,d\mu.$
When $\Sigma$ is compact and connected, if we let $\sigma\rightarrow+\infty$,
and $\varrho\rightarrow 0$, then the area density of $\Sigma$ at $x_{0}$
satisfies
(3.2) $\theta^{2}(\mu,x_{0})\leq\frac{1}{4\pi}W(f).$
If we only let $\varrho\rightarrow 0$, then we get
(3.3)
$\theta^{2}(\mu,x_{0})\leq\frac{\mu(B_{\sigma}(x_{0}))}{\pi\sigma^{2}}+CW(f,B_{\sigma}(x_{0}))+C\left(\frac{\mu(B_{\sigma}(x_{0}))}{\pi\sigma^{2}}\right)^{\frac{1}{2}}W(f,B_{\sigma}(x_{0}))^{\frac{1}{2}}.$
Another useful consequence (cf. [31], [15]) is the following:
(3.4) $\Big{(}\frac{\mu(\Sigma)}{W(f)}\Big{)}^{\frac{1}{2}}\leq{\rm
diam\,}f(\Sigma)\leq C\Big{(}\mu(\Sigma)\,W(f)\Big{)}^{\frac{1}{2}}.$
Proof of Theorem 1. Consider a branched conformal immersion $f_{k}\in
W^{2,2}_{b,c}(\Sigma,h_{k},\mathbb{R}^{n})$, where $h_{k}$ satisfies (1.1).
The following equation clearly holds on $\Sigma$ away from the finitely many
branch points; the singularities at the branch points can be removed by
Theorem 3.1, thus it holds on entire $\Sigma$:
$\Delta_{h_{k}}f_{k}=\frac{1}{2}H_{f_{k}}|\nabla_{h_{k}}f_{k}|^{2}.$
By Remark 3.3, the number of branch points and $\|A_{f_{k}}\|_{L^{2}}$ are
both bounded from above.
By (3.4), $\mbox{\rm diam}f_{k}(\Sigma)\leq{R}$ for some $R>0$. Then
$f_{k}\in\mathcal{F}^{2}_{conf}(\Sigma,h_{k},R+R_{0})$. We only need to prove
that $f_{0}$ in Theorem 2.8 and Theorem 2.12 is also branched conformal
immersion.
In fact, we only need to prove the following: If $f_{k}$ are branched
conformal immersions from $D$ into $\mathbb{R}^{n}$ with a uniform upper bound
on the number of branch points, then the limit $f_{0}$ is either a point or a
branched conformal immersion. Let $P$ be the limit set of the branch points,
and
$\mathcal{S}(\\{f_{k}\\})=\\{z\in D:\lim_{r\rightarrow
0}\varliminf_{k\rightarrow+\infty}\int_{D_{r}(z)}|A_{k}|^{2}\geq\hat{\epsilon}^{2}\\},$
where $\hat{\epsilon}\leq\min\\{\sqrt{4\pi},4\epsilon_{0}\\}$. By Theorem 3.4,
$f_{k}$ will converge weakly in $W^{2,2}_{loc}(D\backslash(\mathcal{S}\cup
P))$ to either a conformal immersion or a point. If the limit is not a single
point, by Theorem 3.1, the limit can be extended across the finite set
${\mathcal{S}}\cup P$ to a branched conformal immersion of $D$, hence
$\Sigma$, in $\mathbb{R}^{n}$. $\hfill\Box$
## 4\. Willmore functional for surfaces in compact manifolds
Let $N$ be a compact Riemannian manifold without boundary. We embed $N$ into
$\mathbb{R}^{n}$ isometrically so that any immersion of $\Sigma$ in $N$ can be
regarded as an immersion in $\mathbb{R}^{n}$. Let
$A_{\Sigma,N},A_{\Sigma,\mathbb{R}^{n}}$ and $A_{N,\mathbb{R}^{n}}$ be the
second fundamental forms of $\Sigma$ in $N$, in $\mathbb{R}^{n}$ and $N$ in
$\mathbb{R}^{n}$ respectively. The $L^{2}$ integrals of these quantities can
be related as in the following simple lemma.
###### Lemma 4.1.
For any $f\in W^{2,2}_{b,c}(\Sigma,h,N)$, we have
(4.1) $\int_{\Sigma}|H_{f,\Sigma,\mathbb{R}^{n}}|^{2}d\mu_{f}\leq
C\mu(f)+\int_{\Sigma}|H_{f,\Sigma,N}|^{2}d\mu_{f},$
and
(4.2) $\int_{\Sigma}|A_{f,\Sigma,\mathbb{R}^{n}}|^{2}d\mu_{f}\leq
C\int_{\Sigma}(1+|H_{f,\Sigma,N}|^{2})d\mu_{f}+C^{\prime},$
where $C$ only depends on $N$ and $C^{\prime}$ only depends on the Euler
characteristic of $\Sigma$.
###### Proof.
Let $e_{1},\dots,e_{n}$ be an orthonormal frame of
$T_{f_{k}(p)}\mathbb{R}^{n}$ with $e_{1}$, $e_{2}\in Tf(\Sigma)$ and
$e_{1},\dots,e_{k}\in TN$. Then for $f_{i}=\frac{\partial f}{\partial x^{i}}$
we have
$\nabla^{N}_{f_{i}}f_{j}=\sum_{l=1}^{k}\langle f_{ij},e_{l}\rangle e_{l}$
and
$A_{\Sigma,N}(f_{i},f_{j})=\sum_{m=3}^{k}\langle\nabla^{N}_{f_{i}}f_{j},e_{m}\rangle
e_{m}=\sum_{m=3}^{k}\langle f_{ij},e_{m}\rangle e_{m}.$
Thus if $F=i\circ f$ where $i:N\to\mathbb{R}^{n}$ is the isometric embedding,
we have
$A_{\Sigma,\mathbb{R}^{n}}(F_{i},F_{j})=\sum_{m=3}^{n}\langle
F_{ij},e_{m}\rangle
e_{m}=A_{\Sigma,N}(f_{i},f_{j})+A_{N,\mathbb{R}^{n}}(F_{i},F_{j}).$
Hence, we have
$H_{\Sigma,\mathbb{R}^{n}}(f)=H_{\Sigma,N}(f)+g^{ij}A_{N,\mathbb{R}^{n}}(F_{i},F_{j}).$
Noting that $H_{\Sigma,N}(f)\perp A_{N,\mathbb{R}^{n}}$, we get
$\left|H_{\Sigma,\mathbb{R}^{n}}(f)\right|^{2}=\left|H_{\Sigma,N}(f)\right|^{2}+\left|g^{ij}A_{N,\mathbb{R}^{n}}(F_{i},F_{j})\right|^{2}\leq\left|H_{\Sigma,\mathbb{R}^{n}}\right|^{2}+\|A_{N,\mathbb{R}^{n}}\|_{L^{\infty}}$
where $\|A_{N,\mathbb{R}^{n}}\|_{L^{\infty}}$ is bounded since $N$ is compact.
Thus, integrating over $\Sigma$ yields (4.1), and by (4.1) and Remark 3.3, we
get
$\int_{\Sigma}|A_{\Sigma,\mathbb{R}^{n}}(f)|^{2}d\mu_{f}<C\left(1+\mu(f)+\int_{\Sigma}|H_{\Sigma,N}(f)|^{2}d\mu_{f}\right).$
$\hfill\Box$
### 4.1. Willmore sphere passing through fixed points
In this subsection, we let
$W_{n}(f)=\int_{S^{2}}\left(1+\frac{1}{4}\left|H_{f}\right|^{2}\right)d\mu_{f}$
where $f$ is a $W^{2,2}$ conformal immersion of $S^{2}$ in the round unit
sphere ${\mathbb{S}}^{n}$ for some $n>2$. We consider the existence of
minimizers of
$\beta_{0}^{n}(y_{1},\dots,y_{m})=\inf\\{W_{n}(f):{y_{1},\dots,y_{m}\in
f(S^{2})}\\}$
where $y_{1},\dots,y_{m}$ are fixed distinct points in ${\mathbb{S}}^{n}$.
When $m\geq 2$, $\beta^{n}_{0}(y_{1},\dots,y_{m})$ is positive by the
conformality of the functions $f$.
###### Proposition 4.2.
When $m\geq 2$, any $\beta_{0}^{n}(y_{1},\dots,y_{m})<8\pi$ is attained by a
$W^{2,2}$-conformally embedded $S^{2}$ in ${\mathbb{S}}^{n}$.
###### Proof.
Let $\\{f_{k}\\}$ be a minimizing sequence of
$\beta_{0}^{n}(y_{1},\dots,y_{m})$. We can consider $f_{k}$ as conformal map
from $S^{2}$ into $\mathbb{R}^{n}$. By Theorem 1, $f_{k}$ will converge to a
mapping $f_{0}$ which is a $W^{2,2}$ branched conformal immersion from a
stratified sphere $\Sigma_{\infty}$ into ${\mathbb{S}}^{n}$ with
$y_{1},\dots,y_{m}\in
f_{0}(\Sigma_{\infty}),\,\,\,\,W_{n}(f_{0})\leq\beta_{0}^{n}(y_{1},\dots,y_{m})<8\pi.$
Composing with a stereographic projection $\Pi$ from ${\mathbb{S}}^{n}$ minus
a point not on $f_{0}(S^{2})$ into $\mathbb{R}^{n}$, we see
$W_{n}(f_{0})=W(\Pi\circ f_{0})$ and
$\theta^{2}_{f_{0}(p)}=\theta^{2}_{\Pi\circ f(p)}$. Now, by (3.2) we have
$\theta^{2}_{f_{0}(p)}\leq\frac{1}{4\pi}W_{n}(f_{0}).$
By Theorem 3.1
$\lambda(p)+1=\theta^{2}_{f(p)}\leq\frac{1}{4\pi}W_{n}(f_{0})<2$
thus $\lambda(p)=0$ which means $f_{0}$ has no branched points. Moreover, that
the area density of $\Sigma_{\infty}$ is one everywhere implies that
$\Sigma_{\infty}$ has only 1 component and $f_{0}$ has no intersection points.
Thus $\Sigma_{\infty}=S^{2}$, and $f_{0}$ is an (Lipschitz) embedding.
$\hfill\Box$
###### Corollary 4.3.
For any $\epsilon>0$, there is a Willmore sphere $f:S^{2}\to{\mathbb{S}}^{n}$
with $W_{n}(f)<4\pi+\epsilon$, which has at least 2 nonremovable singular
points.
###### Proof.
Take five distinct points $y_{1},\dots,y_{5}\in{\mathbb{S}}^{n}$, such that
there is no round 2-sphere passing through all of them. Recall the Willmore
functional $W_{n}$ of a round 2-sphere is $4\pi$. We can choose the five
points to be very closed to a round 2-sphere, such that there is a 2-sphere
$\Sigma$ which is not round and contains $y_{1},\dots,y_{5}$ with
$W_{n}(\Sigma)<4\pi+\epsilon.$
Then we can find a $W^{2,2}$ conformal embedding
$f:S^{2}\rightarrow{\mathbb{S}}^{n}$, such that $f(S^{2})$ passes through
$y_{1},\dots,y_{5}$, and attains $\beta_{0}^{n}(y_{1},\dots,y_{5})$, by
Proposition 4.2.
Choose a point $P\in{\mathbb{S}}^{n}\backslash\Sigma$ as the north pole. Let
$\Pi$ be the stereographic projection from ${\mathbb{S}}^{n}\backslash\\{P\\}$
to $\mathbb{R}^{n}$, and denote $\widetilde{y}_{i}=\Pi(y_{i})$ and
$\widetilde{f}=\Pi(f)$. By the conformal invariance of the Willmore
functional, we have
$W_{n}(f)=\frac{1}{4}\int_{S^{2}}|H_{\widetilde{f}}|^{2}d\mu_{\widetilde{f}}.$
Then $\widetilde{f}$ attains
$\inf\left\\{\frac{1}{4}\int_{S^{2}}|H_{\varphi}|^{2}d\mu_{\varphi}:\varphi\in
W^{2,2}_{conf}(S^{2},\mathbb{R}^{n}),\,\,\,\,\widetilde{y}_{1},\dots,\widetilde{y}_{5}\in\varphi(S^{2})\right\\}.$
Then by results in [27], $\widetilde{f}(S^{2})$ is smooth on
$\widetilde{f}(S^{2})\backslash\\{\widetilde{y}_{1},\dots,\widetilde{y}_{5}\\}$.
However, the Gap Lemma in [14, Theorem 2.7] tells us that there is an
$\epsilon>0$, such that any closed smooth Willmore sphere with Willmore
functional $4\pi+\epsilon$ is a round sphere. Therefore, at least one of
$\widetilde{y}_{1},\dots,\widetilde{y}_{5}$ is a nonremovable singular point.
However, a Willmore sphere cannot have only one singular point, by Lemma 4.2
in [15] (which is true in $\mathbb{R}^{n}$), therefore $\widetilde{f}$ has at
least 2 singular points. $\hfill\Box$
### 4.2. Minimizing Willmore functional subject to area constraint
In this subsection, $N$ stands for a compact closed submanifold of
$\mathbb{R}^{n}$ with induced metric. We say $f\in W^{2,2}_{conf}(\Sigma,h,N)$
if $f\in W^{2,2}_{conf}(\Sigma,h,\mathbb{R}^{n})$ and $f(\Sigma)\subset N$.
For $f\in W^{2,2}_{conf}(\Sigma,h,N)$, we define
$W(f)=W(f,\Sigma,N)=\frac{1}{4}\int_{\Sigma}|H_{f,\Sigma,N}|^{2}d\mu_{f}.$
First, we consider the case of genus zero. Set
$\beta_{0}(N,a)=\inf\\{W(f):\mu(f)=a,\,\,\,\,f\in W^{2,2}_{conf}(S^{2},N)\\}.$
###### Proposition 4.4.
We have
$\lim_{a\rightarrow 0}\beta_{0}(N,a)=4\pi.$
Moreover, when $a$ is sufficiently small, there is an embedding $f\in
W^{2,2}_{conf}(S^{2},N)$, such that
$\mu(f)=a,\,\,\,\,and\,\,\,\,W(f)=\beta_{0}(N,a).$
###### Proof.
First, we show that
(4.3) $\limsup_{a\rightarrow 0}\beta_{0}(N,a)\leq 4\pi.$
Take a point $p\in N$ and a normal coordinate neighborhood $U$ around $p$. Let
$S_{r}=\\{(x^{1},x^{2},x^{3},0,\dots,0)\in
T_{p}N:(x^{1})^{2}+(x^{2})^{2}+(x^{3})^{2}=r^{2}\\}.$
It is easy to check that
$\lim_{r\rightarrow 0}W(\exp_{p}(S_{r}),N)=4\pi.$
For any $a$ which is sufficiently small, we can find $r=r(a)$ such that
$\mu(\exp_{p}(S_{r}))=a$ and $r\rightarrow 0$ as $a\rightarrow 0$. Then (4.3)
follows from $\beta_{0}(N,a)\leq W(\exp_{p}(S_{r}))$.
Next, we prove that $\beta_{0}(N,a)$ can be attained by an embedded 2-sphere.
Let $f_{k}\in W^{2,2}_{conf}(S^{2},N)$ be a minimizing sequence of
$\beta_{0}(N,a)$. By Lemma 4.1 and (4.3), when $a$ is sufficiently small and
$k$ is sufficiently large
$W(f_{k},S^{2},\mathbb{R}^{n})\leq
W(f_{k},S^{2},N)+C\mu(f_{k})<4\pi+\epsilon(a,k)+Ca$
where $\epsilon(a,k)\to 0$ as $a\to 0$ and $k\to\infty$. By Theorem 1,
$\\{f_{k}\\}$ has a limit $f_{0}$, which is a branched conformal immersion
from a stratified sphere $S$ into $N$ with
$\mu(f_{0})=a\,\,\,\,\hbox{and}\,\,\,\,W(f_{0})\leq\beta_{0}(N,a).$
Then by (3.2), for any $p\in S$ it holds
$\theta^{2}(f_{0}(p))<2.$
Thus $S$ is a 2-sphere and $f_{0}$ has no branch points and no self-
intersection points. Hence $f_{0}$ is an embedding. Therefore $f_{0}$ is a
minimizer for $\beta_{0}(N,a)$:
$W(f_{0})=\beta_{0}(N,a).$
Finally, we prove
$\varliminf_{a\rightarrow 0}\beta_{0}(N,a)\geq 4\pi.$
By Lemma 4.1,
$W(f_{0},S^{2},\mathbb{R}^{n})\leq W(f_{0},S^{2},N)+Ca.$
It is well-known that $W(f_{0},S^{2},\mathbb{R}^{n})\geq 4\pi$, which
completes the proof. $\hfill\Box$
We now consider the case of genus larger than 0. Recall a result of Schoen-Yau
[30] and Sacks-Uhlenbeck [29]: If $\varphi:\Sigma\to N$ induces an injection
from the fundamental groups to $\Sigma$ and $N$, then there is a branched
minimal immersion $f:\Sigma\rightarrow N$ so that $f$ induces the same map
between fundamental groups as $\varphi$ and $f$ has least area among all such
maps. If $\pi_{2}(N)=0$ then $f$ is minimizing in its homotopy class. We
denote the area of the branched minimal immersion $f_{\varphi}$ by
$a_{\varphi}$.
Let $g>0$ be the genus of the closed Riemann surface $\Sigma$ and
$\phi:\Sigma\to N$ be a continuous map. Define
$\beta_{g}(N,a,\phi)=\inf\left\\{W(f):f\in\widetilde{W}^{2,2}(\Sigma,N),\,\,\mu(f)=a,\,\,f\sim\phi\right\\},$
where $f\sim\phi$ means that $f$ is homotopic to $\phi$.
###### Proposition 4.5.
Let $\Sigma$ be a closed Riemann surface with genus $g>0$ and $N$ be a compact
Riemannian manifold with $\pi_{2}(N)=0$. Let $\varphi:\Sigma\to N$ be a map
which induces an injective $\varphi_{\\#}:\pi_{1}(\Sigma)\to\pi_{1}(N)$. Then
we can find an $\delta>0$, such that for any
$a\in[a_{\varphi},a_{\varphi}+\delta)$, there is a branched conformal
immersion $f_{0}$ of a smooth Riemann surface $(\Sigma,h)$ of genus $g$ in
$\mathbb{R}^{n}$, such that $\mu(f_{0})=a$ and
$W(f_{0})=\beta_{g}(N,a,\varphi)$ and $f_{0}$ is homotopic to $\varphi$.
Moreover, when $\dim N=3$, we can choose $\delta$ to be small such that
$f_{0}$ is an immersion.
###### Proof.
The proof will be divided into several steps.
Step 1. We prove that $\lim_{a\rightarrow a_{0}}\beta_{g}(N,a,\varphi)=0.$
Let $F\in C^{\infty}(\Sigma\times[0,1],\mathbb{R}^{n})$, such that
$F(\cdot,t)$ is an immersion for each $t$ and
$F(\cdot,0)=f_{\varphi},\,\,\,\,\mu(F(\cdot,1))\geq a_{\varphi}.$
As $F(\cdot,t)\sim\varphi$ and $f_{\varphi}$ is a minimal surface,
$\lim_{a\rightarrow a_{\varphi}}\beta_{p}(N,a,\varphi)\leq\lim_{t\rightarrow
0}W(F(\cdot,t))=W(f_{\varphi})=0.$
Step 2. Smooth convergence of conformal structures.
We take a minimizing sequence $\\{f_{k}\\}$ of $\beta_{g}(N,a,f)$. Recall that
$f_{k}$ are $W^{2,2}$ branched conformal immersions from $(\Sigma,h_{k})$ into
$\mathbb{R}^{n}$, where $h_{k}$ are the smooth metrics with curvature 0 or
$-1$. Because $\pi_{2}(N)=0$ and $f_{k}\sim\varphi$ for each $k$, $f_{k}$
induces the same injective action on the fundamental groups as $\varphi$ does;
hence the conformal structures of $h_{k}$ stay in a compact set of the moduli
space for both the cases $g>1$ and $g=1$, therefore, after passing to a
subsequence if necessary, $\Sigma_{k}=(\Sigma,h_{k})$ converges to a Riemann
surface $(\Sigma,h_{0})$ in $\mathcal{M}_{g}$ (cf. [30]). The results in [30]
applies as $f_{k}$ belong to $W^{1,2}\cap C^{0}$.
Step 3. We prove that $\\{f_{k}\\}$ has no bubbles, i.e. the limit $f_{0}$ is
a map defined on $\Sigma$.
By Remark 2.11, $f_{0}$ is defined on $\Sigma_{\infty}=\Sigma_{0}\cup
S_{1}\cup S_{2}\cdots\cup S_{m}$, where $S_{i}$ are all 2-spheres and
$\Sigma_{0}$ is a smooth surface of genus $g$. We prove $m=0$. Assume $m\geq
1$. By Theorem 1, $\mu(f_{0})=a$ and $W(f_{0})\leq\beta_{p}(N,a,\varphi)$.
Further, $f_{k}(\Sigma)$ converge to $f_{0}(\Sigma_{\infty})$ in Hausdorff
distance and $f_{0}|_{S_{j}}$ is homotopic to a constant map for each
$j=1,\dots,m$ as $\pi_{2}(N)=0$. We conclude that $f_{0}|_{\Sigma_{0}}$ is
homotopic to $\varphi$. Consequently, $\mu(f_{0}(\Sigma_{0}))\geq
a_{\varphi}$. Then we get
$\mu(f_{0},S_{i})\leq\mu(\Sigma)-\mu(\Sigma_{0})\leq
a-a_{\varphi}\,\,\,\,\hbox{and}\,\,\,\,W(f_{0},S_{i},N)\leq\beta_{g}(N,a,\varphi).$
By Lemma 4.1 and Step 1,
$W(f_{0},S_{i},\mathbb{R}^{n})\leq C(a-a_{\varphi})+\beta_{g}(N,a,\varphi)\to
0\,\,\,\,\hbox{as}\,\,\,\,a\to a_{\varphi}.$
This, however, contradicts Proposition 2.2 when the Willmore functional of
$S_{i}$ goes below the gap constant.
Step 4. We consider the case of $\dim N=3$.
We will use the result that there are no branch points for minimal surfaces
[8, 22] to prove that $f_{0}$ has no branch points when $\delta$ is
sufficiently small.
If the claimed result is not true, then there is a sequence of numbers
$a_{k}>a_{\varphi}$ with $a_{k}\rightarrow a_{\varphi}$ and a sequence of
$W^{2,2}$ branched conformal immersions ${f}_{0,k}$ of $(\Sigma,h_{k})$ in $N$
with $\mu(f_{0,k})=a_{k}$, $W(f_{0,k},\Sigma,N)=\beta_{g}(N,a_{k},\varphi)$ by
the first part of the proposition, and each $f_{0,k}$ has at least a branch
point $p_{k}$. By Step 1, $W(f_{0,k},\Sigma,N)\rightarrow 0$.
As in Step 2, $(\Sigma,h_{k})$ converge to a smooth surface $(\Sigma,h_{0})$
in $\mathcal{M}_{g}$. For simplicity, we will still denote
${f}_{0,k}\circ\psi_{k}$ (see Remark 2.11) by ${f}_{0,k}$ which is a branched
conformal immersion from $(\Sigma,\psi^{*}_{k}(h_{k}))$ into $\mathbb{R}^{n}$.
By Theorem 2.8, we may set ${f}_{0,0}$ to be the limit of ${f}_{0,k}$ with
$\mu({f}_{0,0})=a$ and $W({f}_{0,0})=0$. Arguing as in Step 3,
$\\{{f}_{0,k}\\}$ has no bubbles, and ${f}_{0,0}\in
W^{2,2}_{b,c}(\Sigma,h_{0},\mathbb{R}^{n})$ for some smooth $h_{0}$. Moreover,
${f}_{0,0}$ is a minimal surface in $N$. By the result of Gulliver and
Osserman, ${f}_{0,0}$ is a smooth immersion of $\Sigma$ in $N$.
Since $p_{k}$ is a branch point, by Theorem 3.1, the area density
$\theta^{2}_{f_{0,k}(p_{k})}(f_{0,k}(U))\geq 2,$
where $U$ is a neighborhood of $p_{k}\rightarrow p$ in $\Sigma$ for
sufficiently large $k$. As $f_{0,0}$ is immersive, we can take $U$ small so
that ${f}_{0,0}$ is an embedding on $U$ and
$\mu({f}_{0,0}(U))<\epsilon^{\prime}$. Further, by the monotonicity formula
for minimal surfaces, for small $r$ and geodesic balls $B^{N}_{r}(f_{0,0}(p))$
in $N$, it holds
$\mu({f}_{0,0}(U)\cap B^{N}_{r}({f}_{0,0}(p)))\leq(1+\epsilon^{\prime})\pi
r^{2}.$
From the expansion of metric in normal coordinates, for small $r$ and the
Euclidean ball $B_{r}(f_{0,0}(p))$ in $\mathbb{R}^{n}$ we have
$\mu(f_{0,0}(U)\cap B_{r}(f_{0,0}(p)))\leq\mu(f_{0,0}(U)\cap
B^{N}_{r+cr^{2}}(f_{0,0}(p)))\leq(1+\epsilon^{\prime})\pi r^{2}+O(r^{3})$
where $c$ depends on $N$.
In light of Lemma 4.1, $W({f}_{0,k},U,\mathbb{R}^{n})<\epsilon_{0}^{2}$ if we
choose $\epsilon^{\prime}$ to be very small and $k$ large enough. Then
$\\{{f}_{0,k}\\}$ has no blowup points in $U$ by the $\epsilon$-regularity.
Then we have
$\mu({f}_{0,k}(U)\cap B_{r}({f}_{0,k}(p_{k})))\rightarrow\mu({f}_{0,0}(U)\cap
B_{r}({f}_{0,0}(p)))\,\,\,\,\,\,\,\,\hbox{as $k\to\infty$}.$
By Lemma 4.1,
$W({f}_{0,k},U,\mathbb{R}^{n})\leq C\epsilon^{\prime}+W(f_{0,k},U,N).$
Then by (3.3),
$\displaystyle\theta^{2}_{f_{0,k}(p)}(f_{0,k}(U))$ $\displaystyle\leq$
$\displaystyle\frac{\mu(f_{0,k}(U)\cap B_{r}(f_{0,k}(p_{k})))}{\pi
r^{2}}+W(f_{0,k},U,\mathbb{R}^{n})+CW(f_{0,k},U,\mathbb{R}^{n})^{\frac{1}{2}}.$
Hence,
$\displaystyle 2$ $\displaystyle\leq$ $\displaystyle\lim_{U\to
p}\lim_{k\to\infty}\left(\frac{\mu(f_{0,k}(U)\cap B_{r}(f_{0,k}(p_{k})))}{\pi
r^{2}}+W(f_{0,k},U,\mathbb{R}^{n})+CW(f_{0,k},U,\mathbb{R}^{n})^{\frac{1}{2}}\right)$
$\displaystyle\leq$ $\displaystyle 1+\epsilon^{\prime}.$
This is impossible for $\epsilon^{\prime}$ small. $\hfill\Box$
### 4.3. Minimizing Willmore functional of surfaces with a Douglas type
condition
In this subsection, we consider a sufficient condition of Douglas type as in
the minimal surface theory for existence of minimizers of the Willmore
functional.
First, we assume $N$ to be a compact Riemannian manifold with negative
sectional curvatures. In negatively curved $N$, surface area is bounded by the
Willmore functional and the genus of the surface.
###### Lemma 4.6.
Let $N$ be a compact Riemannian manifold with $K\leq-c<0$. Then for any
$f\in\widetilde{W}^{2,2}(\Sigma,N)$,
$\mu(f)\leq c^{-1}\left(W(f,\Sigma,N)-2\pi\chi(\Sigma)\right).$
Especially, when $g(\Sigma)=0$ or $1$,
$\mu(f)\leq c^{-1}W(f,\Sigma,N).$
###### Proof.
From the Gauss equation:
$R^{\Sigma}(X,Y,X,Y)=R^{N}(X,Y,X,Y)+\langle A(X,X),A(Y,Y)\rangle-\langle
A(X,Y),A(X,Y)\rangle.$
we have
$K_{\Sigma}\leq K_{f_{*}(T\Sigma)}+\frac{1}{4}|H_{f,\Sigma,N}|^{2}.$
Then from the generalized Gauss-Bonnet formula - Lemma 3.2, we have
$2\pi\chi(\Sigma)+2\pi b\leq-c\mu_{f}(\Sigma)+W(f,\Sigma,N)$
where $b$ is the number of branch points, in turn
$c\mu_{f}(\Sigma)\leq W(f,\Sigma,N)-2\pi\chi(\Sigma).$
When $g(\Sigma)\leq 1$ the Euler number $\chi(\Sigma)$ is nonnegative, in this
case
$c\mu_{f}(\Sigma)\leq W(f,\Sigma,N).$
Dividing by $c$ yields the desired area bounds. $\hfill\Box$
Recall that any connected stratified surface $\Sigma$ can be written as union
of finitely many connected 2-dimensional components:
$\Sigma=\bigcup_{i}\Sigma_{i}$. Denote the genus of $\Sigma$ and $\Sigma_{i}$
by $g(\Sigma)$ and $g(\Sigma_{i})$, accordingly. We introduce a subset $S(g)$
of all stratified surfaces as follows.
1. (1)
If $g>0$, $S(g)=\left\\{\Sigma:\hbox{$\Sigma=\bigcup_{i}\Sigma_{i}$ with
$g(\Sigma_{i})<g$ for all $i$}\right\\}.$
2. (2)
If $g=0$, $S(0)=\left\\{\Sigma:\hbox{$\Sigma=\bigcup_{i}\Sigma_{i}$ with
$g(\Sigma)=0$ and $i\geq 2$}\right\\}.$
Note that any $\Sigma\in S(g)$ with $g(\Sigma)=g$ must be singular, in the
sense that it has more than one components. Especially,
$S(g)\cap{\mathcal{M}}_{g}=\emptyset$. However, when $g\geq 1$, $S(g)$
contains smooth surfaces of genus $\leq g-1$.
Define
$\displaystyle\alpha^{*}(g)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset N,\Sigma\in S(g)\\}$
$\displaystyle\alpha(g)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset
N,\Sigma\in{\mathcal{M}}_{g}\\}.$
We now state a sufficient condition, similar to the Douglas condition for
minimal surfaces, for existence of minimizers for the Willmore functional.
###### Proposition 4.7.
Let $N$ be a compact Riemannian manifold with negative sectional curvatures.
If $0<\alpha(g)<\alpha^{*}(g)$, then there is a $W^{2,2}$ branched conformal
immersion $f$ from a smooth surface of genus $g$ into $N$ which minimizes the
Willmore functional among all such maps.
###### Proof.
Let $f_{k}:(\Sigma,h_{k})\rightarrow N\hookrightarrow\mathbb{R}^{n}$ be a
minimizing sequence of $\alpha(g)$. By Lemma 4.6, the areas
$\mu(f_{k}(\Sigma))$ are uniformly bounded as well by Lemma 4.6. Since
$\alpha(g)$ is positive, $f_{k}$ cannot converge to a point. Then from Theorem
1, there exists a subsequence of $\\{f_{k}\\}$, still denoted by
$\\{f_{k}\\}$, a limit map $f_{0}\in
W^{2,2}_{b,c}(\Sigma_{\infty},\mathbb{R}^{n})$ from a stratified Riemann
surface $\Sigma_{\infty}$ with $g(\Sigma_{\infty})\leq g$ into
$N\hookrightarrow\mathbb{R}^{n}$, and
$W(f_{0},\Sigma_{\infty},\mathbb{R}^{n})\leq\lim_{k\to\infty}W(f_{k},\Sigma,\mathbb{R}^{n})=\alpha(g).$
We write $\Sigma_{\infty}=\bigcup_{i=1}^{m}\Sigma_{i}$. If
$g(\Sigma_{\infty})=g$, we consider two cases. Case 1: $g(\Sigma_{i})=g$ for
some $i=1,...,m$. In this case,
$W(f_{0}|_{\Sigma_{1}},\Sigma_{1},\mathbb{R}^{n})\leq
W(f_{0},\Sigma_{\infty},\mathbb{R}^{n})=\alpha(g).$
So $f_{0}(\Sigma_{i})$ is a smooth genus $g$ surface attains $\alpha(g)$. Case
2: $g(\Sigma_{i})<g$ for all $i=1,...,m$. Thus $\Sigma_{\infty}\in S(g)$, and
in turn
$\alpha^{*}(g)\leq
W(f_{0},\Sigma_{\infty},\mathbb{R}^{n})\leq\alpha(g)<\alpha^{*}(g).$
This contradiction rules out Case 2. If $g(\Sigma_{\infty})<g$ then
$\Sigma_{g}\in S(g)$. Therefore
$\alpha^{*}(g)\leq
W(f_{0},\Sigma_{\infty},\mathbb{R}^{n})\leq\alpha(g)<\alpha^{*}(g)$
and this is impossible. $\hfill\Box$
Instead of the curvature assumption on $N$, we set, for $0<a<\infty$,
$\displaystyle\gamma^{*}(g,a)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset N,\Sigma\in
S(g),\mu(f(\Sigma))\leq a\\}$ $\displaystyle\gamma(g,a)$ $\displaystyle=$
$\displaystyle\inf\\{W(f,\Sigma,\mathbb{R}^{n}):f\in
W^{2,2}_{b,c}(\Sigma,\mathbb{R}^{n}),f(\Sigma)\subset
N,\Sigma\in{\mathcal{M}}_{g},\mu(f(\Sigma))\leq a\\}.$
Since there is no loss in measures in the limit process, as asserted in
Theorem 1, the same proof above allows us to conclude
###### Proposition 4.8.
Let $N$ be a compact Riemannian manifold. If $0<\gamma(g,a)<\gamma^{*}(g,a)$,
then there is a $W^{2,2}$ branched conformal immersion $f$ from a smooth
surface of genus $g$ into $N$ which minimizes the Willmore functional among
all such maps.
## 5\. Appendix
Wente’s inequality [1, 7, 33] states that if $u\in W^{1,2}_{0}(D)$ solves the
equation
$-\Delta u=\nabla a\,\nabla^{\bot}b,$
then we have
(5.1) $\|u\|_{L^{\infty}(D)}\leq\frac{1}{2\pi}\|\nabla a\|_{L^{2}(D)}\|\nabla
b\|_{L^{2}(D)}.$
and
(5.2) $\|\nabla u\|_{L^{2}(D)}\leq\frac{3}{16\pi}\|\nabla
a\|_{L^{2}(D)}\|\nabla b\|_{L^{2}(D)}.$
###### Lemma 5.1.
Let $u\in W^{1,2}_{0}(D)$ be the unique solution to the equation
$-\Delta u=\nabla a\nabla^{\bot}b,$
where $a,b\in W^{1,2}(D)$. Then $u\in C^{0}(D)$.
###### Proof.
Let $a_{k}\in C^{\infty}(\bar{D})$ with $a_{k}\rightarrow a$ in $W^{1,2}(D)$.
There exist solutions $u_{k}\in W^{2,2}(D)\cap C^{0,\alpha}(D)$ to the
Dirichlet problem
$\displaystyle-\Delta u_{k}$ $\displaystyle=$ $\displaystyle\nabla
a_{k}\nabla^{\bot}b,\,\,\,\,\,\,\,\,\,\,\,\,\mbox{in $D$ }$ $\displaystyle
u_{k}$ $\displaystyle=$ $\displaystyle
0,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\mbox{on
$\partial D$.}$
We have
$-\Delta(u_{k}-u_{m})=\nabla(a_{k}-a_{m})\nabla^{\bot}b.$
Then by (5.1)
$\displaystyle\|u_{k}-u_{m}\|_{C^{0}(D)}$ $\displaystyle=$
$\displaystyle\|u_{k}-u_{m}\|_{L^{\infty}(D)}$ $\displaystyle\leq$
$\displaystyle C\,\|a_{k}-a_{m}\|_{L^{2}(D)}\|b\|_{L^{2}(D)}.$
Hence $u_{k}$ converge in $C^{0}(D)$ to a continuous function $u_{0}$ which
vanished on $\partial D$. By (5.2), we may assume $u_{k}$ converges to $u_{0}$
in $W^{1,2}(D)$ as well. For any smooth $\varphi$,
$\int_{D}\varphi\nabla
a\nabla^{\bot}b=\lim_{k\rightarrow+\infty}\int_{D}\varphi\nabla
a_{k}\nabla^{\bot}b=\lim_{k\rightarrow+\infty}\int_{D}\nabla
u_{k}\nabla\varphi=\int_{D}\nabla u_{0}\nabla\varphi.$
Hence $-\Delta u_{0}=\nabla a\nabla^{\bot}b$, and it follows from the
uniqueness of solution to the Dirichlet problem that $u=u_{0}$, so $u\in
C^{0}(D)$. $\hfill\Box$
Now, we are ready to prove:
###### Proposition 5.2.
If $f\in W^{2,2}_{conf}(D,\mathbb{R}^{n})$ with $df\otimes
df=e^{2u}g_{euclid}$ and $\|u\|_{L^{\infty}}<+\infty$, then $u\in C^{0}(D)$.
###### Proof.
In a complex coordinate $z=x+iy$ on $D$, since $f$ is conformal and the
induced metric is $e^{2u}(dx^{2}+dy^{2})$, we have
$f_{x}\cdot f_{x}=f_{y}\cdot f_{y}=e^{2u}$
and we can take $a=e^{-u}\,f_{x},b=e^{-u}\,f_{y}$ as a local orthonormal frame
for the tangent bundle of $f(D)$. Straight computation shows
$a_{x}=e^{-u}f_{xx}+(e^{-u})_{x}f_{x}=e^{-u}f_{xx}-e^{-3u}(f_{xx}\cdot
f_{x})f_{x}$
which leads to
$a_{x}\cdot a_{x}\leq Ce^{-2u}f_{xx}\cdot f_{xx}\leq
Ce^{2\|u\|_{L^{\infty}}}\left|f_{xx}\right|^{2}$
and similarly
$a_{y}\cdot a_{y}\leq Ce^{2\|u\|_{L^{\infty}}}\left|f_{yy}\right|^{2}.$
Therefore, $a\in W^{1,2}(D)$ and similarly $b\in W^{1,2}(D)$. On the other
hand, we can check
$K_{f}e^{2u}=\nabla a\cdot\nabla^{\perp}b,\,\,\,\,-\Delta u=K_{f}e^{2u}.$
Let $v$ be the harmonic function on $D$ which agrees with $u$ on $\partial D$.
Then
$\displaystyle\Delta(u-v)$ $\displaystyle=$ $\displaystyle\nabla
a\cdot\nabla^{\perp}b\,\,\,\,\,\,\,\,\,\hbox{in $D$}$ $\displaystyle u-v$
$\displaystyle=$ $\displaystyle
0\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\,\hbox{on
$\partial D$}$
and Lemma 5.1 implies $u-v\in C^{0}(D)$, hence $u\in C^{0}(D)$. $\hfill\Box$
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Jingyi Chen | Yuxiang Li
---|---
Department of Mathematics | Department of Mathematical Sciences
University of British Columbia | Tsinghua University
Vancouver,B.C., V6T1Z2, Canada | Beijing 100084, P.R. China
jychen@math.ubc.ca | yxli@math.tsinghua.edu.cn
|
arxiv-papers
| 2011-12-08T12:30:36 |
2024-09-04T02:49:25.089036
|
{
"license": "Public Domain",
"authors": "Jingyi Chen and Yuxiang Li",
"submitter": "Yuxiang Li",
"url": "https://arxiv.org/abs/1112.1818"
}
|
1112.2082
|
2011-204 09.12.2011 Underlying Event measurements with ALICE The ALICE
Collaboration††thanks: See Appendix A for the list of collaboration members
The ALICE Collaboration
We present measurements of Underlying Event observables in pp collisions at
$\sqrt{s}=0.9$ and 7 TeV. The analysis is performed as a function of the
highest charged-particle transverse momentum $p_{\mathrm{T,LT}}$ in the event.
Different regions are defined with respect to the azimuthal direction of the
leading (highest transverse momentum) track: Toward, Transverse and Away. The
Toward and Away regions collect the fragmentation products of the hardest
partonic interaction. The Transverse region is expected to be most sensitive
to the Underlying Event activity. The study is performed with charged
particles above three different $p_{\mathrm{T}}$ thresholds: 0.15, 0.5 and 1.0
${\rm GeV}/c$. In the Transverse region we observe an increase in the
multiplicity of a factor 2-3 between the lower and higher collision energies,
depending on the track $p_{\mathrm{T}}$ threshold considered. Data are
compared to Pythia 6.4, Pythia 8.1 and Phojet. On average, all models
considered underestimate the multiplicity and summed $p_{\mathrm{T}}$ in the
Transverse region by about 10-30%.
## 1 Introduction
The detailed characterization of hadronic collisions is of great interest for
the understanding of the underlying physics. The production of particles can
be classified according to the energy scale of the process involved. At high
transverse momentum transfers ($p_{\mathrm{T}}\gtrsim 2\,\mathrm{\mbox{${\rm
GeV}/c$}}$) perturbative Quantum Chromodynamics (pQCD) is the appropriate
theoretical framework to describe partonic interactions. This approach can be
used to quantify parton yields and correlations, whereas the transition from
partons to hadrons is a non-perturbative process that has to be treated using
phenomenological approaches. Moreover, the bulk of particles produced in high-
energy hadronic collisions originate from low-momentum transfer processes. For
momenta of the order of the QCD scale, $\mathcal{O}$(100 MeV), a perturbative
treatment is no longer feasible. Furthermore, at the center-of-mass energies
of the Large Hadron Collider (LHC), at momentum transfers of a few GeV/$c$,
the calculated QCD cross-sections for 2-to-2 parton scatterings exceed the
total hadronic cross-section [1]. This result indicates that Multiple Partonic
Interactions (MPI) occur in this regime. The overall event dynamics cannot be
derived fully from first principles and must be modeled using phenomenological
calculations. Measurements at different center-of-mass energies are required
to test and constrain these models.
In this paper, we present an analysis of the bulk particle production in pp
collisions at the LHC by measuring the so-called Underlying Event (UE)
activity [2]. The UE is defined as the sum of all the processes that build up
the final hadronic state in a collision excluding the hardest leading order
partonic interaction. This includes fragmentation of beam remnants, multiple
parton interactions and initial- and final-state radiation (ISR/FSR)
associated to each interaction. Ideally, we would like to study the
correlation between the UE and perturbative QCD interactions by isolating the
two leading partons with topological cuts and measuring the remaining event
activity as a function of the transferred momentum scale ($Q^{2}$).
Experimentally, one can identify the products of the hard scattering, usually
the leading jet, and study the region azimuthally perpendicular to it as a
function of the jet energy. Results of such an analysis have been published by
the CDF [2, 3, 4, 5] and STAR [6] collaborations for pp collisions at
$\sqrt{s}=1.8$ and 0.2 TeV, respectively. Alternatively, the energy scale is
given by the leading charged-particle transverse momentum, circumventing
uncertainties related to the jet reconstruction procedure at low
$p_{\mathrm{T}}$. It is clear that this is only an approximation to the
original outgoing parton momentum, the exact relation depends on the details
of the fragmentation mechanism. The same strategy based on the leading charged
particle has recently been applied by the ATLAS [7] and CMS [8]
collaborations.
In the present paper we consider only charged primary particles111Primary
particles are defined as prompt particles produced in the collision and their
decay products (strong and electromagnetic decays), except products of weak
decays of strange particles such as $K^{0}_{S}$ and $\Lambda$., due to the
limited calorimetric acceptance of the ALICE detector systems in azimuth.
Distributions are measured for particles in the pseudorapidity range
$|\eta|<0.8$ with $p_{\mathrm{T}}>p_{\rm T,min}$, where $p_{\rm T,min}=$ 0.15,
0.5 and 1.0 ${\rm GeV}/c$, and are studied as a function of the leading
particle transverse momentum.
Many Monte Carlo (MC) generators for the simulation of pp collisions are
available; see [9] for a recent review discussing for example Pythia [10],
Phojet [11], Sherpa [9] and Herwig [12]. These provide different descriptions
of the UE associated with high energy hadron collisions. A general strategy is
to combine a perturbative QCD treatment of the hard scattering with a
phenomenological approach to soft processes. This is the case for the two
models used in our analysis: Pythia and Phojet. In Pythia the simulation
starts with a hard LO QCD process of the type 2 $\to$ 2\. Multi-jet topologies
are generated with the parton shower formalism and hadronization is
implemented through the Lund string fragmentation model [13]. Each collision
is characterized by a different impact parameter $b$. Small $b$ values
correspond to a large overlap of the two incoming hadrons and to an increased
probability for MPIs. At small $p_{\mathrm{T}}$ values color screening effects
need to be taken into account. Therefore a cut-off $p_{\rm T,0}$ is
introduced, which damps the QCD cross-section for $p_{\mathrm{T}}\ll p_{\rm
T,0}$. This cut-off is one of the main tunable model parameters.
In Pythia version 6.4 [10] MPI and ISR have a common transverse momentum
evolution scale (called interleaved evolution [14]). Version 8.1 [15] is a
natural extension of version 6.4, where the FSR evolution is interleaved with
MPI and ISR and parton rescatterings [16] are considered. In addition initial-
state partonic fluctuations are introduced, leading to a different amount of
color-screening in each event.
Phojet is a two-component event generator, where the soft regime is described
by the Dual Parton Model (DPM) [17] and the high-$p_{\mathrm{T}}$ particle
production by perturbative QCD. The transition between the two regimes happens
at a $p_{\mathrm{T}}$ cut-off value of 3 ${\rm GeV}/c$. A high-energy hadronic
collision is described by the exchange of effective Pomerons. Multiple-Pomeron
exchanges, required by unitarization, naturally introduce MPI in the model.
UE observables allow one to study the interplay of the soft part of the event
with particles produced in the hard scattering and are therefore good
candidates for Monte Carlo tuning. A better understanding of the processes
contributing to the global event activity will help to improve the predictive
power of such models. Further, a good description of the UE is needed to
understand backgrounds to other observables, e.g., in the reconstruction of
high-$p_{\mathrm{T}}$ jets.
The paper is organized in the following way: the ALICE sub-systems used in the
analysis are described in Section 2 and the data samples in Section 3. Section
4 is dedicated to the event and track selection. Section 5 introduces the
analysis strategy. In Sections 6 and 7 we focus on the data correction
procedure and systematic uncertainties, respectively. Final results are
presented in Section 8 and in Section 9 we draw conclusions.
## 2 ALICE detector
Optimized for the high particle densities encountered in heavy-ion collisions,
the ALICE detector is also well suited for the study of pp interactions. Its
high granularity and particle identification capabilities can be exploited for
precise measurements of global event properties [18, 19, 20, 21, 22, 23, 24].
The central barrel covers the polar angle range $45^{\circ}-135^{\circ}$
($|\eta|$ $<$ $1$) and full azimuth. It is contained in the L3 solenoidal
magnet which provides a nominal uniform magnetic field of $0.5\,{\rm T}$. In
this section we describe only the trigger and tracking detectors used in the
analysis, while a detailed discussion of all ALICE sub-systems can be found in
[25].
The V0A and V0C counters consist of scintillators with a pseudorapidity
coverage of $-3.7<\eta<-1.7$ and $2.8<\eta<5.1$, respectively. They are used
as trigger detectors and to reject beam–gas interactions.
Tracks are reconstructed combining information from the two main tracking
detectors in the ALICE central barrel: the Inner Tracking System (ITS) and the
Time Projection Chamber (TPC). The ITS is the innermost detector of the
central barrel and consists of six layers of silicon sensors. The first two
layers, closely surrounding the beam pipe, are equipped with high granularity
Silicon Pixel Detectors (SPD). They cover the pseudorapidity ranges
$|\eta|<2.0$ and $|\eta|<1.4$ respectively. The position resolution is
$12\,{\rm\mu m}$ in $r\phi$ and about $100\,{\rm\mu m}$ along the beam
direction. The next two layers are composed of Silicon Drift Detectors (SDD).
The SDD is an intrinsically 2-dimensional sensor. The position along the beam
direction is measured via collection anodes and the associated resolution is
about $50\,{\rm\mu m}$. The $r\phi$ coordinate is given by a drift time
measurement with a spatial resolution of about $60\,{\rm\mu m}$. Due to drift
field non-uniformities, which were not corrected for in the 2010 data, a
systematic uncertainty of $300\,{\rm\mu m}$ is assigned to the SDD points.
Finally, the two outer layers are made of double-sided Silicon micro-Strip
Detectors (SSD) with a position resolution of $20\,{\rm\mu m}$ in $r\phi$ and
about $800\,{\rm\mu m}$ along the beam direction. The material budget of all
six layers including support and services amounts to 7.7% of a radiation
length.
The main tracking device of ALICE is the Time Projection Chamber that covers
the pseudorapidity range of about $|\eta|<0.9$ for tracks traversing the
maximum radius. In order to avoid border effects, the fiducial region has been
restricted in this analysis to $|\eta|<0.8$. The position resolution along the
$r\phi$ coordinate varies from $1100\,{\rm\mu m}$ at the inner radius to
$800\,{\rm\mu m}$ at the outer. The resolution along the beam axis ranges from
$1250\,{\rm\mu m}$ to $1100\,{\rm\mu m}$.
For the evaluation of the detector performance we use events generated with
the Pythia 6.4 [10] Monte Carlo with tune Perugia-0 [26] passed through a full
detector simulation based on Geant3 [27]. The same reconstruction algorithms
are used for simulated and real data.
## 3 Data samples
The analysis uses two data sets which were taken at the center-of-mass
energies of $\sqrt{s}=$ 0.9 and 7 TeV. In May 2010, ALICE recorded about 6
million good quality minimum-bias events at $\sqrt{s}=$ 0.9 TeV. The
luminosity was of the order of 1026 cm-2 s-1 and, thus, the probability for
pile-up events in the same bunch crossing was negligible. The $\sqrt{s}=$ 7
TeV sample of about 25 million events was collected in April 2010 with a
luminosity of $10^{27}$ cm-2 s-1. In this case the mean number of interactions
per bunch crossing $\mu$ ranges from 0.005 to 0.04. A set of high pile-up
probability runs ($\mu=0.2-2$) was analysed in order to study our pile-up
rejection procedure and determine its related uncertainty. Those runs are
excluded from the analysis.
Corrected data are compared to three Monte Carlo models: Pythia 6.4 (tune
Perugia-0), Pythia 8.1 (tune 1 [15]) and Phojet 1.12.
Collision energy: 0.9 TeV
---
| Events | % of all
Offline trigger | 5,515,184 | 100.0
Reconstructed vertex | 4,482,976 | 81.3
Leading track $p_{\mathrm{T}}>0.15$ ${\rm GeV}/c$ | 4,043,580 | 73.3
Leading track $p_{\mathrm{T}}>0.5$ ${\rm GeV}/c$ | 3,013,612 | 54.6
Leading track $p_{\mathrm{T}}>1.0$ ${\rm GeV}/c$ | 1,281,269 | 23.2
Collision energy: 7 TeV
| Events | % of all
Offline trigger | 25,137,512 | 100.0
Reconstructed vertex | 22,698,200 | 90.3
Leading track $p_{\mathrm{T}}>0.15$ ${\rm GeV}/c$ | 21,002,568 | 83.6
Leading track $p_{\mathrm{T}}>0.5$ ${\rm GeV}/c$ | 17,159,249 | 68.3
Leading track $p_{\mathrm{T}}>1.0$ ${\rm GeV}/c$ | 9,873,085 | 39.3
Table 1: Events remaining after each event selection step.
## 4 Event and track selection
### 4.1 Trigger and offline event selection
Events are recorded if either of the three triggering systems, V0A, V0C or
SPD, has a signal. The arrival time of particles in the V0A and V0C are used
to reject beam–gas interactions that occur outside the nominal interaction
region. A more detailed description of the online trigger can be found in
[20]. An additional offline selection is made following the same criteria but
considering reconstructed information instead of online trigger signals.
For each event a reconstructed vertex is required. The vertex reconstruction
procedure is based on tracks as well as signals in the SPD. Only vertices
within $\pm 10$ cm of the nominal interaction point along the beam axis are
considered. Moreover, we require at least one track with $p_{\mathrm{T}}>$
$p_{\rm T,min}=$ 0.15, 0.5 or 1.0 ${\rm GeV}/c$ in the acceptance
$|\eta|<0.8$.
A pile-up rejection procedure is applied to the set of data taken at
$\sqrt{s}=$ 7 TeV: events with more than one distinct reconstructed primary
vertex are rejected. This cut has a negligible effect on simulated events
without pile-up: only 0.06% of the events are removed. We have compared a
selection of high pile-up probability runs (see Section 3) with a sample of
low pile-up probability runs. The UE distributions differ by 20-25% between
the two samples. After the above mentioned rejection procedure, the difference
is reduced to less than 2%. Therefore, in the runs considered in the analysis,
the effect of pile-up is negligible.
No explicit rejection of cosmic-ray events is applied since cosmic particles
are efficiently suppressed by our track selection cuts [23]. This is further
confirmed by the absence of a sharp enhanced correlation at $\Delta\phi=\pi$
from the leading track which would be caused by almost straight
high-$p_{\mathrm{T}}$ tracks crossing the detector.
Table 1 summarizes the percentage of events remaining after each event
selection step. We do not explicitly select non-diffractive events, although
the above mentioned event selection significantly reduces the amount of
diffraction in the sample. Simulated events show that the event selections
reduce the fraction of diffractive events from 18-33% to 11-16% (Pythia 6.4
and Phojet at 0.9 and 7 TeV). We do not correct for this contribution.
Selection criteria | Value
---|---
Detectors required | ITS,TPC
Minimum number of TPC clusters | 70
Maximum $\chi^{2}$ per TPC cluster | 4
Minimum number of ITS clusters | 3
Minimum number of SPD or $1^{st}$ layer SDD clusters | 1
Maximum $DCA_{Z}$ | 2 cm
Maximum $DCA_{XY}(p_{\mathrm{T}})$ | 7$\sigma$
Table 2: Track selection criteria.
### 4.2 Track cuts
The track cuts are optimized to minimize the contamination from secondary
tracks. For this purpose a track must have at least 3 ITS clusters, one of
which has to be in the first 3 layers. Moreover, we require at least 70 (out
of a maximum of 159) clusters in the TPC drift volume. The quality of the
track fitting measured in terms of the $\chi^{2}$ per space point is required
to be lower than 4 (each space point having 2 degrees of freedom). We require
the distance of closest approach of the track to the primary vertex along the
beam axis (DCAZ) to be smaller than 2 cm. In the transverse direction we apply
a $p_{\mathrm{T}}$ dependent DCAXY cut, corresponding to 7 standard deviations
of its inclusive probability distribution. These cuts are summarized in Table
2.
## 5 Analysis strategy
The Underlying Event activity is characterized by the following observables
[2]:
* •
average charged particle density vs. leading track transverse momentum $p_{\rm
T,LT}$:
$\frac{1}{\Delta\eta\cdot\Delta\Phi}\frac{1}{N_{\rm ev}(p_{\rm T,LT})}N_{\rm
ch}(p_{\rm T,LT})$ (1)
* •
average summed $p_{\mathrm{T}}$ density vs. leading track $p_{\rm T,LT}$:
$\frac{1}{\Delta\eta\cdot\Delta\Phi}\frac{1}{N_{\rm ev}(p_{\rm T,LT})}\sum
p_{\mathrm{T}}(p_{\rm T,LT})$ (2)
* •
$\Delta\phi$-correlation between tracks and the leading track:
$\frac{1}{\Delta\eta}\frac{1}{N_{\rm ev}(p_{\rm T,LT})}\frac{{\rm d}N_{\rm
ch}}{{\rm d}\Delta\phi}$ (3)
(in bins of leading track $p_{\rm T,LT}$).
$N_{\rm ev}$ is the total number of events selected and $N_{\rm ev}(p_{\rm
T,LT})$ is the number of events in a given leading-track transverse-momentum
bin. The first two variables are evaluated in three distinct regions. These
regions, illustrated in Fig. 1, are defined with respect to the leading track
azimuthal angle:
* •
Toward: $|\Delta\phi|<1/3$ $\pi$
* •
Transverse: $1/3$ $\pi<|\Delta\phi|<2/3$ $\pi$
* •
Away: $|\Delta\phi|>2/3$ $\pi$
where $\Delta\phi=\phi_{LT}-\phi$ is defined in $\pm\pi$. In Eq. (1)-(3) the
normalization factor $\Delta\Phi$ is equal to $2/3\pi$, which is the size of
each region. $\Delta\eta=$ 1.6 corresponds to the acceptance in
pseudorapidity. The leading track is not included in the final distributions.
Figure 1: Definition of the regions Toward, Transverse and Away w.r.t. leading
track direction.
## 6 Corrections
We correct for the following detector effects: vertex reconstruction
efficiency, tracking efficiency, contamination from secondary particles and
leading-track misidentification bias. The various corrections are explained in
more detail in the following subsections. We do not correct for the trigger
efficiency since its value is basically 100% for events which have at least
one particle with $p_{\mathrm{T}}>$ 0.15 ${\rm GeV}/c$ in the range
$|\eta|<0.8$. In Table 3 we summarize the maximum effect of each correction on
the measured final observables at the two collision energies for $p_{\rm
T,min}=0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$.
#### Vertex reconstruction
The correction for finite vertex reconstruction efficiency is performed as a
function of the measured multiplicity. Its value is smaller than 0.7% and 0.3%
at $\sqrt{s}=$ 0.9 and $\sqrt{s}=7\,\mathrm{TeV}$, respectively.
#### Tracking efficiency
The tracking efficiency depends on the track level observables $\eta$ and
$p_{\mathrm{T}}$. The projections of the tracking efficiency on the
$p_{\mathrm{T}}$ and $\eta$ axes are shown in Fig. 2. In the pseudorapidity
projection we observe a dip of about 1% at $\eta=0$ due to the central TPC
cathode. The slight asymmetry between positive and negative $\eta$ is due to a
different number of active SPD and SDD modules in the two halves of the
detector. The number of active modules also differs between the data-taking
periods at the two collision energies. Moreover, the efficiency decreases by
5% in the range 1-3 ${\rm GeV}/c$. This is due to the fact that above about 1
${\rm GeV}/c$ tracks are almost straight and can be contained completely in
the dead areas between TPC sectors. Therefore, at high $p_{\mathrm{T}}$ the
efficiency is dominated by geometry and has a constant value of about 80% at
both collision energies. To avoid statistical fluctuations, the estimated
efficiency is fitted with a constant for
$p_{\mathrm{T}}>5\,\mathrm{\mbox{${\rm GeV}/c$}}$ (not shown in the figure).
#### Contamination from secondaries
We correct for secondary tracks that pass the track selection cuts. Secondary
tracks are predominantly produced by weak decays of strange particles (e.g.
$K^{0}_{S}$ and $\Lambda$), photon conversions or hadronic interactions in the
detector material, and decays of charged pions. The relevant track level
observables for the contamination correction are transverse momentum and
pseudorapidity. The correction is determined from detector simulations and is
found to be 15-20% for tracks with $p_{\mathrm{T}}<$ 0.5 ${\rm GeV}/c$ and
saturates at about 2% for higher transverse momenta (see Fig. 3).
Correction | $\sqrt{s}=0.9$ TeV | $\sqrt{s}=7$ TeV
---|---|---
Leading track misidentification | $<5$% | $<8$%
Contamination | $<3$% | $<3$%
Efficiency | $<19$% | $<19$%
Vertex reconstruction | $<0.7$% | $<0.3$%
Table 3: Maximum effect of corrections on final observables for $p_{\rm
T,min}=0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$.
We multiply the contamination estimate by a data-driven coefficient to take
into account the low strangeness yield in the Monte Carlo compared to data
[24]. The coefficient is derived from a fit of the discrepancy between data
and Monte Carlo strangeness yields in the tails of the DCAXY distribution
which are predominantly populated by secondaries. The factor has a maximum
value of 1.07 for tracks with $p_{\mathrm{T}}<$ 0.5 ${\rm GeV}/c$ and is equal
to 1 for $p_{\mathrm{T}}>$ 1.5 ${\rm GeV}/c$. This factor is included in the
Contamination entry in Table 3.
#### Leading-track misidentification
Experimentally, the real leading track can escape detection because of
tracking inefficiency and the detector’s finite acceptance. In these cases
another track (i.e. the sub-leading or sub-sub-leading etc.) will be selected
as the leading one, thus biasing the analysis in two possible ways. Firstly,
the sub-leading track will have a different transverse momentum than the
leading one. We refer to this as leading-track $p_{\mathrm{T}}$ bin migration.
It has been verified with Monte Carlo that this effect is negligible due to
the weak dependence of the final distributions on $p_{\rm T,LT}$. Secondly,
the reconstructed leading track might have a significantly different
orientation with respect to the real one, resulting in a rotation of the
overall event topology. The largest bias occurs when the misidentified leading
track falls in the Transverse region defined by the real leading track.
We correct for leading-track misidentification with a data-driven procedure.
Starting from the measured distributions, for each event the track loss due to
inefficiency is applied a second time to the data (having been applied the
first time naturally by the detector) by rejecting tracks randomly. If the
leading track is considered reconstructed it is used as before to define the
different regions. Otherwise the sub-leading track is used. Since the tracking
inefficiency is quite small (about 20%) applying it on the reconstructed data
a second time does not alter the results significantly. To verify this
statement we compared our results with a two step procedure. In this case the
inefficiency is applied two times on measured data, half of its value at a
time. The correction factor obtained in this way is compatible with the one
step procedure. Furthermore, the data-driven procedure has been tested on
simulated data where the true leading particle is known. We observed a
discrepancy between the two methods, especially at low leading-track
$p_{\mathrm{T}}$ values, which is taken into account in the systematic error.
The maximum leading-track misidentification correction is 8% on the final
distributions.
Figure 2: Tracking efficiency vs. track $p_{\mathrm{T}}$ (left, $|\eta|<0.8$)
and $\eta$ (right, $p_{\mathrm{T}}>0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$) from a
Pythia 6.4 and Geant3 simulation.
#### Two-track effects
By comparing simulated events corrected for single-particle efficiencies with
the input Monte Carlo, we observe a 0.5% discrepancy around $\Delta\phi=0$.
This effect is called non-closure in Monte Carlo (it will be discussed further
in Section 7) and in this case is related to small two-track resolution
effects. Data are corrected for this discrepancy.
## 7 Systematic uncertainties
In Tables 4, 5 and 6 we summarize the systematic uncertainties evaluated in
the analysis for the three track thresholds: $p_{\mathrm{T}}>0.15$, 0.5 and
1.0 ${\rm GeV}/c$. Each uncertainty is explained in more detail in the
following subsections. Uncertainties which are constant as a function of
leading-track $p_{\mathrm{T}}$ are listed in Table 4. Leading-track
$p_{\mathrm{T}}$ dependent uncertainties are summarized in Tables 5 and 6 for
$\sqrt{s}=0.9\,\mathrm{TeV}$ and $7\,\mathrm{TeV}$, respectively. Positive and
negative uncertainties are propagated separately, resulting in asymmetric
final uncertainties.
#### Particle composition
The tracking efficiency and contamination corrections depend slightly on the
particle species mainly due to their decay length and absorption in the
material. To assess the effect of an incorrect description of the particle
abundances in the Monte Carlo, we varied the relative yields of pions,
protons, kaons, and other particles by 30% relative to the default Monte Carlo
predictions. The maximum variation of the final values is 0.9% and represents
the systematic uncertainty related to the particle composition (see Table 4).
Moreover, we have compared our assessment of the underestimation of
strangeness yields with a direct measurement from the ALICE collaboration
[24]. Based on the discrepancy between the two estimates, we assign a
systematic uncertainty of 0-2.3% depending on the $p_{\mathrm{T}}$ threshold
and collision energy, see Tables 5 and 6.
Figure 3: Contamination correction: correction factor vs. track
$p_{\mathrm{T}}$ (left, $|\eta|<0.8$) and $\eta$ (right,
$p_{\mathrm{T}}>0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$) from a Pythia 6.4 and
Geant3 simulation.
#### ITS and TPC efficiency
The tracking efficiency depends on the level of precision of the description
of the ITS and TPC detectors in the simulation and the modeling of their
response. After detector alignment with survey methods, cosmic-ray events and
pp collision events [28], the uncertainty on the efficiency due to the ITS
description is estimated to be below 2% and affects only tracks with
$p_{\mathrm{T}}<$ 0.3 ${\rm GeV}/c$. The uncertainty due to the TPC reaches
4.5% at very low $p_{\mathrm{T}}$ and is smaller than 1.2% for tracks with
$p_{\mathrm{T}}>$ 0.5 ${\rm GeV}/c$. The resulting maximum uncertainty on the
final distributions is below 1.9%. Moreover, an uncertainty of 1% is included
to account for uncertainties in the MC description of the matching between TPC
and ITS tracks (see Table 4).
| $\sqrt{s}=0.9$ TeV
---|---
| $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Particle composition | $\pm$ 0.9% | $\pm$ 0.7% | $\pm$ 0.4%
ITS efficiency | $\pm$ 0.6% | – | –
TPC efficiency | $\pm$ 1.9% | $\pm$ 0.8% | $\pm$ 0.4%
Track cuts | ${}^{+\ 3.0\%}_{-\ 1.1\%}$ | ${}^{+\ 2.0\%}_{-\ 1.1\%}$ | ${}^{+\ 0.9\%}_{-\ 1.5\%}$
ITS/TPC matching | $\pm$ 1.0% | $\pm$ 1.0% | $\pm$ 1.0%
MC dependence | $+$ 1.1% , $+$ 1.1% , $+$ 1.6% | $+$ 0.9% | $+$ 0.9% , $+$ 0.9% , $+$ 1.3%
Material budget | $\pm$ 0.6% | $\pm$ 0.2% | $\pm$ 0.2%
| $\sqrt{s}=7$ TeV
| $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Particle composition | $\pm$ 0.9% | $\pm$ 0.7% | $\pm$ 0.5%
ITS efficiency | $\pm$ 0.5% | – | –
TPC efficiency | $\pm$ 1.8% | $\pm$ 0.8% | $\pm$ 0.5%
Track cuts | ${}^{+\ 2.1\%}_{-\ 2.3\%}$ | ${}^{+\ 1.6\%}_{-\ 3.2\%}$ | ${}^{+\ 2.5\%}_{-\ 3.5\%}$
ITS/TPC matching | $\pm$ 1.0% | $\pm$ 1.0% | $\pm$ 1.0%
MC dependence | $+$ 0.8% , $+$ 0.8% , $+$ 1.2% | $+$ 0.8% | $+$ 1.0%
Material budget | $\pm$ 0.6% | $\pm$ 0.2% | $\pm$ 0.2%
Table 4: Constant systematic uncertainties at both collision energies. When
more than one number is quoted, separated by a comma, the first value refers
to the number density distribution, the second to the summed $p_{\mathrm{T}}$
and the third to the azimuthal correlation. Some of the uncertainties are
quoted asymmetrically.
#### Track cuts
By applying the efficiency and contamination corrections we correct for those
particles which are lost due to detector effects and for secondary tracks
which have not been removed by the selection cuts. These corrections rely on
detector simulations and therefore, one needs to estimate the systematic
uncertainty introduced in the correction procedure by one particular choice of
track cuts. To do so, we repeat the analysis with different values of the
track cuts, both for simulated and real data. The variation of the final
distributions with different track cuts is a measure of the systematic
uncertainty. The overall effect, considering all final distributions, is
smaller than 3.5% at both collision energies (see Table 4).
#### Misidentification bias
The uncertainty on the leading-track misidentification correction is estimated
from the discrepancy between the data-driven correction used in the analysis
and that based on simulations. The effect influences only the first two
leading-track $p_{\mathrm{T}}$ bins at both collision energies. The maximum
uncertainty ($\sim 18\%$) affects the first leading-track $p_{\mathrm{T}}$ bin
for the track $p_{\mathrm{T}}$ cut-off of 0.15 GeV/$c$. In all other bins this
uncertainty is of the order of few percent. As summarized in Tables 5 and 6,
the uncertainty has slightly different values for the various UE
distributions.
#### Vertex-reconstruction efficiency
The analysis accepts reconstructed vertices with at least one contributing
track. We repeat the analysis requiring at least two contributing tracks. The
systematic uncertainty related to the vertex reconstruction efficiency is
given by the maximum variation in the final distributions between the cases of
one and two contributing tracks. Its value is 2.4% for $p_{\rm T,min}=$ 0.15
GeV/$c$ and below 1% for the other cut-off values (see Tables 5 and 6). The
effect is only visible in the first leading-track $p_{\mathrm{T}}$ bin.
#### Non-closure in Monte Carlo
By correcting a Monte Carlo prediction after full detector simulation with
corrections extracted from the same generator, we expect to obtain the input
Monte Carlo prediction within the statistical uncertainty. This consideration
holds true only if each correction is evaluated with respect to all the
variables to which the given correction is sensitive. Any statistically
significant difference between input and corrected distributions is referred
to as non-closure in Monte Carlo.
The overall non-closure effect is sizable ($\sim 17\%$) in the first leading-
track $p_{\mathrm{T}}$ bin and is 0.6-5.3% in all other bins at both collision
energies.
#### Monte-Carlo dependence
The difference in final distributions when correcting the data with Pythia 6.4
or Phojet generators is of the order of 1% and equally affects all the
leading-track $p_{\mathrm{T}}$ bins.
#### Material budget
The material budget has been measured by reconstructing photon conversions
which allows a precise $\gamma$-ray tomography of the ALICE detector. For the
detector regions important for this analysis the remaining uncertainty on the
extracted material budget is less than 7%. Varying the material density in the
detector simulation, the effect on the observables presented is determined to
be 0.2-0.6% depending on the $p_{\mathrm{T}}$ threshold considered.
| | $\sqrt{s}=0.9$ TeV
---|---|---
| | Number density
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ (17.8, 16.3, 16.3)% | $+$ (4.6, 3.5, 3.5)% | $+$ (4.2, 2.9, 1.7)%
| $2^{nd}$ bin | $+$ 2.9% | $+$ 1.3% | –
MC non closure | $1^{st}$ bin | $-$ 17.2% | $-$ 3.6% | $-$ 1.2%
| $2^{nd}$ bin | $-$ 3.2% | $-$ 0.8% | $-$ 1.2%
| others | $-$ 0.6% | $-$ 0.8% | $-$ 1.2%
Strangeness | $1^{st}$ bin | $\pm$ 1.9% | $\pm$ 0.2% | –
| others | $\pm$ 1.0% | $\pm$ 0.2% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.7% | $-$ 0.5%
| | Summed $p_{\mathrm{T}}$
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ (20.0, 18.1, 18.1)% | $+$ (5.3, 4.1, 4.1)% | $+$ (4.8, 3.4, 3.4)%
| $2^{nd}$ bin | $+$ 3.7% | $+$ 1.6% | –
MC non closure | $1^{st}$ bin | $-$ 17.0% | $-$ 2.8% | $-$ 1.1%
| $2^{nd}$ bin | $-$ 3.0% | $-$ 1.0% | $-$ 1.1%
| others | $-$ 0.7% | $-$ 1.0% | $-$ 1.1%
Strangeness | $1^{st}$ bin | $\pm$ 1.9% | $\pm$ 0.2% | –
| others | $\pm$ 1.0% | $\pm$ 0.2% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.7% | $-$ 0.5%
| | Azimuthal correlation
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ 12.0% | $+$ 3.9% | $+$ 2.5%
| $2^{nd}$ bin | $+$ 2.6% | $+$ 1.1% | –
MC non closure | $1^{st}$ bin | $-$ 17.1% | $-$ 3.3% | $-$ 1.6%
| $2^{nd}$ bin | $-$ 3.5% | $-$ 3.0% | $-$ 1.6%
| others | $-$ 2.4% | $-$ 3.0% | $-$ 1.6%
Strangeness | $1^{st}$ bin | $\pm$ 1.9% | $\pm$ 0.2% | –
| others | $\pm$ 1.0% | $\pm$ 0.2% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.4% | –
| others | $-$ 0.5% | $-$ 0.4% | –
Table 5: Systematic uncertainties vs. leading track $p_{\mathrm{T}}$ at $\sqrt{s}=0.9\,\mathrm{TeV}$. When more than one number is quoted, separated by a comma, the first value refers to the Toward, the second to the Transverse and the third to the Away region. The second column denotes the leading track $p_{\mathrm{T}}$ bin for which the uncertainty applies. The numbering starts for each case from the first bin above the track $p_{\mathrm{T}}$ threshold. | | $\sqrt{s}=7$ TeV
---|---|---
| | Number density
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ (17.9, 16.3, 16.3)% | $+$ (4.0, 3.2, 3.2)% | $+$ (2.5, 1.2, 1.2)%
| $2^{nd}$ bin | $+$ 2.7% | – | $+$ 0.7%
MC non closure | $1^{st}$ bin | $-$ 16.8% | $-$ 2.6% | $-$ 1.9%
| $2^{nd}$ bin | $-$ 2.9% | $-$ 1.4% | $-$ 1.9%
| others | $-$ 0.6% | $-$ 1.0% | $-$ 1.9%
Strangeness | $1^{st}$ bin | $\pm$ 1.8% | $\pm$ 2.3% | –
| others | $\pm$ 1.0% | $\pm$ 2.3% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.7% | $-$ 0.5%
| | Summed $p_{\mathrm{T}}$
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ (20.0, 17.9, 17.9)% | $+$ (4.9, 3.8, 3.8)% | $+$ (3.4, 1.9, 1.9)%
| $2^{nd}$ bin | $+$ 3.4% | $+$ 0.8% | $+$ 1.1%
MC non closure | $1^{st}$ bin | $-$ 16.7% | $-$ 2.7% | $-$ 1.5%
| $2^{nd}$ bin | $-$ 2.6% | $-$ 1.2% | $-$ 1.5%
| others | $-$ 0.8% | $-$ 1.0% | $-$ 1.5%
Strangeness | $1^{st}$ bin | $\pm$ 1.8% | $\pm$ 2.3% | –
| others | $\pm$ 1.0% | $\pm$ 2.3% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.7% | $-$ 0.5%
| | Azimuthal correlation
| $p_{\rm T,LT}$ | $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$ | $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Lead. track misid. | $1^{st}$ bin | $+$ 16.8% | $+$ 3.4% | $+$ 0.9%
| $2^{nd}$ bin | $+$ 2.5% | – | –
MC non closure | $1^{st}$ bin | $-$ 25.3% | $-$ 4.3% | $-$ 1.2%
| $2^{nd}$ bin | $-$ 5.3% | $-$ 2.1% | $-$ 1.2%
| others | $-$ 2.1% | $-$ 2.1% | $-$ 1.2%
Strangeness | $1^{st}$ bin | $\pm$ 1.8% | $\pm$ 2.3% | –
| others | $\pm$ 1.0% | $\pm$ 2.3% | –
Vertex reco. | $1^{st}$ bin | $-$ 2.4% | $-$ 0.4% | –
| others | $-$ 0.5% | $-$ 0.4% | –
Table 6: Systematic uncertainties vs. leading track $p_{\mathrm{T}}$ at
$\sqrt{s}=7\,\mathrm{TeV}$. When more than one number is quoted, separated by
a comma, the first value refers to the Toward, the second to the Transverse
and the third to the Away region. The second column denotes the leading track
$p_{\mathrm{T}}$ bin for which the uncertainty applies. The numbering starts
for each case from the first bin above the track $p_{\mathrm{T}}$ threshold.
## 8 Results
In this section we present and discuss the corrected results for the three UE
distributions in all regions at the two collision energies. The upper part of
each plot shows the relevant measured distribution (black points) compared to
a set of Monte Carlo predictions (coloured curves). Shaded bands represent the
systematic uncertainty only. Error bars along the $x$ axis indicate the bin
width. The lower part shows the ratio between Monte Carlo and data. In this
case the shaded band is the sum in quadrature of statistical and systematic
uncertainties.
The overall agreement of data and simulations is of the order of 10-30% and we
were not able to identify a preferred model that can reproduce all measured
observables. In general, all three generators underestimate the event activity
in the Transverse region. Nevertheless, an agreement of the order of 20% has
to be considered a success, considering the complexity of the system under
study. Even though an exhaustive comparison of data with the latest models
available is beyond the scope of this paper, in the next sections we will
indicate some general trends observed in the comparison with the chosen
models.
In the following discussion we define the leading track $p_{\mathrm{T}}$ range
from 4 to 10 ${\rm GeV}/c$ at $\sqrt{s}=0.9\,\mathrm{TeV}$ and from 10 to 25
${\rm GeV}/c$ at $\sqrt{s}=7\,\mathrm{TeV}$ as the plateau.
| $\sqrt{s}=0.9$ TeV
---|---
| Number density | Summed $p_{\mathrm{T}}$
| Slope $(\mbox{${\rm GeV}/c$})^{-1}$ | Mean | Slope | Mean $(\mbox{${\rm GeV}/c$})$
$p_{\mathrm{T}}>$ 0.15 GeV/c | 0.00 $\pm$ 0.02 | 1.00 $\pm$ 0.04 | 0.00 $\pm$ 0.01 | 0.62 $\pm$ 0.02
$p_{\mathrm{T}}>$ 0.5 GeV/c | 0.00 $\pm$ 0.01 | 0.45 $\pm$ 0.02 | 0.01 $\pm$ 0.01 | 0.45 $\pm$ 0.02
$p_{\mathrm{T}}>$ 1.0 GeV/c | 0.003 $\pm$ 0.003 | 0.16 $\pm$ 0.01 | 0.006 $\pm$ 0.005 | 0.24 $\pm$ 0.01
| $\sqrt{s}=7$ TeV
| Number density | Summed $p_{\mathrm{T}}$
| Slope $(\mbox{${\rm GeV}/c$})^{-1}$ | Mean | Slope | Mean $(\mbox{${\rm GeV}/c$})$
$p_{\mathrm{T}}>$ 0.15 GeV/c | 0.00 $\pm$ 0.01 | 1.82 $\pm$ 0.06 | 0.01 $\pm$ 0.01 | 1.43 $\pm$ 0.05
$p_{\mathrm{T}}>$ 0.5 GeV/c | 0.005 $\pm$ 0.007 | 0.95 $\pm$ 0.03 | 0.01 $\pm$ 0.01 | 1.15 $\pm$ 0.04
$p_{\mathrm{T}}>$ 1.0 GeV/c | 0.001 $\pm$ 0.003 | 0.41 $\pm$ 0.01 | 0.008 $\pm$ 0.006 | 0.76 $\pm$ 0.03
| $\sqrt{s}=1.8$ TeV (CDF)
| Number density (at leading charged jet
$p_{\mathrm{T}}=20\,\mathrm{\mbox{${\rm GeV}/c$}}$)
$p_{\mathrm{T}}>$ 0.5 GeV/c | 0.60
Table 7: Saturation values in the Transverse region for the two collision
energies. The result from CDF is also given, for details see text.
### 8.1 Number density
In Fig. 4-6 we show the multiplicity density as a function of leading track
$p_{\mathrm{T}}$ in the three regions: Toward, Transverse and Away. Toward and
Away regions are expected to collect the fragmentation products of the two
back-to-back outgoing partons from the elementary hard scattering. We observe
that the multiplicity density in these regions increases monotonically with
the $p_{\rm T,LT}$ scale. In the Transverse region, after a monotonic increase
at low leading track $p_{\mathrm{T}}$, the distribution tends to flatten out.
The same behaviour is observed at both collision energies and all values of
$p_{\rm T,min}$.
The rise with $p_{\rm T,LT}$ has been interpreted as evidence for an impact
parameter dependence in the hadronic collision [29]. More central collisions
have an increased probability for MPI, leading to a larger transverse
multiplicity. Nevertheless, we must be aware of a trivial effect also
contributing to the low $p_{\rm T,LT}$ region. For instance for any
probability distribution, the maximum value per randomized sample averaged
over many samples rises steadily with the sample size $M$. In our case, the
conditional probability density $\mathcal{P}(p_{\rm T,LT}|M)$ shifts towards
larger $p_{\rm T,LT}$ with increasing $M$. Using Bayes’ theorem one expects
the conditional probability density $\mathcal{P}(M|p_{\rm T,LT})$ to shift
towards larger $M$ with rising $p_{\rm T,LT}$:
$\mathcal{P}(M|p_{\rm T,LT})\sim\mathcal{P}(p_{\rm T,LT}|M)\mathcal{P}(M).$
(4)
The saturation of the distribution at higher values of $p_{\rm T,LT}$
indicates the onset of the event-by-event partitioning into azimuthal regions
containing the particles from the hard scattering and the UE region. The bulk
particle production becomes independent of the hard scale.
The plateau range is fitted with a line. The fit slopes, consistent with zero,
and mean values for the three $p_{\mathrm{T}}$ thresholds are reported in
Table 7. In the fit, potential correlations of the systematic uncertainties in
different $p_{\mathrm{T}}$ bins are neglected.
ATLAS has published a UE measurement where the hard scale is given by the
leading track $p_{\mathrm{T}}$, with a $p_{\mathrm{T}}$ threshold for
particles of 0.5 ${\rm GeV}/c$ and an acceptance of $|\eta|<2.5$ [7]. Given
the different acceptance with respect to our measurement, the results in the
Toward and Away regions are not comparable. On the other hand the mean values
of the Transverse plateaus from the two measurements are in good agreement,
indicating an independence of the UE activity on the pseudorapidity range. The
CDF collaboration measured the UE as a function of charged particle jet
$p_{\mathrm{T}}$ at a collision energy of 1.8 TeV[2]. The particle
$p_{\mathrm{T}}$ threshold is 0.5 ${\rm GeV}/c$ and the acceptance $|\eta|<1$.
In the Transverse region CDF measures 3.8 charged particles per unit
pseudorapidity above the $p_{\mathrm{T}}$ threshold at leading-jet
$p_{\mathrm{T}}=20\,\mathrm{\mbox{${\rm GeV}/c$}}$. This number needs to be
divided by $2\pi$ in order to be compared with the average number of particles
in the plateau from Table 7 at the same threshold value. The scaled CDF result
is 0.60, also shown in Table 7 for comparison. As expected it falls between
our two measurements at $\sqrt{s}=0.9\,\mathrm{\mbox{${\rm TeV}$}}$ and
$\sqrt{s}=7\,\mathrm{\mbox{${\rm TeV}$}}$. The values do not scale linearly
with the collision energy, in particular the increase is higher from 0.9 to
1.8 ${\rm TeV}$ than from 1.8 to 7 ${\rm TeV}$. Interpolating between our
measurements assuming a logarithmic dependence on $\sqrt{s}$ results in 0.62
charged particles per unit area at 1.8 ${\rm TeV}$, consistent with the CDF
result.
For illustration, Figure 7 presents the number density in the plateau of the
Transverse region for $p_{\mathrm{T}}>0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$ (our
measurement as well as the value measured by CDF at 1.8 TeV) compared with
$dN_{\rm ch}/d\eta|_{\eta=0}$ of charged particles with $p_{\mathrm{T}}>$ 0.5
${\rm GeV}/c$ in minimum-bias events [30] (scaled by $1/2\pi$).222These data
are for events that have at least one charged particle in $|\eta|<2.5$. The UE
activity in the plateau region is more than a factor 2 larger than the
$dN_{\rm ch}/d\eta$. Both can be fitted with a logarithmic dependence on $s$
($a+b\ln{s}$). The relative increase from 0.9 to 7 TeV for the UE is larger
than that for the $dN_{\rm ch}/d\eta$: about 110% compared to about 80%,
respectively.
In Fig. 8 (left) we show the ratio between the number density distribution at
$\sqrt{s}=7\,\mathrm{\mbox{${\rm TeV}$}}$ and
$\sqrt{s}=0.9\,\mathrm{\mbox{${\rm TeV}$}}$. Most of the systematic
uncertainties are expected to be correlated between the two energies,
therefore we consider only statistical uncertainties. The ratio saturates for
leading track $p_{\mathrm{T}}>4\,\mathrm{\mbox{${\rm GeV}/c$}}$. The results
of a constant fit in the range $4<p_{\mathrm{T,LT}}<10\,\mathrm{\mbox{${\rm
GeV}/c$}}$ are reported in Table 8. The measured scaling factor for a
$p_{\mathrm{T}}$ threshold of 0.5 ${\rm GeV}/c$ is in agreement with the
observations of ATLAS [7, 31] and CMS [32].
For the track threshold $p_{\mathrm{T}}>0.15\,\mathrm{\mbox{${\rm GeV}/c$}}$
all models underestimate the charged multiplicity in the Transverse and Away
regions. In particular at $\sqrt{s}=7\,\mathrm{\mbox{${\rm TeV}$}}$ PHOJET
predictions largely underestimate the measurement in the Transverse region (up
to $\sim 50\%$), the discrepancy being more pronounced with increasing
$p_{\mathrm{T}}$ cut-off value. Pythia 8 correctly describes the Toward region
at both collision energies and Phojet only at
$\sqrt{s}=0.9\,\mathrm{\mbox{${\rm TeV}$}}$. For track $p_{\mathrm{T}}>$ 1
${\rm GeV}/c$, Pythia 8 systematically overestimates the event activity in the
jet fragmentation regions (Toward and Away).
| Number density | Summed $p_{\mathrm{T}}$
---|---|---
$p_{\mathrm{T}}>0.15\,\mathrm{\mbox{${\rm GeV}/c$}}$ | 1.76 $\pm$ 0.02 | 2.00 $\pm$ 0.03
$p_{\mathrm{T}}>0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$ | 1.97 $\pm$ 0.03 | 2.16 $\pm$ 0.03
$p_{\mathrm{T}}>1.0\,\mathrm{\mbox{${\rm GeV}/c$}}$ | 2.32 $\pm$ 0.04 | 2.48 $\pm$ 0.05
Table 8: Constant fit in $4<p_{\mathrm{T,LT}}<10\,\mathrm{GeV/\textit{c}}$ to
the ratio between $\sqrt{s}=0.9\,\mathrm{TeV}$ and $\sqrt{s}=7\,\mathrm{TeV}$
for number density (left) and summed $p_{\mathrm{T}}$ (right) distributions in
the Transverse region. The shown uncertainties are based on statistical and
systematic uncertainties summed in quadrature.
### 8.2 Summed $p_{\mathrm{T}}$
In Fig. 9-11 we show the summed $p_{\mathrm{T}}$ density as a function of
leading track $p_{\mathrm{T}}$ in the three topological regions. The shape of
the distributions follows a trend similar to that discussed above for the
number density.
The general trend of Pythia 8 is to overestimate the fragmentation in the
Toward region at all $p_{\mathrm{T}}$ cut-off values. Also in this case at
$\sqrt{s}=7\,\mathrm{\mbox{${\rm TeV}$}}$ PHOJET largely underestimates the
measurement in the Transverse region (up to $\sim 50\%$), especially at higher
values of $p_{\mathrm{T}}$ cut-off. Other systematic trends are not very
pronounced.
In Table 7 we report the mean value of a linear fit in the plateau range. Our
results agree with the ATLAS measurement in the Transverse plateau.
In Fig. 8 (right) we show the ratio between the distribution at
$\sqrt{s}=7\,\mathrm{\mbox{${\rm TeV}$}}$ and
$\sqrt{s}=0.9\,\mathrm{\mbox{${\rm TeV}$}}$, considering as before only
statistical errors. The results of a constant fit in the range
$4<p_{\mathrm{T,LT}}<10\,\mathrm{\mbox{${\rm GeV}/c$}}$ are reported in Table
8. Also in this case the scaling factor is in agreement with ATLAS and CMS
results.
The summed $p_{\mathrm{T}}$ density in the Transverse region can be
interpreted as a measurement of the UE activity in a given leading track
$p_{\mathrm{T}}$ bin. Therefore, its value in the plateau can be used, for
example, to correct jet spectra.
### 8.3 Azimuthal correlation
In Fig. 12-22 azimuthal correlations between tracks and the leading track are
shown in different ranges of leading track $p_{\mathrm{T}}$. The range
$1/3\pi<|\Delta\phi|<2/3\pi$ corresponds to the Transverse region. The regions
$-1/3\pi<\Delta\phi<1/3\pi$ (Toward) and $2/3\pi<|\Delta\phi|<\pi$ (Away)
collect the fragmentation products of the leading and sub-leading jets. In
general, all Monte Carlo simulations considered fail to reproduce the shape of
the measured distributions. Pythia 8 provides the best prediction for the
Transverse activity in all leading track $p_{\mathrm{T}}$ ranges considered.
Unfortunately the same model significantly overestimates the jet fragmentation
regions.
## 9 Conclusions
We have characterized the Underlying Event in pp collisions at $\sqrt{s}=$ 0.9
and 7 ${\rm TeV}$ by measuring the number density, the summed $p_{\mathrm{T}}$
distribution and the azimuthal correlation of charged particles with respect
to the leading particle. The analysis is based on about $6\cdot 10^{6}$
minimum bias events at $\sqrt{s}=$ 0.9 ${\rm TeV}$ and $25\cdot 10^{6}$ events
at $\sqrt{s}=$ 7 ${\rm TeV}$ collected during the data taking periods from
April to July 2010. Measured data have been corrected for detector related
effects; in particular we applied a data-driven correction to account for the
misidentification of the leading track. The fully corrected final
distributions are compared with Pythia 6.4, Pythia 8 and Phojet, showing that
pre-LHC tunes have difficulties describing the data. These results are an
important ingredient in the required retuning of those generators.
Among the presented distributions, the Transverse region is particularly
sensitive to the Underlying Event. We find that the ratio between the
distributions at $\sqrt{s}=$ 0.9 and 7 ${\rm TeV}$ in this region saturates at
a value of about 2 for track $p_{\mathrm{T}}>$ 0.5 ${\rm GeV}/c$. The summed
$p_{\mathrm{T}}$ distribution rises slightly faster as a function of
$\sqrt{s}$ than the number density distribution, indicating that the available
energy tends to increase the particle’s transverse momentum in addition to the
multiplicity. This is in qualitative agreement with an increased relative
contribution of hard processes to the Underlying Event with increasing
$\sqrt{s}$. Moreover, the average number of particles at large $p_{\rm T,LT}$
in the Transverse region seems to scale logarithmically with the collision
energy. In general our results are in good qualitative and quantitative
agreement with measurements from other LHC experiments (ATLAS and CMS) and
show similar trends to that of the Tevatron (CDF).
Our results show that the activity in the Transverse region increases
logarithmically and faster than $dN_{\rm ch}/d\eta$ in minimum-bias events.
Models aiming to correctly reproduce these minimum-bias and underlying event
distributions need a precise description of the interplay of the hard process,
the associated initial and final-state radiation and multiple parton
interactions.
## Number Density - track $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$
Figure 4: Number density in Toward (top), Transverse (middle) and Away
(bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV (right).
Right and left vertical scales differ by a factor 2. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
## Number Density - track $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$
Figure 5: Number density in Toward (top), Transverse (middle) and Away
(bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV (right).
Right and left vertical scales differ by a factor 2. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
## Number Density - track $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Figure 6: Number density in Toward (top), Transverse (middle) and Away
(bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV (right).
Right and left vertical scales differ by a factor 2. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 7: Comparison of number density in the plateau of the Transverse region
(see Table 8) and $dN_{\rm ch}/d\eta$ in minimum-bias events (scaled by
$1/2\pi$) [30]. Both are for charged particles with
$p_{\mathrm{T}}>0.5\,\mathrm{\mbox{${\rm GeV}/c$}}$. For this plot,
statistical and systematic uncertainties have been summed in quadrature. The
lines show fits with the functional form $a+b\ln{s}$.
Figure 8: Ratio between $\sqrt{s}=0.9\,\mathrm{TeV}$ and
$\sqrt{s}=7\,\mathrm{TeV}$ for number density (left) and summed
$p_{\mathrm{T}}$ (right) distributions in the Transverse region. Statistical
uncertainties only.
## Summed $p_{\mathrm{T}}$ \- track $p_{\mathrm{T}}>0.15\,\mbox{${\rm
GeV}/c$}$
Figure 9: Summed $p_{\mathrm{T}}$ in Toward (top), Transverse (middle) and
Away (bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV
(right). Right and left vertical scales differ by a factor 4 (2) in the top
(middle and bottom) panel. Shaded area in upper plots: systematic
uncertainties. Shaded areas in bottom plots: sum in quadrature of statistical
and systematic uncertainties. Horizontal error bars: bin width.
## Summed $p_{\mathrm{T}}$ \- track $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$
Figure 10: Summed $p_{\mathrm{T}}$ in Toward (top), Transverse (middle) and
Away (bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV
(right). Right and left vertical scales differ by a factor 4 (2) in the top
(middle and bottom) panel. Shaded area in upper plots: systematic
uncertainties. Shaded areas in bottom plots: sum in quadrature of statistical
and systematic uncertainties. Horizontal error bars: bin width.
## Summed $p_{\mathrm{T}}$ \- track $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Figure 11: Summed $p_{\mathrm{T}}$ in Toward (top), Transverse (middle) and
Away (bottom) regions at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$ TeV
(right). Right and left vertical scales differ by a factor 4 (3) in the top
(middle and bottom) panel. Shaded area in upper plots: systematic
uncertainties. Shaded areas in bottom plots: sum in quadrature of statistical
and systematic uncertainties. Horizontal error bars: bin width.
## Azimuthal correlations - track $p_{\mathrm{T}}>0.15\,\mbox{${\rm GeV}/c$}$
Figure 12: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 0.5 $<p_{T,LT}<$ 1.5 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 13: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 2.0 $<p_{T,LT}<$ 4.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 14: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 4.0 $<p_{T,LT}<$ 6.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 15: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 6.0 $<p_{T,LT}<$ 10.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
## Azimuthal correlations - track $p_{\mathrm{T}}>0.5\,\mbox{${\rm GeV}/c$}$
Figure 16: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 0.5 $<p_{T,LT}<$ 1.5 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 17: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 2.0 $<p_{T,LT}<$ 4.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 18: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 4.0 $<p_{T,LT}<$ 6.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 19: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 6.0 $<p_{T,LT}<$ 10.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
## Azimuthal correlations - track $p_{\mathrm{T}}>1.0\,\mbox{${\rm GeV}/c$}$
Figure 20: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 2.0 $<p_{T,LT}<$ 4.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 21: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 4.0 $<p_{T,LT}<$ 6.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
Figure 22: Azimuthal correlation at $\sqrt{s}=0.9$ TeV (left) and $\sqrt{s}=7$
TeV (right). Leading-track: 6.0 $<p_{T,LT}<$ 10.0 GeV/$c$. For visualization
purposes the $\Delta\phi$ axis is not centered around 0. Shaded area in upper
plots: systematic uncertainties. Shaded areas in bottom plots: sum in
quadrature of statistical and systematic uncertainties. Horizontal error bars:
bin width.
## 10 Acknowledgements
The ALICE collaboration would like to thank all its engineers and technicians
for their invaluable contributions to the construction of the experiment and
the CERN accelerator teams for the outstanding performance of the LHC complex.
The ALICE collaboration acknowledges the following funding agencies for their
support in building and running the ALICE detector:
Calouste Gulbenkian Foundation from Lisbon and Swiss Fonds Kidagan, Armenia;
Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq),
Financiadora de Estudos e Projetos (FINEP), Fundação de Amparo à Pesquisa do
Estado de São Paulo (FAPESP);
National Natural Science Foundation of China (NSFC), the Chinese Ministry of
Education (CMOE) and the Ministry of Science and Technology of China (MSTC);
Ministry of Education and Youth of the Czech Republic;
Danish Natural Science Research Council, the Carlsberg Foundation and the
Danish National Research Foundation;
The European Research Council under the European Community’s Seventh Framework
Programme;
Helsinki Institute of Physics and the Academy of Finland;
French CNRS-IN2P3, the ‘Region Pays de Loire’, ‘Region Alsace’, ‘Region
Auvergne’ and CEA, France;
German BMBF and the Helmholtz Association;
General Secretariat for Research and Technology, Ministry of Development,
Greece;
Hungarian OTKA and National Office for Research and Technology (NKTH);
Department of Atomic Energy and Department of Science and Technology of the
Government of India;
Istituto Nazionale di Fisica Nucleare (INFN) of Italy;
MEXT Grant-in-Aid for Specially Promoted Research, Japan;
Joint Institute for Nuclear Research, Dubna;
National Research Foundation of Korea (NRF);
CONACYT, DGAPA, México, ALFA-EC and the HELEN Program (High-Energy physics
Latin-American–European Network);
Stichting voor Fundamenteel Onderzoek der Materie (FOM) and the Nederlandse
Organisatie voor Wetenschappelijk Onderzoek (NWO), Netherlands;
Research Council of Norway (NFR);
Polish Ministry of Science and Higher Education;
National Authority for Scientific Research - NASR (Autoritatea Naţională
pentru Cercetare Ştiinţifică - ANCS);
Federal Agency of Science of the Ministry of Education and Science of Russian
Federation, International Science and Technology Center, Russian Academy of
Sciences, Russian Federal Agency of Atomic Energy, Russian Federal Agency for
Science and Innovations and CERN-INTAS;
Ministry of Education of Slovakia;
Department of Science and Technology, South Africa;
CIEMAT, EELA, Ministerio de Educación y Ciencia of Spain, Xunta de Galicia
(Consellería de Educación), CEADEN, Cubaenergía, Cuba, and IAEA (International
Atomic Energy Agency);
Swedish Reseach Council (VR) and Knut $\&$ Alice Wallenberg Foundation (KAW);
Ukraine Ministry of Education and Science;
United Kingdom Science and Technology Facilities Council (STFC);
The United States Department of Energy, the United States National Science
Foundation, the State of Texas, and the State of Ohio.
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## Appendix A The ALICE Collaboration
B. Abelevorg1234&A. Abrahantes Quintanaorg1197&D. Adamováorg1283&A.M.
Adareorg1260&M.M. Aggarwalorg1157&G. Aglieri Rinellaorg1192&A.G.
Agocsorg1143&A. Agostinelliorg1132&S. Aguilar Salazarorg1247&Z.
Ahammedorg1225&N. Ahmadorg1106&A. Ahmad Masoodiorg1106&S.U.
Ahnorg1160,org1215&A. Akindinovorg1250&D. Aleksandrovorg1252&B.
Alessandroorg1313&R. Alfaro Molinaorg1247&A. Aliciorg1133,org1192,org1335&A.
Alkinorg1220&E. Almaráz Aviñaorg1247&T. Altorg1184&V. Altiniorg1114,org1192&S.
Altinpinarorg1121&I. Altsybeevorg1306&C. Andreiorg1140&A. Andronicorg1176&V.
Anguelovorg1200&C. Ansonorg1162&T. Antičićorg1334&F. Antinoriorg1271&P.
Antonioliorg1133&L. Aphecetcheorg1258&H. Appelshäuserorg1185&N.
Arbororg1194&S. Arcelliorg1132&A. Arendorg1185&N. Armestoorg1294&R.
Arnaldiorg1313&T. Aronssonorg1260&I.C. Arseneorg1176&M. Arslandokorg1185&A.
Asryanorg1306&A. Augustinusorg1192&R. Averbeckorg1176&T.C. Awesorg1264&J.
Äystöorg1212&M.D. Azmiorg1106&M. Bachorg1184&A. Badalàorg1155&Y.W.
Baekorg1160,org1215&R. Bailhacheorg1185&R. Balaorg1313&R. Baldini
Ferroliorg1335&A. Baldisseriorg1288&A. Balditorg1160&F. Baltasar Dos Santos
Pedrosaorg1192&J. Bánorg1230&R.C. Baralorg1127&R. Barberaorg1154&F.
Barileorg1114&G.G. Barnaföldiorg1143&L.S. Barnbyorg1130&V. Barretorg1160&J.
Bartkeorg1168&M. Basileorg1132&N. Bastidorg1160&B. Bathenorg1256&G.
Batigneorg1258&B. Batyunyaorg1182&C. Baumannorg1185&I.G. Beardenorg1165&H.
Beckorg1185&I. Belikovorg1308&F. Belliniorg1132&R. Bellwiedorg1205&E. Belmont-
Morenoorg1247&S. Beoleorg1312&I. Berceanuorg1140&A. Bercuciorg1140&Y.
Berdnikovorg1189&D. Berenyiorg1143&C. Bergmannorg1256&D. Berzanoorg1313&L.
Betevorg1192&A. Bhasinorg1209&A.K. Bhatiorg1157&N. Bianchiorg1187&L.
Bianchiorg1312&C. Bianchinorg1270&J. Bielčíkorg1274&J. Bielčíkováorg1283&A.
Bilandzicorg1109&F. Blancoorg1242&F. Blancoorg1205&D. Blauorg1252&C.
Blumeorg1185&M. Boccioliorg1192&N. Bockorg1162&A. Bogdanovorg1251&H.
Bøggildorg1165&M. Bogolyubskyorg1277&L. Boldizsárorg1143&M. Bombaraorg1229&J.
Bookorg1185&H. Borelorg1288&A. Borissovorg1179&C. Bortolinorg1270,Dipartimento
di Fisica dell’Universita, Udine, Italy&S. Boseorg1224&F.
Bossúorg1192,org1312&M. Botjeorg1109&S. Böttgerorg27399&B. Boyerorg1266&P.
Braun-Munzingerorg1176&M. Bregantorg1258&T. Breitnerorg27399&M. Brozorg1136&R.
Brunorg1192&E. Brunaorg1260,org1312,org1313&G.E. Brunoorg1114&D.
Budnikovorg1298&H. Bueschingorg1185&S. Bufalinoorg1312,org1313&K.
Bugaievorg1220&O. Buschorg1200&Z. Butheleziorg1152&D. Caffarriorg1270&X.
Caiorg1329&H. Cainesorg1260&E. Calvo Villarorg1338&P. Cameriniorg1315&V. Canoa
Romanorg1244,org1279&G. Cara Romeoorg1133&F. Carenaorg1192&W. Carenaorg1192&N.
Carlin Filhoorg1296&F. Carminatiorg1192&C.A. Carrillo Montoyaorg1192&A.
Casanova Díazorg1187&M. Caselleorg1192&J. Castillo Castellanosorg1288&J.F.
Castillo Hernandezorg1176&E.A.R. Casulaorg1145&V. Catanescuorg1140&C.
Cavicchioliorg1192&J. Cepilaorg1274&P. Cerelloorg1313&B.
Changorg1212,org1301&S. Chapelandorg1192&J.L. Charvetorg1288&S.
Chattopadhyayorg1224&S. Chattopadhyayorg1225&M. Cherneyorg1170&C.
Cheshkovorg1192,org1239&B. Cheynisorg1239&E. Chiavassaorg1313&V. Chibante
Barrosoorg1192&D.D. Chinellatoorg1149&P. Chochulaorg1192&M.
Chojnackiorg1320&P. Christakoglouorg1109,org1320&C.H. Christensenorg1165&P.
Christiansenorg1237&T. Chujoorg1318&S.U. Chungorg1281&C. Cicaloorg1146&L.
Cifarelliorg1132,org1192&F. Cindoloorg1133&J. Cleymansorg1152&F.
Coccettiorg1335&J.-P. Coffinorg1308&F. Colamariaorg1114&D. Colellaorg1114&G.
Conesa Balbastreorg1194&Z. Conesa del Valleorg1192,org1308&P.
Constantinorg1200&G. Continorg1315&J.G. Contrerasorg1244&T.M.
Cormierorg1179&Y. Corrales Moralesorg1312&P. Corteseorg1103&I. Cortés
Maldonadoorg1279&M.R. Cosentinoorg1125,org1149&F. Costaorg1192&M.E.
Cotalloorg1242&E. Crescioorg1244&P. Crochetorg1160&E. Cruz Alanizorg1247&E.
Cuautleorg1246&L. Cunqueiroorg1187&A. Daineseorg1270,org1271&H.H.
Dalsgaardorg1165&A. Danuorg1139&I. Dasorg1224,org1266&K. Dasorg1224&D.
Dasorg1224&A. Dashorg1127,org1149&S. Dashorg1254,org1313&S. Deorg1225&A. De
Azevedo Moregulaorg1187&G.O.V. de Barrosorg1296&A. De Caroorg1290,org1335&G.
de Cataldoorg1115&J. de Cuvelandorg1184&A. De Falcoorg1145&D. De
Gruttolaorg1290&H. Delagrangeorg1258&E. Del Castillo Sanchezorg1192&A.
Delofforg1322&V. Demanovorg1298&N. De Marcoorg1313&E. Dénesorg1143&S. De
Pasqualeorg1290&A. Deppmanorg1296&G. D Erasmoorg1114&R. de Rooijorg1320&D. Di
Bariorg1114&T. Dietelorg1256&C. Di Giglioorg1114&S. Di Libertoorg1286&A. Di
Mauroorg1192&P. Di Nezzaorg1187&R. Diviàorg1192&Ø. Djuvslandorg1121&A.
Dobrinorg1179,org1237&T. Dobrowolskiorg1322&I. Domínguezorg1246&B.
Dönigusorg1176&O. Dordicorg1268&O. Drigaorg1258&A.K. Dubeyorg1225&L.
Ducrouxorg1239&P. Dupieuxorg1160&M.R. Dutta Majumdarorg1225&A.K. Dutta
Majumdarorg1224&D. Eliaorg1115&D. Emschermannorg1256&H. Engelorg27399&H.A.
Erdalorg1122&B. Espagnonorg1266&M. Estienneorg1258&S. Esumiorg1318&D.
Evansorg1130&G. Eyyubovaorg1268&D. Fabrisorg1270,org1271&J. Faivreorg1194&D.
Falchieriorg1132&A. Fantoniorg1187&M. Faselorg1176&R. Fearickorg1152&A.
Fedunovorg1182&D. Fehlkerorg1121&L. Feldkamporg1256&D. Feleaorg1139&G.
Feofilovorg1306&A. Fernández Téllezorg1279&A. Ferrettiorg1312&R.
Ferrettiorg1103&J. Figielorg1168&M.A.S. Figueredoorg1296&S.
Filchaginorg1298&R. Finiorg1115&D. Finogeevorg1249&F.M. Fiondaorg1114&E.M.
Fioreorg1114&M. Florisorg1192&S. Foertschorg1152&P. Fokaorg1176&S.
Fokinorg1252&E. Fragiacomoorg1316&M. Fragkiadakisorg1112&U.
Frankenfeldorg1176&U. Fuchsorg1192&C. Furgetorg1194&M. Fusco
Girardorg1290&J.J. Gaardhøjeorg1165&M. Gagliardiorg1312&A. Gagoorg1338&M.
Gallioorg1312&D.R. Gangadharanorg1162&P. Ganotiorg1264&C. Garabatosorg1176&E.
Garcia-Solisorg17347&I. Garishviliorg1234&J. Gerhardorg1184&M.
Germainorg1258&C. Geunaorg1288&A. Gheataorg1192&M. Gheataorg1192&B.
Ghidiniorg1114&P. Ghoshorg1225&P. Gianottiorg1187&M.R. Girardorg1323&P.
Giubellinoorg1192&E. Gladysz-Dziadusorg1168&P. Glässelorg1200&R.
Gomezorg1173&E.G. Ferreiroorg1294&L.H. González-Truebaorg1247&P. González-
Zamoraorg1242&S. Gorbunovorg1184&A. Goswamiorg1207&S. Gotovacorg1304&V.
Grabskiorg1247&L.K. Graczykowskiorg1323&R. Grajcarekorg1200&A.
Grelliorg1320&C. Grigorasorg1192&A. Grigorasorg1192&V. Grigorievorg1251&A.
Grigoryanorg1332&S. Grigoryanorg1182&B. Grinyovorg1220&N. Grionorg1316&P.
Grosorg1237&J.F. Grosse-Oetringhausorg1192&J.-Y. Grossiordorg1239&R.
Grossoorg1192&F. Guberorg1249&R. Guernaneorg1194&C. Guerra Gutierrezorg1338&B.
Guerzoniorg1132&M. Guilbaudorg1239&K. Gulbrandsenorg1165&T. Gunjiorg1310&A.
Guptaorg1209&R. Guptaorg1209&H. Gutbrodorg1176&Ø. Haalandorg1121&C.
Hadjidakisorg1266&M. Haiducorg1139&H. Hamagakiorg1310&G. Hamarorg1143&B.H.
Hanorg1300&L.D. Hanrattyorg1130&A. Hansenorg1165&Z. Harmanovaorg1229&J.W.
Harrisorg1260&M. Hartigorg1185&D. Haseganorg1139&D. Hatzifotiadouorg1133&A.
Hayrapetyanorg1192,org1332&S.T. Heckelorg1185&M. Heideorg1256&H.
Helstruporg1122&A. Herghelegiuorg1140&G. Herrera Corralorg1244&N.
Herrmannorg1200&K.F. Hetlandorg1122&B. Hicksorg1260&P.T. Hilleorg1260&B.
Hippolyteorg1308&T. Horaguchiorg1318&Y. Horiorg1310&P. Hristovorg1192&I.
Hřivnáčováorg1266&M. Huangorg1121&S. Huberorg1176&T.J. Humanicorg1162&D.S.
Hwangorg1300&R. Ichouorg1160&R. Ilkaevorg1298&I. Ilkivorg1322&M.
Inabaorg1318&E. Incaniorg1145&P.G. Innocentiorg1192&G.M. Innocentiorg1312&M.
Ippolitovorg1252&M. Irfanorg1106&C. Ivanorg1176&M. Ivanovorg1176&V.
Ivanovorg1189&A. Ivanovorg1306&O. Ivanytskyiorg1220&A. Jachołkowskiorg1192&P.
M. Jacobsorg1125&L. Jancurováorg1182&H.J. Jangorg20954&S. Jangalorg1308&R.
Janikorg1136&M.A. Janikorg1323&P.H.S.Y. Jayarathnaorg1205&S. Jenaorg1254&R.T.
Jimenez Bustamanteorg1246&L. Jirdenorg1192&P.G. Jonesorg1130&W. Jungorg1215&H.
Jungorg1215&A. Juskoorg1130&A.B. Kaidalovorg1250&V. Kakoyanorg1332&S.
Kalcherorg1184&P. Kaliňákorg1230&M. Kaliskyorg1256&T. Kalliokoskiorg1212&A.
Kalweitorg1177&K. Kanakiorg1121&J.H. Kangorg1301&V. Kaplinorg1251&A. Karasu
Uysalorg1192,org15649&O. Karavichevorg1249&T. Karavichevaorg1249&E.
Karpechevorg1249&A. Kazantsevorg1252&U. Kebschullorg1199,org27399&R.
Keidelorg1327&P. Khanorg1224&M.M. Khanorg1106&S.A. Khanorg1225&A.
Khanzadeevorg1189&Y. Kharlovorg1277&B. Kilengorg1122&J.H. Kimorg1300&D.J.
Kimorg1212&D.W. Kimorg1215&J.S. Kimorg1215&M. Kimorg1301&S.H. Kimorg1215&S.
Kimorg1300&B. Kimorg1301&T. Kimorg1301&S. Kirschorg1184,org1192&I.
Kiselorg1184&S. Kiselevorg1250&A. Kisielorg1192,org1323&J.L. Klayorg1292&J.
Kleinorg1200&C. Klein-Bösingorg1256&M. Kliemantorg1185&A. Klugeorg1192&M.L.
Knichelorg1176&K. Kochorg1200&M.K. Köhlerorg1176&A. Kolojvariorg1306&V.
Kondratievorg1306&N. Kondratyevaorg1251&A. Konevskikhorg1249&A.
Korneevorg1298&C. Kottachchi Kankanamge Donorg1179&R. Kourorg1130&M.
Kowalskiorg1168&S. Koxorg1194&G. Koyithatta Meethaleveeduorg1254&J.
Kralorg1212&I. Králikorg1230&F. Kramerorg1185&I. Krausorg1176&T.
Krawutschkeorg1200,org1227&M. Kretzorg1184&M. Krivdaorg1130,org1230&F.
Krizekorg1212&M. Krusorg1274&E. Kryshenorg1189&M. Krzewickiorg1109,org1176&Y.
Kucheriaevorg1252&C. Kuhnorg1308&P.G. Kuijerorg1109&P. Kurashviliorg1322&A.B.
Kurepinorg1249&A. Kurepinorg1249&A. Kuryakinorg1298&S. Kushpilorg1283&V.
Kushpilorg1283&H. Kvaernoorg1268&M.J. Kweonorg1200&Y. Kwonorg1301&P. Ladrón de
Guevaraorg1246&I. Lakomovorg1266,org1306&R. Langoyorg1121&C. Laraorg27399&A.
Lardeuxorg1258&P. La Roccaorg1154&C. Lazzeroniorg1130&R. Leaorg1315&Y. Le
Bornecorg1266&S.C. Leeorg1215&K.S. Leeorg1215&F. Lefèvreorg1258&J.
Lehnertorg1185&L. Leistamorg1192&M. Lenhardtorg1258&V. Lentiorg1115&H.
Leónorg1247&I. León Monzónorg1173&H. León Vargasorg1185&P. Lévaiorg1143&X.
Liorg1118&J. Lienorg1121&R. Lietavaorg1130&S. Lindalorg1268&V.
Lindenstruthorg1184&C. Lippmannorg1176,org1192&M.A. Lisaorg1162&L.
Liuorg1121&P.I. Loenneorg1121&V.R. Logginsorg1179&V. Loginovorg1251&S.
Lohnorg1192&D. Lohnerorg1200&C. Loizidesorg1125&K.K. Looorg1212&X.
Lopezorg1160&E. López Torresorg1197&G. Løvhøidenorg1268&X.-G. Luorg1200&P.
Luettigorg1185&M. Lunardonorg1270&J. Luoorg1329&G. Luparelloorg1320&L.
Luquinorg1258&C. Luzziorg1192&R. Maorg1260&K. Maorg1329&D.M. Madagodahettige-
Donorg1205&A. Maevskayaorg1249&M. Magerorg1177,org1192&D.P.
Mahapatraorg1127&A. Maireorg1308&M. Malaevorg1189&I. Maldonado
Cervantesorg1246&L. Malininaorg1182,M.V.Lomonosov Moscow State University,
D.V.Skobeltsyn Institute of Nuclear Physics, Moscow, Russia&D.
Mal’Kevichorg1250&P. Malzacherorg1176&A. Mamonovorg1298&L. Manceauorg1313&L.
Mangotraorg1209&V. Mankoorg1252&F. Mansoorg1160&V. Manzariorg1115&Y.
Maoorg1194,org1329&M. Marchisoneorg1160,org1312&J. Marešorg1275&G.V.
Margagliottiorg1315,org1316&A. Margottiorg1133&A. Marínorg1176&C.
Markertorg17361&I. Martashviliorg1222&P. Martinengoorg1192&M.I.
Martínezorg1279&A. Martínez Davalosorg1247&G. Martínez Garcíaorg1258&Y.
Martynovorg1220&A. Masorg1258&S. Masciocchiorg1176&M. Maseraorg1312&A.
Masoniorg1146&L. Massacrierorg1239&M. Mastromarcoorg1115&A.
Mastroserioorg1114,org1192&Z.L. Matthewsorg1130&A. Matyjaorg1258&D.
Mayaniorg1246&C. Mayerorg1168&J. Mazerorg1222&M.A. Mazzoniorg1286&F.
Meddiorg1285&A. Menchaca-Rochaorg1247&J. Mercado Pérezorg1200&M.
Meresorg1136&Y. Miakeorg1318&A. Michalonorg1308&J. Midoriorg1203&L.
Milanoorg1312&J. Milosevicorg1268,Institute of Nuclear Sciences, Belgrade,
Serbia&A. Mischkeorg1320&A.N. Mishraorg1207&D. Miśkowiecorg1176,org1192&C.
Mituorg1139&J. Mlynarzorg1179&A.K. Mohantyorg1192&B. Mohantyorg1225&L.
Molnarorg1192&L. Montaño Zetinaorg1244&M. Montenoorg1313&E. Montesorg1242&T.
Moonorg1301&M. Morandoorg1270&D.A. Moreira De Godoyorg1296&S.
Morettoorg1270&A. Morschorg1192&V. Mucciforaorg1187&E. Mudnicorg1304&S.
Muhuriorg1225&H. Müllerorg1192&M.G. Munhozorg1296&L. Musaorg1192&A.
Mussoorg1313&B.K. Nandiorg1254&R. Naniaorg1133&E. Nappiorg1115&C.
Nattrassorg1222&N.P. Naumovorg1298&S. Navinorg1130&T.K. Nayakorg1225&S.
Nazarenkoorg1298&G. Nazarovorg1298&A. Nedosekinorg1250&M. Nicassioorg1114&B.S.
Nielsenorg1165&T. Niidaorg1318&S. Nikolaevorg1252&V. Nikolicorg1334&V.
Nikulinorg1189&S. Nikulinorg1252&B.S. Nilsenorg1170&M.S. Nilssonorg1268&F.
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Zhouorg1320&D. Zhouorg1329&X. Zhuorg1329&A. Zichichiorg1132,org1335&A.
Zimmermannorg1200&G. Zinovjevorg1220&Y. Zoccaratoorg1239&M. Zynovyevorg1220
## Affiliation notes
0Deceased
Dipartimento di Fisica dell’Universita, Udine, ItalyAlso at: Dipartimento di
Fisica dell’Universita, Udine, Italy
M.V.Lomonosov Moscow State University, D.V.Skobeltsyn Institute of Nuclear
Physics, Moscow, RussiaAlso at: M.V.Lomonosov Moscow State University,
D.V.Skobeltsyn Institute of Nuclear Physics, Moscow, Russia
Institute of Nuclear Sciences, Belgrade, SerbiaAlso at: ”Vinča” Institute of
Nuclear Sciences, Belgrade, Serbia
## Collaboration Institutes
org1279Benemérita Universidad Autónoma de Puebla, Puebla, Mexico
org1220Bogolyubov Institute for Theoretical Physics, Kiev, Ukraine
org1262Budker Institute for Nuclear Physics, Novosibirsk, Russia
org1292California Polytechnic State University, San Luis Obispo, California,
United States
org14939Centre de Calcul de l’IN2P3, Villeurbanne, France
org1197Centro de Aplicaciones Tecnológicas y Desarrollo Nuclear (CEADEN),
Havana, Cuba
org1242Centro de Investigaciones Energéticas Medioambientales y Tecnológicas
(CIEMAT), Madrid, Spain
org1244Centro de Investigación y de Estudios Avanzados (CINVESTAV), Mexico
City and Mérida, Mexico
org1335Centro Fermi – Centro Studi e Ricerche e Museo Storico della Fisica
“Enrico Fermi”, Rome, Italy
org17347Chicago State University, Chicago, United States
org1118China Institute of Atomic Energy, Beijing, China
org1288Commissariat à l’Energie Atomique, IRFU, Saclay, France
org1294Departamento de Física de Partículas and IGFAE, Universidad de Santiago
de Compostela, Santiago de Compostela, Spain
org1106Department of Physics Aligarh Muslim University, Aligarh, India
org1121Department of Physics and Technology, University of Bergen, Bergen,
Norway
org1162Department of Physics, Ohio State University, Columbus, Ohio, United
States
org1300Department of Physics, Sejong University, Seoul, South Korea
org1268Department of Physics, University of Oslo, Oslo, Norway
org1132Dipartimento di Fisica dell’Università and Sezione INFN, Bologna, Italy
org1315Dipartimento di Fisica dell’Università and Sezione INFN, Trieste, Italy
org1145Dipartimento di Fisica dell’Università and Sezione INFN, Cagliari,
Italy
org1270Dipartimento di Fisica dell’Università and Sezione INFN, Padova, Italy
org1285Dipartimento di Fisica dell’Università ‘La Sapienza’ and Sezione INFN,
Rome, Italy
org1154Dipartimento di Fisica e Astronomia dell’Università and Sezione INFN,
Catania, Italy
org1290Dipartimento di Fisica ‘E.R. Caianiello’ dell’Università and Gruppo
Collegato INFN, Salerno, Italy
org1312Dipartimento di Fisica Sperimentale dell’Università and Sezione INFN,
Turin, Italy
org1103Dipartimento di Scienze e Tecnologie Avanzate dell’Università del
Piemonte Orientale and Gruppo Collegato INFN, Alessandria, Italy
org1114Dipartimento Interateneo di Fisica ‘M. Merlin’ and Sezione INFN, Bari,
Italy
org1237Division of Experimental High Energy Physics, University of Lund, Lund,
Sweden
org1192European Organization for Nuclear Research (CERN), Geneva, Switzerland
org1227Fachhochschule Köln, Köln, Germany
org1122Faculty of Engineering, Bergen University College, Bergen, Norway
org1136Faculty of Mathematics, Physics and Informatics, Comenius University,
Bratislava, Slovakia
org1274Faculty of Nuclear Sciences and Physical Engineering, Czech Technical
University in Prague, Prague, Czech Republic
org1229Faculty of Science, P.J. Šafárik University, Košice, Slovakia
org1184Frankfurt Institute for Advanced Studies, Johann Wolfgang Goethe-
Universität Frankfurt, Frankfurt, Germany
org1215Gangneung-Wonju National University, Gangneung, South Korea
org1212Helsinki Institute of Physics (HIP) and University of Jyväskylä,
Jyväskylä, Finland
org1203Hiroshima University, Hiroshima, Japan
org1329Hua-Zhong Normal University, Wuhan, China
org1254Indian Institute of Technology, Mumbai, India
org1266Institut de Physique Nucléaire d’Orsay (IPNO), Université Paris-Sud,
CNRS-IN2P3, Orsay, France
org1277Institute for High Energy Physics, Protvino, Russia
org1249Institute for Nuclear Research, Academy of Sciences, Moscow, Russia
org1320Nikhef, National Institute for Subatomic Physics and Institute for
Subatomic Physics of Utrecht University, Utrecht, Netherlands
org1250Institute for Theoretical and Experimental Physics, Moscow, Russia
org1230Institute of Experimental Physics, Slovak Academy of Sciences, Košice,
Slovakia
org1127Institute of Physics, Bhubaneswar, India
org1275Institute of Physics, Academy of Sciences of the Czech Republic,
Prague, Czech Republic
org1139Institute of Space Sciences (ISS), Bucharest, Romania
org27399Institut für Informatik, Johann Wolfgang Goethe-Universität Frankfurt,
Frankfurt, Germany
org1185Institut für Kernphysik, Johann Wolfgang Goethe-Universität Frankfurt,
Frankfurt, Germany
org1177Institut für Kernphysik, Technische Universität Darmstadt, Darmstadt,
Germany
org1256Institut für Kernphysik, Westfälische Wilhelms-Universität Münster,
Münster, Germany
org1246Instituto de Ciencias Nucleares, Universidad Nacional Autónoma de
México, Mexico City, Mexico
org1247Instituto de Física, Universidad Nacional Autónoma de México, Mexico
City, Mexico
org23333Institut of Theoretical Physics, University of Wroclaw
org1308Institut Pluridisciplinaire Hubert Curien (IPHC), Université de
Strasbourg, CNRS-IN2P3, Strasbourg, France
org1182Joint Institute for Nuclear Research (JINR), Dubna, Russia
org1143KFKI Research Institute for Particle and Nuclear Physics, Hungarian
Academy of Sciences, Budapest, Hungary
org18995Kharkiv Institute of Physics and Technology (KIPT), National Academy
of Sciences of Ukraine (NASU), Kharkov, Ukraine
org1199Kirchhoff-Institut für Physik, Ruprecht-Karls-Universität Heidelberg,
Heidelberg, Germany
org20954Korea Institute of Science and Technology Information
org1160Laboratoire de Physique Corpusculaire (LPC), Clermont Université,
Université Blaise Pascal, CNRS–IN2P3, Clermont-Ferrand, France
org1194Laboratoire de Physique Subatomique et de Cosmologie (LPSC), Université
Joseph Fourier, CNRS-IN2P3, Institut Polytechnique de Grenoble, Grenoble,
France
org1187Laboratori Nazionali di Frascati, INFN, Frascati, Italy
org1232Laboratori Nazionali di Legnaro, INFN, Legnaro, Italy
org1125Lawrence Berkeley National Laboratory, Berkeley, California, United
States
org1234Lawrence Livermore National Laboratory, Livermore, California, United
States
org1251Moscow Engineering Physics Institute, Moscow, Russia
org1140National Institute for Physics and Nuclear Engineering, Bucharest,
Romania
org1165Niels Bohr Institute, University of Copenhagen, Copenhagen, Denmark
org1109Nikhef, National Institute for Subatomic Physics, Amsterdam,
Netherlands
org1283Nuclear Physics Institute, Academy of Sciences of the Czech Republic,
Řež u Prahy, Czech Republic
org1264Oak Ridge National Laboratory, Oak Ridge, Tennessee, United States
org1189Petersburg Nuclear Physics Institute, Gatchina, Russia
org1170Physics Department, Creighton University, Omaha, Nebraska, United
States
org1157Physics Department, Panjab University, Chandigarh, India
org1112Physics Department, University of Athens, Athens, Greece
org1152Physics Department, University of Cape Town, iThemba LABS, Cape Town,
South Africa
org1209Physics Department, University of Jammu, Jammu, India
org1207Physics Department, University of Rajasthan, Jaipur, India
org1200Physikalisches Institut, Ruprecht-Karls-Universität Heidelberg,
Heidelberg, Germany
org1325Purdue University, West Lafayette, Indiana, United States
org1281Pusan National University, Pusan, South Korea
org1176Research Division and ExtreMe Matter Institute EMMI, GSI
Helmholtzzentrum für Schwerionenforschung, Darmstadt, Germany
org1334Rudjer Bošković Institute, Zagreb, Croatia
org1298Russian Federal Nuclear Center (VNIIEF), Sarov, Russia
org1252Russian Research Centre Kurchatov Institute, Moscow, Russia
org1224Saha Institute of Nuclear Physics, Kolkata, India
org1130School of Physics and Astronomy, University of Birmingham, Birmingham,
United Kingdom
org1338Sección Física, Departamento de Ciencias, Pontificia Universidad
Católica del Perú, Lima, Peru
org1146Sezione INFN, Cagliari, Italy
org1115Sezione INFN, Bari, Italy
org1313Sezione INFN, Turin, Italy
org1133Sezione INFN, Bologna, Italy
org1155Sezione INFN, Catania, Italy
org1316Sezione INFN, Trieste, Italy
org1286Sezione INFN, Rome, Italy
org1271Sezione INFN, Padova, Italy
org1322Soltan Institute for Nuclear Studies, Warsaw, Poland
org1258SUBATECH, Ecole des Mines de Nantes, Université de Nantes, CNRS-IN2P3,
Nantes, France
org1304Technical University of Split FESB, Split, Croatia
org1168The Henryk Niewodniczanski Institute of Nuclear Physics, Polish Academy
of Sciences, Cracow, Poland
org17361The University of Texas at Austin, Physics Department, Austin, TX,
United States
org1173Universidad Autónoma de Sinaloa, Culiacán, Mexico
org1296Universidade de São Paulo (USP), São Paulo, Brazil
org1149Universidade Estadual de Campinas (UNICAMP), Campinas, Brazil
org1239Université de Lyon, Université Lyon 1, CNRS/IN2P3, IPN-Lyon,
Villeurbanne, France
org1205University of Houston, Houston, Texas, United States
org20371University of Technology and Austrian Academy of Sciences, Vienna,
Austria
org1222University of Tennessee, Knoxville, Tennessee, United States
org1310University of Tokyo, Tokyo, Japan
org1318University of Tsukuba, Tsukuba, Japan
org21360Eberhard Karls Universität Tübingen, Tübingen, Germany
org1225Variable Energy Cyclotron Centre, Kolkata, India
org1306V. Fock Institute for Physics, St. Petersburg State University, St.
Petersburg, Russia
org1323Warsaw University of Technology, Warsaw, Poland
org1179Wayne State University, Detroit, Michigan, United States
org1260Yale University, New Haven, Connecticut, United States
org1332Yerevan Physics Institute, Yerevan, Armenia
org15649Yildiz Technical University, Istanbul, Turkey
org1301Yonsei University, Seoul, South Korea
org1327Zentrum für Technologietransfer und Telekommunikation (ZTT),
Fachhochschule Worms, Worms, Germany
|
arxiv-papers
| 2011-12-09T12:18:32 |
2024-09-04T02:49:25.126448
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/4.0/",
"authors": "ALICE collaboration",
"submitter": "Alice Publications",
"url": "https://arxiv.org/abs/1112.2082"
}
|
1112.2160
|
On the role of enrichment and statical admissibility of recovered fields in
a-posteriori error estimation for enriched finite element methods
Octavio A. González-Estrada1, Juan José Ródenas2, Stéphane P.A. Bordas1, Marc
Duflot3, Pierre Kerfriden1, Eugenio Giner2
1Institute of Mechanics and Advanced Materials. Cardiff School of Engineering,
Cardiff University, The Parade, Cardiff CF24 3AA Wales, UK.
2Centro de Investigación de Tecnología de Vehículos(CITV),
Universitad Politècnica de València, E-46022-Valencia, Spain.
3CENAERO, Rue des Frères Wright 29, B-6041 Gosselies, Belgium
###### Abstract
Purpose – This paper aims at assessing the effect of (1) the statical
admissibility of the recovered solution; (2) the ability of the recovered
solution to represent the singular solution; on the accuracy, local and global
effectivity of recovery-based error estimators for enriched finite element
methods (e.g. the extended finite element method, XFEM).
Design/methodology/approach – We study the performance of two recovery
techniques. The first is a recently developed superconvergent patch recovery
procedure with equilibration and enrichment (SPR-CX). The second is known as
the extended moving least squares recovery (XMLS), which enriches the
recovered solutions but does not enforce equilibrium constraints. Both are
extended recovery techniques as the polynomial basis used in the recovery
process is enriched with singular terms for a better description of the
singular nature of the solution.
Findings – Numerical results comparing the convergence and the effectivity
index of both techniques with those obtained without the enrichment
enhancement clearly show the need for the use of extended recovery techniques
in Zienkiewicz-Zhu type error estimators for this class of problems. The
results also reveal significant improvements in the effectivities yielded by
statically admissible recovered solutions.
Originality/value – This work shows that both extended recovery procedures and
statical admissibility are key to an accurate assessment of the quality of
enriched finite element approximations.
Keywords extended finite element method; error estimation; linear elastic
fracture mechanics; statical admissibility; extended recovery
Paper Type Research Paper
## 1 Introduction
Engineering structures, in particular in aerospace engineering, are intended
to operate with flawless components, especially for safety critical parts.
However, there is always a possibility that cracking will occur during
operation, risking catastrophic failure and associated casualties.
The mission of Damage Tolerance Assessment (DTA) is to assess the influence of
these defects, cracks and damage on the ability of a structure to perform
safely and reliably during its service life. An important goal of DTA is to
estimate the fatigue life of a structure, i.e. the time during which it
remains safe given a pre-existing flaw.
Damage tolerance assessment relies on the ability to accurately predict crack
paths and growth rates in complex structures. Since the simulation of three-
dimensional crack growth is either not supported by commercial software, or
requires significant effort and time for the analysts and is generally not
coupled to robust error indicators, reliable, quality-controlled, industrial
damage tolerance assessment is still a major challenge in engineering
practice.
The extended finite element method (XFEM) (Moës et al., 1999) is now one of
many successful numerical methods to solve fracture mechanics problems. The
particular advantage of the XFEM, which relies on the partition of unity (PU)
property (Melenk and Babuška, 1996) of finite element shape functions, is its
ability to model cracks without the mesh conforming to their geometry. This
allows the crack to split the background mesh arbitrarily, thereby leading to
significantly increased freedom in simulating crack growth.
This feat is achieved by adding new degrees of freedom to (i) describe the
discontinuity of the displacement field across the crack faces, within a given
element, and (ii) reproduce the asymptotic fields around the crack tip. Thanks
to the advances made in the XFEM during the last years, the method is now
considered to be a robust and accurate means of analysing fracture problems,
has been implemented in commercial codes (ABAQUS, 2009, CENAERO, 2011) and is
industrially in use for damage tolerance assessment of complex three
dimensional structures (Bordas and Moran, 2006, Wyart et al., 2007, Duflot,
2007).
While XFEM allows to model cracks without (re)meshing the crack faces as they
evolve and yields exceptionally accurate stress intensity factors (SIF) for 2D
problems, Bordas and Moran (2006), Wyart et al. (2007), Duflot (2007) show
that for realistic 3D structures, a “very fine” mesh is required to accurately
capture the complex three-dimensional stress field and obtain satisfactory
stress intensity factors (SIFs), the main drivers of linear elastic crack
propagation.
Consequently, based on the mesh used for the stress analysis, a new mesh
offering sufficient refinement throughout the whole potential path of the
crack must be constructed _a priori_ , i.e. before the crack path is known.
Practically, this is done by running preliminary analyses on coarse meshes to
obtain an approximative crack path and heuristically refining the mesh around
this path to increase accuracy.
Typically, this refinement does not rely on sound error measures, thus the
heuristically chosen mesh is in general inadequately suited and can cause
large inaccuracies in the crack growth path, especially around holes, which
can lead to non-conservative estimates of the safe life of the structure.
Thus, it is clear that although XFEM simplifies the treatment of cracks
compared to the standard FEM by lifting the burden of a geometry conforming
mesh, it still requires iterations and associated user intervention and it
employs heuristics which are detrimental to the robustness and accuracy of the
simulation process.
It would be desirable to minimise the changes to the mesh topology, thus user
intervention, while ensuring the stress fields and SIFs are accurately
computed at each crack growth step. This paper is one more step in attempting
to control the discretization error committed by enriched FE approximations,
decrease human intervention in damage tolerance assessment of complex
industrial structures and enhance confidence in the results by providing
enriched FEMs with sound error estimators which will guarantee a predetermined
accuracy level and suppress recourse to manual iterations and heuristics.
The error assessment procedures used in finite element analysis are well known
and can be usually classified into different families (Ainsworth and Oden,
2000): residual based error estimators, recovery based error estimators, dual
techniques, etc. Residual based error estimators were substantially improved
with the introduction of the _residual equilibration_ by Ainsworth and Oden
(2000). These error estimators have a strong mathematical basis and have been
frequently used to obtain lower and upper bounds of the error (Ainsworth and
Oden, 2000, Díez et al., 2004). Recovery based error estimates were first
introduced by Zienkiewicz and Zhu (1987) and are often preferred by
practitioners because they are robust and simple to use and provide an
enhanced solution. Further improvements were made with the introduction of new
recovery processes such as the superconvergent patch recovery (SPR) technique
proposed by Zienkiewicz and Zhu (1992a, b). Dual techniques based on the
evaluation of two different fields, one compatible in displacements and
another equilibrated in stresses, have also been used by Pereira et al. (1999)
to obtain bounds of the error. Herein we are going to focus on recovery based
techniques which follow the ideas of the Zienkiewicz-Zhu (ZZ) error estimator
proposed by Zienkiewicz and Zhu (1987).
The literature on error estimation techniques for mesh based partition of
unity methods is still scarce. One of the first steps in that direction was
made in the context of the Generalized Finite Element Method (GFEM) by
Strouboulis et al. (2001). The authors proposed a recovery-based error
estimator which provides good results for h-adapted meshes. In a later work
two new a posteriori error estimators for GFEM were presented (Strouboulis et
al., 2006). The first one was based on patch residual indicators and provided
an accurate theoretical upper bound estimate, but its computed version
severely underestimated the exact error. The second one was an error estimator
based on a recovered displacement field and its performance was closely
related to the quality of the GFEM solution. This recovery technique
constructs a recovered solution on patches using a basis enriched with
handbook functions.
In order to obtain a recovered stress field that improves the accuracy of the
stresses obtained by XFEM, Xiao and Karihaloo (2006) proposed a moving least
squares fitting adapted to the XFEM framework which considers the use of
statically admissible basis functions. Nevertheless, the recovered stress
field was not used to obtain an error indicator.
Pannachet et al. (2009) worked on error estimation for mesh adaptivity in XFEM
for cohesive crack problems. Two error estimates were used, one based on the
error in the energy norm and another that considers the error in a local
quantity of interest. The error estimation was based on solving a series of
local problems with prescribed homogeneous boundary conditions.
Panetier et al. (2010) presented an extension to enriched approximations of
the constitutive relation error (CRE) technique already available to evaluate
error bounds in FEM (Ladevèze and Pelle, 2005). This procedure has been used
to obtain local error bounds on quantities of interest for XFEM problems.
Bordas and Duflot (2007) and Bordas et al. (2008) proposed a recovery based
error estimator for 2D and 3D XFEM approximations for cracks known as the
extended moving least squares (XMLS). This method intrinsically enriched an
MLS formulation to include information about the singular fields near the
crack tip. Additionally, it used a diffraction method to introduce the
discontinuity in the recovered field. This error estimator provided accurate
results with effectivity indices close to unity (optimal value) for 2D and 3D
fracture mechanics problems. Later, Duflot and Bordas (2008) proposed a global
derivative recovery formulation extended to XFEM problems. The recovered
solution was sought in a space spanned by the near tip strain fields obtained
from differentiating the Westergaard asymptotic expansion. Although the
results provided by this technique were not as accurate as those in Bordas and
Duflot (2007), Bordas et al. (2008), they were deemed by the authors to
require less computational power.
Ródenas et al. (2008) presented a modification of the superconvergent patch
recovery (SPR) technique tailored to the XFEM framework called
$\textrm{SPR}_{\textrm{XFEM}}$. This technique was based on three key
ingredients: (1) the use of a _singular_ +_smooth_ stress field decomposition
procedure around the crack tip similar to that described by Ródenas et al.
(2006) for FEM; (2) direct calculation of recovered stresses at integration
points using the partition of unity, and (3) use of different stress
interpolation polynomials at each side of the crack when a crack intersects a
patch. In order to obtain an equilibrated smooth recovered stress field, a
simplified version of the SPR-C technique presented by Ródenas et al. (2007)
was used. This simplified SPR-C imposed the fulfilment of the boundary
equilibrium equation at boundary nodes but did not impose the satisfaction of
the internal equilibrium equations.
The numerical results presented by Bordas and Duflot (2007), Bordas et al.
(2008) and Ródenas et al. (2008) showed promising accuracy for both the XMLS
and the $\textrm{SPR}_{\textrm{XFEM}}$ techniques. It is apparent in those
papers that $\textrm{SPR}_{\textrm{XFEM}}$ led to effectivity indices
remarkably close to unity. Yet, these papers do not shed any significant light
on the respective roles played by two important ingredients in those error
estimators, namely:
1. 1.
The statistical admissibility of the recovered solution;
2. 2.
The enrichment of the recovered solution.
Moreover, the differences in the test cases analysed and in the quality
measures considered in each of the papers makes it difficult to objectively
compare the merits of both methods.
The aim of this paper is to assess the role of statistical admissibility and
enrichment of the recovered solution in recovery based error estimation of
enriched finite element approximation for linear elastic fracture. To do so,
we perform a systematic study of the results obtained when considering the
different features of two error estimation techniques: XMLS and SPR-CX. SPR-CX
is an enhanced version of the $\textrm{SPR}_{\textrm{XFEM}}$ presented by
Ródenas et al. (2008) for 2D problems. Advantages and disadvantages of each
method are also provided.
The outline of the paper is as follows. In Section 2, the XFEM is briefly
presented. Section 3 deals with error estimation and quality assessment of the
solution. In Sections 4 and 5 we introduce the error estimators used in XFEM
approximations: the SPR-CX and the XMLS, respectively. Some numerical examples
analysing both techniques and the effect of the enrichment functions in the
recovery process are presented in Section 6. Finally some concluding remarks
are provided in Section 7.
## 2 Reference problem and XFEM solution
Let us consider a 2D linear elastic fracture mechanics (LEFM) problem on a
bounded domain $\Omega\subset\mathbb{R}^{2}$. The unknown displacement field
$\bm{\mathrm{u}}$ is the solution of the boundary value problem
$\displaystyle\nabla\cdot\boldsymbol{\upsigma}(\bm{\mathrm{u}})+\bm{\mathrm{b}}$
$\displaystyle=\mathbf{0}$ $\displaystyle\textrm{in }\Omega$ (1)
$\displaystyle\boldsymbol{\upsigma}(\bm{\mathrm{u}})\cdot\bm{\mathrm{n}}$
$\displaystyle=\bm{\mathrm{t}}$ $\displaystyle\textrm{on }\Gamma_{N}$ (2)
$\displaystyle\boldsymbol{\upsigma}(\bm{\mathrm{u}})\cdot\bm{\mathrm{n}}$
$\displaystyle=\mathbf{0}$ $\displaystyle\textrm{on }\Gamma_{C}$ (3)
$\displaystyle\bm{\mathrm{u}}$ $\displaystyle=\bar{\bm{\mathrm{u}}}$
$\displaystyle\textrm{on }\Gamma_{D}$ (4)
where $\Gamma_{N}$ and $\Gamma_{D}$ are the Neumann and Dirichlet boundaries
and $\Gamma_{C}$ represents the crack faces, with
$\partial\Omega=\Gamma_{N}\cup\Gamma_{D}\cup\Gamma_{C}$ and
$\Gamma_{N}\cap\Gamma_{D}\cap\Gamma_{C}=\varnothing$. $\bm{\mathrm{b}}$ are
the body forces per unit volume, $\bm{\mathrm{t}}$ the tractions applied on
$\Gamma_{N}$ (being $\bm{\mathrm{n}}$ the normal vector to the boundary) and
$\bar{\bm{\mathrm{u}}}$ the prescribed displacements on $\Gamma_{D}$. The weak
form of the problem reads: Find $\bm{\mathrm{u}}\in V$ such that:
$\forall\bm{\mathrm{v}}\in V\qquad
a(\bm{\mathrm{u}},\bm{\mathrm{v}})=l(\bm{\mathrm{v}}),$ (5)
where $V$ is the standard test space for elasticity problems such that
$V=\\{\bm{\mathrm{v}}\;|\;\bm{\mathrm{v}}\in
H^{1}(\Omega),\bm{\mathrm{v}}|_{\Gamma_{D}}(\bm{\mathrm{x}})=\mathbf{0},\bm{\mathrm{v}}\;\textrm{discontinuous
on }\Gamma_{C}\\}$, and
$\displaystyle a(\bm{\mathrm{u}},\bm{\mathrm{v}})$
$\displaystyle:=\int_{\Omega}\boldsymbol{\upsigma}(\bm{\mathrm{u}}):\bm{\mathrm{\epsilon}}(\bm{\mathrm{v}})d\Omega=\int_{\Omega}\boldsymbol{\upsigma}(\bm{\mathrm{u}}):\boldsymbol{\mathsf{S}}:\boldsymbol{\upsigma}(\bm{\mathrm{v}})d\Omega$
(6) $\displaystyle l(\bm{\mathrm{v}})$
$\displaystyle:=\int_{\Omega}\bm{\mathrm{b}}\cdot\bm{\mathrm{v}}d\Omega+\int_{\Gamma_{N}}\bm{\mathrm{t}}\cdot\bm{\mathrm{v}}d\Gamma,$
(7)
where $\boldsymbol{\mathsf{S}}$ is the compliance tensor,
$\boldsymbol{\upsigma}$ and $\bm{\mathrm{\epsilon}}$ represent the stress and
strain operators.
LEFM problems are denoted by the singularity present at the crack tip. The
following expressions represent the first term of the asymptotic expansion
which describes the displacements and stresses for combined loading modes I
and II in 2D. These expressions can be found in the literature (Szabó and
Babuška, 1991, Ródenas et al., 2008) and are reproduced here for completeness:
$\displaystyle u_{1}(r,\phi)$ $\displaystyle=\frac{K_{\rm
I}}{2\mu}\sqrt{\frac{r}{2\pi}}\cos\frac{\phi}{2}\left(\kappa-\cos\phi\right)+\frac{K_{\rm
II}}{2\mu}\sqrt{\frac{r}{2\pi}}\sin\frac{\phi}{2}\left(2+\kappa+\cos\phi\right)$
(8) $\displaystyle u_{2}(r,\phi)$ $\displaystyle=\frac{K_{\rm
I}}{2\mu}\sqrt{\frac{r}{2\pi}}\sin\frac{\phi}{2}\left(\kappa-\cos\phi\right)+\frac{K_{\rm
II}}{2\mu}\sqrt{\frac{r}{2\pi}}\cos\frac{\phi}{2}\left(2-\kappa-\cos\phi\right)$
$\displaystyle\sigma_{11}(r,\phi)$ $\displaystyle=\frac{K_{\rm I}}{\sqrt{2\pi
r}}\cos\frac{\phi}{2}\left(1-\sin\frac{\phi}{2}\sin\frac{3\phi}{2}\right)-\frac{K_{\rm
II}}{\sqrt{2\pi
r}}\sin\frac{\phi}{2}\left(2+\cos\frac{\phi}{2}\cos\frac{3\phi}{2}\right)$ (9)
$\displaystyle\sigma_{22}(r,\phi)$ $\displaystyle=\frac{K_{\rm I}}{\sqrt{2\pi
r}}\cos\frac{\phi}{2}\left(1+\sin\frac{\phi}{2}\sin\frac{3\phi}{2}\right)+\frac{K_{\rm
II}}{\sqrt{2\pi r}}\sin\frac{\phi}{2}\cos\frac{\phi}{2}\cos\frac{3\phi}{2}$
$\displaystyle\sigma_{12}(r,\phi)$ $\displaystyle=\frac{K_{\rm I}}{\sqrt{2\pi
r}}\sin\frac{\phi}{2}\cos\frac{\phi}{2}\cos\frac{3\phi}{2}+\frac{K_{\rm
II}}{\sqrt{2\pi
r}}\cos\frac{\phi}{2}\left(1-\sin\frac{\phi}{2}\sin\frac{3\phi}{2}\right)$
where $r$ and $\phi$ are the crack tip polar coordinates, $K_{\rm I}$ and
$K_{\rm II}$ are the stress intensity factors for modes I and II, $\mu$ is the
shear modulus, and $\kappa$ the Kolosov’s constant, defined in terms of the
parameters of material $E$ (Young’s modulus) and $\upsilon$ (Poisson’s ratio),
according to the expressions:
$\mu=\frac{E}{2\left(1+\upsilon\right)},\qquad\kappa=\left\\{\begin{array}[]{l
c r}{\displaystyle 3-4\upsilon\qquad\textrm{plane strain}}\\\ \vskip 6.0pt
plus 2.0pt minus
2.0pt\cr{\displaystyle\frac{3-\upsilon}{1+\upsilon}\,\qquad\textrm{plane
stress}}\end{array}\right.$
This type of problem is difficult to model using a standard FEM approximation
as the mesh needs to explicitly conform to the crack geometry. With the XFEM
the discontinuity of the displacement field along the crack faces is
introduced by adding degrees of freedom to the nodes of the elements
intersected by the crack. This tackles the problem of adjusting the mesh to
the geometry of the crack (Moës et al., 1999, Stolarska et al., 2001).
Additionally, to describe the solution around the crack tip, the numerical
model introduces a basis that spans the near tip asymptotic fields. The
following expression is generally used to interpolate the displacements at a
point of coordinates $\bm{\mathrm{x}}$ accounting for the presence of a crack
tip in a 2D XFEM approximation:
$\bm{\mathrm{u}}_{h}(\bm{\mathrm{x}})=\sum_{i\in\mathcal{I}}N_{i}(\bm{\mathrm{x}})\textbf{a}_{i}+\sum_{j\in\mathcal{J}}N_{j}(\bm{\mathrm{x}})H(\bm{\mathrm{x}})\textbf{b}_{j}+\sum_{m\in\mathcal{M}}N_{m}(\bm{\mathrm{x}})\left(\sum_{\ell=1}^{4}F_{\ell}(\bm{\mathrm{x}})\textbf{c}_{m}^{\ell}\right)$
(10)
where $N_{i}$ are the shape functions associated with node $i$,
$\textbf{a}_{i}$ represent the conventional nodal degrees of freedom,
$\textbf{b}_{j}$ are the coefficients associated with the discontinuous
enrichment functions, and $\bm{\mathrm{c}}_{m}$ those associated with the
functions spanning the asymptotic field. In the above equation, $\mathcal{I}$
is the set of all the nodes in the mesh, $\mathcal{M}$ is the subset of nodes
enriched with crack tip functions, and $\mathcal{J}$ is the subset of nodes
enriched with the discontinuous enrichment (see Figure 1). In (10), the
Heaviside function $H$, with unitary modulus and a change of sign on the crack
face, describes the displacement discontinuity if the finite element is
intersected by the crack. The $F_{\ell}$ are the set of branch functions used
to represent the asymptotic expansion of the displacement field around the
crack tip seen in (8). For the 2D case, the following functions are used
(Belytschko and Black, 1999):
$\left\\{F_{\ell}\left(r,\phi\right)\right\\}\equiv\sqrt{r}\left\\{\sin\frac{\phi}{2},\cos\frac{\phi}{2},\sin\frac{\phi}{2}\sin\phi,\cos\frac{\phi}{2}\sin\phi\right\\}$
(11)
Figure 1: Classification of nodes in XFEM. Fixed enrichment area of radius
$r_{e}$
The main features of the XFEM implementation considered to evaluate the
numerical results is described in detail in Ródenas et al. (2008) and can be
summarized as follows:
* •
Use of bilinear quadrilaterals.
* •
Decomposition of elements intersected by the crack into integration subdomains
that do not contain the crack. Alternatives which do not required this
subdivision are proposed by Ventura (2006), Natarajan et al. (2010).
* •
Use of a quasi-polar integration with a $5\times 5$ quadrature rule in
triangular subdomains for elements containing the crack tip.
* •
No correction for blending elements. Methods to address blending errors are
proposed by Chessa et al. (2003), Gracie et al. (2008), Fries (2008), Tarancón
et al. (2009).
### 2.1 Evaluation of stress intensity factors
The stress intensity factors (SIFs) in LEFM represent the amplitude of the
singular stress fields and are key quantities of interest to simulate crack
growth in LEFM. Several post–processing methods, following local or global
(energy) approaches, are commonly used to extract SIFs (Banks-Sills, 1991) or
to calculate the energy release rate $G$. Energy or global methods are
considered to be the most accurate and efficient methods (Banks-Sills, 1991,
Li et al., 1985). Global methods based on the equivalent domain integral (EDI
methods) are specially well-suited for FEM and XFEM analyses as they are easy
to implement and can use information far from the singularity. In this paper,
the interaction integral as described in Shih and Asaro (1988), Yau et al.
(1980) has been used to extract the SIFs. This technique provides $K_{\rm I}$
and $K_{\rm II}$ for problems under mixed mode loading conditions using
auxiliary fields. Details on the implementation of the interaction integral
can be found for example in Moës et al. (1999), Ródenas et al. (2008), Giner
et al. (2005).
## 3 Error estimation in the energy norm
The approximate nature of the FEM and XFEM approximations implies a
discretization error which can be quantified using the error in the energy
norm for the solution
$\lVert\bm{\mathrm{e}}\rVert=\lVert\bm{\mathrm{u}}-\bm{\mathrm{u}}^{h}\rVert$.
To obtain an estimate of the discretization error
$\lVert\bm{\mathrm{e}}_{es}\rVert$, in the context of elasticity problems
solved using the FEM, the expression for the ZZ estimator is defined as
(Zienkiewicz and Zhu, 1987) (in matrix form):
$\lVert\bm{\mathrm{e}}\rVert^{2}\approx\lVert\bm{\mathrm{e}}_{es}\rVert^{2}=\int_{\Omega}\left(\bm{\mathrm{\sigma}}^{*}-\bm{\mathrm{\sigma}}^{h}\right)^{T}\bm{\mathrm{D}}^{-1}\left(\bm{\mathrm{\sigma}}^{*}-\bm{\mathrm{\sigma}}^{h}\right)d\Omega$
(12)
or alternatively for the strains
$\lVert\bm{\mathrm{e}}_{es}\rVert^{2}=\int_{\Omega}\left(\bm{\mathrm{\varepsilon}}^{*}-\bm{\mathrm{\varepsilon}}^{h}\right)^{T}\bm{\mathrm{D}}\left(\bm{\mathrm{\varepsilon}}^{*}-\bm{\mathrm{\varepsilon}}^{h}\right)d\Omega,$
where the domain $\Omega$ refers to the full domain of the problem or a local
subdomain (element), $\bm{\mathrm{\sigma}}^{h}$ represents the stress field
provided by the FEM, $\bm{\mathrm{\sigma}}^{*}$ is the recovered stress field,
which is a better approximation to the exact solution than
$\bm{\mathrm{\sigma}}^{h}$ and $\bm{\mathrm{D}}$ is the elasticity matrix of
the constitutive relation
$\bm{\mathrm{\sigma}}=\bm{\mathrm{D}}\bm{\mathrm{\varepsilon}}$.
The recovered stress field $\bm{\mathrm{\sigma}}^{*}$ is usually interpolated
in each element using the shape functions $\bm{\mathrm{N}}$ of the underlying
FE approximation and the values of the recovered stress field calculated at
the nodes $\bar{\bm{\mathrm{\sigma}}}^{*}$
$\bm{\mathrm{\sigma}}^{*}(\bm{\mathrm{x}})=\sum_{i=1}^{n_{e}}N_{i}(\bm{\mathrm{x}})\bar{\bm{\mathrm{\sigma}}}^{*}_{i}(\bm{\mathrm{x}}_{i}),$
(13)
where $n_{e}$ is the number of nodes in the element under consideration and
$\bar{\bm{\mathrm{\sigma}}}^{*}_{i}(\bm{\mathrm{x}}_{i})$ are the stresses
provided by the least squares technique at node $i$. The components of
$\bar{\bm{\mathrm{\sigma}}}^{*}_{i}$ are obtained using a polynomial
expansion, $\bar{\sigma}_{i,j}=\bm{\mathrm{p}}\bm{\mathrm{a}}$ (with
$j=xx,yy,xy$), defined over a set of contiguous elements connected to node $i$
called patch, where $\bm{\mathrm{p}}$ is the polynomial basis and
$\bm{\mathrm{a}}$ are the unknown coefficients.
The ZZ error estimator is asymptotically exact if the recovered solution used
in the error estimation converges at a higher rate than the finite element
solution (Zienkiewicz and Zhu, 1992b).
In this paper, we are interested in the role played by statical admissibility
and enrichment of the recovered solution in estimating the error committed by
XFEM. To do so, we study the performance of two recovery-based error
estimators, which exhibit different features, that have been recently
developed for XFEM:
* •
The SPR-CX derived from the error estimator developed by Ródenas et al. (2008)
and summarized in Section 4 (two other versions of the SPR-CX technique have
been also considered);
* •
The XMLS proposed by Bordas and Duflot (2007) and summarized in Section 5.
Both estimators provide a recovered stress field in order to evaluate the
estimated error in the energy norm by means of the expression shown in (12).
## 4 SPR-CX error estimator
The SPR-CX error estimator is an enhancement of the error estimator first
introduced by Ródenas et al. (2008), which incorporates the ideas proposed in
Ródenas et al. (2007) to guarantee the exact satisfaction of the equilibrium
locally on patches. In Ródenas et al. (2008) a set of key ideas are proposed
to modify the standard SPR by Zienkiewicz and Zhu (1992a), allowing its use
for singular problems.
The recovered stresses $\bm{\mathrm{\sigma}}^{*}$ are directly evaluated at an
integration point $\bm{\mathrm{x}}$ through the use of a partition of unity
procedure, properly weighting the stress interpolation polynomials obtained
from the different patches formed at the vertex nodes of the element
containing $\bm{\mathrm{x}}$:
$\bm{\mathrm{\sigma}}^{*}(\bm{\mathrm{x}})=\sum_{i=1}^{n_{v}}N_{i}(\bm{\mathrm{x}})\bm{\mathrm{\sigma}}^{*}_{i}(\bm{\mathrm{x}}),$
(14)
where $N_{i}$ are the shape functions associated to the vertex nodes $n_{v}$.
One major modification is the introduction of a splitting procedure to perform
the recovery. For singular problems the exact stress field
$\bm{\mathrm{\sigma}}$ is decomposed into two stress fields, a smooth field
$\bm{\mathrm{\sigma}}_{smo}$ and a singular field
$\bm{\mathrm{\sigma}}_{sing}$:
$\bm{\mathrm{\sigma}}=\bm{\mathrm{\sigma}}_{smo}+\bm{\mathrm{\sigma}}_{sing}.$
(15)
Then, the recovered field $\bm{\mathrm{\sigma}}^{*}$ required to compute the
error estimate given in (12) can be expressed as the contribution of two
recovered stress fields, one smooth $\bm{\mathrm{\sigma}}^{*}_{smo}$ and one
singular $\bm{\mathrm{\sigma}}^{*}_{sing}$:
$\bm{\mathrm{\sigma}}^{*}=\bm{\mathrm{\sigma}}^{*}_{smo}+\bm{\mathrm{\sigma}}^{*}_{sing}.$
(16)
For the recovery of the singular part, the expressions in (9) which describe
the asymptotic fields near the crack tip are used. To evaluate
$\bm{\mathrm{\sigma}}^{*}_{sing}$ from (9) we first obtain estimated values of
the stress intensity factors $K_{\rm I}$ and $K_{\rm II}$ using the
interaction integral as indicated in Section 2.1. The recovered part
$\bm{\mathrm{\sigma}}^{*}_{sing}$ is an equilibrated field as it satisfies the
internal equilibrium equations.
Once the field $\bm{\mathrm{\sigma}}^{*}_{sing}$ has been evaluated, an FE
approximation to the smooth part $\bm{\mathrm{\sigma}}^{h}_{smo}$ can be
obtained subtracting $\bm{\mathrm{\sigma}}^{*}_{sing}$ from the raw FE field:
$\bm{\mathrm{\sigma}}^{h}_{smo}=\bm{\mathrm{\sigma}}^{h}-\bm{\mathrm{\sigma}}^{*}_{sing}.$
(17)
Then, the field $\bm{\mathrm{\sigma}}^{*}_{smo}$ is evaluated applying an
SPR-C recovery procedure over the field $\bm{\mathrm{\sigma}}^{h}_{smo}$.
For patches intersected by the crack, the recovery technique uses different
stress interpolation polynomials on each side of the crack. This way it can
represent the discontinuity of the recovered stress field along the crack
faces, which is not the case for SPR that smoothes out the discontinuity (see,
e.g. Bordas and Duflot (2007)).
In order to obtain an equilibrated recovered stress field
$\bm{\mathrm{\sigma}}^{*}_{smo}$, the SPR-CX enforces the fulfilment of the
equilibrium equations locally on each patch. The constraint equations are
introduced via Lagrange multipliers into the linear system used to solve for
the coefficients of the polynomial expansion of the recovered stresses on each
patch. These include the satisfaction of the:
* •
Internal equilibrium equations.
* •
Boundary equilibrium equation: A point collocation approach is used to impose
the satisfaction of a second order approximation to the tractions along the
Neumann boundary.
* •
Compatibility equation: This additional constraint is also imposed to further
increase the accuracy of the recovered stress field.
To evaluate the recovered field, quadratic polynomials have been used in the
patches along the boundary and crack faces, and linear polynomials for the
remaining patches. As more information about the solution is available along
the boundary, polynomials one degree higher are useful to improve the quality
of the recovered stress field.
The enforcement of equilibrium equations provides an equilibrated recovered
stress field locally on patches. However, the process used to obtain a
continuous field $\bm{\mathrm{\sigma}}^{*}$ shown in (14) introduces a small
lack of equilibrium as explained in Ródenas et al. (2010a). The reader is
referred to Ródenas et al. (2010a), Díez et al. (2007) for details.
## 5 XMLS error estimator
In the XMLS the solution is recovered through the use of the _moving least
squares_ (MLS) technique, developed by mathematicians to build and fit
surfaces. The XMLS technique extends the work of Tabbara et al. (1994) for FEM
to enriched approximations. The general idea of the XMLS is to use the
displacement solution provided by XFEM to obtain a recovered strain field
(Bordas and Duflot, 2007, Bordas et al., 2008). The smoothed strains are
recovered from the derivative of the MLS-smoothed XFEM displacement field:
$\displaystyle\bm{\mathrm{u}}^{*}(\bm{\mathrm{x}})$
$\displaystyle=\sum_{i\in\mathcal{N}_{x}}\psi_{i}(\bm{\mathrm{x}})\bm{\mathrm{u}}_{i}^{h}$
(18) $\displaystyle\bm{\mathrm{\varepsilon}}^{*}(\bm{\mathrm{x}})$
$\displaystyle=\sum_{i\in\mathcal{N}_{x}}\nabla_{s}(\psi_{i}(\bm{\mathrm{x}})\bm{\mathrm{u}}_{i}^{h}),$
(19)
where the $\bm{\mathrm{u}}_{i}^{h}$ are the raw nodal XFEM displacements and
$\psi_{i}(\bm{\mathrm{x}})$ are the MLS shape function values associated with
a node $i$ at a point $\bm{\mathrm{x}}$. $\mathcal{N}_{x}$ is the set of
$n_{x}$ nodes in the domain of influence of point $\bm{\mathrm{x}}$,
$\nabla_{s}$ is the symmetric gradient operator, $\bm{\mathrm{u}}^{*}$ and
$\bm{\mathrm{\varepsilon}}^{*}$ are the recovered displacement and strain
fields respectively. At each point $\bm{\mathrm{x}}_{i}$ the MLS shape
functions $\psi_{i}$ are evaluated using weighting functions $\omega_{i}$ and
an enriched basis $\bm{\mathrm{p}}(\bm{\mathrm{x}}_{i})$. The total $n_{x}$
non-zero MLS shape functions at point $\bm{\mathrm{x}}$ are evaluated as:
$(\psi_{i}(\bm{\mathrm{x}}))_{1\leq i\leq
n_{x}}=(\bm{\mathrm{A}}^{-1}(\bm{\mathrm{x}})\bm{\mathrm{p}}(\bm{\mathrm{x}}))^{T}\bm{\mathrm{p}}(\bm{\mathrm{x}}_{i})\omega_{i}(\bm{\mathrm{x}})$
(20)
where
$\bm{\mathrm{A}}(\bm{\mathrm{x}})=\sum_{i=1}^{n_{x}}\omega_{i}(\bm{\mathrm{x}})\bm{\mathrm{p}}(\bm{\mathrm{x}}_{i})\bm{\mathrm{p}}(\bm{\mathrm{x}}_{i})^{T}$
is a matrix to be inverted at every point $\bm{\mathrm{x}}$ (see Bordas and
Duflot (2007), Bordas et al. (2008) for further details).
For each supporting point $\bm{\mathrm{x}}_{i}$, the weighting function
$\omega_{i}$ is defined such that:
$\omega_{i}(s)=f_{4}(s)=\begin{cases}1-6s^{2}+8s^{3}-3s^{4}&\text{if
}\left|s\right|\leq 1\\\ 0&\text{if }\left|s\right|>1\end{cases}$ (21)
where $s$ is the normalized distance between the supporting point
$\bm{\mathrm{x}}_{i}$ and a point $\bm{\mathrm{x}}$ in the computational
domain. In order to describe the discontinuity, the distance $s$ in the weight
function defined for each supporting point is modified using the _diffraction
criterion_ (Belytschko et al., 1996). The basic idea of this criterion is
depicted in Figure 2. The weight function is continuous except across the
crack faces since the points at the other side of the crack are not considered
as part of the support. Near the crack tip the weight of a node $i$ over a
point of coordinates $\bm{\mathrm{x}}$ diminishes as the crack hides the
point. When the point $\bm{\mathrm{x}}$ is hidden by the crack the following
expression is used:
$s=\frac{\left\|\bm{\mathrm{x}}-\bm{\mathrm{x}}_{C}\right\|+\left\|\bm{\mathrm{x}}_{C}-\bm{\mathrm{x}}_{i}\right\|}{d_{i}}$
(22)
where $d_{i}$ is the radius of the support.
Figure 2: Diffraction criteria to introduce the discontinuity in the XMLS
approximation.
The MLS shape functions can reproduce any function in their basis. The basis
$\bm{\mathrm{p}}$ used is a linear basis enriched with the functions that
describe the first order asymptotic expansion at the crack tip as indicated in
(11):
$\bm{\mathrm{p}}=\left[1,x,y,\left[F_{1}(r,\phi),F_{2}(r,\phi),F_{3}(r,\phi),F_{4}(r,\phi)\right]\right]$
(23)
Note that although the enriched basis can reproduce the singular behaviour of
the solution around the crack tip, the resulting recovered field not
necessarily satisfies the equilibrium equations.
## 6 Numerical results
In this section, numerical experiments are performed to verify the behaviour
of both XFEM recovery based error estimators considered in this paper.
Babuška et al. (1994a, b, 1997) proposed a robustness patch test for quality
assessment of error estimators. However, the use of this test is not within
the scope of this paper and furthermore, to the authors’ knowledge, it has not
been used in the context of XFEM. Therefore, the more traditional approach of
using benchmark problems is considered here to analyse the response of the
different error estimators.
The accuracy of the error estimators is evaluated both locally and globally.
This evaluation has been based on the effectivity of the error in the energy
norm, which is quantified using the _effectivity index_ $\theta$ defined as:
$\theta=\frac{\lVert\textbf{e}_{es}\rVert}{\lVert\textbf{e}\rVert}.$ (24)
When the enhanced or recovered solution is close to the analytical solution
the effectivity approaches the theoretical value of 1, which indicates that it
is a good error estimator, i.e. the approximate error is close to the exact
error.
To assess the quality of the estimator at a local level, the local effectivity
$D$, inspired on the _robustness index_ found in Babuška et al. (1994a), is
used. For each element $e$, $D$ represents the variation of the effectivity
index in this element, $\theta^{e}$, with respect to the theoretical value
(the error estimator can be considered to be of good quality if it yields $D$
values close to zero). $D$ is defined according to the following expression,
where superscripts e indicate the element $e$:
$\begin{array}[]{c}{\displaystyle D=\theta^{e}-1\qquad{\rm
if}\qquad\theta^{e}\geq 1}\\\ {\displaystyle
D=1-\frac{1}{\theta^{e}}\qquad{\rm
if}\qquad\theta^{e}<1}\end{array}\qquad\qquad{\rm
with}\qquad\theta^{e}=\frac{\left\|\textbf{e}_{es}^{e}\right\|}{\left\|\textbf{e}^{e}\right\|}.$
(25)
To evaluate the overall quality of the error estimator we use the global
effectivity index $\theta$, the mean value $m\left(|D|\right)$ and the
standard deviation $\sigma\left(D\right)$ of the local effectivity. A good
quality error estimator yields values of $\theta$ close to one and values of
$m\left(|D|\right)$ and $\sigma\left(D\right)$ close to zero.
The techniques can be used in practical applications. However, in order to
properly compare their performance we have used an academic problem with exact
solution. In the analysis we solve the Westergaard problem (Gdoutos, 1993) as
it is one of the few problems in LEFM under mixed mode with an analytical
solution. In Giner et al. (2005), Ródenas et al. (2008) we can find explicit
expressions for the stress fields in terms of the spatial coordinates. In the
next subsection we show a description of the Westergaard problem and XFEM
model, taken from Ródenas et al. (2008) and reproduced here for completeness.
### 6.1 Westergaard problem and XFEM model
The Westergaard problem corresponds to an infinite plate loaded at infinity
with biaxial tractions $\sigma_{x\infty}=\sigma_{y\infty}=\sigma_{\infty}$ and
shear traction $\tau_{\infty}$, presenting a crack of length $2a$ as shown in
Figure 3. Combining the externally applied loads we can obtain different
loading conditions: pure mode I, II or mixed mode.
Figure 3: Westergaard problem. Infinite plate with a crack of length $2a$
under uniform tractions $\sigma_{\infty}$ (biaxial) and $\tau_{\infty}$.
Finite portion of the domain $\Omega_{0}$, modelled with FE.
In the numerical model only a finite portion of the domain ($a=1$ and $b=4$ in
Figure 3) is considered. The projection of the stress distribution
corresponding to the analytical Westergaard solution for modes I and II, given
by the expressions below, is applied to its boundary:
$\begin{array}[]{r@{\hspace{1ex}}c@{\hspace{1ex}}l}{\sigma_{x}^{I}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip
4.30554pt&{\displaystyle\frac{\sigma_{\infty}}{\sqrt{\left|t\right|}}\bigg{[}\left(x\cos\frac{\phi}{2}-y\sin\frac{\phi}{2}\right)+y\frac{a^{2}}{\left|t\right|^{2}}\left(m\sin\frac{\phi}{2}-n\cos\frac{\phi}{2}\right)\bigg{]}}\\\
\vskip 6.0pt plus 2.0pt minus 2.0pt\cr{\sigma_{y}^{I}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip
4.30554pt&{\displaystyle\frac{\sigma_{\infty}}{\sqrt{\left|t\right|}}\bigg{[}\left(x\cos\frac{\phi}{2}-y\sin\frac{\phi}{2}\right)-y\frac{a^{2}}{\left|t\right|^{2}}\left(m\sin\frac{\phi}{2}-n\cos\frac{\phi}{2}\right)\bigg{]}}\\\
\vskip 6.0pt plus 2.0pt minus 2.0pt\cr{\tau_{xy}^{I}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip 4.30554pt&{\displaystyle
y\frac{a^{2}\sigma_{\infty}}{\left|t\right|^{2}\sqrt{\left|t\right|}}\left(m\cos\frac{\phi}{2}+n\sin\frac{\phi}{2}\right)}\end{array}$
(26)
$\begin{array}[]{r@{\hspace{1ex}}c@{\hspace{1ex}}l}{\sigma_{x}^{II}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip
4.30554pt&{\displaystyle\frac{\tau_{\infty}}{\sqrt{\left|t\right|}}\bigg{[}2\left(y\cos\frac{\phi}{2}+x\sin\frac{\phi}{2}\right)-y\frac{a^{2}}{\left|t\right|^{2}}\left(m\cos\frac{\phi}{2}+n\sin\frac{\phi}{2}\right)\bigg{]}}\\\
\vskip 6.0pt plus 2.0pt minus 2.0pt\cr{\sigma_{y}^{II}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip 4.30554pt&{\displaystyle
y\frac{a^{2}\tau_{\infty}}{\left|t\right|^{2}\sqrt{\left|t\right|}}\left(m\cos\frac{\phi}{2}+n\sin\frac{\phi}{2}\right)}\\\
\vskip 6.0pt plus 2.0pt minus 2.0pt\cr{\tau_{xy}^{II}}(x,y)\hskip
4.30554pt&{=}\hfil\hskip
4.30554pt&{\displaystyle\frac{\tau_{\infty}}{\sqrt{\left|t\right|}}\bigg{[}\left(x\cos\frac{\phi}{2}-y\sin\frac{\phi}{2}\right)+y\frac{a^{2}}{\left|t\right|^{2}}\left(m\sin\frac{\phi}{2}-n\cos\frac{\phi}{2}\right)\bigg{]}}\end{array}$
(27)
where the stress fields are expressed as a function of $x$ and $y$, with
origin at the centre of the crack. The parameters $t$, $m$, $n$ and $\phi$ are
defined as
$\begin{split}t&=(x+iy)^{2}-a^{2}=(x^{2}-y^{2}-a^{2})+i(2xy)=m+in\\\
m&=\textrm{Re}(t)=\textrm{Re}(z^{2}-a^{2})=x^{2}-y^{2}-a^{2}\\\
n&=\textrm{Im}(t)=(z^{2}-a^{2})=2xy\\\
\phi&=\textrm{Arg}(\bar{t})=\textrm{Arg}(m-in)\qquad\textrm{with
}\phi\in\left[-\pi,\pi\right],\;i^{2}=-1\end{split}$ (28)
For the problem analysed, the exact value of the SIF is given as
$K_{{\rm I},ex}=\sigma\sqrt{\pi a}\qquad\qquad K_{{\rm II},ex}=\tau\sqrt{\pi
a}$ (29)
Three different loading configurations corresponding to the _pure mode I_
($\sigma_{\infty}=100,\;\tau_{\infty}=0$), _pure mode II_
($\sigma_{\infty}=0,\;\tau_{\infty}=100$), and _mixed mode_
($\sigma_{\infty}=30,\;\tau_{\infty}=90$) cases of the Westergaard problem are
considered. The geometric models and boundary conditions are shown in Figure
4.
Figure 4: Model for an infinite plate with a crack subjected to biaxial
tractions $\sigma_{\infty},\;\tau_{\infty}$ in the infinite.
Bilinear elements are used in the models, with a _singular_ +_smooth_
decomposition area of radius $\rho=0.5$ equal to the radius $r_{e}$ used for
the fixed enrichment area (note that the splitting radius $\rho$ should be
greater or equal to the enrichment radius $r_{e}$, Ródenas et al. (2008)). The
radius for the Plateau function used to extract the SIF is $r_{q}=0.9$.
Young’s modulus is $E=10^{7}$ and Poisson’s ratio $\upsilon=0.333$. For the
numerical integration of standard elements we use a $2\times 2$ Gaussian
quadrature rule. For split elements we use 7 Gauss points in each triangular
subdomain, and a $5\times 5$ quasipolar integration (Béchet et al., 2005) in
the subdomains of the element containing the crack tip.
Regarding global error estimation, the evolution of global parameters in
sequences of uniformly refined structured (Figure 5) and unstructured (Figure
6) meshes is studied. In the first case, the mesh sequence is defined so that
the crack tip always coincided with a node. For a more general scenario, this
condition has not been applied for the unstructured meshes.
Figure 5: Sequence of structured meshes. Figure 6: Sequence of unstructured
meshes.
### 6.2 Mode I and structured meshes
The first set of results presented in this study are for the Westergaard
problem under mode I load conditions. Figure 7 shows the values for the local
effectivity index $D$ for the first mesh in the sequence analysed. In the
figure, the size of the enrichment area is denoted by a circle. It can be
observed that far from the enrichment area both techniques yield similar
results, however, the XMLS technique exhibits a higher overestimation of the
error close to the singularity.
Figure 7: Mode I, structured mesh 1. Local effectivity index $D$ (ideal value
$D=0$). The circles denote the enriched zones around the crack tips. The
superiority of the statically admissible recovery (SPR-CX) compared to the
standard XMLS is clear.
The same behaviour is also observed for more refined meshes in the sequence.
Figure 8 shows a zoom in the enriched area of a finer mesh where, as before, a
higher overestimation of the error can be observed for the results obtained
with the XMLS error estimator. In this example and in order to study the
evolution of the accuracy of the recovered stress field when considering
different features in the recovery process, we consider two additional
versions:
* •
The first considers a recovery procedure which enforces both internal and
boundary equilibrium, but does not include the singular+smooth splitting
technique (SPR-C). This approach is very similar in form to conventional SPR-
based recovery techniques widely used in FEM. It is well-known that this type
of recovery process will produce unreliable results when the stress field
contains a singularity because the polynomial representation of the recovered
stresses is not able to describe the singular field (e.g. Bordas and Duflot
(2007)).
* •
The second version of the error estimator performs the splitting but does not
equilibrate the recovered field (SPR-X). As previously commented,
$\bm{\mathrm{\sigma}}^{*}_{sing}$ is an equilibrated field but
$\bm{\mathrm{\sigma}}^{*}_{smo}$ is not equilibrated. The aim of this second
version is to assess the influence of enforcing the equilibrium constraints.
Table 1 summarizes the main features of the different recovery procedures we
considered.
| Singular | Equilibrated sin- | Locally equili-
---|---|---|---
| functions | gular functions | brated field
SPR-X | yes | yes | no
SPR-C | no | not applicable | yes
SPR-CX | yes | yes | yes
XMLS | yes | no | no
Table 1: Comparison of features for the different recovery procedures Figure
8: Mode I, structured mesh 4. Local effectivity index $D$ (ideal value $D=0$).
The circles denote the enriched zones around the crack tips. The results show
the need for the inclusion of the near-tip fields in the recovery process
(SPR-X is far superior to SPR-C). It is also clear that enforcing statical
admissibility (SPR-CX) greatly improves the accuracy along the crack faces and
the Neumann boundaries. Also notice that XMLS leads to slightly improved
effectivities compared to SPR-CX around the crack faces, between the enriched
zone and the left boundary. SPR-C is clearly not able to reproduce the
singular fields, which is shown by large values of $D$ inside the enriched
region.
The performance of SPR-X in Figure 8 is poor, especially along the Neumann
boundaries, where the equilibrium equations are not enforced, as they are in
the SPR-CX. Particularly interesting are the results for the SPR-C technique,
where considerable overestimation of the exact error is observed _in the whole
enriched region_. This is due to the inability of polynomial functions to
reproduce the near-tip fields, even relatively far from the crack tip.
The above compared the relative performance of the methods locally. It is also
useful to have a measure of the global accuracy of the different error
estimators, Figure 9(left) shows the convergence of the estimated error in the
energy norm $\lVert\textbf{e}_{es}\rVert$ (see equation (12)) evaluated with
the proposed techniques, the convergence of the exact error
$\lVert\textbf{e}\rVert$111The exact error measures the error of the raw XFEM
compared to the exact solution. The theoretical convergence rate is
$\mathcal{O}(h^{p})$, $h$ being the element size and $p$ the polynomial degree
(1 for linear elements,…), or $\mathcal{O}(-\text{dof}^{p/2})$ if we consider
the number of degrees of freedom. is shown for comparison. It can be observed
that SPR-CX leads to the best approximation to the exact error in the energy
norm. The best approach is that which captures as closely as possible the
exact error. Non-singular formulations such as SPR-C greatly overestimate the
error, especially close to the singular point. Only SPR-C suffers from the
characteristic suboptimal convergence rate (factor 1/2, associated with the
strength of the singularity)
In order to evaluate the quality of the recovered field
$\bm{\mathrm{\sigma}}^{*}$ obtained with each of the recovery techniques, the
convergence in energy norm of the approximate error on the recovered field
$\lVert\textbf{e}^{*}\rVert=\lVert\bm{\mathrm{u}}-\bm{\mathrm{u}}^{*}\rVert$
is compared for the error estimators in Figure 9(right). It can be seen that
the XMLS, the SPR-CX and the SPR-X provide convergence rates higher than the
convergence rate for the exact error $\lVert\textbf{e}\rVert$, which is an
indicator of the asymptotic exactness of the error estimators based on these
recovery techniques (Zienkiewicz and Zhu, 1992b).
Figure 9: Mode I, structured meshes: Convergence of the estimated error in the
energy norm $\|\mathbf{e}_{es}\|$ (left). All methods which include the near-
tip fields do converge with the optimal convergence rate of 1/2. The best
method is SPR-CX. XMLS performs worse than SPR-X for coarse meshes, but it
becomes equivalent to SPR-CX for fine meshes whilst the SPR-X error increases
slightly. Convergence of the error for the recovered field
$\|\mathbf{e}_{es}^{*}\|$ (right). For the recovered field, SPR-CX is the best
method closely followed by SPR-X. XMLS errors are half an order of magnitude
larger than SPR-CX/SPR-X, but superior to the non-enriched recovery method
SPR-C. For all methods, except SPR-C, the error on the recovered solution
converges faster than the error on the raw solution, which shows that the
estimators are asymptotically exact.
Figure 10 shows the evolution of the parameters $\theta$, $m\left(|D|\right)$
y $\sigma\left(D\right)$ with respect to the number of degrees of freedom for
the error estimators. It can be verified that although the SPR-C technique
includes the satisfaction of the equilibrium equations, the effectivity of the
error estimator does not converge to the theoretical value ($\theta=1$). This
is due to the absence of the near-tip fields in the recovered solution. The
influence of introducing singular functions in the recovery process can be
observed in the results provided by SPR-X. In this case, the convergence is
obtained both locally and globally. This shows the importance of introducing a
description of the singular field in the recovery process, i.e. of using
_extended recovery techniques_ in XFEM.
Figure 10: Global indicators $\theta$, $m\left(|D|\right)$ and
$\sigma\left(D\right)$ for mode I and structured meshes.
Regarding the enforcement of the equilibrium equations, we can see that an
equilibrated formulation (SPR-CX) leads to better effectivities compared to
other non-equilibrated configurations (SPR-X, XMLS), which is an indication of
the advantages associated with equilibrated recoveries in this context.
Still in Figure 10, the results for the SPR-CX and XMLS show that the values
of the global effectivity index are close to (and less than) unity for the
SPR-CX estimator, whilst the XMLS is a lot less effective and shows
effectivities larger than 1.
The next step in our analysis is the verification of the convergence of the
mean value and standard deviation of the local effectivity index $D$. From a
practical perspective, one would want the local effectivity to be equally
good, on average, everywhere inside the domain; one would also want this
property to improve with mesh refinement. In other words, the average local
effectivity $m$ and its standard deviation $\sigma$ should decrease with mesh
refinement: $m$ and $\sigma$ measure the average behaviour of the method and
how far the results deviate from the mean, i.e. how spatially consistent they
are. The idea is to spot areas where there may be compensation between
overestimated and underestimated areas that could produce an apparently
’accurate’ global estimation. It can be confirmed that, although the curves
for $m(|D|)$ and $\sigma(D)$ for both estimators tend to zero with mesh
refinement, the SPR-CX is clearly superior to the XMLS. These results can be
explained by the fact that the SPR-CX enforces the fulfilment of the
equilibrium equations in patches, and evaluates the singular part of the
recovered field using the equilibrated expression that represents the first
term of the asymptotic expansion, whereas the XMLS enrichment uses only a set
of singular functions that would be able to reproduce this term but are not
necessarily equilibrated.
Although the results for the SPR-CX proved to be more accurate, the XMLS-type
techniques will prove useful to obtain error bounds. There is an increasing
interest in evaluating upper and lower bounds of the error for XFEM
approximations. Some work has been already done to obtain upper bounds using
recovery techniques as indicated in Ródenas et al. (2010a), where the authors
showed that an upper bound can be obtained if the recovered field is
continuous and equilibrated. However, they demonstrated that due to the use of
the conjoint polynomials process to enforce the continuity of
$\bm{\mathrm{\sigma}}^{*}$, some residuals in the equilibrium equations
appear, and the recovered field is only nearly equilibrated. Then, correction
terms have to be evaluated to obtain the upper bound. An equilibrated version
of the XMLS could provide a recovered stress field which would be continuous
and equilibrated, thus facilitating the evaluation of the upper bound. Initial
results for this class of techniques can be found in Ródenas et al. (2009,
2010b).
In this first example we have shown the effect of using the SPR-C and SPR-X
recovery techniques in the error estimation. Considering that these two
techniques are special cases of the SPR-CX with inferior results, for further
examples we will focus only on the SPR-CX and XMLS techniques.
### 6.3 Mode II and structured meshes
Figure 11 presents the results considering mode II loading conditions for the
local effectivity index $D$ on the fourth mesh of the sequence. Similarly to
the results for mode I, the XMLS estimator presents a higher overestimation of
the error near the enrichment area. This same behaviour is observed for the
whole set of structured meshes analysed under pure mode II.
Figure 11: Mode II, structured mesh 4. Local effectivity index $D$ (ideal
value $D=0$). SPR-CX leads to much better local effectivities than XMLS. It is
also remarkable that the effectivities are clearly worse in the enriched
region (circle) for both the XMLS and the SPR-CX, and that this effect is more
pronounced in the former method. Note the slightly worse results obtained by
SPR-CX around the crack faces between the boundary of the enriched region and
the left boundary, as in mode I.
As for the mode I case, the evolution of global accuracy parameters in mode II
exhibits the same behaviour seen for mode I loading conditions. Figure 12
shows the results for the convergence of the error in the energy norm. Figure
13 shows the evolution of global indicators $\theta$, $m\left(|D|\right)$ and
$\sigma\left(D\right)$. Again, SPR-CX provides the best results for the
Westergaard problem under pure mode II, using structured meshes.
Figure 12: Mode II, structured meshes: Convergence of the estimated error in
the energy norm $\|\mathbf{e}_{es}\|$ (left). The convergence rates are very
similar to those obtained in the mode I case, and, SPR-CX is still the most
accurate method, i.e. that for which the estimated error is closest to the
exact error. Convergence of the error in the recovered field
$\|\mathbf{e}_{es}^{*}\|$ (right). It can be noticed that the SPR-CX error
still converges faster than the exact error, but with a lower convergence rate
(0.59 versus 0.73) than in the mode I case. There is no such difference for
the XMLS results, which converge at practically the same rate (0.71 versus
0.72). The error level difference between SPR-CX and XMLS is about half an
order of magnitude, as in the mode I case. Figure 13: Global indicators
$\theta$, $m\left(|D|\right)$ and $\sigma\left(D\right)$ for mode II and
structured meshes. The results are qualitatively and quantitatively similar to
those obtained in mode I.
### 6.4 Mixed mode and unstructured meshes
Considering a more general problem, Figure 14 shows the local effectivity
index $D$ for the fourth mesh in a sequence of unstructured meshes, having a
number of degrees of freedom (dof) similar to the mesh represented previously
for load modes I and II. The XMLS recovered field presents higher
overestimation of the error around the crack tip which corroborates the
results found in previous load cases.
Figure 14: Mixed mode, unstructured mesh 4. Local effectivity index $D$ (ideal
value $D=0$). Note the improved results compared to the structured case with
$D$ values ranging from -1.5 to 1.5 as opposed to -4 to 4.
Figures 15 and 16 represent the evolution of global parameters for
unstructured meshes and mixed mode. Once more, the best results for the error
estimates are obtained using the SPR-CX technique.
Figure 15: Mixed mode, unstructured meshes: Convergence of the estimated
error in the energy norm $\|\mathbf{e}_{es}\|$ (left). Convergence of the
error for the recovered field $\|\mathbf{e}_{es}^{*}\|$ (right). The results
are almost identical to the structured mesh case, except for the faster
convergence obtained for SPR-CX in the unstructured compared to the structured
case (0.77 versus 0.59), keeping in mind that optimal convergence rates are
only formally obtained for structured meshes. Figure 16: Global indicators
$\theta$, $m\left(|D|\right)$ and $\sigma\left(D\right)$ for mixed mode and
unstructured meshes. Notice that the XMLS performs more closely to SPR-CX
using those indicators, for unstructured than for structured meshes,
especially when measuring the mean value of the effectivities
$m\left(|D|\right)$.
## 7 Conclusions
The aim of this paper was to assess the accuracy gains provided by
1. 1.
Including relevant enrichment functions in the recovery process;
2. 2.
Enforcing statical admissibility of the recovered solution.
We focused on two recovery-based error estimators already available for LEFM
problems using the XFEM. The first technique called SPR-CX is an enhancement
of the SPR-based error estimator presented by Ródenas et al. (2008), where the
stress field is split into two parts (singular and smooth) and equilibrium
equations are enforced locally on patches. The second technique is the XMLS
proposed by Bordas and Duflot (2007), Bordas et al. (2008) which enriches the
basis of MLS shape functions, and uses a diffraction criterion, in order to
capture the discontinuity along the crack faces and the singularity at the
crack tip.
To analyse the behaviour of the ZZ error estimator using both techniques and
to assess the quality of the recovered stress field, we have evaluated the
effectivity index considering problems with an exact solution. Convergence of
the estimated error in the energy norm and other local error indicators are
also evaluated. To analyse the influence of the special features introduced in
the recovery process we have also considered two additional configurations:
SPR-C and SPR-X.
The results indicate that both techniques, SPR-CX and XMLS, provide error
estimates that converge to the exact value and can be considered as
asymptotically exact. Both techniques could be effectively used to estimate
the error in XFEM approximations, while other conventional recovery procedures
which do not include the enrichment functions in the recovery process have
proved not to converge to the exact error. This shows (albeit only
numerically) the need for the use of extended recovery techniques for accuracy
assessment in the XFEM context.
Better results are systematically obtained when using the SPR-CX to recover
the stress field, specially in the areas close to the singular point. For all
the different test cases analysed, the XMLS produced higher values of the
effectivity index in the enriched area, where the SPR-CX proved to be more
accurate. This can be ascribed to the fact that the SPR-CX technique recovers
the singular part of the solution using the known equilibrated exact
expressions for the asymptotic fields around the crack tip and enforces the
fulfilment of the equilibrium equations on patches. Further work currently in
progress will include the development of an equilibrated XMLS formulation,
which could provide a continuous and globally equilibrated recovered stress
field which can then be used to obtain upper bounds of the error in energy
norm.
The aim of our project is to tackle practical engineering problems such as
those presented in Bordas and Moran (2006), Wyart et al. (2007) where the
accuracy of the stress intensity factor is the target. The superiority of SPR-
CX may then be particularly relevant. To minimise the error on the stress
intensity factors we will target this error directly, through goal-oriented
error estimates. This will be reported in a forthcoming publication. However,
although the SPR-CX results are superior in general, an enhanced version of
the XMLS technique presented in this paper where we enforce equilibrium
conditions could result useful to evaluate upper error bounds when considering
goal-oriented error estimates, as it directly produces a continuous
equilibrated recovered stress field.
Our next step will be the comparison of available error estimators in three
dimensional settings in terms of accuracy versus computational cost to
minimise the error on the crack path and damage tolerance of the structure.
## 8 Acknowledgements
Stéphane Bordas would like to thank the support of the Royal Academy of
Engineering and of the Leverhulme Trust for his Senior Research Fellowship
entitled “Towards the next generation surgical simulators” as well as the
support of EPSRC under grant EP/G042705/1 Increased Reliability for
Industrially Relevant Automatic Crack Growth Simulation with the eXtended
Finite Element Method.
This work has been carried out within the framework of the research projects
DPI 2007-66773-C02-01, DPI2010-20542 and DPI2010-20990 of the Ministerio de
Ciencia e Innovación (Spain). Funding from FEDER, Universitat Politècnica de
València and Generalitat Valenciana is also acknowledged.
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Corresponding author
Octavio Andrés González-Estrada can be contacted at: ocgones@upv.es
|
arxiv-papers
| 2011-12-09T17:41:45 |
2024-09-04T02:49:25.142079
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Octavio A. Gonz\\'alez-Estrada, Juan Jos\\'e R\\'odenas, St\\'ephane P.A.\n Bordas, Marc Duflot, Pierre Kerfriden, Eugenio Giner",
"submitter": "Octavio Andres Gonzalez Estrada",
"url": "https://arxiv.org/abs/1112.2160"
}
|
1112.2196
|
# Ultra-high-Q wedge-resonator on a silicon chip
Hansuek Lee, Tong Chen, Jiang Li, Ki Youl Yang, Seokmin Jeon, Oskar Painter,
and Kerry J. Vahala
T. J. Watson Laboratory of Applied Physics, California Institute of
Technology, Pasadena, California 91125, USA
Ultra-high-Q optical resonators are being studied across a wide range of
research subjects including quantum information, nonlinear optics, cavity
optomechanics, and telecommunications vahala_nature ; review_vahala ;
review_vahala2 ; review_ilchenko ; review_ilchenko2 ; review_comb ; cqed .
Here, we demonstrate a new, resonator on-a-chip with a record Q factor of 875
million, surpassing even microtoroids toroid . Significantly, these devices
avoid a highly specialized processing step that has made it difficult to
integrate microtoroids with other photonic devices and to also precisely
control their size. Thus, these devices not only set a new benchmark for Q
factor on a chip, but also provide, for the first time, full compatibility of
this important device class with conventional semiconductor processing. This
feature will greatly expand the possible kinds of system on a chip functions
enabled by ultra-high-Q devices.
Figure 1: Micrographs and mode renderings of the wedge resonator from top and
side views. (a) An optical micrograph shows a top view of a $1\,$mm diameter
wedge resonator. (b) A scanning electron micrograph shows the side-view of a
resonator. The insets here give slightly magnified micrographs of resonators
in which the wedge angle is $12$ degrees (upper inset) and $27\,$ degrees
(lower inset). (c) A rendering shows calculated fundamental mode intensity
profiles in 10 degree and 30 degree wedge angle resonators at two wavelengths.
As a guide, the center-of-motion of the mode is provided to illustrate how the
wedge profile introduces normal dispersion that is larger for smaller wedge
angles.
Long photon storage time (high Q factor) in microcavities relies critically
upon use of low absorption dielectrics and creation of very smooth (low
scattering) dielectric interfaces. For chip-compatible devices, silica has by
far the lowest intrinsic material loss. Microtoroid resonators combine this
low material loss with a reflow technique in which surface tension is used to
smooth lithographic and etch-related blemishes toroid . At the same time,
reflow smoothing makes it very challenging to fabricate larger diameter UHQ
resonators and likewise to leverage the full range of integration tools and
devices available on silicon. The devices reported here attain ultra-high-Q
performance using only conventional semiconductor processing methods on a
silicon wafer. Moreover, the best Q performance occurs for diameters greater
than $500$ microns, a size range that is difficult to access for microtoroids
on account of the limitations of the reflow process. Microcombs will benefit
from such a combination of UHQ and larger diameter resonators (microwave-rate
free-spectral-range) to create combs that are both efficient in turn-on power
and that can be self-referenced review_comb . Moreover, integrated reference
cavities and ring gyroscopes are two other applications that can benefit from
larger ($1-50\,$mm diameter) UHQ resonators. Fabrication control of the free-
spectral range to $1:20,000$ is also demonstrated here, opening the
possibility of precision repetition rate control in microcombs or precision
spectral placement of modes in certain nonlinear oscillators SBS_carmon ;
SBS_ivan .
Earlier work considered the Q factor in a wedge-shaped resonator fabricated of
silica on a silicon wafer. Q factors as high as $50$ million were obtained PRA
. That approach isolated the mode from the lithographic blemishes near the
outer rim of the resonator by using a shallow wedge angle. In the current
work, we have boosted the optical Q by about $20$X beyond these earlier
results through a combination of process improvements. These improvements make
it unnecessary to isolate the mode from the resonator rim. Indeed, the highest
Q factor demonstrated uses the largest wedge angles. A top-view optical
micrograph is provided in figure $1$ to illustrate the basic geometry. The
process flow begins with thermal oxide on silicon, followed by lithography and
oxide etching with buffered hydrofluoric acid. In the insets to figure $1$,
scanning electron micrographs of devices featuring $12$-degree and $27$-degree
wedge angles are imaged. Empirically, the angle can be controlled through
adjustment of the photoresist adhesion using commercially available adhesion
promoters. The oxide disk structures function as an etch mask for an isotropic
dry etch of the silicon using XeF2. During the dry etch, the silicon undercut
is set so as to reduce coupling of the optical mode to the silicon support
pillar. This value is typically set to about $100$ microns for $1\,$mm
diameter structures and over $150$ microns for $7.5\,$mm diameter disks,
however, smaller undercuts are possible while preserving ultra-high-Q
performance. Further information on the processing is given in the Methods
section.
To measure intrinsic Q factor, devices were coupled to SMF-$28$ optical fiber
using a fiber taper taper ; ideality_PRL and spectral lineshape data were
obtained by tuning an external cavity semiconductor laser across the resonance
while monitoring transmission on an oscilloscope. To accurately calibrate the
laser scan in this measurement, a portion of the laser output was also
monitored after transmission through a calibrated Mach-Zehnder interferometer
having a free spectral range of $7.75$ MHz. The inset in figure $2$ shows a
spectral scan obtained on a device having a record Q factor of $875$ million.
In these measurements, the taper coupling was applied on the upper surface of
the resonator near the center of the wedge region. Modeling shows that the
fundamental mode has its largest field amplitude in this region. Moreover,
this mode is expected to feature the lowest overall scattering loss resulting
from the three, dielectric-air interfaces as well as from the silicon support
pillar. An additional test that can be performed to verify the fundamental
mode is to measure the mode index by monitoring the free-spectral-range (FSR).
The fundamental mode features the largest mode index and hence smallest FSR.
Figure 2: Data showing the measured Q factor plotted versus resonator diameter
with oxide thickness as a parameter. The solid lines show the predicted Q
factor from a model that accounts for surface roughness induced scattering
loss and also material loss. The rms roughness is measured using an AFM (see
Methods section for values) and the fitted bulk material loss corresponds to a
Q value of $2.5$ billion. The red data points correspond to a wedge angle of
$27$ degrees. All other data are obtained using a wedge angle of approximately
$10$ degrees. The inset shows a spectral scan for the case of a record Q
factor of $875$ million. The sinusoidal curve accompanying the spectrum is a
calibration scan performed using a fiber interferometer.
The typical coupled power in all measurements was maintained around $1$
microWatt to minimize thermal effects. However, there was little or no
evidence of thermal effects in the optical spectrum. Typically, these appear
as an asymmetry in the lineshape and also a scan-direction dependent (to
higher or lower frequency) spectral linewidth. As a further check that thermal
effects were negligible, ring down measurements ring_down were also performed
on a range of devices for comparison to the spectral-based Q measurement. For
these, the laser was tuned into resonance with the cavity and a lithium
niobate modulator was used to abruptly switch off the input. The output cavity
decay rate was then monitored to ascertain the cavity lifetime. Ring-down data
and spectral linewidths were consistently in good agreement. This
insensitivity to thermal effects is a result of the larger mode volumes of
these devices in comparison to earlier work on microtoroids (for which thermal
effects must be carefully monitored). The mode volumes in the present devices
are typically $100-1000$X larger.
Figure 3: Plot of measured free spectral range (FSR) versus the target design-
value resonator diameter on a lithographic mask. The plot shows one device at
each size and five different sizes. The rms variance is $2.4\,$MHz (relative
variance of less than $1:4,500$). The inset shows the FSR data measured on
four devices having the same target FSR. An improved variance of $0.45\,$MHz
is obtained (a relative variance of $1:20,000$).
Measurements showing the effects of oxide thickness and device diameter on Q
factor are presented in the figure $2$ main panel. Four, oxide thicknesses are
shown ($2$, $4$, $7.5$ and $10$ microns) over diameters ranging from $0.2\,$mm
to $7.5\,$mm. All data points, with the exception of the red points,
correspond to a wedge angle of approximately $10$ degrees. The upper most
(highest Q at a given diameter) data correspond to a wedge angle of $27$
degrees. The solid curves are a model of optical loss caused by surface
scattering on the upper, wedge, and lower oxide-air interfaces and by bulk-
oxide loss. In the model, the surface roughness was measured independently on
each of these surfaces using an atomic force microscope (AFM) (r.m.s.
roughness values are given in the Methods section). The bulk optical loss of
the thermal silica corresponds to a Q value of $2.5$ billion by fit to the
data. The data corresponding to the $10$ degree wedge angle show that Q
increases for thicker oxides and also larger diameters. Using the model, this
trend can be understood to result from loss that is caused primarily by
scattering at the oxide-air interfaces. Specifically, both thicker oxides and
larger diameter structures feature a reduced field amplitude at the
dielectric-air interface, leading to reduced scattered power. A slight,
overall boost to the Q factor is possible by increasing the wedge angle. In
this case, the mode experiences reduced upper and lower surface scattering as
compared to the smaller angle case. A record Q factor of $875$ million for any
chip based resonators is obtained under these conditions. In general, there is
reasonably good agreement between the model and the data, except in the case
of the thinner oxides. For these thinner structures, there is a tendency for
stress-induced buckling to occur at larger radii. This is believed to create
the discrepancy with the model.
Figure 4: Data plot showing the effect of etch time on appearance of the
“foot” region in etching of a $10$ micron thick silica layer. The foot region
is a separate etch front produced by wet etch of silica that is empirically
observed to adversely affect the optical Q factor. The data show that by
control of the etch time the “foot” region can be eliminated. The upper-left
inset is an image of the foot region and the lower right inset shows the foot
region eliminated by increase of the wet etch time.
The ability to lithographically define ultra-high-Q resonators as opposed to
relying upon the reflow process enables a multi-order-of-magnitude improvement
in precision control of resonator diameter and FSR. This feature is especially
important in microcombs and also certain nonlinear sources SBS_carmon ;
SBS_ivan . As a preliminary test of the practical limits of FSR control, two
studies were conducted. In the first, a series of resonator diameters were set
in a CAD file used to create a photo mask. A plot of the measured FSR
(fundamental mode) versus CAD file target diameter is provided in figure $3$
(main panel). The variance from ideal linear behavior is $2.4\,$MHz, giving a
relative variance of better than $1:4,500$ (FSR $\approx$ $11$ GHz). The inset
to figure $3$ shows that for separate devices having the same target CAD file
diameter, the variance is further improved to a value of $0.45\,$MHz or
$1:20,000$.
The Q factor for these new resonators is not only higher in an absolute sense
than what has been possible with microtoroids, but it also accesses an
important regime of resonator FSR that has not been possible using
microtoroids. To date, the smallest FSR achieved with the toroid reflow
process has been $86\,$GHz (D $=750\,\mu m$) and the corresponding Q factor
was $20$ million kippenberg_PRL_2008 . The present structures attain their
best Q factors for FSRs that are complementary to microtoroids (FSRs less than
$100\,$GHz). This range has become increasingly important in applications like
microcombs where self-referencing is important. Specifically, low turn-on
power and microwave-rate repetition are conflicting requirements in these
devices on account of the inverse dependence of threshold power on FSR.
However, such increases can be compensated using ultra-high-Q because turn-on
power depends inverse quadratically on Q OPO_vahala . The ability to
manipulate normal dispersion through the wedge angle (see figure 1) can be
shown to provide control over the zero dispersion point in spectral regions
where silica exhibits anomalous dispersion. Ultra-high-Q performance in large
area resonators is also important in rotation sensing rotation and for on-
chip frequency references freq_ref ; freq_ref2 . In the former case, the
larger resonator area enhances the Sagnac effect. In the latter, the larger
mode volume lowers the impact of thermal fluctuations on the frequency noise
of the resonator noise . The precision control of FSR is important to
determine repetition rate in microcombs, and also in applications such as
stimulated Brillouin lasers where a precise match of FSR to the Brillouin
shift is a prerequisite for oscillation. Application of these devices to low
turn-on power, microwave-rate microcombs and to high-efficiency SBS lasers
will be reported elsewhere. Finally, an upper bound to the material loss of
thermal silica was established in this work. The value of $2.5$ billion bodes
well for further application of thermal silica to photonic devices.
Methods
Disks were fabricated on ($100$) prime grade float zone silicon wafers. Photo-
resist was patterned using a GCA $6300$ stepper on thermally grown oxide of
thickness in the range of $2-10$ microns. Post exposure bake followed in order
to cure the surface roughness of photo-resist pattern which acted as an etch
mask during immersion in buffered hydrofluoric solution (Transene, buffer-HF
improved). Careful examination of the wet etch revealed that the vertex formed
by the lower oxide and upper surface contains an etch front that is distinct
from that associated with the upper surface (see “foot” region in figure $4$
inset). This region has a roughness level that is higher than any other
surface and is a principle contributor to Q degradation. By extending the etch
time beyond what is necessary to reach the silicon substrate, this foot region
can be eliminated as shown in figure $4$. With elimination of the foot etch
front, the isotropic and uniform etching characteristic of buffered
hydrofluoric solution results in oxide disks and waveguides having very smooth
wedge-profiles which enhance Q factors. After the conventional cleaning
process to remove photo-resist and organics, silicon was isotropically etched
by xenon difluoride to create an air-cladding whispering gallery resonator. An
atomic force microscope was used to measure the surface roughness of the
three, silica-air dielectric surfaces. For the lower surface, the resonators
were detached by first etching the silicon pillar to a few microns in diameter
and then removing the resonator using tape. The r.m.s. roughness values on
$10$-degree wedge-angle devices are: $0.15\,$nm (upper), $0.46\,$nm (wedge),
$0.70\,$nm (lower); and for $27$-degree wedge-angle devices are: $0.15\,$ nm
(upper), $0.75\,$nm (wedge), $0.70\,$nm (lower). The correlation length is
approximately a few hundred nm. The difference in the wedge surface roughness
obtained for the large and small wedge angle cases is not presently
understood.
Acknowledgments We gratefully acknowledge the Defense Advanced Research
Projects Agency under the iPhod and Orchid programs and also the Kavli
Nanoscience Institute at Caltech. H. L. thanks the Center for the Physics of
Information.
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* (20) Gorodetsky, M. L. & Grudinin, I. S. Fundamental thermal fluctuations in microspheres. _J. Opt. Soc. Am. B_ 21, 697–705 (2004).
|
arxiv-papers
| 2011-12-09T20:10:40 |
2024-09-04T02:49:25.153919
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Hansuek Lee, Tong Chen, Jiang Li, Ki Youl Yang, Seokmin Jeon, Oskar\n Painter, and Kerry J. Vahala",
"submitter": "Tong Chen",
"url": "https://arxiv.org/abs/1112.2196"
}
|
1112.2343
|
# Spatial confinement effects on quantum harmonic oscillator I: Nonlinear
coherent state approach
M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi Quantum Optics Group,
Department of physics, University of Isfahan, Isfahan, Iran
m-baghreri@phys.ui.ac.ir rokni@sci.ui.ac.ir mhnaderi@phys.ui.ac.ir , ,
###### Abstract
In this paper we study some basic quantum confinement effects through
investigation of a deformed harmonic oscillator algebra. We show that spatial
confinement effects on a quantum harmonic oscillator can be represented by a
deformation function within the framework of nonlinear coherent states theory.
We construct the coherent states associated with the spatially confined
quantum harmonic oscillator in a one-dimensional infinite well and examine
some of their quantum statistical properties, including sub-poissonian
statistics and quadrature squeezing.
## 1 Introduction
The harmonic oscillator is one of the models most extensively used in both
classical and quantum mechanics. The usefulness and simplicity make this model
a subject of lots of studies. One of the most important aspects of quantum
harmonic oscillator (QHO) is its dynamic algebra i.e. Weyl-Heisenberg algebra.
This algebra appears in many areas of modern theoretical physics, as an
example we notice that the one-dimensional quantum harmonic oscillator was
successfully used in second quantization formalism [1].
Due to the relevance of Weyl-Heisenberg algebra, some efforts have been
devoted to studying possible deformations of the QHO algebra [2]. A deformed
algebra is a nontrivial generalization of a given algebra through the
introduction of one or more deformation parameters, such that, in a certain
limit of parameters the non-deformed algebra is recovered. A particular
deformation of Heisenberg algebra has led to the notion of $f$-oscillator [3].
An $f$-oscillator is a non-harmonic system, that from mathematical point of
view its dynamical variables (creation and annihilation operators) are
constructed from a non canonical transformation through
$\hat{A}=\hat{a}f(\hat{n})\hskip 28.45274pt,\hskip
28.45274pt\hat{A}^{{\dagger}}=f(\hat{n})\hat{a}^{{\dagger}},$ (1)
where $\hat{a}$ and $\hat{a}^{{\dagger}}$ are the usual (non-deformed)
harmonic oscillator operators with $[\hat{a},\hat{a}^{{\dagger}}]=1$ and
$\hat{n}=\hat{a}^{{\dagger}}\hat{a}$. The function $f(\hat{n})$ is called
deformation function which depends on the number of excitation quanta and some
physical parameters. The presence of the operator-valued deformation function
causes the Heisenberg algebra of the standard QHO to transform into a deformed
Heisenberg algebra. The nonlinearity in $f$-oscillators means dependence of
the oscillation frequency on the intensity [4]. On the other hand, in contrast
to the standard QHO, $f$-oscillators have not equal spaced energy spectrum.
For example, if we confine a simple QHO inside an infinite well, due to the
spatial confinement, the energy levels constitute a spectrum that is not equal
spaced. Therefore, in this case it is reasonable to investigate the
corresponding $f$-oscillator.
The confined QHO can be used to describe confinement effects on physical
properties of confined systems. Physical size and shape of the materials
strongly affect the nature, dynamics of the electronic excitations, lattice
vibrations, and dynamics of carriers. For example, in the mesoscopic systems,
the dimension of system is comparable with the coherence length of carriers
and this leads to some new phenomena that they do not appear in a bulk
semiconductor, such as quantum interference between carrier’s motion [5].
Recent progress in growth techniques and development of micromachinig
technology in designing mesoscopic systems and nanostructures, have led to
intensive theoretical [6] and experimental investigations [7] on electronic
and optical properties of those systems. The most important point about the
nanoscale structures is that the quantum confinement effects play the center-
stone role. One can even say, in general, that recent success in
nanofabrication technique have resulted in great interest in various
artificial physical systems (quantum dots, quantum wires and quantum wells)
with new phenomena driven by the quantum confinement. A number of recent
experiments have demonstrated that isolated semiconductor quantum dots are
capable of emitting light [8]. It becomes possible to combine high-Q optical
microcavities with quantum dot emitters as the active medium [9]. Furthermore,
there are many theoretical attempts for understanding the optical and
electronic properties of nanostructures especially semiconductor quantum dots
[10]. On the other hand, a nanostructure such as quantum dot, is a system that
carrier’s motion is confined inside a small region, and during the interaction
with other systems, the generated excitations such as phonons, excitons and
plasmons are confined in small region. In order to describe the physical
properties of these excitations one can consider them as harmonic oscillator.
As another application of deformed algebra we can refer to the notion of
parastatistics [11]. The concept of parastatistics has found many application
in fractal statistics and anyon theory [12]. In addition to the anyon theory,
the parastatistics has found many interesting application in supersymmetry and
non-commutative quantum mechanics [13].
The construction of generalized deformed oscillators corresponding to well-
known potentials and study of the correspondence between the properties of the
conventional potential picture and the algebraic one has been done [14].
Recently, the generalized deformed algebra and its associated generalized
operators have been considered [15]. By looking at the classical
correspondence of the Hamiltonian, the potential energy and the effective mass
function is obtained. In this contribution we derive the generalized operators
associated with a definite potential by comparing the physical properties of
system and physical results of generalized algebra.
One of the most interesting features of the QHO is the construction of its
coherent states as the eigenfunctions of the annihilation operator. As is well
known [3], one can introduce nonlinear coherent states (NLCSs) or $f$-coherent
states as the right-hand eigenstates of the deformed annihilation operator
$\hat{A}$. It has been shown [16] that these families of generalized coherent
states exhibit various non-classical properties. Due to these properties and
their applications, generation of these states is a very important issue in
the context of quantum optics. The $f$-coherent states may appear as
stationary states of the center-of-mass motion of a trapped ion [17].
Furthermore, a theoretical scheme for generation of these states in a
coherently pumped micromaser within the frame-work of intensity-dependent
Jaynes-Cummings model has been proposed [18].
One of the most important questions is the physical meaning of the deformation
in the NLCSs theory. It has been shown [19] that there is a close connection
between the deformation function appeared in the algebraic structure of NLCSs
and the non-commutative geometry of the configuration space. Furthermore, it
has been shown recently [20], that a two-mode QHO confined on the surface of a
sphere, can be interpreted as a single mode deformed oscillator, whose quantum
statistics depends on the curvature of sphere.
Motivated by the above-mentioned studies, in the present contribution we are
intended to investigate the spatial confinement effects on physical properties
of a standard QHO. It will be shown that the spatial confinement leads to
deformation of standard QHO. We consider a QHO confined in a one-dimensional
infinite well without periodic boundary conditions, and we find its energy
levels, as well as associated ladder operators. We show that the ladder
operators can be interpreted as a special kind of the so-called $f$-deformed
creation and annihilation operators [3].
This paper is organized as follows: In section 2, we review some physical
properties of $f$-oscillator and its coherent states. In section 3 we consider
the spatially confined QHO in a one-dimensional infinite well and construct
its associated coherent states. We shall also examine some of their quantum
statistical properties, including sub-Poissonian statistics and quadrature
squeezing. Finally, we summarize our conclusions in section 4.
## 2 $f$-oscillator and nonlinear coherent states
In this section, we review the basics of the $f$-deformed quantum oscillator
and the associated coherent states known in the literature as nonlinear
coherent states. In the first step, to investigate one of the sources of
deformation we consider an eigenvalue problem for a given quantum physical
system and we focus our attention on the properties of creation and
annihilation operators, that allow to make transition between the states of
discrete spectrum of the system Hamiltonian [21]. As usual, we expand the
Hamiltonian in its eigenvectors
$\hat{H}=\sum_{i=0}^{\infty}E_{i}|i\rangle\langle i|\>,$ (2)
where we have choosed $E_{0}=0$. We introduce the creation (raising) and
annihilation (lowering) operators as follows
$\hat{a}^{{\dagger}}=\sum_{i=0}^{\infty}\sqrt{E_{i+1}}|i+1\rangle\langle
i|,\hskip 56.9055pt\hat{a}=\sum_{i=0}^{\infty}\sqrt{E_{i}}|i-1\rangle\langle
i|\>,$ (3)
so that $\hat{a}|0\rangle=0$. These ladder operators satisfy the following
commutation relation
$[\hat{a},\hat{a}^{{\dagger}}]=\sum_{i=1}^{\infty}(E_{i+1}-E_{i})|i\rangle\langle
i|\>.$ (4)
Obviously if the energy spectrum is equally spaced that is, it should be
linear in quantum numbers, as in the case of ordinary QHO, then
$E_{i+1}-E_{i}=c$, where $c$ is a constant and in this condition the
commutator of $\hat{a}$ and $\hat{a}^{{\dagger}}$ becomes a constant (a
rescaled Weyl-Heisenberg algebra). On the other hand, if the energy spectrum
is not equally spaced, the ladder operators of the system satisfy a deformed
Heisenberg algebra, i.e. their commutator depends on the quantum numbers that
appear in the energy spectrum. This is one of the most important properties of
the quantum $f$-oscillators [3].
An $f$-oscillator is a non-harmonic system characterized by a Hamiltonian of
the harmonic oscillator form
$\hat{H}_{D}=\frac{\Omega}{2}(\hat{A}\hat{A}^{{\dagger}}+\hat{A}^{{\dagger}}\hat{A})\hskip
28.45274pt(\hbar=1)\>,$ (5)
($\hat{A}=\hat{a}f(\hat{n})$) with a specific frequency $\Omega$ and deformed
boson creation and annihilation operators defined in (1). The deformed
operators obey the commutation relation
$[\hat{A}\>,\>\hat{A}^{{\dagger}}]=(\hat{n}+1)f^{2}(\hat{n}+1)-\hat{n}f^{2}(\hat{n})\>.$
(6)
The $f$-deformed Hamiltonian $\hat{H}_{D}$ is diagonal on the eingenstates
$|n\rangle$ in the Fock space and its eigenvalues are
$E_{n}=\frac{\Omega}{2}[(n+1)f^{2}(n+1)+nf^{2}(n)].$ (7)
In the limit $f\rightarrow 1$, the ordinary expression
$E_{n}=\Omega(n+\frac{1}{2})$ and the usual (non-deformed) commutation
relation $[\hat{a}\>,\>\hat{a}^{{\dagger}}]=1$ are recovered.
Furthermore, by using the Heisenberg equation of motion with Hamiltonian (5)
$i\frac{d\hat{A}}{dt}=[\hat{A}\>,\>\hat{H}_{D}],$ (8)
we obtain the following solution for the $f$-deformed operators $\hat{A}$ and
$\hat{A}^{{\dagger}}$
$\hat{A}(t)=e^{-i\Omega G(\hat{n})t}\hat{A}(0),\hskip
28.45274pt\hat{A}^{{\dagger}}(t)=\hat{A}^{{\dagger}}(0)e^{i\Omega
G(\hat{n})t},$ (9)
where
$G(\hat{n})=\frac{1}{2}\left((\hat{n}+2)f^{2}(\hat{n}+2)-\hat{n}f^{2}(\hat{n})\right).$
(10)
In this sense, the $f$-deformed oscillator can be interpreted as a nonlinear
oscillator whose frequency of vibrations depends explicitly on its number of
excitation quanta [4]. It is interesting to point out that recent studies have
revealed strictly physical relationship between the nonlinearity concept
resulting from the $f$-deformation and some nonlinear optical effects, e.g.,
Kerr nonlinearity, in the context of atom-field interaction [22].
The nonlinear transformation of the creation and annihilation operators leads
naturally to the notion of nonlinear coherent states or $f$-coherent states.
The nonlinear coherent states $|\alpha\rangle_{f}$ are defined as the right-
hand eigenstates of the deformed operator
$\hat{A}|\alpha\rangle_{f}=\alpha|\alpha\rangle_{f}\>.$ (11)
From Eq.(11) one can obtain an explicit form of the nonlinear coherent states
in a number state representation
$|\alpha\rangle_{f}=C\sum_{n=0}^{\infty}\alpha^{n}d_{n}|n\rangle,$ (12)
where the coefficients $d_{n}$’s and normalization constant $C$ are,
respectively, given by
$\displaystyle d_{0}$ $\displaystyle=$ $\displaystyle 1\hskip
14.22636pt,\hskip 14.22636ptd_{n}=\left(\sqrt{n!}[f(n)]!\right)^{-1}\hskip
14.22636pt,\hskip 14.22636pt[f(n)]!=\prod_{j=1}^{n}f(j),$ $\displaystyle C$
$\displaystyle=$
$\displaystyle\left(\sum_{n=0}^{\infty}d_{n}^{2}|\alpha|^{2n}\right)^{\frac{-1}{2}}.$
(13)
In recent years the nonlinear coherent states have been paid much attentions
because they exhibit nonclassical features [16] and many quantum optical
states, such as squeezed states, phase states, negative binomial states and
photon-added coherent states can be viewed as a sort of nonlinear coherent
states [23].
## 3 Quantum harmonic oscillator in a one dimensional infinite well
### 3.1 $f$-deformed oscillator description of confined QHO
In this section we consider a quantum harmonic oscillator confined in a one
dimensional infinite well. Many attempts have been done for solving this
problem (see [24],[25], and references therein). In most of those works,
authors tried to solve the problem numerically. But in our consideration we
try to solve the problem analytically, to reveal the relationship between the
confinement effect and given deformation function. We start from the
Schrödinger equation ($\hbar=1$)
$\left[-\frac{1}{2m}\frac{d^{2}}{dx^{2}}+\frac{1}{2}kx^{2}+V(x)\right]\psi(x)=E\psi(x),$
(14)
where
$V(x)=\left\\{\begin{array}[]{ll}0&\textrm{$-a\leq x\leq a$}\\\
\infty&\textrm{elsewhere}.\end{array}\right.$
According to the approach introduced in previous section, we can obtain
raising and lowering operators from the spectrum of Schrödinger operator. On
the other hand, by comparing the energy spectrum of particular system with
energy spectrum of general $f$-deformed oscillator (7), one could obtain
deformed raising and lowering operators. Hence, we need an analytical
expression for energy spectrum of the system which explicitly shows dependence
on special quantum numbers. The original problem, confined QHO (14), can be
solved only by using the approximation methods. When applying perturbation
theory, one is usually concern with a small perturbation of an exactly
solvable Hamiltonian system. In the case of confined QHO we deal with three
limits. Inside the well, for small values of position we have harmonic
oscillator, for large values we have an infinite well and at the positions of
the boundaries the two potentials have the same power. Hence the approximation
method can not lead to acceptable results. Therefore, we model the original
problem by a model potential that has mathematical behavior such as confined
QHO. Instead of solving the Schrödinger equation for the QHO confined between
infinite rectangular walls in positions $\pm a$, we propose to solve the
eigenvalue problem for the potential
$V(x)=\frac{1}{2}k\left(\frac{\tan(\delta x)}{\delta}\right)^{2}\>,$ (15)
where $\delta=\frac{\pi}{2a}$, is a scaling factor depending on the width of
the well. This potential models a QHO placed in the center of the rectangular
infinite well [26]. The potential $V(x)$ (15) fulfills two asymptotic
requirements: 1) $V(x)\rightarrow\frac{1}{2}kx^{2}$ when $a\rightarrow\infty$
(free harmonic oscillator limit). 2) $V(x)$ at equilibrium position has the
same curvature as a free QHO, $\left[\frac{d^{2}V}{dx^{2}}\right]_{x=0}=k$.
This model potential belongs to the exactly solvable trigonometric Pöschl-
Teller potentials family [27]. Stationary coherent states for special kind of
this potential have been considered [28].
Now we consider the following equation
$\left[-\frac{1}{2m}\frac{d^{2}}{dx^{2}}+\frac{1}{2}k\left(\frac{\tan(\delta
x)}{\delta}\right)^{2}-E\right]\psi(x)=0\;.$ (16)
To solve analytically this equation, we use the factorization method [29]. By
changing the variable and some mathematical manipulation, the corresponding
energy eigenvalues are found as
$E_{n}=\gamma(n+\frac{1}{2})^{2}+\sqrt{\gamma^{2}+\omega^{2}}(n+\frac{1}{2})+\frac{\gamma}{4}\>,$
(17)
where $\gamma=\frac{4\pi^{2}}{32a^{2}m}$,and $\omega=\sqrt{\frac{k}{m}}$ is
the frequency of the QHO. The first term in the energy spectrum can be
interpreted as the energy of a free particle in a well, the second term
denotes the energy spectrum of the QHO, and the last term shifts the energy
spectrum by a constant amount. It is evident that if $a\rightarrow\infty$ then
$\gamma\rightarrow 0$ and the energy spectrum (17) reduces to the spectrum of
a free QHO. As is clear from (17), different energy levels are not equally
spaced. Hence, confining a free QHO leads to deformation of its dynamical
algebra and we can interpret the parameter $\gamma$ as the corresponding
deformation parameter. In Table 1 the numerical results associated with the
original potential, given in Ref. [24], are compared with the generated
results from the model potential under consideration. As is seen, the results
are in a good agreement when boundary size is of order of characteristic
length of the harmonic oscillator. The original oscillator potential when
approaches to the boundaries of the well becomes infinite suddenly, while the
model potential is smooth and approaches to the infinity asymptotically.
Therefore, the model potential (15) is more appropriate for the physical
systems.
If we normalize Eq.(17) to energy quanta of the simple harmonic oscillator and
introduce the new variables $n+\frac{1}{2}=h$,
$\sqrt{\frac{\gamma^{2}}{\omega^{2}}+1}=\eta$, and
$\gamma^{\prime}=\frac{\gamma}{\omega}$ then it takes the following form
$E_{l}=\gamma^{\prime}h^{2}+\eta h+\frac{\gamma^{\prime}}{4}.$ (18)
By comparing this spectrum with the energy spectrum of an $f$-deformed
oscillator, given by (7), we find the corresponding deformation function as
$f(\hat{n})=\sqrt{\gamma^{\prime}\hat{n}+\eta}.$ (19)
Furthermore, the ladder operators associated with the confined oscillator
under consideration can be written in terms of the conventional (non-deformed)
operators $\hat{a}$ , $\hat{a}^{{\dagger}}$ as follows
$\hat{A}=\hat{a}\sqrt{\gamma^{\prime}\hat{n}+\eta}\hskip 28.45274pt,\hskip
28.45274pt\hat{A}^{{\dagger}}=\sqrt{\gamma^{\prime}\hat{n}+\eta}\,\,\hat{a}^{{\dagger}}.$
(20)
These two operators satisfy the following commutation relation
$[\hat{A},\hat{A}^{{\dagger}}]=\gamma^{\prime}(2\hat{n}+1)+\eta.$ (21)
It is obvious that in the limiting case $a\rightarrow\infty$
($\gamma^{\prime}\rightarrow 0$,$\eta\rightarrow 1$), the right hand side of
the above commutation relation becomes independent of $\hat{n}$, and the
deformed algebra reduces to a the conventional Weyl-Heisenberg algebra for a
free QHO.
Classically, harmonic oscillator is a particle that attached to an ideal
spring, and can oscillate with specific amplitude. When that particle be
confined, boundaries can affect particle’s motion if the boundaries position
be in a smaller distance in comparison with a characteristic length that
particle oscillates within it. This characteristic length for the QHO is given
by $\frac{\hbar}{m\omega}$ $(\hbar=1)$ , and if $2a\leq\frac{1}{m\omega}$,
then the presence of the boundaries affects the behavior of QHO, otherwise it
behaves like a free QHO. Therefore, one can interpret
$l_{0}=\frac{1}{m\omega}$ as a scale length where the deformation effects
become relevant.
### 3.2 Coherent states of confined oscillator
Now, we focus our attention on the coherent states associated with the QHO
under consideration. As usual, we define coherent states as the right-hand
eigenstates of the deformed annihilation operator
$\hat{A}|\beta\rangle_{f}=\beta|\beta\rangle_{f}.$ (22)
From (22) and using the NLCS formalism introduced in (11)-(2) the explicit
form of the corresponding NLCS of the confined QHO is written as
$|\beta\rangle_{f}=\mathcal{N}\sum_{n}\frac{\beta^{n}}{\sqrt{n!(\gamma^{\prime}n+\eta)!}}|n\rangle,$
(23)
where
$\mathcal{N}=\left(\sum_{n}\frac{|\beta|^{2n}}{[f(n)!]^{2}n!}\right)^{-\frac{1}{2}}$
is the normalization factor, $\beta$ is a complex number, and the deformation
function $f(n)$ is given by Eq.(19). The ensemble of states
$|\beta\rangle_{f}$ labeled by the single complex number $\beta$ is called a
set of coherent states if the following conditions are satisfied [30]:
* •
normalizability
$_{f}\langle\beta|\beta\rangle_{f}=1,$ (24)
* •
continuity in the label $\beta$
$|\beta-\beta^{\prime}|\rightarrow 0\hskip 14.22636pt\Rightarrow\hskip
14.22636pt\|\;|\beta\rangle_{f}-|\beta^{\prime}\rangle_{f}\|\rightarrow 0,$
(25)
* •
resolution of the identity
$\int_{c}d^{2}\beta|\beta\rangle_{f}{}_{f}\langle\beta|w(|\beta|^{2})=\hat{I},$
(26)
where $w(|\beta|^{2})$ is a proper measure that ensures the completeness and
the integration is restricted to the part of the complex plane where
normalization converges.
The first two conditions can be proved easily. For the third condition, we
choose the normalization constant as
$\mathcal{N}^{2}=\frac{|\beta|^{\eta}}{I_{\eta}^{\gamma^{\prime}}(2|\beta|)},$
(27)
where
$I_{\eta}^{\gamma^{\prime}}(x)=\sum_{s=0}^{\infty}\frac{1}{s!(\gamma^{\prime}s+\eta)!}(\frac{x}{2})^{2s+\eta},$
(28)
is similar to the modified Bessel function of the first kind of the order
$\eta$ with the series expansion
$I_{\eta}(x)=\sum_{s=0}^{\infty}\frac{1}{s!(s+\eta)!}(\frac{x}{2})^{2s+\eta}$.
Resolution of the identity of the deformed coherent states $|\beta\rangle_{f}$
can be written as
$\displaystyle\int d^{2}\beta|\beta\rangle_{f}\langle\beta|w(|\beta|)=$
$\displaystyle\pi\sum_{n}|n\rangle\langle
n|\frac{1}{n!(\gamma^{\prime}n+\eta)!}\int_{0}^{\infty}d|\beta||\beta||\beta|^{2n}$
$\displaystyle\times\frac{|\beta|^{\eta}}{I_{\eta}^{\gamma^{\prime}}(2|\beta|)}w(|\beta|).$
Now we introduce the new variable $|\beta|^{2}=x$ and the measure
$w(\sqrt{x})=\frac{8}{\pi}I_{\eta}^{\gamma^{\prime}}(2\sqrt{x})K_{m}(2\sqrt{x})x^{\frac{l}{2}},$
(30)
where $K_{m}(x)$ is the modified Bessel function of the second kind of the
order $m$, $m=(\gamma^{\prime}-1)n+\alpha$, and $l=(\gamma^{\prime}-1)n+1$.
Using the integral relation
$\int_{0}^{\infty}K_{\nu}(t)t^{\mu-1}dt=2^{\mu-2}\Gamma\left(\frac{\mu-\nu}{2}\right)\Gamma\left(\frac{\mu+\nu}{2}\right)$
[31], we obtain
$\int
d^{2}\beta|\beta\rangle_{f}{}_{f}\langle\beta|w(|\beta|)=\sum_{n}|n\rangle\langle
n|=\hat{I}.$ (31)
We therefore conclude that the states $|\beta\rangle_{f}$ qualify as coherent
states in the sense described by the conditions (24)-(26).
We now proceed to examine some nonclassical properties of the nonlinear
coherent states $|\beta\rangle_{f}$. As an important quantity, we consider the
variance of the number operator $\hat{n}$. Since for the coherent states the
variance of number operator is equal to its average, deviation from the
Poissonian statistics can be measured with the Mandel parameter [32]
$M=\frac{(\Delta n)^{2}-\langle\hat{n}\rangle}{\langle\hat{n}\rangle}.$ (32)
This parameter vanishes for the Poisson distribution, is positive for the
super-Poissonian distribution (bunching effect), and is negative for the sub-
Poissonian distribution (antibunchig effect). Fig. 1 shows the size dependence
of the Mandel parameter for different values of $|\beta|^{2}$. As is seen, the
Mandel parameter exhibits the sub-Poissonian statistics and with further
increasing values of $a$ it is finally stabilized at an asymptotical zero
value corresponding to the Poissonian statistics. In addition, the smaller the
parameter $|\beta|^{2}$ is, the more rapidly the Mandel parameter tends to the
Poissonian statistics.
As another important nonclassical property we examine the quadrature
squeezing. For this purpose we first consider the conventional quadrature
operators $\hat{X}_{a}$ and $\hat{Y}_{a}$ defined in terms of nondeformed
operators $\hat{a}$ and $\hat{a}^{\dagger}$ as [33]
$\hat{X}_{a}=\frac{1}{2}(\hat{a}e^{i\phi}+\hat{a}^{{\dagger}}e^{-i\phi})\hskip
28.45274pt\hat{Y}_{a}=\frac{1}{2i}(\hat{a}e^{i\phi}-\hat{a}^{{\dagger}}e^{-i\phi}).$
(33)
In this equation, $\phi$ is the phase of quadrature operators which can
effectivly affect the squeezing properties. The commutation relation for
$\hat{a}$ and $\hat{a}^{{\dagger}}$ leads to the following uncertainty
relation
$(\Delta\hat{X}_{a})^{2}(\Delta\hat{Y}_{a})^{2}\geq\frac{1}{4}|\langle[\hat{X}_{a},\hat{Y}_{a}]\rangle|^{2}=\frac{1}{16}.$
(34)
For the vacuum state $|0\rangle$, we have
$(\Delta\hat{X}_{a})^{2}=(\Delta\hat{Y}_{a})^{2}=\frac{1}{4}$ and hence
$(\Delta\hat{X}_{a})^{2}(\Delta\hat{Y}_{a})^{2}=\frac{1}{16}$. A given quantum
state of the QHO is said to be squeezed when the variance of one of the
quadrature components $\hat{X}_{a}$ and $\hat{Y}_{a}$ satisfies the relation
$(\Delta\hat{O}_{a})^{2}<(\Delta\hat{O}_{a})^{2}_{vacuum}=\frac{1}{4}\hskip
14.22636pt(\hat{O}_{a}=\hat{X}_{a}\hskip 8.5359ptor\hskip
8.5359pt\hat{Y}_{a}).$ (35)
The degree of quadrature squeezing can be measured by the squeezing parameter
$s_{\hat{O}}$ defined by
$s_{\hat{O}}=4(\Delta\hat{O}_{a})^{2}-1.$ (36)
Then, the condition for squeezing in the quadrature component can be simply
written as $s_{\hat{O}}<0$. In Fig. 2 we have plotted the parameter
$s_{\hat{X}_{a}}$ corresponding to the squeezing of $\hat{X}_{a}$ with respect
to the phase angle $\phi$ for three different values of $a$. As is seen, the
state $|\beta\rangle_{f}$ exhibits squeezing for different values of the
confinement size, and when $a_{l}=\frac{a}{l_{0}}=2.5$, the quadrature
$\hat{X}_{a}$ exhibits squeezing for all values of the phase angle $\phi$.
Fig. 3 shows the plot of $s_{\hat{X}_{a}}$ versus the dimensionless parameter
$a_{l}=\frac{a}{l_{0}}$ for different values of the phase $\phi$. As is seen,
with the increasing value of $a_{l}\;(\frac{a}{l_{0}})$, the quadrature
component tends to the zero according to the vacuum fluctuation. Let us also
consider the deformed quadrature operators $\hat{X}_{A}$ and $\hat{Y}_{A}$
defined in terms of the deformed operators $\hat{A}$ and $\hat{A}^{{\dagger}}$
as
$\hat{X}_{A}=\frac{1}{2}(\hat{A}e^{i\phi}+\hat{A}^{{\dagger}}e^{-i\phi}),\hskip
28.45274pt\hat{Y}_{A}=\frac{1}{2i}(\hat{A}e^{i\phi}-\hat{A}^{{\dagger}}e^{-i\phi}).$
(37)
By considering the commutation relation (6) for the deformed operators
$\hat{A}$ and $\hat{A}^{{\dagger}}$, the squeezing condition for the deformed
quadrature operators $\hat{O}_{A}$ ($=\hat{X}_{A}$, $\hat{Y}_{A}$)can be
written as
$S=4(\Delta\hat{O}_{A})^{2}-\langle(\hat{n}+1)f^{2}(\hat{n}+1)\rangle+\langle\hat{n}f^{2}(\hat{n})\rangle<0.$
(38)
In Fig. 4 we have plotted the parameter $S_{\hat{X}_{A}}$ versus the
dimensionless parameter $\frac{a}{l_{0}}$ for three different values of
$|\beta|^{2}$. As is seen, the deformed quadrature operator exhibits squeezing
for all values of $a$. Furthermore, with the increasing value of $|\beta|^{2}$
the squeezing of the quadrature $\hat{X}_{A}$ is enhanced.
## 4 Conclusion
In this paper, we have considered the relation between the spatial confinement
effects and a special kind of $f$-deformed algebra. We have found that the
confined simple harmonic oscillator can be interpreted as an $f$-oscillator,
and we have obtained the corresponding deformation function. By constructing
the associated NLCSs, we have examined the effects of confinement size on non-
classical statistical properties of those states. The result show that the
stronger confinement leads to the strengthening of non-classical properties.
We hope that our approach may be used in description of phonons in the strong
excitation regimes, photons in a microcavity and different elementary
excitations in confined systems. The work on this direction is in progress.
Acknowledgment The authors wish to thank the Office of Graduate Studies of the
University of Isfahan and Iranian Nanotechnology initiative for their support.
Table 1: Calculated energy levels of the confined QHO in a one dimensional infinite well by using our model potential in comparison with the numerical results given in Ref.[24] state | boundary size | model potential | numerical results
---|---|---|---
0 | a=0.5 | 4.98495312 | 4.95112932
0 | 1 | 1.41089325 | 1.29845983
0 | 2 | 0.67745392 | 0.53746120
0 | 3 | 0.57321464 | 0.50039108
0 | 4 | 0.54003728 | 0.50000049
1 | a=0.5 | 19.88966157 | 19.77453417
1 | 1 | 5.46638033 | 5.07558201
1 | 2 | 2.34078691 | 1.76481643
1 | 3 | 1.85672176 | 1.50608152
1 | 4 | 1.69721813 | 1.50001461
2 | a=0.5 | 44.66397441 | 44.45207382
2 | 1 | 11.98926850 | 11.25882578
2 | 2 | 4.62097017 | 3.39978824
2 | 3 | 3.41438455 | 2.54112725
2 | 4 | 3.00861155 | 2.50020117
3 | a=0.5 | 79.30789166 | 78.99692115
3 | 1 | 20.97955777 | 19.89969649
3 | 2 | 7.51800371 | 5.58463907
3 | 3 | 5.24620303 | 3.66421964
3 | 4 | 4.47421754 | 3.50169153
4 | a=0.5 | 123.82141330 | 123.41071050
4 | 1 | 32.43724814 | 31.00525450
4 | 2 | 11.03188752 | 8.36887442
4 | 3 | 7.35217718 | 4.95418047
4 | 4 | 6.09403610 | 4.50964099
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Figure 1: Plots of the Mandel parameter versus the dimensionless parameter
$a_{l}=\frac{a}{l_{0}}$. For $|\beta|^{2}=0.5$ (dashed curve), for
$|\beta|^{2}=1$ (longdashed curve), for $|\beta|^{2}=1.5$ (solid curve) and
for $|\beta|^{2}=4.0$ (bold curve). Figure 2: Plot of $s_{\hat{X}_{a}}$ versus
$\phi$ for $|\beta|^{2}=4$. The dashed, longdashed and solid curves
respectively relate to $a=2.5$, $a=1$, $a=0.5$ (the values of $a$ are
renormalized to $l_{0}$). Figure 3: Plots of $s_{\hat{X}_{a}}$ versus the
dimensionless parameter $a_{l}=\frac{a}{l_{0}}$ for different phases and
$|\beta|^{2}=1$. Dashed curve, solid curve and bold curve ,respectively,
correspond to $\phi=100$, $\phi=110$ and $\phi=90$. Figure 4: Plots of
deformed squeezing parameter $S_{X_{A}}$ versus the dimensionless parameter
$a_{l}=\frac{a}{l_{0}}$. The dashed curve, longdashed curv, solid curve and
bold curve are respectively, correspond to $|\beta|^{2}=1$, $|\beta|^{2}=1.5$
$|\beta|^{2}=2.5$ and $|\beta|^{2}=4$.
|
arxiv-papers
| 2011-12-11T10:52:14 |
2024-09-04T02:49:25.163343
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi",
"submitter": "Malek Bagheri",
"url": "https://arxiv.org/abs/1112.2343"
}
|
1112.2344
|
# Nonlinear coherent state of an exciton in a wide quantum dot
M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi Quantum Optics Group, Physics
Department, University of Isfahan m-baghreri@phys.ui.ac.ir rokni@sci.ui.ac.ir
mhnaderi@phys.ui.ac.ir , ,
###### Abstract
In this paper, we derive the dynamical algebra of a particle confined in an
infinite spherical well by using the $f$-deformed oscillator approach. We
consider an exciton with definite angular momentum in a wide quantum dot
interacting with two laser beams. We show that under the weak confinement
condition, and quantization of the center-of-mass motion of exciton, the
stationary state of it can be considered as a special kind of nonlinear
coherent states which exhibits the quadrature squeezing.
## 1 Introduction
The conventional coherent states of the quantum harmonic oscillator, defined
by Glauber [1] as the right-hand eigenstates of non-hermitian annihilation
operator $\hat{a}([\hat{a},\hat{a}^{{\dagger}}]=1)$, have found many
interesting applications in different areas of physics such as quantum optics,
condensed matter physics, statistical physics and atomic physics [2]. These
states play an important role in the quantum theory of coherence, are
considered as the most classical ones among the pure quantum states, and laser
light can be supposed as a physical realization of them. Due to the vast
application of these states, there have been many attempts to generalize them
[3]. Among the all generalizations, nonlinear coherent states (NLCS) [4] have
been paid attention in recent years because they exhibit nonclassical features
such as quadrature squeezing and sub-poissonian statistics [5]. These states
are defined as the right-hand eigenstates of a deformed operator $\hat{A}$
$\hat{A}=\hat{a}f(\hat{n})\hskip
56.9055pt\hat{A}|\alpha,f\rangle=\alpha|\alpha,f\rangle,$ (1)
where the deformation function $f(\hat{n})$ is an operator-valued function of
the number operator $\hat{n}$. From (1) one can obtain an explicit form of
NLCS in the number state representation
$\displaystyle|\alpha,f\rangle=\mathcal{N}_{f}\sum_{n}\frac{\alpha^{n}}{\sqrt{n!}f(n)!}|n\rangle,$
$\displaystyle\mathcal{N}_{f}=\left[\sum_{n}\frac{|\alpha|^{2n}}{[f(n)!]^{2}n!}\right]^{-\frac{1}{2}}.$
(2)
A class of NLCS can be realized physically as the stationary state of the
center-of-mass motion of a laser driven trapped ion [6, 7]. Furthermore, it
has been proposed a theoretical scheme to show the possibility of generating
various families of NLCS [8] of the radiation field in a lossless coherently
pumped micromaser within the frame work of the intensity-dependent Jaynes-
Cummings model.
Recently, the influences of the spatial confinement [9] and the curvature of
physical space [10] on the algebraic structure of the coherent states of the
quantum harmonic oscillator have been investigated within the frame work of
nonlinear coherent states approach. It has been shown that if a quantum
harmonic oscillator be confined within a small region of order of its
characteristic length [9] or its physical space to be a sphere [10], then it
can be regarded as a deformed oscillator, i.e., an oscillator that its
creation and annihilation operators are deformed operator $\hat{A}$ and
$\hat{A}^{{\dagger}}$ given by Eq.(1).
On the other hand, we can consider nanostructures as systems whose physical
properties are related to the confinement effects. Thus, we expect that it is
possible to realize some natural deformations in these systems [9, 11]. In
addition, in nanostructures different kinds of quantum states can be prepared.
One of the most applicable of these states is exciton state. Exciton is an
elementary excitation in semiconductors interacting with light, electron in
conduction band which is bounded to hole in valance band that can easily move
through the sample. In one of the nano size systems, quantum dot (QD), due to
the confinement in three dimensions, energy bands reduce to quasi energy
levels. Therefore, in order to describe the interaction of QD with light we
can consider it as a few level atom [12]. These Exciton states can be used in
quantum information processes. It has been shown that excitons in coupled QDs
are ideal for preparation of entangled state in solid-state systems [13].
Entanglement of the exciton states in a single QD or in a QD molecule has been
demonstrated experimentally [14]. Entanglement of the coherent states of the
excitons in a system of two coupled QDs has been considered [15]. Recently,
coherent exciton states of excitonic nano-crystal-molecules has been
considered [16]. Theoretical approach for generating Dick states of excitons
in optically driven QD has been proposed in Ref.[17]. In a QD, the effects of
exciton-phonon interaction, exciton-impurity interaction and exciton-exciton
interaction play an important role. These effects are the main sources for the
decoherence of exciton states [18]. Furthermore, these effects cause the
exciton has the spontaneous recombination or scattered to other exciton modes
[19, 20].
In this paper we propose a theoretical scheme for generating excitonic NLCS.
We will show that under certain conditions the quantized motion of wave packet
of center-of-mass of exciton can be consider as a special kind of NLCSs. Our
scheme is based on the interaction of a quantum dot with two laser beams. By
using the approach considered in Ref.[6], we propose a theoretical scheme for
generation of NLCS of an exciton in a wide QD.
In section 2, we consider different confinement regimes in a QD, and the
explicit forms of the creation and annihilation operators for a particle
confined in an infinite well are derived by using the deformed quantum
oscillator approach. In section 3, we consider an exciton in a wide QD which
interacts with two laser beams. We shall show that under the weak confinement
condition, the stationary state of the exciton center-of-mass motion can be
considered as a NLCS.
## 2 Algebraic approach for a particle in an infinite spherical well
In nanostructures and confined systems, there are three different confinement
regimes. The criteria for this classification is based on the comparison
between excitation Bohr radius and the spatial dimensions of the system under
consideration. In the case of a QD, these regimes are defined as follows [21].
We first introduce three quantities $\Delta E_{c}$, $\Delta E_{v}$ and
$V_{exc}$ which, respectively denote: the electron energy due to the
confinement, the hole energy due to the confinement and Coulomb energy between
correlated electron-hole (exciton).
1) $V_{exc}>\Delta E_{c}-\Delta E_{v}$: In this case, the exciton energy is
much greater than the confinement energies of electron and hole. If we show
the system size by $L$ and the exciton Bohr radius by $a$, then in this regime
$L>a$. This regime corresponds to the weak confinement (in some literature the
weak confinement is characterized by the situation in which the electron and
the hole are not in the same matter, for example, hole be in QD and excited
electron in host matter. In this paper, by the weak confinement regime we mean
$L>a$ and the excitations in the same matter). In this regime due to the
confinement, the center-of-mass motion of the exciton is quantized and the
confinement do not affect electron and hole separately. Hence, the confinement
affect the exciton motion as a whole [22].
2) $V_{exc}<\Delta E_{c}\;,\Delta E_{v}$: This regime, in contrast to the
previous one, is associated with the cases where $L<a$. In this regime the
exciton is completely localized, and the confinement affects both the electron
and the hole independently and their states become quantized in conduction and
valance bands. This regime is called strong confinement.
3) $\Delta E_{c}>V_{exc}\;,\Delta E_{v}$: This condition is equivalent to the
situation $a_{c}<a<a_{v}$, where $a_{c}$ and $a_{v}$ are, respectively, the
Bohr radii of electron and hole. Here, due to the different effective masses
of electron and hole, the hole which has heavier effective mass is localized
and the electron motion will be quantized. This regime is called intermediate
confinement.
In the first case (weak confinement), in a wide QD, an exciton can move due to
its center-of-mass momentum, and because of the presence of the barriers, its
center-of-mass motion is quantized. Therefore, it moves as a whole between
energy levels of an infinite well. We consider a wide spherical QD whose
energy levels are equivalent to the energy levels of a spherical well
$E_{nl}=\frac{\hbar^{2}}{2M}\frac{\alpha_{nl}^{2}}{R^{2}},$ (3)
where $\alpha_{nl}$ is the n’th zero of the first kind Bessel function of
order $l$, $j_{l}(x)$. In this energy spectrum according to the azimuthal
symmetry around $z$ axis, we have a degenerate spectrum. As mentioned before,
in the weak confinement regime, the Coulomb potential plays an essential role
and its spectrum is given by
$E_{k}^{b}=\frac{\mu e^{4}}{2\hbar^{2}\varepsilon^{2}}\frac{1}{k^{2}}\hskip
14.22636pt,\hskip 14.22636pt\mu=\frac{m_{e}m_{h}}{m_{e}+m_{h}},$ (4)
where superscript $b$ shows binding energy related to the Coulomb interaction
and $\varepsilon$ shows dielectric constant of the system. As is usual, we
interpret the Coulomb part as an exciton and another degree of freedom (motion
between energy levels of the well) as the exciton center-of-mass motion.
Therefore, in a wide QD an exciton has two different kinds of degrees of
freedom: internal degrees of freedom due to the Coulomb potential and external
degrees of freedom related to the quantum confinement. Here we consider the
lowest exciton state, $1s$ exciton, because this exciton state has the largest
oscillator strength among other exciton state. Then the energy of the exciton
in a wide QD can be written as
$E_{nlm}=E_{g}-E_{1}^{b}+\frac{\hbar^{2}}{2MR^{2}}\alpha_{nl}^{2},$ (5)
where $E_{g}$ is the energy gap of QD, $E_{1}^{b}=E_{k}^{b}|_{k=1}$ is the
exciton binding energy, $M=m_{e}+m_{h}$ is the total mass of exciton, and $R$
is the radius of QD. Due to the relation of quantum numbers $l$ and $m$ with
the angular momentum and the selection rules for optical transitions, we can
fix $l$ and $m$ (by choosing a certain condition), and hence the energy of
exciton depends only on a single quantum number
$E_{n}=E_{g}-E_{1}^{b}+\frac{\hbar^{2}}{2MR^{2}}\alpha_{nl}^{2}.$ (6)
Therefore, we can prepare the conditions under which the exciton center-of-
mass motion has a one-dimensional degree of freedom. Due to the quantization
of the exciton center-of-mass motion, we can describe the exciton motion
between the energy levels by the action of a special kind of ladder operators.
In order to find these operators we use the $f$-deformed oscillator approach
[4].
As mentioned elsewhere [9], if the energy spectrum of the system is equally
spaced, such as harmonic oscillator, its creation and annihilation operators
satisfy the ordinary Weyl-Heisenberg algebra, otherwise we can interpret them
as the generators of a generalized Weyl-Heisenberg algebra.
The energy spectrum of a particle with mass $M$ confined in an infinite
spherical well can be written as (3). According to the conservation of angular
momentum, we assume that particle has been prepared with definite angular
momentum (for example by measuring its angular momentum). Then $l$ becomes
completely determined, i.e., in the energy spectrum the number $l$ is a
constant. By determining the number $l$ and considering the rotational
symmetry of the system around the $z$ axis, the angular part of the spectrum
becomes completely determined, and the radius part is described by (3). Now we
use a factorization method and write the Hamiltonian of the center-of-mass
motion of the system as follows
$\hat{H}=\frac{1}{2}(\hat{A}\hat{A}^{{\dagger}}+\hat{A}^{{\dagger}}\hat{A}),$
(7)
where $\hat{A}$ and $\hat{A}^{{\dagger}}$ are defined through the relation
(1). Therefore the spectrum of $\hat{H}$, after straightforward calculation,
is obtained as
$E_{n}=\frac{1}{2}[(n+1)f^{2}(n+1)+nf^{2}(n)].$ (8)
By comparing (8) with Eq.(3) we arrive at the following expression for the
corresponding deformation function $f_{1}(\hat{n})$
$f_{1}(n)=\sqrt{\frac{\hbar^{2}}{MR^{2}}\frac{(-1)^{n}}{n}\sum_{i=1}^{n}(-1)^{i}\alpha_{i-1l}^{2}}.$
(9)
Then, the ladder operators associated with the radial motion of a confined
particle in a spherical infinite well is given by
$\displaystyle\hat{A}=\hat{a}\sqrt{\frac{\hbar^{2}}{MR^{2}}\frac{(-1)^{n}}{n}\sum_{i=1}^{n}(-1)^{i}\alpha_{i-1l}^{2}},$
(10)
$\displaystyle\hat{A}^{{\dagger}}=\sqrt{\frac{\hbar^{2}}{MR^{2}}\frac{(-1)^{n}}{n}\sum_{i=1}^{n}(-1)^{i}\alpha_{i-1l}^{2}}\,\,\hat{a}^{{\dagger}}.$
These two deformed operators obey the following commutation relation
$[\,\hat{A}\,,\,\hat{A}^{{\dagger}}]=-nf_{1}^{2}(n)+\frac{\hbar^{2}}{MR^{2}}\alpha_{nl}^{2}.$
(11)
As is usual in the $f$-deformation approach, for a particular limit of the
corresponding deformation parameter, the deformed algebra should be reduced to
the conventional oscillator algebra. However, in this treatment we note that
there is no thing in common between the harmonic oscillator potential and an
infinite spherical well. Only in the limit $R\rightarrow\infty$, the system
reduces to a free particle which has continuous spectrum.
As a result, in this section we conclude that the radial motion of a particle
confined in a three-dimensional infinite spherical well can be interpreted by
an $f$-deformed Weyl-Heisenberg algebra.
## 3 Exciton dynamics in QD
Now we consider the formation of an exciton and its dynamics in a wide QD
during the exciton lifetime. As mentioned before, in this situation the
center-of-mass motion of the exciton is quantized. The exciton is created
during the interaction of a QD with light, and because of the angular momentum
conservation, the exciton has a well-defined angular momentum. The exciton is
a quasiparticle composed of an electron and a hole and thus the exciton spin
state can be in a singlet state or a triplet state. According to the optical
transition selection rules, the triplet state is optically inactive and is
called dark exciton [23]. By adding spin and angular momentum of absorbed
photons, the angular momentum of the exciton state can be determined. Hence,
the exciton behaves like a particle in a spherical well with the definite
angular momentum. According to the previous section, the center-of-mass motion
of the exciton in the QD and the barriers of QD can be described by an
oscillator-like Hamiltonian expressed in terms of the $f$-deformed
annihilation and creation operators given by Eq.(10)
$H_{well}=\frac{1}{2}(\hat{A}\hat{A}^{{\dagger}}+\hat{A}^{{\dagger}}\hat{A}),$
(12)
where we interpret the operator $\hat{A}$ $(\hat{A}^{{\dagger}})$ as the
operator whose action causes the transition of exciton center-of-mass motion
to a lower (an upper) energy state. In fact the Hamiltonian (12) is related to
the external degree of freedom of exciton. On the other hand, one can imagine
QD as a two-level system with the ground state $|g\rangle$ and the excited
state $|e\rangle$ (associated with the presence of exciton). Thus, for the
internal degree of freedom we can consider the following Hamiltonian
$H_{ex}=\hbar\omega_{ex}\hat{S}_{22},$ (13)
where $\hat{S}_{22}=|e\rangle\langle e|-|g\rangle\langle g|$ and
$\hbar\omega_{ex}=E_{g}-E_{1}^{b}$ is the exciton energy.
We consider a single exciton of frequency $\omega_{ex}$ confined in a wide QD
interacting with two laser fields, respectively, tuned to the internal degree
of freedom of the frequency $\omega_{ex}$ and to the non-equal spaced energy
levels of the infinite well. It is necessary that the second laser has special
conditions, because it should give rise to the transitions between energy
levels whose frequencies depend on intensity. The interacting system can be
described by the Hamiltonian
$\hat{H}=\hat{H}_{0}+\hat{H}_{int},$ (14)
where $\hat{H}_{0}=\hat{H}_{well}+\hat{H}_{ex}$ and
$H_{int}=g[E_{0}e^{-i(k_{0}\hat{x}-\omega_{ex}t)}+E_{1}e^{-i(k_{1}\hat{x}-(\omega_{ex}-\omega_{\overline{n}})t)}]\hat{S}_{12}+H.c.,$
(15)
in which $g$ is the coupling constant, $k_{0}$ and $k_{1}$ are the wave
vectors of the laser fields, $\hat{S}_{12}=|g\rangle\langle e|$ is the exciton
annihilation operator, and $\omega_{\overline{n}}$ is the frequency of exciton
transition between energy levels of QD due to the spatial confinement. Here,
we consider transition between specific side-band levels hence, we show the
frequency transition with definite dependence to $n$. We show this by a
c-number quantity $\overline{n}$.
The exciton has a finite lifetime that in systems with small dimension, is
increased [24]. The interaction with phonons is the main reason of damping of
the exciton [25]. Phonons in bulk matter have a continuous spectrum while in a
confined system such as QD their spectrum becomes discrete. Hence in a QD, the
resonant interaction between the exciton and phonons decreases and in this
system the exciton lifetime will increase. Therefore during the lifetime of an
exciton, its dynamics is under influence of a bath reservoir, and its damping
play an important role. We assume that during the presence of the exciton in
QD, it interacts with the reservoir and hence we can consider its steady
state. We consider an exciton in dark state. Experimental preparation methods
of such exciton has been described in [23]. In this situation lifetime of
exciton will increase and exciton has not spontaneously recombination
radiation. However, its interaction with phonons causes a finite lifetime for
it.
The operator of the center-of-mass motion position $\hat{x}$ of the exciton in
a spherical QD may be defined as
$\hat{x}=\frac{\kappa}{k_{ex}}(\hat{A}+\hat{A}^{{\dagger}}),$ (16)
where $\kappa$ being a parameter similar to the Lamb-Dick parameter in ion
trapped systems and is defined as the ratio of QD radius to the wavelength of
the driving laser (because of the spatial confinement of exciton, its wave
function width is determined by the barriers of QD), and we assume
$k_{0}\simeq k_{1}\simeq k_{ex}$ ($k_{ex}$ is the wavevector of the exciton).
The operators $\hat{A}$ and $\hat{A}^{{\dagger}}$ are defined in Eq.(10). The
interaction Hamiltonian (15) can be written as
$H_{int}=\hbar
e^{i\omega_{ex}t}\Omega_{1}\left[\frac{\Omega_{0}}{\Omega_{1}}+e^{-i\omega_{\overline{n}}t}\right]e^{i\kappa(\hat{A}+\hat{A}^{{\dagger}})}\hat{S}_{12}+H.c.,$
(17)
where $\Omega_{0}=\frac{gE_{0}}{\hbar}$ and $\Omega_{1}=\frac{gE_{1}}{\hbar}$
are the Rabi frequencies of the lasers, respectively, tuned to the electronic
transition of QD (internal degree of freedom) and the first center-of-mass
motion transition of exciton. Since the external degree of freedom is
definite, then $\omega_{\overline{n}}$ depends on a special value of $n$ such
that it can be consider as a c-number quantity. The frequency
$\omega_{\overline{n}}$ is depend on the number of quanta for each transition
and hence the laser tuned to the center-of-mass motion must be so strong that
causes transition. This allows us to treat the interaction of the confined
exciton in a wide QD with two lasers separately, by using a nonlinear Jaynes-
Cummings Hamiltonian [26] for each coupling. The interaction Hamiltonian in
the interaction picture can be written as
$H_{I}=\hbar\Omega_{1}\hat{S}_{12}\left[\frac{\Omega_{0}}{\Omega_{1}}+e^{i\omega_{\overline{n}}t}\right]\exp[i\kappa(e^{-i\omega_{\hat{n}}t}\hat{A}+\hat{A}^{{\dagger}}e^{i\omega_{\hat{n}}t})]+H.c.,$
(18)
where
$\omega_{\hat{n}}=\frac{1}{2\hbar}[(\hat{n}+2)f_{1}(\hat{n}+2)-\hat{n}f_{1}(\hat{n})]$.
By using the vibrational rotating wave approximation [6], applying the
disentangling formula introduced in [27], and using the fact that in the
present case the Lamb-Dick parameter is small, the interaction Hamiltonian
(18) is simplified to
$H_{I}^{(1)}=\hbar\Omega_{1}\hat{S}_{12}\left[F_{0}(\hat{n},\kappa)\frac{\Omega_{0}}{\Omega_{1}}+i\kappa
F_{1}(\hat{n},\kappa)\hat{a}\right]+H.c.,$ (19)
where the function $F_{i}(\hat{n},\kappa)\;(i=0,1)$ is defined by
$\displaystyle F_{i}(\hat{n},\kappa)$ $\displaystyle=$ $\displaystyle
e^{-\frac{\kappa^{2}}{2}((n+1+i)f_{1}^{2}(n+1+i)-(n+i)f_{1}^{2}(n+i))}\times$
$\displaystyle\sum_{l=0}^{n}\frac{\left(i\kappa\right)^{2l}}{l!(l+i)!}\frac{f_{1}(\hat{n})f_{1}(\hat{n}+i)}{[f_{1}(\hat{n}-l)!]^{2}}(\hat{a}^{{\dagger}})^{l}\hat{a}^{l}.$
It should be noted that this function in the limit $f_{1}(\hat{n})\rightarrow
1$ (which is equivalent to the harmonic confinement) is proportional to the
associated Laguerre polynomials
$F_{i}(\hat{n},\kappa)|_{f_{1}(\hat{n})\rightarrow
1}=\frac{e^{-\frac{\kappa^{2}}{2}}}{\hat{n}+i}L_{\hat{n}}^{i}\left(\kappa^{2}\right).$
(21)
Now we write the function $F_{i}(\hat{n},\kappa)$ (3)
$F_{i}(\hat{n},\kappa)=\frac{e^{-\frac{\kappa^{2}}{2}((n+1+i)f_{1}^{2}(n+1+i)-(n+i)f_{1}^{2}(n+i))}}{\hat{n}+i}f_{1}(\hat{n})!f_{1}(\hat{n}+1)!L_{f,\hat{n}}^{i}\left(\kappa^{2}\right),$
(22)
where the function $L_{f,\hat{n}}^{i}(x)$ is defined as
$L_{f,\hat{n}}^{i}(x)=\sum_{l=0}^{n}\frac{1}{[f_{1}(\hat{n}-l)!]^{2}}\frac{(\hat{n}+i)!}{(\hat{n}-l)!l!(l+i)!}(-x)^{l}.$
(23)
This function is similar to the associated Laguerre function.
The time evolution of the system under consideration is characterized by the
master equation
$\frac{d\hat{\rho}}{dt}=-\frac{i}{\hbar}[\hat{H}_{I}^{(1)},\hat{\rho}]+\mathfrak{L}\hat{\rho},$
(24)
where $\mathfrak{L}\hat{\rho}$ defines damping of the system due to the
different kinds of interactions which lead to annihilation of exciton. We
assume a bosonic reservoir that causes damping of exciton system. Due to the
properties of dark exciton, the rate of spontaneous recombination and hence
spontaneous emission is decrease. On the other hand, interactions of exciton-
phonon and exciton-impurities cause the exciton to be damped. In fact in low
temperatures it is possible to ignore the phonon effects and by assuming a
pure system we neglect the impurity effects. Hence we can write
$\mathfrak{L}\hat{\rho}=\frac{\Gamma}{2}(2\hat{b}\hat{\rho}\hat{b}^{{\dagger}}-\hat{b}^{{\dagger}}\hat{b}\hat{\rho}-\hat{\rho}\hat{b}^{{\dagger}}\hat{b}),$
(25)
where $\Gamma$ is the energy relaxation rate, $\hat{b}$ and
$\hat{b}^{{\dagger}}$ are the annihilation and creation operators of the
reservoir. Due to the confinement and dark state properties, spontaneous
recombination of exciton decreases and hence the lifetime of exciton becomes
so long that we can consider the stationary solution of Eq.(24). We assume a
finite lifetime for exciton, and during this time we neglect damping effects.
The stationary solution of the master equation (24) in the time scales of our
interest is
$\hat{\rho}=|e\rangle|\psi\rangle\langle\psi|\langle e|,$ (26)
where $|e\rangle$ is the electronic excited state correspond to the presence
of exciton and $|\psi\rangle$ is the center-of-mass motion state of the
exciton, which can be considered as a right-hand eigenstate of the deformed
operator $\hat{A}=\frac{F_{1}(\hat{n},\kappa)}{F_{0}(\hat{n},\kappa)}\hat{a}$
$\frac{F_{1}(\hat{n},\kappa)}{F_{0}(\hat{n},\kappa)}\hat{a}|\psi\rangle=\frac{i\Omega_{0}}{\Omega_{1}\kappa}|\psi\rangle.$
(27)
According to Eq.(22) the corresponding deformation function reads as
$\displaystyle f(\hat{n})$ $\displaystyle=$
$\displaystyle\frac{F_{1}(\hat{n}-1,\kappa)}{F_{0}(\hat{n}-1,\kappa)}$
$\displaystyle=$
$\displaystyle\frac{f_{1}(\hat{n})L_{f,\hat{n}-1}^{1}(\kappa^{2})}{nL_{f,\hat{n}-1}^{0}(\kappa^{2})}e^{-\frac{\kappa^{2}}{2}\left((n+1)f_{1}^{2}(n+1)-(n-1)f_{1}^{2}(n-1)\right)}.$
Hence, we can express the state $|\psi\rangle$ in the Fock space
representation as
$|\psi\rangle=\mathcal{N}_{f}\sum_{n}\frac{\chi^{n}}{\sqrt{n!}f(n)!}|n\rangle,$
(29)
where $\chi=\frac{i\Omega_{0}}{\kappa\Omega_{1}}$. According to the definition
(1), it is evident that the state $|\psi\rangle$ can be regarded as a special
kind of NLCS. As is seen from equation (27), the eigenvalues of the deformed
operator $\hat{A}$ depends on some physical parameters such as the Rabi
frequencies, the parameter $\kappa$ and radius of QD.
As is clear from equation (3), the deformation function $f(\hat{n})$ depends
on the quantum number $\hat{n}$ and physical parameters such as QD radius and
$\kappa$ which characterizes the confinement regime. In the limit
$f_{1}(\hat{n})\rightarrow 1$, (harmonic confinement), which corresponds, for
example, to a QD in lens shape [28], the function $L_{f,\hat{n}}^{i}$ reduces
to the ordinary associated Laguerre polynomials, its argument tends to
$\kappa^{2}$ and therefore, the deformation function (3) takes the following
form
$f(\hat{n})=e^{-\kappa^{2}}L_{\hat{n}}^{1}(\kappa^{2})[(\hat{n}+1)L_{\hat{n}}^{0}(\kappa^{2})]^{-1}.$
(30)
This is the deformation function that appears in the center-of-mass motion of
a trapped ion confined in a harmonic trap [6].
In order to investigate the nonclassical behavior of the NLCS $|\psi\rangle$
we consider the quadrature squeezing of the center-of-mass motion. For this
purpose, we define the deformed quadratures operators as follows
$\hat{X}_{1}=\frac{1}{2}(\hat{A}e^{i\phi}+\hat{A}^{{\dagger}}e^{-i\phi}),\hskip
28.45274pt\hat{X}_{2}=\frac{1}{2i}(\hat{A}e^{i\phi}-\hat{A}^{{\dagger}}e^{-i\phi}).$
(31)
In the limiting case $f(\hat{n})\rightarrow 1$, these two operators reduce to
the conventional (non-deformed) quadrature operators [29]. The commutation
relation of $\hat{X}_{1}$ and $\hat{X}_{2}$ is
$[\hat{X}_{1},\hat{X}_{2}]=\frac{i}{2}[(\hat{n}+1)f^{2}(\hat{n}+1)-\hat{n}f^{2}(\hat{n})].$
(32)
The variances
$\langle(\Delta\hat{X}_{i})^{2}\rangle\equiv\langle\hat{X}_{i}^{2}\rangle-\langle\hat{X}_{i}\rangle^{2}(i=1,2)$
satisfy the uncertainty relation
$\langle(\Delta\hat{X}_{1})^{2}\rangle\langle(\Delta\hat{X}_{2})^{2}\rangle\geq\frac{1}{16}(\langle(\hat{n}+1)f^{2}(\hat{n}+1)-\hat{n}f^{2}(\hat{n})\rangle)$
(33)
A quantum state is said to be squeezed when one of the quadratures components
$\hat{X}_{1}$ and $\hat{X}_{2}$ satisfies the relation
$\langle(\Delta\hat{X}_{i})^{2}\rangle<\frac{1}{4}\langle(\hat{n}+1)f^{2}(\hat{n}+1)-\hat{n}f^{2}(\hat{n})\rangle\hskip
28.45274pti=1\;or\;2$ (34)
The degree of squeezing can be measured by the squeezing parameter
$s_{i}(i=1,2)$ defined by
$s_{i}=4\langle(\Delta\hat{X}_{i})^{2}\rangle-\frac{1}{4}\langle(\hat{n}+1)f^{2}(\hat{n}+1)-\hat{n}f^{2}(\hat{n})\rangle.$
(35)
Then the condition for squeezing in the quadrature component can be simply
written as $s_{i}<0$. In Fig.(1) we plot the squeezing parameter $s_{1}$
versus the parameter $\frac{R}{a_{B}}$ defined as the ratio of the QD radius
to the Bohr radius of exciton for two different values of ratio
$\frac{\Omega_{1}}{\Omega_{0}}$. As is clear from Fig.(1) for small values of
the parameter $\frac{R}{a_{B}}$ the state shows quadrature squeezing and by
increasing this parameter the quadrature squeezing disappears.
## 4 Conclusion
In this paper, we first considered a particle confined in a spherical infinite
well and we found the explicit forms of its creation and annihilation
operators by using the $f$-deformed oscillator approach. Then we considered an
exciton in a wide QD interacts with two laser beams. We showed that under the
weak confinement condition, the exciton is influenced as a whole and its
center-of-mass motion will be quantized. Within the framework of the
$f$-deformed oscillator approach, we found that under certain circumstances of
exciton-laser interaction the stationary state of the exciton center-of-mass
is a nonlinear coherent state which exhibits the quadrature squeezing.
Acknowledgment The authors wish to thank the Office of Graduate Studies of the
University of Isfahan and Iranian Nanotechnology initiative for their support.
## References
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Figure 1: Plots of squeezing versus ration $\frac{r}{a_{B}}$. Solid line
correspond to $\frac{\Omega_{0}}{\Omega_{1}}=0.5$, dash line correspond to
$\frac{\Omega_{0}}{\Omega_{1}}=0.2$. In both plots Lamb-Dick parameter is
chosen as $\kappa=0.3$.
|
arxiv-papers
| 2011-12-11T11:05:03 |
2024-09-04T02:49:25.171054
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi",
"submitter": "Malek Bagheri",
"url": "https://arxiv.org/abs/1112.2344"
}
|
1112.2346
|
# $Q$-deformed description of excitons and associated physical results
M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi Quantum Optics Group, Physics
Department, University of Isfahan, Iran m-baghreri@phys.ui.ac.ir
rokni@sci.ui.ac.ir mhnaderi@phys.ui.ac.ir , ,
###### Abstract
We consider excitons in a quantum dot as $q$-deformed systems. Interaction of
some excitonic systems with one cavity mode is considered. Dynamics of the
system is obtained by diagonalizing total Hamiltonian and emission spectrum of
quantum dot is derived. Physical consequences of $q$-deformed exciton on
emission spectrum of quantum dot is given. It is shown that when the exciton
system deviates from Bose statistics, emission spectra will become multi peak.
With our investigation we try to find the origin of the $q$-deformation of
exciton. The optical response of excitons, which affected by the nonlinear
nature of $q$-deformed systems, up to the second order of approximation is
calculated and absorption spectra of the system is given.
###### pacs:
73.20.La, 71.35.Cc, 03.65.Fd
††: J. Phys. B: At. Mol. Phys.
## 1 Introduction
Exciton is an elementary excitation of a semiconductor which consists of a
pair of two correlated fermions, the electron and the hole. Analogous to the
Hydrogen atom, it is characterized by a binding energy $E_{b}$ and a Bohr
radius $a_{B}$. Because an exciton is composed of two fermions, it is a
composite boson. Particularly, in a bulk semiconductor, when the excitation of
system is dilute, i.e. $n_{ex}a_{B}^{3}\ll 1$, where $n_{ex}$ is the exciton
density, a bosonic description of system is convenient [1]. Also Bose-Einstein
condensation of excitons, which is an essential characteristic of boson
systems has been considered theoretically [2]. When density of excitons
increase, the above condition is violated. In this situation the statistics of
excitons deviates from Bose statistics.
In low dimensional semiconductor systems such as quantum well (QW), quantum
wire and quantum dot (QD), due to the small dimensions and loss of
translational symmetry, exciton excitation differs from exciton in bulk
materials. In semiconductor nanostructures the size of the system strongly
affects exciton properties. For example, in the case of quantum well it is
shown [3], if the well width is larger (smaller) than the Bohr radius of
exciton, the spectrum of quantum well has properties similar to the situation
in which excitons are boson (fermion). Hence, the size of the system directly
affects the quantum statistics of excitons in that system. Recently similar
results has been obtained for QD [4]. In Ref.[4] the effects of different
statistics of excitons on emission spectra of a QD is investigated, and the
origin of different statistics of excitons is considered. The same results
have also been obtained for the quantum well. If the size of QD is smaller
(larger) than the exciton Bohr radius, excitons behave like fermion (boson).
Real statistics of excitons in the interaction is considered in [5] and
references therein. As mentioned before in high density regime, exciton
statistics deviates from Bose statistics. This is due to the increase of
mutual forces between the excitations of the system and then the Pauli
exclusion principle plays a dominant role [6]. Appearance of Bose statistics
of exciton-biexciton system and Pauli exclusion effects in superlattice has
been considered experimentally [7].
Bosons and fermions are the only two kinds of particles realized in nature.
The conditions mentioned for excitons (in one regime they are like bosons and
in another one like fermions) are property of a special kind of statistics
called intermediate statistics [8]. Bose and Fermi statistics are two limiting
cases of this statistics. Properties of this statistics have been considered
by many authors [9]$-$[11]. Operator realization of intermediate statistics is
similar to $q$-deformed operators [12]. Bosonic $q$-deformed operators [13]
are a generalization of the Heisenberg algebra obtained by introducing a
deformation parameter $q$. Deviation of this parameter from 1 shows deviation
of algebra from the Heisenberg algebra. It is shown that it is possible to
describe correlated fermion pairs with $q$-deformed bosons [14]. Therefore it
is reasonable to consider an exciton system as a $q$-deformed system. We
assume the creation and annihilation operators of excitons obey a $q$-deformed
algebra. A $q$-deformed description of Frenkel exciton has been considered
recently [15].
The algebra generated by $q$-deformed operators are given by
$\displaystyle[\hat{b}_{q},\hat{b}_{q}^{\dagger}]_{q}$ $\displaystyle=$
$\displaystyle\hat{b}_{q}\hat{b}_{q}^{\dagger}-q^{-1}\hat{b}_{q}^{\dagger}\hat{b}_{q}=q^{\hat{n}},$
(1) $\displaystyle[\hat{n},\hat{b}_{q}^{\dagger}]$ $\displaystyle=$
$\displaystyle\hat{b}_{q}^{\dagger},\hskip
42.67912pt[\hat{n},\hat{b}_{q}]=-\hat{b}_{q}.$
where $\hat{n}=\hat{b}^{\dagger}\hat{b}$ is the usual particle number
operator. Representation of this algebra is given in [16]. In the case of
excitons, $q$-parameter can depend on excitation number and physical size of
system.
In this paper we consider the interaction of light with a QD embedded in a
microcavity. By considering excitons in QD as $q$-deformed bosons (case of
$q$-deformed fermion is straightforward) we study the emission spectrum of the
system. As is clear, the commutator (1) explicitly depends on the number of
excitons. Hence, this is a system in which light interacts with a nonlinear
active medium. Therefore, we shall obtain the linear and nonlinear response of
a $q$-deformed exciton system. Knowledge of interaction of light with a
nonlinear medium ($q$-deformed excitons) and its optical response is important
for the interpretation of experimental results such as [30]. On the other
hand, we compared the obtained results with some experimental ones and in this
manner we investigate the physical origin of $q$-deformation of excitons. In
section 2 we derive the spectrum of QD when one exciton mode interacts with a
single mode cavity-field. In section 3 we consider the interaction of two
exciton modes with a single cavity mode. In section 4 the nonlinear response
of QD is derived up to second order of approximation. Finally we summarize our
conclusions in section 5.
## 2 Model Hamiltonian
We consider a QD embedded in a microcavity which interacts with a single mode
cavity-field. We assume the excitations in QD have an intermediate statistics
[4], and their creation and annihilation operators obey $q$-deformed algebra.
We can express the $q$-deformed operators in terms of ordinary boson operators
by the following maps
$\hat{b}_{q}=\hat{b}\sqrt{\frac{q^{\hat{n}}-q^{-\hat{n}}}{\hat{n}(q-q^{-1})}},\hskip
42.67912pt\hat{b}_{q}^{\dagger}=\sqrt{\frac{q^{\hat{n}}-q^{-\hat{n}}}{\hat{n}(q-q^{-1})}}\hat{b}^{\dagger},$
(2)
where $\hat{b}$ and $\hat{b}^{\dagger}$ are the ordinary boson operators and
$\hat{n}=\hat{b}^{\dagger}\hat{b}$. Ordinary commutator of $q$-deformed
exciton operators is
$[\hat{b}_{q},\hat{b}_{q}^{\dagger}]=\frac{q}{q+1}[q^{n}+q^{-(n+1)}]\equiv
k(\hat{n}).$ (3)
Deviation of this commutator from ordinary boson algebra
($[\hat{b},\hat{b}^{\dagger}]=1$) relates to deviation of $q$-parameter from
1. It is clear, this generalized commutator depends on the number of
excitations. It seems that by using this algebra we can consider some
nonlinear phenomena in the system related to the population of excitons. For
example, biexciton effects can be considered in this manner as an effective
approach. So that the deformation parameter $q$ can represent some physical
parameters such as the ratio of size of system to the Bohr radius of exciton.
Interaction of QD with single mode cavity-field in rotating wave approximation
can be described by the following Hamiltonian
$\hat{H}=\hbar\omega\hat{a}^{\dagger}\hat{a}+\hbar\omega_{ex}\hat{b}_{q}^{\dagger}\hat{b}_{q}+\hbar
g(\hat{a}\hat{b}_{q}^{\dagger}+\hat{a}^{\dagger}\hat{b}_{q}),$ (4)
where $\hat{a}$ and $\hat{a}^{\dagger}$ are creation and annihilation
operators of cavity field and
$[\hat{a}_{i},\hat{a}_{j}^{\dagger}]=\delta_{ij}$. We shall consider a
phenomenological damping for the system which relates to both subsystems:
photon and exciton. As is clear from the Hamiltonian (4), the exciton number
is not a constant of motion. Because of the dependence of exciton operator,
$\hat{b}_{q}$, on the exciton number, resulting equations of motion become a
nontrivial set of coupled equations. On the other hand, since the total number
of excitation (exciton and photon) is conserved we can diagonalize the
Hamiltonian in the subspace of a definite excitation. To consider this
dynamics we propose an approach based on diagonalization of the Hamiltonian by
using the polariton transformation [17]. This procedure depends on some
unitary transformations which diagonalize the model Hamiltonian. As is usual
in this procedure [18], new operators have the same commutation relation as
the original operators (free operators). Here, there are two distinct sets of
operators, the cavity mode operators which obey the usual boson commutation
relation and exciton operators that are $q$-deformed boson. Therefore, with
the presence of these two different statistics, mixed operators (polariton
operators) do not have specific statistics. They can be considered as ordinary
boson operators or $q$-deformed operators. We consider both situations and we
study the physical results associated with each situation in the resonance
fluorescence spectrum of QD.
### 2.1 Boson polaritons
In order to solve the dynamical system, we perform the following
transformation
$\hat{p}_{k}=u_{k}\hat{b}_{q}+v_{k}\hat{a}.$ (5)
Due to the presence of $q$-deformed operator $\hat{b}_{q}$, we call this
transformation a polariton-like transformation. As mentioned before,
$\hat{b}_{q}$ depends on the number of excitons explicitly and this causes the
Hopfield coefficients $u_{k}$ and $v_{k}$ will depend on the number of
excitons. Hence, the transformation (5) can be considered as a nonlinear
polariton transformation. This kind of transformation has been considered
recently for the case of Bogoliubov transformation [19, 20]. We assume
polariton-like operators obey the usual boson commutation relations
$[\hat{p}_{k},\hat{p}^{\dagger}_{k^{\prime}}]=\delta_{kk^{\prime}}\Rightarrow[\hat{p}_{k},\hat{p}^{\dagger}_{k}]=|u_{k}|^{2}k(\hat{n})+|v_{k}|^{2}=1,$
(6)
where the operator valued function $k(\hat{n})$ was introduced by Eq.(3). We
choose unknown coefficients $u_{k}$ and $v_{k}$ so that the Hamiltonian (4)
becomes diagonal in terms of the polariton-like operators
$\hat{H}=\hbar\sum_{k}\Omega_{k}\hat{p}_{k}^{\dagger}\hat{p}_{k},$ (7)
where $\Omega_{k}$ is the polariton spectrum and $k$ refers to different
polariton branches. By taking into account a phenomenological damping for
exciton and photon systems separately, the unknown parameters satisfy
following set of equations
$[\omega_{ex}k(\hat{n})-\Omega_{k}-i\gamma_{ex}]u_{k}+v_{k}g=0,\hskip
31.2982ptu_{k}gk(\hat{n})+(\omega-\Omega_{k}-i\gamma_{ph})v_{k}=0.$ (8)
In this set of equations, $\gamma_{ex}$ and $\gamma_{ph}$ are the exciton and
photon damping constants, respectively. From these equations the polariton
spectrum can be obtained as
$\displaystyle\Omega_{k}$ $\displaystyle=$
$\displaystyle\frac{\omega_{ex}k(\hat{n})+\omega-i(\gamma_{ex}+\gamma_{ph})}{2}$
(9) $\displaystyle\pm$
$\displaystyle\frac{1}{2}\sqrt{[\omega_{ex}k(\hat{n})-\omega-i(\gamma_{ex}-\gamma_{ph})]^{2}+4g^{2}k(\hat{n}))}.$
It is apparent that $q$-deformed description of excitons causes the splitting
between these energy eigenvalues be increased in compare to the case of
bosonic description of exciton. Using the set of equations (8) and the
polariton spectrum (9) we find the coefficients for two polariton branches
$\displaystyle
u_{k}=\sqrt{\frac{\omega-i\gamma_{ph}-\Omega_{k}}{k(\hat{n})[\omega-2\Omega_{k}+\omega_{ex}k(\hat{n})-i(\gamma_{ex}+\gamma_{ph})]}},$
(10) $\displaystyle
v_{k}=-\sqrt{\frac{\omega_{ex}k(\hat{n})-i\gamma_{ex}-\Omega_{k}}{\omega-2\Omega_{k}+\omega_{ex}k(\hat{n})-i(\gamma_{ex}+\gamma_{ph})}}.$
By employing these coefficients all necessary parameters for the polariton
Hamiltonian are determined.
Now we can consider the dynamics of polariton operators. The time evolution of
polariton operators is governed by the polariton Hamiltonian (7)
$\hat{\dot{p}}_{k}=\frac{-i}{\hbar}[\hat{p}_{k},\hat{H}]=-i\Omega_{k}\hat{p}_{k}.$
(11)
Let us consider damping effects by taking into account a phenomenological
damping term and noise operator in the dynamical equations of polariton
operators. Hence, the time evolution of polariton operator is given by
$\hat{\dot{p}}_{k}=-i\Omega_{k}\hat{p}_{k}-\Gamma_{k}\hat{p}_{k}+\hat{F}_{\hat{p}_{k}}(t),$
(12)
where $\hat{F}_{\hat{p}_{k}}(t)$ is the Langevin noise operator which depends
on the reservoir variables and $\Gamma_{k}$ is the damping constant of $k$th
polariton branch given by $\Gamma_{k}=\frac{\gamma_{ex}+\gamma_{ph}}{2}$.
Correlation functions of the noise operators determine physical properties of
the system. The Langevin noise operator are such that their expectation values
$\langle\hat{F}_{x}\rangle$ vanishes, but their second order moments do not
[21]. They are intimately linked up with the global dissipation and in a
Markovian environment they take the form
$\langle\hat{F}^{\dagger}_{\hat{p}_{k}}(t)\hat{F}_{\hat{p}_{k}}(t^{\prime})\rangle=2\Gamma_{k}\delta(t-t^{\prime}).$
(13)
With neglecting the phonon effects by decreasing the temperature, other
sources of damping like spontaneous recombination of exciton and photon loss
are considered as Markovian procedures. It follows, on solving Eq.(12), that
$\hat{p}_{k}(t)=\hat{p}_{k}(0)e^{(-i\Omega_{k}-\Gamma_{k})t}+\int_{0}^{t}e^{(-i\Omega_{k}-\Gamma_{k})(t-t^{\prime})}\hat{F}_{\hat{p}_{k}}(t^{\prime})dt^{\prime}.$
(14)
In this equation we set initial time equal zero.
The power spectrum of the scattered light for statistical stationary fields is
given by [22]
$S(r,\omega)=\frac{1}{\pi}Re\int_{0}^{\infty}\langle\hat{E}^{-}(r,t)\hat{E}^{+}(r,t+\tau)\rangle
e^{i\omega\tau}d\tau,$ (15)
where $\hat{E}^{\pm}$ are the positive and negative frequency parts of the
electric field operator. Expressing field operators in terms of creation and
annihilation operators we have
$S(r,\omega)=\frac{A(r)}{\pi}Re\int_{0}^{\infty}\langle\hat{a}^{\dagger}(0)\hat{a}(\tau)\rangle
e^{i\omega\tau}d\tau.$ (16)
Here, we set $t=0$, and $A(r)$ depends on mode function of the cavity-field.
Now we can express, the field and exciton creation and annihilation operators
in terms of polariton ones:
$\hat{a}=v_{1}^{\ast}\hat{p}_{1}+v_{2}^{\ast}\hat{p}_{2},\hskip
42.67912pt\hat{b}_{q}=k(\hat{n})(u_{1}^{\ast}\hat{p}_{1}+u_{2}^{\ast}\hat{p}_{2}),$
(17)
and at the time $t$ we have
$\hat{a}(t)=v_{1}^{\ast}\hat{p}_{1}(t)+v_{2}^{\ast}\hat{p}_{2}(t).$ (18)
Now to calculate the resonance fluorescence spectrum we have to determine the
initial state of system. we assume at $t=0$, the cavity-field is in a coherent
state $|\alpha\rangle$, and the exciton subsystem in its vacuum state. Under
this condition, by using Eq.(14) the resonance fluorescence spectrum is
obtained as
$S(r,\omega)=\frac{A(r)|\alpha|^{2}}{\pi}\left[|v_{1}|^{2}\frac{\Gamma_{1}}{(\omega-\Omega_{1})^{2}+\Gamma_{1}^{2}}+|v_{2}|^{2}\frac{\Gamma_{2}}{(\omega-\Omega_{2})^{2}+\Gamma_{2}^{2}}\right].$
(19)
In deriving this result we implicitly assume that at $t=0$ the noise operator
and polariton operators are uncorrelated. Fig.(1) shows the plot of
$S(r,\omega)$ versus $\omega$ for different values of deformation parameter q.
Material parameters are chosen as $\omega=1.75\;eV$, $\omega_{ex}=1.75\;eV$,
$g=200\;\mu eV$, $\gamma_{ex}=20\;\mu eV$, $\gamma_{ph}=40\;\mu eV$ [23],
$n=100$ and $|\alpha|^{2}=9$ . As is clear when $q=1$, spectrum has similar
variation as experimental results [23]. This figure shows that when $q=1$
(nondeformed case) the power spectrum of the fluorescence light is a double
peak centered at $\omega=\Omega_{1}$ and $\omega=\Omega_{2}$. By increasing
deviation of q from 1, it is apparent from the different plots in this figure
that splitting between two peaks increases and the height of one of peaks
decreases. This result has been reported in resonance fluorescence of excitons
when the biexcitonic interaction is taken into account. It has been shown [3]
that biexcitonic effects are a red shift of the transition frequencies,
emergence of sidebands due to the switch-on forbidden transitions and
asymmetry of the emission spectrum. The binding energy of biexciton in QD
causes a shift in the spectrum of the system. In the present model the
splitting of spectrum (Rabi splitting) depends on the $q$-parameter. Hence,
changing this parameter affects the spectrum. Then as a one reason of
deviation of excitons from ideal Bose system we can consider Coulomb
interaction between them. On the other hand, $q$-deformed exciton operators
depend on the total number of exciton, and biexciton interaction occurs when
there are more than one exciton. This similarity makes this clue that the
$q$-deformation can be consider as an effective approach to take into account
the biexciton effects. As mentioned before, the $q$-parameter can depend on
the size of sample. The plotted resonance fluorescence spectrum in Fig.(1)
makes clear some differences of optical properties of different size QD. For
large values of q, compare with 1, spectrum will reduce to one peak. This case
is a characteristic of the weak coupling regime.
### 2.2 $Q$-deformed polaritons
In this subsection we assume that the polariton operators are $q$-deformed
operators. According to the $q$-deformed nature of the exciton system we
assume the following algebra for polariton operators
$[\hat{p}_{k},\hat{p}^{\dagger}_{k}]_{s}=\hat{p}_{k}\hat{p}^{\dagger}_{k}-s^{-1}\hat{p}^{\dagger}_{k}\hat{p}_{k}=s^{\hat{n}_{k}},$
(20)
where $s$ denotes the deformation parameter corresponding to the polariton
system and $\hat{n}_{k}$ shows the number operator for $k$th polariton branch.
Ordinary commutator for these operators is
$[\hat{p}_{k},\hat{p}^{\dagger}_{k}]=|u_{k}|^{2}k(\hat{n})+|v_{k}|^{2}=\frac{s}{s+1}[s^{\hat{n}_{k}}+s^{-(\hat{n}_{k}+1)}]=M(\hat{n}_{k}).$
(21)
Using the same approach of the previous subsection we obtain the following set
of equations for the coefficients of transformation
$\displaystyle[(\omega_{ex}k(\hat{n})-i\gamma_{ex}-\Omega^{\prime}_{k}M(n_{k})]u_{k}+v_{k}g$
$\displaystyle=$ $\displaystyle 0,$ $\displaystyle
u_{k}gk(\hat{n})+[\omega-i\gamma_{ph}-\Omega^{\prime}_{k}M(n_{k})]v_{k}$
$\displaystyle=$ $\displaystyle 0.$ (22)
From this set of equations we derive the deformed polariton spectrum as
$\displaystyle\Omega^{\prime}_{k}$ $\displaystyle=$
$\displaystyle\frac{\omega_{ex}k(\hat{n})+\omega-i(\gamma_{ex}+\gamma_{ph})}{2M(n_{k})}$
(23) $\displaystyle\pm$
$\displaystyle\frac{\sqrt{[\omega_{ex}k(\hat{n})-\omega-i(\gamma_{ex}-\gamma_{ph})]^{2}+4g^{2}k(\hat{n}))}}{2M(n_{k})}.$
and the transformation coefficients read as
$\displaystyle
u_{k}=-\sqrt{\frac{M(n_{k})[\omega-i\gamma_{ph}-\Omega^{\prime}_{k}M(n_{k})]}{k(\hat{n})[\omega-2\Omega^{\prime}_{k}M(n_{k})+\omega_{ex}k(\hat{n})-i(\gamma_{ex}+\gamma_{ph})]}},$
(24) $\displaystyle
v_{k}=\sqrt{\frac{M(n_{k})[\omega_{ex}k(\hat{n})-i\gamma_{ex}-\Omega^{\prime}_{k}M(n_{k})]}{\omega-2\Omega^{\prime}_{k}M(n_{k})+\omega_{ex}k(\hat{n})-i(\gamma_{ex}+\gamma_{ph})}}.$
By determining all the variables, polariton Hamiltonian (diagonal Hamiltonian)
will be determined. By applying the same procedure as before we derive the
resonance fluorescence spectrum in this case as follows
$S(r,\omega)=\frac{A(r)|\alpha|^{2}(|v_{1}|^{2}+|v_{2}|^{2})}{\pi}\sum_{i=1,2}|v_{i}|^{2}\frac{\Gamma_{i}}{(\omega-\Omega^{\prime}_{i}M(n_{k}))^{2}+\Gamma_{i}^{2}}.$
(25)
Fig. (2) shows the plot of $S(r,\omega)$ versus $\omega$ for different values
of polariton deformation parameter $s$. This figure shows that changes of
s-parameter (deformation parameter of polariton) does not cause any shift in
transition frequencies, but causes strengths of peaks increase.
## 3 Interaction of light with two exciton modes
We now consider the interaction of one cavity mode with QD when two exciton
modes are coupled to the field mode. As before, we assume exciton system is
expressed by the $q$-deformed operators. The total Hamiltonian of the system
under consideration can be written as follows
$\hat{H}=\hbar\omega\hat{a}^{\dagger}\hat{a}+\hbar\sum_{i=1,2}\omega_{{ex}_{i}}\hat{b}_{q_{i}}^{\dagger}\hat{b}_{q_{i}}+\hbar
g\sum_{i=1,2}(\hat{a}\hat{b}_{q_{i}}^{\dagger}+\hat{a}^{\dagger}\hat{b}_{q_{i}}).$
(26)
We assume both excitons have the same coupling constant with the cavity mode.
We solve this system as before by diagonalizing the Hamiltonian. For this
purpose we perform the following transformation
$\hat{p}_{k}=u_{k}\hat{b}_{q_{1}}+x_{k}\hat{b}_{q_{2}}+v_{k}\hat{a}.$ (27)
We consider the situation in which the polariton operators obey the
nondeformed Bose statistics
$[\hat{p}_{k},\hat{p}^{\dagger}_{k}]=|u_{k}|^{2}k(\hat{n}_{1})+|x_{k}|^{2}k(\hat{n}_{2})+|v_{1}|^{2}=1,$
(28)
where $\hat{n}_{i}$ ($i=1,2$) represents the number operator for each
excitonic mode. As is clear in this case there are three polariton branches.
Assuming the transformation (27) diagonalizes the Hamiltonian (26), this
polariton Hamiltonian takes the following form
$\hat{H}=\hbar\sum_{k}\Omega_{k}\hat{p}^{\dagger}_{k}\hat{p}_{k},$ (29)
where summation is over all polariton branches. The following equation
determines the polariton spectrum
$(c-\Omega_{k})[(d-\Omega_{k})(\omega-i\gamma_{ph}-\Omega_{k})-g^{2}k(\hat{n}_{1})]-g^{2}k(\hat{n}_{2})(d-\Omega_{k})=0,$
(30)
where $c=\omega_{ex_{1}}k(\hat{n}_{1})-i\gamma_{ex_{1}}$ and
$d=\omega_{ex_{2}}k(\hat{n}_{2})-i\gamma_{ex_{2}}$. By deriving the polariton
spectrum the transformation parameters are obtained as
$\displaystyle
u_{k}=\frac{g[(d-\Omega_{k})(\omega-i\gamma_{ph}-\Omega_{k})-g^{2}k(\hat{n}_{2})]}{A},$
(31) $\displaystyle x_{k}=\frac{g^{3}k(\hat{n}_{1})}{A}$ $\displaystyle
v_{k}=-\frac{(c-\Omega_{k})[(d-\Omega_{k})(\omega-i\gamma_{ph}-\Omega_{k})-g^{2}k(\hat{n}_{2})]}{A},$
where the parameter $A$ is given by
$\displaystyle A$ $\displaystyle=$
$\displaystyle([g^{2}k(\hat{n}_{1})+(c-\Omega_{k})^{2}][(d-\Omega_{k})(\omega-i\gamma_{ph}-\Omega_{k})-g^{2}k(\hat{n}_{2})]^{2}$
(32) $\displaystyle+$ $\displaystyle
g^{6}k^{2}(\hat{n}_{1})k(\hat{n}_{2}))^{\frac{1}{2}}.$
In this manner, all the parameters which appear in the polariton Hamiltonian
are determined. By repeating the approach of previous section the resonance
fluorescence spectrum of system with different initial conditions can be
determined. If we assume at $t=0$ the cavity mode is in the coherent state
$|\alpha\rangle$ and QD in vacuum state $|0\rangle$, the resonance
fluorescence spectrum is given by
$S(r,\omega)=\frac{|\alpha|^{2}A(r)}{\pi}\sum_{k}\frac{|v_{k}|^{2}\Gamma_{k}}{\Gamma_{k}^{2}+(\omega-\Omega_{k})^{2}}.$
(33)
To show complex structure (multi-peak structure) of this spectrum Fig. (3)
presents the spectra on a logarithmic scale. For the sake of clarity, we have
powered some peaks compare to other ones in this figure. In the case of $q=1$
(nondeformed exciton) the spectrum has two peaks. Increasing the $q$-parameter
causes that splitting between peaks be increased and spectrum becomes multi-
peaks. Multi-peaks structure in emission of exciton such as Mollow triplet was
predicted when excitons obey statistics different from Bose statistics [3, 4].
When, $q$-parameter is changed, the energy and intensities of emission change.
Effects of exciton number on absorption spectrum of QD is considered . Due to
the relation of absorption spectrum and resonance fluorescence, similar result
is obtain in [24].
## 4 Nonlinear response of excitons in $q$-deformed regime
In previous sections we considered some physical results of $q$-deformed
description of excitons. The $q$-deformed description can be served as a
nonlinear description of excitons. It is well-known that different kinds of
nonlinearity in an exciton system lead to different orders of nonlinear
response of the system [25, 26]. Therefore, we try to obtain optical response
of a driven quantum dot, which its optical excitations are considered as
$q$-deformed systems. For this purpose we will calculate the coefficient
absorption of a QD in this regime. In this section we neglect all damping
effects and we consider the Hamiltonian of the system as follows
$\hat{H}=\hbar\omega\hat{a}^{\dagger}\hat{a}+\hbar\omega_{ex}\hat{b}_{q}^{\dagger}\hat{b}_{q}+\hbar
g(\hat{a}\hat{b}_{q}^{\dagger}+\hat{a}^{\dagger}\hat{b}_{q}).$ (34)
In the electron picture, the induced dipole moment by transition of an
electron is described by
$\hat{\mu}=\hat{a}^{\dagger}_{v}\hat{a}_{c}+\hat{a}^{\dagger}_{c}\hat{a}_{v}$
[27]. The operator $\hat{a}^{\dagger}_{v}\;(\hat{a}_{v})$ is the creation
(annihilation) operator for an electron in the valance band (level in the case
of QD), and $\hat{a}^{\dagger}_{c}\;(\hat{a}_{c})$ is the creation
(annihilation) operator for an electron in the conduction band. Hence,
creation of an exciton is denoted by
$\hat{a}^{\dagger}_{c}\hat{a}_{v}=\hat{b}_{q}^{\dagger}$. Therefore we can
write the dipole operator of QD as
$\hat{\mu}=\hat{b}^{\dagger}_{q}+\hat{b}_{q}$. The macroscopic polarization is
expectation value of polarization operator. The optical response function
represents the reaction of the system to an external classic field $E(t)$
coupled to the variables of system [28], i.g., the dipole operator. Hence, we
consider an external field as a pump source and we treat the reaction of QD to
it. Then the total Hamiltonian of system is then given by
$\hat{H}=\hbar\omega\hat{a}^{\dagger}\hat{a}+\hbar\omega_{ex}\hat{b}_{q}^{\dagger}\hat{b}_{q}+\hbar
g(\hat{a}\hat{b}_{q}^{\dagger}+\hat{a}^{\dagger}\hat{b}_{q})-[\vec{d}_{vc}\cdot\vec{E}(t)\hat{b}_{q}+\vec{d}_{cv}\cdot\vec{E}(t)\hat{b}^{\dagger}_{q}],$
(35)
where $\vec{d}_{vc}$ denotes the dipole matrix element. The Hamiltonian in the
interaction picture has the form
$\displaystyle\hat{H}_{int}$ $\displaystyle=$ $\displaystyle\hbar
g\left[\hat{a}\hat{b}^{\dagger}_{q}e^{-i[\omega-\omega_{ex}k(\hat{n})]t}+\hat{a}^{\dagger}e^{i[\omega-\omega_{ex}k(\hat{n})]t}\hat{b}_{q}\right]$
$\displaystyle-$
$\displaystyle\left[\vec{d}_{cv}\cdot\vec{E}(t)e^{-i\omega_{ex}k(\hat{n})t}\hat{b}_{q}+\vec{d}_{vc}\cdot\vec{E}(t)\hat{b}^{\dagger}_{q}e^{i\omega_{ex}k(\hat{n})t}\right].$
The observable of interest for the optical response is the time-dependent
dipole density
$\mu(t)=\langle\hat{b}_{q}(t)\rangle+h.c.=Tr_{ex}(\hat{b}_{q}\rho_{ex}(t))+h.c.$,
where $Tr_{ex}$ means trace over the exciton system and
$\rho_{ex}(t)=Tr_{f}\rho(t)$, which $\rho(t)$ is the total time dependent
density matrix of the system and $\rho_{ex}(t)$ is the time dependent density
matrix of exciton system. The total time dependent density matrix is given by
$\hat{\rho}(t)=\hat{U}(t,t_{0})\hat{\rho}(t_{0})\hat{U}^{-1}(t,t_{0}),$ (37)
where
$U(t,t_{0})=\hat{T}\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{int}(t^{\prime})dt^{\prime}]$
is the time ordered evolution operator and $\hat{\rho}(t_{0})$ is the total
density matrix of system at initial time. We assume that the quantum field and
exciton system are both in vacuum state. Therefore, the time dependent density
matrix of excitons is given by
$\hat{\rho}_{ex}(t)=\sum_{n}\langle
n|\hat{U}(t,t_{0})(|0\rangle_{f}|0\rangle_{ex})(_{ex}\langle 0|_{f}\langle
0|)\hat{U}^{-1}(t,t_{0})|n\rangle,$ (38)
where summation is carried on field state and the matrix elements of the time
evolution operator are in the basis of field states. By using the Feynman
disentanglement theorem [29] the matrix elements of the time evolution
operator $\hat{U}(t,t_{0})$ can be evaluated. We can write Hamiltonian in (4)
as $\hat{H}_{int}=\hat{H}_{1}(t)+\hat{H}_{2}(t)$, where
$\displaystyle\hat{H}_{1}(t)=\hbar
g\left[\hat{a}\hat{b}^{\dagger}_{q}e^{-i[\omega-\omega_{ex}k(\hat{n})]t}+\hat{a}^{\dagger}e^{i[\omega-\omega_{ex}k(\hat{n})]t}\hat{b}_{q}\right]$
(39)
$\displaystyle\hat{H}_{2}(t)=-\left[\vec{d}_{cv}\cdot\vec{E}(t)e^{-i\omega_{ex}k(\hat{n})t}\hat{b}_{q}+\vec{d}_{vc}\cdot\vec{E}(t)\hat{b}^{\dagger}_{q}e^{i\omega_{ex}k(\hat{n})t}\right].$
As is clear $\hat{H}_{2}(t)$ depends only on exciton operators. The time
evolution operator can be written as
$\displaystyle\hat{U}(t,t_{0})$ $\displaystyle=$
$\displaystyle\hat{T}\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}(\hat{H}_{1}(t^{\prime})+\hat{H}_{2}(t))dt^{\prime}]$
(40) $\displaystyle=$
$\displaystyle\hat{T}\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{2}(t^{\prime})dt^{\prime}]\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{1}(s)ds].$
In this equation we use Feynman notation [29]. These two exponential terms are
not disentangle from each other. They are correlated and in doing integration,
we have to take into account ordering of operators. In calculation of matrix
element of this operator in the basis of field states, second exponential can
be considered as a ordinary c-number function of $t^{\prime}$, because it is
independent of field operators:
$\langle i|\hat{U}(t,t_{0})|j\rangle=\langle
i|\hat{T}\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{1}(t^{\prime})dt^{\prime}]|j\rangle\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{2}(s)ds].$
(41)
On the other hand we consider all the exciton operators in $\hat{H}_{1}(t)$ as
ordinary c-number functions, and we can write
$\displaystyle\langle i|\hat{U}(t,t_{0})|j\rangle$ $\displaystyle=$
$\displaystyle\langle
i|\hat{T}\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}(\hat{a}^{\dagger}_{t^{\prime}}e^{i\omega
t^{\prime}}B(t^{\prime})+\hat{a}_{t^{\prime}}e^{-i\omega
t^{\prime}}B^{\ast}(t^{\prime}))dt^{\prime}]|j\rangle$ (42)
$\displaystyle\times$
$\displaystyle\exp[-\frac{i}{\hbar}\int_{t_{0}}^{t}\hat{H}_{2}(s)ds],$
where $B(t)$ is a ordinary function corresponding to exciton operators. As is
clear this matrix element is a function of exciton operators. The influence of
exciton system is completely contained in this operator functional and
factored term in (42). By using Feynman theorem, above matrix element can
written as
$\langle i|\exp[-ig\int_{t_{0}}^{t}\hat{a}^{\dagger}_{t^{\prime}}e^{i\omega
t^{\prime}}B(t^{\prime})dt^{\prime}]\exp[-ig\int_{t_{0}}^{t}\hat{a}^{\prime}_{t^{\prime}}e^{-i\omega
t^{\prime}}B^{\ast}(t^{\prime})dt^{\prime}]|j\rangle,$ (43)
where in this equation
$\hat{a}^{\prime}_{t^{\prime}}=\hat{V}^{-1}(t)\hat{a}_{t}\hat{V}(t)$, and
$\hat{V}(t)=\exp[-ig\hat{a}^{\dagger}\int_{t_{0}}^{t}B(t^{\prime})e^{i\omega
t^{\prime}}].$ (44)
In this manner the density matrix of exciton system takes the following form
$\hat{\rho}_{ex}(t)=\sum_{n}\frac{1}{n!}S_{1}(\hat{n},\hat{b}_{q},\hat{b}^{\dagger}_{q})|0\rangle_{ex}\langle
0|S_{2}(\hat{n},\hat{b}_{q},\hat{b}^{\dagger}_{q}),$ (45)
where
$\displaystyle S_{1}(\hat{n},\hat{b}_{q},\hat{b}^{\dagger}_{q})$
$\displaystyle=$
$\displaystyle\left[-g\hat{b}^{\dagger}_{q}\frac{e^{i\omega_{ex}k(\hat{n}_{ex})(t-t_{0})}}{\omega_{ex}k(\hat{n}_{ex})}\right]^{n}e^{-\frac{g^{2}}{2}L(\hat{n}_{ex})}f(\hat{b}_{q},\hat{b}_{q}^{\dagger}),$
$\displaystyle S_{2}(\hat{n},\hat{b}_{q},\hat{b}^{\dagger}_{q})$
$\displaystyle=$
$\displaystyle\left[-g\hat{b}_{q}\frac{e^{-i\omega_{ex}k(\hat{n}_{ex}+1)(t-t_{0})}}{\omega_{ex}k(\hat{n}_{ex}+1)}\right]^{n}e^{-\frac{g^{2}}{2}L(\hat{n}_{ex})}f^{-1}(\hat{b}_{q},\hat{b}_{q}^{\dagger}),$
$\displaystyle L(\hat{n}_{ex})$ $\displaystyle=$
$\displaystyle\hat{b}_{q}^{\dagger}\hat{b}_{q}[\frac{e^{-i\omega_{ex}[k(\hat{n}_{ex}+1)-k(\hat{n}_{ex}-1)](t-t_{0})}}{[\omega-\omega_{ex}k(\hat{n}_{ex}-1)][\omega-\omega_{ex}k(\hat{n}_{ex}+1)]}$
(46) $\displaystyle+$
$\displaystyle\frac{e^{-i\omega_{ex}[k(\hat{n}_{ex}+2)-k(\hat{n}_{ex}))](t-t_{0})}}{[\omega-\omega_{ex}k(\hat{n}_{ex}+2)][\omega-\omega_{ex}k(\hat{n}_{ex})]}],$
(47)
and
$\displaystyle f(\hat{b}_{q},\hat{b}_{q}^{\dagger})=\hat{T}\exp$
$\displaystyle[$
$\displaystyle\frac{i}{\hbar}\int_{t_{0}}^{t}dt^{\prime}(\vec{d}_{cv}\cdot\vec{E}(t^{\prime})\hat{b}_{q}e^{-i\omega_{ex}k(\hat{n}_{ex}+1)t^{\prime}}$
(48) $\displaystyle+$
$\displaystyle\vec{d}_{vc}\cdot\vec{E}(t^{\prime})\hat{b}^{\dagger}_{q}e^{i\omega_{ex}k(\hat{n}_{ex})t^{\prime}})],$
By expanding the function $f(\hat{b}_{q},\hat{b}^{\dagger}_{q})$ up to second
order in $E(t)$ and using (45) we obtain the time-dependent dipole density as
follows
$\displaystyle
p(t)=\sum_{n}\frac{g^{2n}}{n!}h_{1}(n)!\sqrt{f_{q}(n)!}e^{-\frac{g^{2}}{2}L(n)}\times$
$\displaystyle[\frac{i}{\hbar}h_{0}(n+1)!e^{-\frac{g^{2}}{2}L(1)}\sqrt{f_{q}(n+1)!f_{q}(n+1)}\int_{t_{0}}^{t}\vec{d}_{cv}\cdot\vec{E}(t^{\prime})e^{i\omega_{ex}t^{\prime}}dt^{\prime}$
$\displaystyle-\frac{i}{\hbar}\sqrt{f_{q}(n)}h_{0}(n)!e^{-\frac{g^{2}}{2}L(0)}\sqrt{f_{q}(n)!f_{q}(n)}\int_{t_{0}}^{t}\vec{d}_{cv}\cdot\vec{E}(t^{\prime})e^{i\omega_{ex}k(n-1)t^{\prime}}dt^{\prime}$
$\displaystyle+\frac{i}{2\hbar^{3}}\sqrt{f_{q}(n+1)}h_{0}(n+2)!e^{-\frac{g^{2}}{2}L(0)}\sqrt{f_{q}(n+2)!f_{q}(n+2)}\times$
$\displaystyle\int_{t_{0}}^{t}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\vec{d}_{vc}\cdot\vec{E}(t^{\prime})\vec{d}_{cv}\cdot\vec{E}(r)\vec{d}_{cv}\cdot\vec{E}(s)e^{i\omega_{ex}[s+r-k(n+1)t^{\prime}]}dt^{\prime}drds$
$\displaystyle+\frac{i}{2\hbar^{3}}\sqrt{f_{q}(n)}h_{0}(n)!e^{-\frac{g^{2}}{2}L(0)}\sqrt{f_{q}(n)!f_{q}(n)}\times$
$\displaystyle\int_{t_{0}}^{t}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\vec{d}_{cv}\cdot\vec{E}(t^{\prime})\vec{d}_{vc}\cdot\vec{E}(r)\vec{d}_{cv}\cdot\vec{E}(s)e^{i\omega_{ex}[k(n-1)t^{\prime}-(r-s)]}dt^{\prime}drds$
$\displaystyle-\frac{i}{2\hbar^{3}}\sqrt{f_{q}(n+1)}h_{0}(n+1)!e^{-\frac{g^{2}}{2}L(1)}\sqrt{f_{q}(n+1)!f_{q}(n+1)}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\times$
$\displaystyle\vec{d}_{vc}\cdot\vec{E}(r)\vec{d}_{cv}\cdot\vec{E}(s)\vec{d}_{cv}\cdot\vec{E}(t^{\prime})e^{-i\omega_{ex}[k(n+1)(r-s)-t^{\prime}]}dt^{\prime}drds$
$\displaystyle-\frac{i}{2\hbar^{3}}\sqrt{f_{q}(n)}h_{0}(n+1)!e^{-\frac{g^{2}}{2}L(1)}\sqrt{f_{q}(n+1)!f_{q}(n+1)}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\int_{t_{0}}^{t}\times$
$\displaystyle\vec{d}_{cv}\cdot\vec{E}(r)\vec{d}_{vc}\cdot\vec{E}(s)\vec{d}_{cv}\cdot\vec{E}(t^{\prime})e^{i\omega_{ex}[k(n+1)(r-s)+t^{\prime}]}dt^{\prime}drds],$
(49)
where
$h_{i}(n)=\frac{e^{(-1)^{i}i\omega_{ex}k(n+i)(t-t_{0})}}{\omega_{ex}k(n+i)}$
and $f_{q}(n)=\sqrt{\frac{q_{n}-q^{-n}}{q-q^{-1}}}$. These equation shows that
in this conditions second order response function is equal zero. Now we can
calculate linear and nonlinear electric susceptibility of this exciton system
from this equation. Generalized linear and nonlinear absorption spectra of
this system is shown in figures (4)-(6) for different values of $q$-parameter.
In these plots, $1s$-exciton is considered. In these figures we choose
$\hbar=e=1$, $g=200\;\mu ev$ and $\omega_{ex}=1574mev$. Fig.(4) shows plots of
linear absorption spectra and Fig.(6) shows plots of nonlinear spectra. On the
other hand, 3-dimension plot of linear absorption coefficients is given in
figure (5). It is clear that changes of $q$-parameter strongly affects
absorption spectra of the system. These figures show in the presence of
$q$-values absorption of probe beam shows a complex structure: a multiple-like
absorption pattern appears with one strong peak and some side bands. Presence
of these side bands is a signature of the optical generation of an nonlinear
exciton (an exciton which expresses with $q$-deformed operator). Negative part
of the absorption spectrum demonstrates gain of the probe beam. Due to the
resonance interaction of pump with exciton transition, the gain effect comes
from the coherent energy exchange between the pump and probe beams through the
QD nonlinearity. The obtained absorption spectra are very similar to
experimental results [30]. In Ref. [30] absorption spectra of a driven charged
QD is derived experimentally. Charged QD is a nonlinear medium and is similar
to our model. Then It can be consider as a experimental test of our model.
## 5 Conclusion
$Q$-deformed description of excitons in a QD and its physical consequences was
considered. We showed that increasing the $q$-parameter will lead to increase
of splitting between peaks in the spectrum and asymmetry of spectrum. Similar
effects were observe when biexciton effects taken into account. In experiments
of QD it is shown [23] the same results are obtained in different
temperatures. Then we can associate this physical parameter as source of
$q$-deformation. The temperature dependence of emission energy of system can
be attributed to the change in the refractive index of its active medium with
temperature. We have derived the optical response of QD with $q$-deformed
exciton. As mentioned before $q$-deformed description of excitons will lead to
dependence of optical response on $q$ parameter. Hence, due to the wide range
of $q$ parameter and its effects on optical response we can consider some
parameters like temperature and interaction between the excitons which affects
the optical response of QD as sources of $q$-deformation of excitons. As
mentioned, the relation of quantum statistics of excitons in the QD and the
size of QD has been considered. Then we can consider the ratio of exciton Bohr
radius to dimension of system and exciton population as two main sources of
$q$-deformation. $Q$-deformed operator depends on total number of associated
particles of system. Therefore we can interpret $q$-deformed operator as an
operator which consists of effects of other excitations of system implicitly.
Then it is reasonable to consider this description as an effective description
which takes into account some nonlinearity in exciton system. As we saw, in
the case of interaction of light with two excitons, when $q=1$ this system
showed a two peaks spectrum. While by increasing deviation of exciton from
Bose statistics, spectrum becomes multi peak. Due to the nonlinear nature of
$q$-deformed exciton we showed that different orders of nonlinear response
function of this system can be calculated. From coincidence of obtained
results and experimental results, we can conclude that $q$-deformed
description of excitons can be a considerable model for excitons. With
comparing the obtained results in this paper with experimental ones we can
investigate the origin of this description of excitons. As pointed out the
ratio of system dimension to the Bohr radius of exciton is one of the sources
of deviation of excitons from usual boson. The obtained results are very
similar to the effects of the exciton-exciton interaction [3],[31] which is
relates to exciton population and biexciton binding energy. On the other hand,
it is shown that [1] exciton density is another source of their deviation from
ordinary bosons. To sum up we attribute the origin of $q$-deformation of the
excitons to their density, their mutual interactions, confinement size and
other parameters which cause fluctuation of optical response of the system.
$Q$-deformed description of an active medium causes that the optical
properties of system depend on the $q$-parameter. Then, it is seem that
parameters which can affect optical properties of the active medium (like
refractive index) their effects can be considered by this formulation. The
$q$-parameter can be considered as a variation parameter which its values can
be obtained from comparison of theoretical and experimental results.
Acknowledgment The authors wish to thank the Office of Graduate Studies of the
University of Isfahan and Iranian Nanotechnology initiative for their support.
## References
## References
* [1] A. S. Davydov, Theory of Molecular Excitons, (Plenum Press, NewYork, 1971).
* [2] E. Hanamura and H. Haug, Phys. Rep. 33, 209 (1977).
* [3] Y. Yamamoto, F. Tassone and H. Cao, Semiconductor Cavity Electrodynamics, (Springer-Verlag Berlin Heidelberg, 2000).
* [4] F. P. Laussy, M. M. Glazov, A. Kavokin, D. M. Whittaker and G. Malpuech, Phys. Rev. B 73, 115343 (2006).
* [5] B. Laikhtman, J. Phys: Condens. Matter 19, 295214 (2007).
* [6] M. Combescot, O. Betbeder-Matibet and F. Dubin, Phys. Rev. A 76, 033601 (2007).
* [7] H. Ichida, M. Nakayama and J. Lumin, J. Lumin. 94-95, 379 (2001).
* [8] A. Khare, Fractional Statistics and Quantum Theory, (World Scientific, Singapore, 1997).
* [9] M. P. Blencowe and N. C-Koshnick, J. Math. Phys. 42, 5713 (2001).
* [10] W. S. Dai and M. Xie, Physica A 331, 497 (2204).
* [11] Y. Shen, W. S. Dai and M. Xie, Phys. Rev. A 75, 042111(2007).
* [12] S. Chaturvedi, V. Srinivasan, Phys. Rev. A 44, 8024 (1991).
* [13] A. J. Macfarlane, J. Phys. A: Math. Gen. 22, 4581 (1989); L. C. Biedenharn, J. Phys. A: Math. Gen. 22, L873 (1989).
* [14] D. Bonatsos, J. Phys. A: Math. Gen. 25, L101 (1992).
* [15] Y. X. Liu, C. P. Sun, S. X. Yu and D. L. Zhou, Phys. Rev. A 63, 023802 (2001).
* [16] G. Rideau, Lett. Math. Phys. 24, 147(1992).
* [17] J. J. Hopfield, Phys. Rev. 112, 1555 (1958).
* [18] U. Fano. Phys. Rev. 103, 1202 (1956).
* [19] J. Katriel, Phys. Lett. A 307, 1(2003).
* [20] M. H. Naderi, R. Roknizadeh and M. Soltanolkotabi, Prog. Theor. Phys. 112, 797 (2004); ibid 112, 811 (2004).
* [21] M. Lax, Phys. Rev. 145, 110 (1966).
* [22] M. O. Scully, M. S. Zubairy, Quantum Optics, (Cambridge University Press 1997).
* [23] E. Peter, P. Senellart, D. Martrou, A. Lemaître, J. Hours, J. M. Gérard, and J. Bloch, Phys. Rev. Lett. 95, 067401 (2005).
* [24] A. Franceschetti and Y. Zhang, Phys. Rev. Lett. 100, 136805 (2008).
* [25] V. M. Axt, A. Stahl, Z. Phys. B 93, 205 (1994).
* [26] Th. Östreich, K. Schönhammer, L. J. Sham, Phys. Rev. Lett. 74, 4698 (1995).
* [27] H. Haug, S. W. Koch, Quantum Theory of The Optical and Electronic Properties of Semiconductor, 4rd edition (World Scientific, Singapore, 2004).
* [28] S. Mukamel, Principles of Nonlinear Optical Spectroscopy, (Oxford, NewYork, 1995).
* [29] R. P. Feynman, Phys. Rev. 84, 108 (1951).
* [30] X. Xu et al., arxiv:0803,0734 (Cond-Mat. mes-hall).
* [31] U. Hohenester and E. Molinari, Phys. Stat. Sol. (b) 221, 19 (2000).
Figure 1: Plots of $S(\omega)$ versus $\omega$. Parameter are choose as
$\omega=1.75\;eV$, $\omega_{ex}=1.75\;eV$, $g=200\;\mu eV$,
$\gamma_{ex}=20\;\mu eV$, $\gamma_{ph}=40\;\mu eV$, $n=1$ and
$|\alpha|^{2}=9$. Solid plot corresponds to $q=1$, nondeformed case. Dotted
one corresponds to $q=1.01$, and for dash line $q$ is equal $1.015$. Figure 2:
Plots of $S(\omega)$ versus $\omega$. Parameter are choose as
$\omega=1.75\;eV$, $\omega_{ex}=1.75\;eV$, $g=200\;\mu eV$,
$\gamma_{ex}=20\;\mu eV$, $\gamma_{ph}=40\;\mu eV$, $n=1$ and
$|\alpha|^{2}=9$. In all figures we have $q=1$. Solid line corresponds to case
$s=1$. In dotted one we have $s=1.007$ and for dash line $s=1.01$. Figure 3:
Plots of $S(\omega)$ versus $\omega$. Parameter are choose as
$\omega=1.75\;eV$, $\omega_{ex_{1}}=1.75\;eV$, $\omega_{ex_{2}}=1.77\;eV$,
$g=200\;\mu eV$, $\gamma_{ex_{1}}=\gamma_{ex_{2}}=200\;\mu eV$,
$\gamma_{ph}=45\;\mu eV$, $n_{1}=1$, $n_{2}=1$ and $|\alpha|^{2}=9$. Dotted
line corresponds to nondeformed case $q_{1},q_{2}=1$. For solid line
$q_{1},q_{2}=1.04$. In the case of dashed line $q_{1},q_{2}=1.08$. Figure 4:
Plots of spectrum absorption versus $\omega$. We consider 1s-exciton and
Parameter are choose as $\hbar=e=1$, $g=200\;\mu ev$ and
$\omega_{ex}=1574\;mev$. Solid plot corresponds to nondeformed case $q=1$. For
dotted one $q=1.01$ and in dash one $q=0.99$. Figure 5: 3D-Plots of spectrum
absorption versus $\omega$ and deformation parameter $q$. Physical parameter
are the same as Fig.(4). Figure 6: Plots of nonlinear spectrum absorption
versus $\omega$. We consider 1s-exciton and Parameters are choose as
$\hbar=e=1$, $g=200\;\mu ev$ and $\omega_{ex}=1574\;mev$ Solid plot
corresponds to nondeformed case $q=1$. In dotted plot $q=1.01$. In dash plot
$q=0.99$.
|
arxiv-papers
| 2011-12-11T11:12:54 |
2024-09-04T02:49:25.179194
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "M. Bagheri Harouni, R. Roknizadeh, M. H. Naderi",
"submitter": "Malek Bagheri",
"url": "https://arxiv.org/abs/1112.2346"
}
|
1112.2531
|
# Double-relativistic-electron-shell laser proton acceleration
Yongsheng Huang http://www.anianet.com/adward huangyongs@gmail.com China
Institute of Atomic Energy, Beijing 102413, China. Naiyan Wang Xiuzhang Tang
Yijin Shi China Institute of Atomic Energy, Beijing 102413, China. Yan
Xueqing Institute of Heavy Ion Physics, Peking University, Beijing 100871,
China Zhang Shan Beijing Normal University, Beijing 100875, China
###### Abstract
A new laser-proton acceleration structure combined by two relativistic
electron shells, a suprathermal electron shell and a thermal electron cloud is
proposed for $a\gtrapprox 80\sigma_{0}$, where a is the normalized laser field
and $\sigma_{0}$ is the normalized plasma surface density. In the new region,
a uniform energy distribution of several GeV and a monoenergetic hundreds-of-
MeV proton beam have been obtained for $a=39.5$. The first relativistic
electron shell maintains opaque for incident laser pulse in the whole process.
A monoenergetic electron beam has been generated with energy hundreds of MeV
and charge of hundreds of pC. It is proposed a stirring solution for
relativistic laser-particle acceleration.
###### pacs:
52.38.Kd,41.75.Jv,52.40.Kh,52.65.-y
Laser-ion acceleration has been an international research focusMakoTajima ;
Machnisms ; Esirkepov , however it is still a challenge to obtain mono-
energetic proton beams larger than $100\mathrm{MeV}$. Although the field in
the laser-plasma acceleration is three to four orders higher than that of the
classic accelerators, it decreases to zero quickly in several pulse durations
for target normal sheath acceleration (TNSA)Machnisms for
$a\ll\sigma_{0}=\frac{n_{e}l_{e}}{n_{c}\lambda}$, where $a=eE_{l}/m\omega c$,
$E_{l}$ is the electric field of the laser pulse, $\omega$ is the laser
frequency, $e$ is the elementary charge, $m$ is the electron mass, $c$ is the
light velocity, $n_{0}$ is the initial plasma density, $n_{c}$ is the critical
density, $\lambda$ is the wave length, $n_{e}$ is the electron density,
$l_{e}$ is the thickness of the electron shell. As a promising method to
generate relativistic mono-energetic protons, radiation pressure acceleration
(RPA) has attracted more attentionEsirkepov ; Henig ; YanxqPRL ; unlimitedRPA
and becomes dominant in the interaction of the ultra-intense laser pulse with
thin foils if $a\approx\sigma_{0}$. Even in the unlimited ion
accelerationunlimitedRPA , only the ions trapped in the electron shell can
obtain efficient acceleration, therefore, the total charge of the ion beam is
quite limited due to the transverse expansionunlimitedRPA . Although in RPA
region, the energy dispersion will become worse with time.
Fortunately, for $a\gtrapprox 80\sigma_{0}$, in the relativistic case, a new
acceleration region appears: double relativistic electron shells come into
being. The ions between the two electron shells will be accelerated most
efficiently and obtain a uniform energy distribution. The first electron shell
is ultra-relativistic and is totally separated from the ions. The second
electron shell comes into being in the potential well induced by the electron
recirculation and will also be relativistic. It is in the ion region and
follows the ion front and forms a potential well which traps energetic ions
and accelerates them to be quasi-mono-energetic and relativistic. Following
the second electron shell, a suprathermal electron beam comes into being and
induces another potential well which can also trap lots of ions and accelerate
them to obtain a monoenergetic relativistic one. On the whole, the maximum ion
energy can reach several GeV and a relativistic monoenergetic ion beam with
relative energy dispersion smaller then $5\%$ can be obtained.
As a ultraintense laser pulse is shot on a ultra-thin plasma foil for
$a\gtrapprox 80\sigma_{0}$, the electron shell is compressed to ultra-high
density and pushed forward to be separated from the ion shell totally, and
gains ultra-relativistic energy that can make sure it opaque for the laser
pulse. In the whole process, the first electron shell keeps opaque for the
incident laser pulse and is pushed by it continuously. According to Eq.
(22)unlimitedRPA , it can be satisfied that the opaqueness condition of
electron shell for laser in the acceleration:
$a_{0}\leq\pi(\gamma_{e}+p_{e})\hat{n}_{e}\hat{l}_{e},$ (1)
since $d\ln(p\hat{n}_{e}\hat{l}_{e})/dt>0$, as pointed by Bulanov, $p\propto
t^{1/3}$ is the normalized electron momentum, and $n_{e}l_{e}=n_{0}l_{0}$ in
the no transverse expansion case, where the electron density $n_{e},n_{0}$ are
normalized by $n_{c}$ and $l_{e}$ is normalized by $\lambda$,
$a_{0}=eE_{0}/m\omega_{0}c$, $E_{0}$ is the electric field of the laser pulse,
$\omega_{0}$ is the laser frequency, $e$ is the elementary charge, $m$ is the
electron mass, $c$ is the light velocity, $n_{0}$ is the initial plasma
density, $n_{c}$ is the critical density, $\lambda$ is the wave length. In the
ultra-relativistic case, the r.h.s. of Eq. (1) is approximate $2\pi
p\hat{n}_{e}\hat{l}_{e}$.
For $a=39.5$ and $\hat{n}_{e}\hat{l}_{e}=49\times 0.01=0.49$, $a\approx
81\hat{n}_{e}\hat{l}_{e}$. When the laser pulse interacts with the plasma
shell, the electrons are compressed to high density and pushed forward to be
separated from the ion shell. As shown by Figure 1 (c) and (d), at
$t=25\mathrm{fs}$, the normalized momentum of the electron shell reaches
$20-100$. With Eq. (1), the compressed high-density electron shell is opaque
for the laser pulse as shown by Figure 1 (f). The electrons will be
accelerated efficiently and continuously by the radiation pressure of the
laser.
Figure 1: (Color online) Simulation results by one-dimensional VORPAL at
$t=25$fs: the first opaque relativistic electron shell forms and is totally
separated from ions. (a) and (c): the phase space distribution of ions and
electrons. The electron shell is totally separated from the ions and is
relativistic and opaque for laser pulse. (b) and (d): the density distribution
of the ions and electrons with normalized momentum. The ion energy
distribution is uniform. The electron energy distribution contains several
monoenergetic ones. (e): the longitudinal space charge separation field, which
is similar to the field in a capacity, is uniform and is $9.95\times
10^{12}$V/m. (f): the potential and laser field. The electron shell is
relativistic and opaque for the laser pulse.
This high-density relativistic electron shell is called the first relativistic
electron shell. Between the first electron shell and the ions, a uniform
space-charge separation field forms and accelerates the ions at the ion front
and drags the electrons as shown by Figure 1 (e). The ”capacity” field is
decided by the surface density of the electron shell:
$E_{cap}=\frac{en_{e}l_{e}}{\epsilon_{0}},$ (2)
For $n_{e}=5.5\times 10^{22}\mathrm{/cm^{3}}$, $l_{e}=10\mathrm{nm}$, the
stable field is $9.95\times 10^{12}\mathrm{V/m}$, which accelerates the ions
at the rear of the ion shell continuously. Therefore the maximum ion energy is
proportional to the acceleration length, $d_{acc}$,
$E_{i}=E_{cap}d_{acc}(\mathrm{eV}),$ (3)
before the electron shell breaks up.
After several hundreds of femtoseconds, some electrons leak out from the
electron shell continuously and move backward and round again and follow the
ion front, however, they can not catch the ion front. Figure 2 (c), (e) and
(f) shows the electron recirculation, the decrease of the separation field,
and the formation of the potential well for electrons at $t=450\mathrm{fs}$
respectively. The normalized maximum electron momentum reaches about $500$ as
shown by Figure 2 (d). Figure 2 (f) shows that the deepness of the potential
well for electrons increases about to $0.3\mathrm{GeV}$. From Figure 2 (a) and
(b), the maximum ion energy is about $500\mathrm{MeV}$. The acceleration
length is about $65\mathrm{\mu m}$. A monoenergetic electron beam of
$186\mathrm{MeV}$ and $132$pC is obtained as shown by Figure 2 (d). It is
obvious that the first relativistic electron shell is still opaque for the
laser pulse.
Figure 2: (Color online) Simulation results by one-dimensional VORPAL at
$t=450$fs: electron recirculation begins and generates a potential well for
electrons, where $p_{i},p_{e}$ is momentum of ion and electron respectively,
$M$ is the proton mass. (a) and (c): the phase space of ions and electrons
respectively. The electron recirculation happens and some of them move in the
opposite direction relativistically. (b) and (d): the energy distribution of
ions and electrons respectively. The first electron shell is accelerated most
efficiently and continuously. (e) the longitudinal field decreases due to the
electron recirculation. However it is still uniform between the ion front and
the first electron shell. (f): a potential hill for ions and a potential well
for electrons forms due to the electron recirculation. The first relativistic
electron shell maintains opaque for the laser pulse.
As shown by Figure 3 (d), the recirculating electrons cumulate and drive up
the electron potential. In front of and behind the accumulating electrons, two
potential wells are forming for electrons. The potential well I still traps
and accelerates the accumulating electrons to generate the second relativistic
high-density electron shell. In the potential well II, some electrons at the
end of the second relativistic electron shell drop into it and will be trapped
and accelerated to form a suprathermal electron shell as shown by Figure 5
(d). At the same time and at the local position of the second relativistic
electron shell, potential well III for ions traps lots of ions and accelerates
them to be relativistic and monoenergetic.
Figure 3: (Color online) Simulation results by one-dimensional VORPAL at
$t=1.25$ps: the second relativistic electron shell forms and traps ions to be
accelerated to relativistic. (a): the second relativistic electron shell forms
in the potential well. It drives up the potential and forms two potential well
I and II for electrons in front of and behind itself, a potential well III for
ions at the local position of itself. It is shown clearly in (d). (b): the
number density distribution of electrons. (c): the electron recirculation
decreases the longitudinal field continuously. It is still uniform between the
ion front and the first electron shell. Figure 4: (Color online) Simulation
results by one-dimensional VORPAL at $t=3.075$ps: the suprathermal electron
shell forms and traps ions to obtain a flap-top $3-47$MeV energy distribution.
(a) and (c): the phase space of ions and electrons. The suprathermal electron
shell forms in the potential well II and then potential V for electrons and
potential IV are generated. (b) and (d) the energy distribution of ions and
electrons. A monoenergy ion beam with energy of $627$MeV is obtained in the
second relativistic electron shell. (e): the longitudinal field induced by the
double electron shell. (f): the potential IV for ions is induced by the
suprathermal electron shell and traps ions and improves the energy dispersion
of the $3-47$MeV ion beam. The potential V can trap slow electrons and
thermalize them to obtain thermal electron cloud. Potential III is nearly
filled to be flat and the energy dispersion of the 627MeV monoenergetic ion
beam will become worse.
With time, the slow recirculating electrons can also be trapped in potential
well V as shown in Figure 4(f) and the third suprathermal electron shell forms
in the potential well II as shown by Figure 5. At the position of the shell,
potential well IV traps the ions and accelerates them to obtain a quasi-
monoenergetic distribution of $171\pm 10\mathrm{MeV}$. Behind the shell, a
potential well for electrons traps them and a thermal electron cloud is
generated. The ions between the double relativistic electron shell have a
uniform distribution from $1\mathrm{GeV}$ to $2.18\mathrm{GeV}$. Trapped by
the second relativistic electron shell, the maximum energy reach
$981\mathrm{MeV}$. As shown in Figure 4 (f), the potential well III has been
filled up nearly, then the energy dispersion will become worse. The ions with
larger energy will coast down the following potential slope and get into the
ion beam between the double relativistic electron shell. The ion number has a
steep descent for the energy larger than $981\mathrm{MeV}$ and has a slow drop
for the energy smaller than $981\mathrm{MeV}$. In the electron energy
distribution, there is a monoenergetic one of $385\pm 10\mathrm{MeV}$ and
$163\mathrm{pC}$, a ultra-relativistic one of $1\mathrm{GeV}$ and a Maxwellian
one which contains the thermal electron cloud and the suprathermal electron
shell. As shown by Figure 5(c) and (f), the ion front is between the double
electron shell. The ions between the double electron shell coast down the
potential slope and obtain relativistic energy as shown in Figure 5(h) and
(b).
Figure 5: (Color online) Simulation results by one-dimensional VORPAL at
$t=4.025$ps: the thermal electron cloud and the suprathermal electron shell
come into being. (a) and (d): the phase-space of ions and electrons
respectively. The electrons contain four main parts: the double relativistic
electron shells, the suprathermal electron shell, the thermal electron cloud.
(b) and (e): the energy distribution of ions and electrons respectively. A
monoenergetic ion beam with energy of $171\pm 10$MeV is obtained by the
suprathermal electron shell. Trapped and accelerated by the second
relativistic electron shell, the ion energy distribution drops down at
$981$MeV. (c) and (f) the number density of ions and electrons respectively.
(g) the longitudinal field. (h) the potential for ions and electrons. (i) the
laser pulse field. The first electron shell maintains opaque for laser pulse.
In conclusion, the double relativistic electron shells, the suprathermal
electron shell and the thermal electron cloud induce a new region of laser
particle acceleration. In the process, several potential wells for ions and
electrons are generated. On the whole, the double relativistic electron shells
induce two relativistic platforms of the ion energy distribution. The
suprathermal electron shell traps and accelerates a monoenergetic ion beam
with several hundreds of MeV, whose relative energy dispersion is near $5\%$.
Together with the thermal electron cloud, a thermal Maxwellian ion beam has
been obtained.
###### Acknowledgements.
The authors would like to thank Dr. Hong-Yu Wang for useful discussion. The
computation was carried out at the HSCC of Beijing Normal University. This
work was supported by the Key Project of Chinese National Programs for
Fundamental Research (973 Program) under contract No. $2011CB808104$ and the
Chinese National Natural Science Foundation under contract No. $10834008$.
## References
* (1) F. Mako and T. Tajima, Phys. Fluids 27, 1815 (1984).
* (2) H. Schwoerer, S. Pfotenhauer, O. Jackel, et al., Nature 439, 445 (2006). electronacc
* (3) T. Esirkepov, M. Borghesi, S. V. Bulanov, G. Mourou, and T. Tajima, Phys. Rev. Lett. 92, 175003 (2004).
* (4) A. Henig, S. Steinke, M. Schn rer, et al., Phys. Rev. Lett. 103, 245003 (2009).
* (5) H. Schwoerer, S. Pfotenhauer, O. Jackel, et al., Nature 439, 445 (2006).
* (6) X. Q. Yan, C. Lin, Z.M. Sheng, et al., Phys. Rev. Lett. 100, 135003 (2008).
* (7) S. V. Bulanov, E. Yu. Echkina, T. Zh. Esirkepov, et al., Phys. Rev. Lett. 104, 135003 (2010).
* (8) X. Q. Yan, T. Tajima, M. Hegelich, et al., Appl. Phys. B, 98 711-721 (2010).
* (9) Y. S. Huang,
* (10) T. P. Yu, A. Pukhov, G. Shvets, and M. Chen, Phys. Rev. Lett. 105, 065002 (2010).
* (11) A. Macchi, F. Cattani, T. V. Liseykina and F. Cornolti, Phys. Rev. Lett. 94, 165003 (2005).
|
arxiv-papers
| 2011-12-12T12:41:42 |
2024-09-04T02:49:25.190632
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Yongsheng Huang, Naiyan Wang, Xiuzhang Tang, Yijin Shi, Zhang Shan",
"submitter": "Yongsheng Huang",
"url": "https://arxiv.org/abs/1112.2531"
}
|
1112.2570
|
# Minkowski momentum of an MHD wave
Tadas K Nakamura CFAAS, Fukui Prefectural University, 910-1195 Fukui, Japan
###### Abstract
The momentum of an MHD wave has been examined from the view point of the
electromagnetic momentum expression derived by Minkowski. Basic calculations
show that the Minkowski momentum is the sum of electromagnetic momentum and
the momentum of the medium, as proposed in some of the past literature. The
result has been explicitly confirmed by an example of an MHD wave, whose
dynamics can be easily and precisely calculated from basic equations. The
example of MHD wave also demonstrates the possiblility to construct a
symmetric energy-momentum tensor based on the Minkowski momentum.
###### pacs:
03.50.De,42.20.Jb,52.35.Bj
## 1 Introduction
The Minkowski-Abraham controversy has been discussed by a number of authors
over a hundred years. Minkowski [1] proposed the electromagnetic momentum
density in a dielectric medium must be $\mathbf{D}\times\mathbf{B}$, and
Abraham [2, 3] proposed $\mathbf{E}\times\mathbf{H}$ for that (in the present
letter symbols have conventional meanings, e.g., $\mathbf{E}$ = electric
field, unless otherwise stated). There have been published numerous papers on
this problem both theoretically and experimentally, but the final conclusion
is still yet to come; papers are still beeing published in this century (see,
e.g., [4] for a review).
Several authors [5, 4, 6] pointed that the electromagnetic field inevitably
affect the dynamics of the medium to change its energy-momentum, and
therefore, the energy-momentum of an electromagnetic wave must include the
contribution of the medium. In the present letter we show that the Minkowski
momentum is the sum of electromagnetic momentum and the momentum of the
medium. Feigel [5] obtained a similar result based on the Noether’s theorem
using Lagrangian formulation. Compared to his elegant approach, the
calculation here is rather a down-to-eath type, which is more closer to the
Minkowski’s original derivation. This approach is less elegant, however,
easier to understand its meaning intuitively.
Perhaps the largest weak point of the Minkowski momentum is the fact that the
four dimensional energy-momentum tensor does not become symmetric with this
momentum, which means the violation of angular momentum conservation (see,
e.g. [7]). Most of the past literature argued the legitimacy of the momentum
part of the tensor in this point. Here, in contrast, we elucidate the
possibility to alter the energy part to make a symmetric tensor; provided the
momentum part of the Minkowski energy-momentum tensor includes the momentum of
medium, the same should be true for the energy part. To treat it in a
relativistically consistent way, the mass flux must be included in the energy
flux even in the non-relativistic regime. The energy-momentum tensor with
Minkowski momentum can become symmetric when the mass flux is taken into
account.
The consideration stated above is confirmed by an example of an MHD wave in a
collisionless magnetized plasma. Usually the behavior of an ordinary medium is
complicated and need to calculate microscopic states of molecules, which is
difficult to solve exactly. A collisionless plasma is, in contrast, easy to
calculate its response to the electromagnetic field from the classical basic
equations (Maxwell equations and Newtonian mechanics). Here in this short
letter we use the MHD approximation, however, if one wishes it is possible to
derive an exact solution of the basic equation system to confirm the result.
The result agree with the “frozen-in” of a magnetized plasma, which has been
confirmed by a wide variety of experimental and observational facts.
## 2 Basics
Microscopic Ampere’s equation in a medium is
$-\varepsilon_{0}\frac{\partial\mathbf{E}}{\partial
t}+\mu_{0}^{-1}\nabla\times\mathbf{B}=\mathbf{J}$ (1)
Suppose there is no external current, and $\mathbf{J}$ consists of the
polarization current $\mathbf{J}_{P}$ and magnetization current
$\mathbf{J}_{M}$, which are generated in response to the electric field
$\mathbf{E}$ and magnetic field $\mathbf{B}$ respectively. We introduce the
polarization vector $\mathbf{P}$ and magnetization vector $\mathbf{M}$ such
that
$\frac{\partial}{\partial
t}\mathbf{P=\mathbf{J}}_{P}\,,\;\;\nabla\times\mathbf{M}=\mathbf{J}_{M}$ (2)
In this context, $\mathbf{P}$ and $\mathbf{M}$ should be understood as
convenient mathematical expressions to represent the response of the medium to
the electromagnetic field, rather than real physical entities.
When averaged over a microscopically large but macroscopically small volume,
$\bar{\mathbf{P}}$ and $\bar{\mathbf{M}}$ are assumed to have simple linear
relations to the electromagnetic fields as
$\bar{\mathbf{P}}=\chi_{P}\bar{\mathbf{E}}\,,\;\;\bar{\mathbf{M}}=\chi_{M}\bar{\mathbf{B}}\,.$
(3)
The linear coefficients $\chi_{P}$ and $\chi_{M}$ are matrices in general
because the medium may not be isotropic (as in our example of magnetized
plasmas). The fields $\bar{\mathbf{D}}$ and $\bar{\mathbf{H}}$ are then
defined as macroscopic quantities as
$\bar{\mathbf{D}}=\varepsilon_{0}\bar{\mathbf{E}}+\bar{\mathbf{P}}=\varepsilon\bar{\mathbf{E}}\,,\;\;\bar{\mathbf{H}}=\mu_{0}^{-1}\bar{\mathbf{B}}+\bar{\mathbf{M}}=\mu^{-1}\bar{\mathbf{B}}$
(4)
## 3 Minkowski Momentum
The momentum of microscopic electromagnetic field is
$\varepsilon_{0}\mathbf{E}\times\mathbf{B}$ and its conservation law is
$\frac{\partial}{\partial
t}(\varepsilon_{0}\mathbf{E}\times\mathbf{B})+\nabla\cdot
T+(\mathbf{J}_{P}+\mathbf{J}_{M})\times\mathbf{B}+(Q_{P}+Q_{M})\mathbf{E}=0\,,$
(5)
where $T$ is the Maxwell stress tensor and we denote $\partial T_{ij}/\partial
x_{i}=(\nabla\cdot T)_{j}$ in short. The polarization/magnetization charge
$Q_{P}$ and $Q_{M}$ are the result of polarization/magnetization current
($\partial Q_{P,M}/\partial t=\nabla\mathbf{J}_{P,M}$). The charge due to
magnetization current vanishes when averaged, $\bar{Q}_{M}=0$ since
$\bar{\mathbf{J}}_{M}$ satisfies (2).
The third and fourth term of (5) are the Lorentz and Coulomb force acting on
the medium, and therefore, it can expressed by the momentum change of the
medium.
$(\mathbf{J}_{p}+\mathbf{J}_{M})\times\mathbf{B}+(Q_{P}+Q_{M})\mathbf{E}=\frac{\partial}{\partial
t}\mathbf{g}+\nabla\cdot T_{M}\,,$ (6)
where $\mathbf{g}$ and $T_{M}$ are the momentum density and stress tensor of
the medium. It should be noted that the right hand side of the above
expression has mathematical ambiguity. If we define new values of
momentum/stress by $\mathbf{g}^{\prime}=\mathbf{g}+\mathbf{a}$ and
$T^{\prime}=T+G$ with arbitrary vector $\mathbf{a}$ and tensor $G$ that
satisfy $\partial\mathbf{a}/\partial t=\nabla G=0$, they also satisfy the
above equation. Therefore, $\mathbf{g}$ and $T$ do not necessarily have to be
the total momentum/stress of the medium. For example, the medium may contain a
part that does not interact with the electromagnetic field, and such part
causes this ambiguity.
From (4) we obtain
$\displaystyle\bar{\mathbf{J}}_{P}\times\bar{\mathbf{B}}$ $\displaystyle=$
$\displaystyle\frac{\partial}{\partial
t}(\bar{\mathbf{P}}\times\bar{\mathbf{B}})+\chi_{P}\left[(\bar{\mathbf{E}}\cdot\nabla)\bar{\mathbf{E}}+\frac{1}{2}\nabla\bar{\mathbf{E}}^{2}+\bar{\mathbf{E}}(\nabla\bar{\mathbf{E}})\right]\,,$
(7) $\displaystyle\bar{\mathbf{J}}_{M}\times\bar{\mathbf{B}}$ $\displaystyle=$
$\displaystyle\chi_{M}\left[(\bar{\mathbf{B}}\bigtriangledown)\bar{\mathbf{B}}+\frac{1}{2}\nabla\bar{\mathbf{B}}^{2}\right]\,.$
Here we neglected cross terms of fluctuation in averaging as usually done in
this kind of calculation, e.g.,
$\overline{\mathbf{P}\times\mathbf{B}}=\bar{\mathbf{P}}\times\bar{\mathbf{B}}$.
Combining (5), (6) and (7) we obtain
$\frac{\partial}{\partial
t}(\varepsilon_{0}\bar{\mathbf{E}}\times\bar{\mathbf{B}}+\bar{\mathbf{g}})+\nabla\cdot(\bar{T}+\bar{T}_{M})=\frac{\partial}{\partial
t}(\bar{\mathbf{D}}\times\bar{\mathbf{B}})+\nabla\cdot\bar{T}^{\prime}=0\,,$
(8)
where $\bar{T}^{\prime}$ is the stress tensor in a dielectric medium defined
as
$\bar{T}^{\prime}_{ij}=\bar{E}_{i}\bar{D}_{j}+\mu_{0}^{-1}\bar{B}_{i}\bar{H}_{j}-\frac{1}{2}\delta_{ij}(\bar{\mathbf{E}}\bar{\cdot\mathbf{D}}+\bar{\mathbf{B}}\cdot\bar{\mathbf{H}})\,,$
(9)
which is the sum of the fluxes of electromagnetic momentum and momentum
carried by the medium.
Now that we understand the Minkowski momentum includes the part of the medium
then so should be for the energy. The energy is equivalent to mass in
relativity, thus the energy flux of the medium must include mass, and then the
flux may take the form of $\bar{\mathbf{D}}\times\bar{\mathbf{B}}$ to make the
energy momentum symmetric. We will check this for the case of an Alfven wave
in the following.
## 4 MHD wave
Let us confirm the above discussion with an example of an MHD wave. Suppose a
linearly polarized one dimensional ($\partial/\partial x=\partial/\partial
y=0$) MHD wave (Alfven wave in this case) propagating in the $z$ direction,
which is the direction of the background magnetic field: $B_{0}=B_{Z}$. The
wave amplitude is small enough for linear approximation, and the plasma
velocity is low enough for non-relativistic approximation. Also we assume
“cold plasma limit”, which means the thermal energy of plasma particles is
negligibly small.
The wave has an electric field perpendicular to its propagation direction, and
we take the $x$ axis in this electric field direction. Then the current is
also in the $x$ direction, whereas the magnetic perturbation and the plasma
velocity is in the $y$ direction (see Appendix). The magnetic field created by
cyclotron motion of plasma particles is negligible in a cold plasma limit, and
thus we treat the case with $\mu=\mu_{0}$ hereafter.
The current is the polarization current due to the temporal change of the
electric field, which is
$J_{z}=\frac{\mu_{0}}{V_{A}^{2}}\frac{\partial}{\partial t}E_{x}$ (10)
where $V_{A}$ is th Alfven speed defined by $V_{A}=B_{0}/\sqrt{\mu_{0}\rho}$
with $\rho$ being the mass density of the plasma. From (2) and the above
expression we obtain
$D_{x}=\varepsilon_{0}\left(1+\frac{c^{2}}{V_{A}^{2}}\right)E_{x}$ (11)
The plasma frozen-in condition $\mathbf{E}+\mathbf{v}\times\mathbf{B}=0$ means
that the plasma is moving in the $y$ direction with the $E\times B$ drift
speed as
$v_{y}=\frac{E_{x}}{B_{0}}\,.$ (12)
Then the momentum carried by the plasma particles is
$\rho v_{y}=\frac{\mu_{0}B_{0}}{V_{A}^{2}}E_{x}\,.$ (13)
The $y$ component of the Minkowski momentum can be calculated from (11) and
(13), which is
$(\mathbf{D}\times\mathbf{B})_{y}=D_{x}B{}_{0}=\varepsilon_{0}E_{x}B_{0}+\rho
v_{y}\,.$ (14)
The above expression means the Minkowski momentum is the sum of
electromagnetic momentum and momentum of the plasma particles as long as the
frozen-in condition is satisfied.
When multiplied by $c^{2}$, the first term of the right hand side of (14)
becomes the electromagnetic energy flux in the $y$ direction. The second term
becomes the relativistic energy comes from the rest mass, which is the
predominant energy flux in the non-relativistic limit here; thermal or kinetic
energy flux is negligible. The balance equation of the energy-momentum tensor
is in a derivative form, and therefore, it has an ambiguity as discussed below
(6) for the momentum. The same is true for the energy, and its conservation
also holds when we add the mass density and mass flux ($\rho,\rho\mathbf{v})$,
since $\partial\rho/\partial t+\nabla(\rho\mathbf{v})=0$. The energy momentum
tensor becomes symmetric with these terms.
## 5 Concluding Remarks
Momentum carried by an electromagnetic wave has been examined with an example
of an MHD wave. It has been shown that the total momentum (electromagnetic
momentum plus momentum of the medium) is expressed by
$\mathbf{D}\times\mathbf{B}$ as proposed by Minkowski. This result is based on
the very simple and basic two properties of an MHD plasma, the frozen-in
condition (12) and polarization current (10), namely. It would not be
exaggeration if one says the whole kingdom of MHD plasma physics would fall if
these two basic properties were wrong.
Here in this letter we examined a simplest case of an parallel (to the
$\mathbf{B}$ field) propagating MHD wave, but similar calculations can be done
for more complicated plasma waves to confirm the result here. A collisionless
plasma contains a wide variety of wave phenomena, and the properties of waves
can be precisely calculated at least in the linear limit. Calculation of the
Minkowski momentum for various plasma waves would be a good exercise to
understand the Abraham-Minkowski controversy.
The drawback of the Minkowski momentum has been believed that the momentum
fails to form a symmetric four dimensional energy-momentum tensor when coupled
with the Poyinting flux; an asymmetric energy-momentum tensor means the
violation of angular momentum conservation. This difficulty can be overcome
when we include the mass flux as a part of energy flux, which is reasonable
from the relativistic point of view. The energy-momentum tensor can be
symmetric as we have examined with an MHD wave here.
What we have shown in the present letter is that the Minkowski momentum can be
self consistent description of the total momentum of an electromagnetic wave
in a polarizable medium. This does not necessarily mean the Abraham momentum
is wrong and inconsistent; it might be possible to give Abraham momentum
another appropriate meaning to make it consistent. For example, Barnett [8]
recently argued both Abraham and Minkowski momentum can be consistent when we
interpret the former as kinetic momentum and latter as canonical momentum. It
is out of our scope here to examine this argument, however, it should be noted
the legitimacy of the Minkowski momentum does not automatically exclude the
validity of the Abraham momentum.
## Appendix
This appendix is to derive (10) and (12) in a shortest way for a physicist not
familiar with plasma physics. For further information, see any textbook on
plasma physics, e.g., [9, 10]. Note that many books derive Alfven waves from
the MHD equations, which is different from the derivation here; of course the
result is the same.
Suppose a plasma consists of equal number of protons and electrons in a
uniform magnetic field, which is in the $z$ direction of Cartesian
coordinates. The plasma response to the electromagnetic field can be expressed
by $\mathbf{P}$ only and we do not need the magnetization current for our
calculation. Therefore we can set $\mathbf{H}=\mu_{0}^{-1}\mathbf{B}$ here. We
assume an MHD wave described above is propagating in this plasma.
Let us denote a vector in the $xy$ plane by a complex number as
$A=A_{x}+iA_{y}$. Then the wave electric filed in the $x$ direction is denoted
as
$E(t)=\frac{E_{0}}{2}(e^{-i\omega t}+e^{\omega t})\,.$ (15)
The equation of motion of a plasma particle in the $xy$ plane is written as
$\frac{dv}{dt}=i\Omega v+\frac{e}{m}E(t)\,,$ (16)
where $\Omega=eB/m$ is the gyro frequency. We include the sign of the charge
in $\Omega$, thus $\Omega$ is positive/negative for a proton/electron. The
above equation can be directly solved as
$v=v_{0}e^{i\Omega t}+\frac{eE_{0}}{2m}\left(\frac{e^{-i\omega
t}}{\Omega-\omega}+\frac{e^{i\omega t}}{\Omega+\omega}\right)\,,$ (17)
where $v_{0}$ is the integration constant.
Now we assume the wave frequency is much smaller than the gyro frequency
($\omega\ll\Omega$), which is true for most of MHD waves. Then the effect of
the first term in (17) will be averaged out for MHD time scale; we do not pay
attention to this term hereafter.
The rest of the motion is called “drift” in plasma physics. The drift velocity
$v_{d}$ can be expanded as
$v_{d}=\frac{E_{0}}{B}\left(i\cos\omega t+\frac{\omega}{\Omega}\sin\omega
t+\cdots\right)\,.$ (18)
The first term of (18) is called $E\times B$ drift; protons and electrons
drift in the same direction with the same speed with this drift. This term is
pure imaginary, which means the drift is in the $y$ direction. This drift
gives the predominant motion of the bulk plasma as in (12), however, it does
not cause a current because both protons and electrons have the same drift
velocity. What contribute to a current is the second term of (18), which is
called the polarization drift. Since this term contains $\Omega^{-1}$ factor,
protons and electrons moves in the opposite direction with different speed.
The electron gyro frequency is much larger than that of protons, therefore,
protons predominantly carry currents. The drift direction is the same as the
electric field since it is pure real, and the drift speed is proportional to
time derivative of the field because of the factor $\omega$ and $\cos\omega
t\rightarrow\sin\omega t$. Multiplying the proton’s second term of (18) with
the number density and charge, and replacing the factor $\omega$ and
$\cos\omega t\rightarrow\sin\omega t$ by the time derivative, we obtain (18).
From the Maxwell’s equation we have
$\nabla\times\nabla\times\mathbf{E}=c^{-2}\partial^{2}\mathbf{E}/\partial^{2}t+\mu_{0}\mathbf{J}\,.$
(19)
When we assume the wave propagation is in the $z$ direction
($\partial/\partial x=\partial/\partial y=0$) and use (10), we obtain the
propagation equation of an MHD wave (Alfven wave) as
$\left(\frac{1}{c^{2}}+\frac{1}{V_{A}^{2}}\right)\frac{\partial^{2}}{\partial
t^{2}}E_{x}-\frac{\partial^{2}}{\partial z^{2}}E_{x}=0\,.$ (20)
The Alfven speed $V_{A}$ is often much smaller than the speed of light in
space and laboratory plasmas. We obtain an wave propagating with the Alfven
speed $V_{A}$ when we ignore the $1/c^{2}$ term in the above expression.
## References
## References
* [1] H. Minkowski. Math. Ann, 68:472, 1910.
* [2] M. Abraham. Rend. Pal, 28:1, 1909.
* [3] M. Abraham. Rend. Pal, 30:33, 1910.
* [4] R.N.C. Pfeifer, T.A. Nieminen, N.R. Heckenberg, and H. Rubinsztein-Dunlop. Colloquium: Momentum of an electromagnetic wave in dielectric media. Reviews of Modern Physics, 79(4):1197, 2007.
* [5] A. Feigel. Quantum vacuum contribution to the momentum of dielectric media. Physical review letters, 92(2):20404, 2004.
* [6] R N C Pfeifer, T A Nieminen, N R Heckenberg, and H Rubinsztein-Dunlop. Constraining Validity of the Minkowski Energy-Momentum Tensor. Physical Review A, 79(2):023813, 2009.
* [7] J. D. Jackson. Cassical Electrodynamics. Wiley & Sons, 1962.
* [8] Stephen M Barnett. Resolution of the Abraham-Minkowski Dilemma. Physical Review Letters, 104(7):070401, 2010.
* [9] F. Chen. Introduction to Plasma Physics. Plenum Press, 1974.
* [10] R. Dendy. Plasma Dynamics. Oxford University Press, 1990.
|
arxiv-papers
| 2011-12-09T10:21:33 |
2024-09-04T02:49:25.197093
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Tadas K Nakamura",
"submitter": "Tadas Nakamura",
"url": "https://arxiv.org/abs/1112.2570"
}
|
1112.2678
|
# Self-similar solutions of viscous and resistive ADAFs with thermal
conduction
Kazem Faghei
###### Abstract
We have studied the effects of thermal conduction on the structure of viscous
and resistive advection-dominated accretion flows (ADAFs). The importance of
thermal conduction on hot accretion flow is confirmed by observations of hot
gas that surrounds Sgr A∗ and a few other nearby galactic nuclei. In this
research, thermal conduction is studied by a saturated form of it, as is
appropriated for weakly-collisional systems. It is assumed the viscosity and
the magnetic diffusivity are due to turbulence and dissipation in the flow.
The viscosity also is due to angular momentum transport. Here, the magnetic
diffusivity and the kinematic viscosity are not constant and vary by position
and $\alpha$-prescription is used for them. The govern equations on system
have been solved by the steady self-similar method. The solutions show the
radial velocity is highly subsonic and the rotational velocity behaves sub-
Keplerian. The rotational velocity for a specific value of the thermal
conduction coefficient becomes zero. This amount of conductivity strongly
depends on magnetic pressure fraction, magnetic Prandtl number, and viscosity
parameter. Comparison of energy transport by thermal conduction with the other
energy mechanisms implies that thermal conduction can be a significant energy
mechanism in resistive and magnetized ADAFs. This property is confirmed by
non-ideal magnetohydrodynamics (MHD) simulations.
00footnotetext: School of Physics, Damghan University, Damghan, Iran
e-mail: kfaghei@du.ac.ir
Received: 14 October 2011 / Accepted: 30 November 2011
Keywords accretion, accretion discs – conduction – magnetohydrodynamics: MHD
## 1 Introduction
The observational features in active galactic nuclei (AGN) and X-ray binaries
can be successfully explained by the standard geometrically thin, optically
thick accretion disc model (Shakura & Sunyaev 1973). The motion of the matter
in the standard thin model of the accretion disc is approximately Keplerian,
and the energy released in the accreting gas is radiated away locally. In the
past two decades, another type of accretion flow has been considered that the
energy released due to heating processes in the flow may be trapped within
accreting gas. As, only the small fraction of the energy released in the
accretion flow is radiated away due to inefficient cooling, and most of the
energy is stored in the accretion flow and advected to the central object.
This type of accretion flow is called as advection-dominated accretion flow
(ADAF). The models of ADAF have been studied by a number of researchers (e.g.
Ichimaru 1977; Rees et al. 1982; Narayan & Yi 1994; Abramowicz et al. 1995;
Blandford & Begelman 1999; Ogilvie 1999).
The observations of black holes confirm existence of hot accretion flow that
contrasted with classical cold and thin accretion disc model (Shakura &
Sunyaev 1973). Hot accretion flows can be seen in supermassive black holes of
galactic nuclei and during quiescent of accretion onto stellar-mass black
holes in X-ray transients (e.g., Lasota et al. 1996; Esin et al. 1997, 2001;
Narayan et al. 1998a; Menou et al. 1999; Di Matteo et al. 2000; see Narayan et
al. 1998b; Melia & Falcke 2001; Narayan 2002; Narayan & Quataert 2005 for
reviews). Chandra observations provide constraints on the density and
temperature of gas at or near the the Bondi capture radius in Sgr A∗ and
several nearby galactic nuclei (Loewenstein et al. 2001; Baganoff et al. 2003;
Di Matteo et al. 2003; Ho et al. 2003). Tanaka & Menou (2006) exploited these
constraints to calculate mean free path of the observed gas. They suggested
accretion in such systems are under weakly collisional condition. Moreover,
they suggested thermal conduction as a possible mechanism by which the
sufficient extra heating is provided in hot accretion flows. In following,
Johnson & Quataert (2007) studied the effects of electron thermal conduction
on the properties of a spherical hot accretion flows. Their model is
applicable for Sgr A∗ in the Galactic centre. Because, electron heat
conduction is important for low accretion rate systems. They also found a
supervirial temperature in the presence of thermal conduction. Similar to
Tanaka & Menou (2006), they assumed a steady state model, but they solved
their equations numerically. Abbassi et al. (2008) presented a set of self-
similar solutions for ADAFs with a toroidal magnetic field in which the
saturated thermal conduction has a important role in the energy transport in
the radial direction.
Since the observational evidences and magnetohydrodynamics (MHD) simulations
have expressed the toroidal magnetic field and the magnetic diffusivity are
important in accretion flows (see Faghei 2011a and references therein), Faghei
(2011a) examined the self-similar solutions of viscous and resistive ADAFs in
the presence of a toroidal magnetic field. But, he did not consider the
effects of thermal conduction in his model. While, the recent studies of
resistive accretion flows have represented that thermal conduction can play an
importance role in such systems (e.g. Sharma et al. 2008; Ghanbari et al.
2009). Sharma et al. (2008) studied the effects of thermal conduction on
magnetized spherical accretion flows using global axisymmetric MHD
simulations. In their model, when the magnetic energy density becomes
comparable to the gravitational potential energy density, the plasma due to
resistivity is heated to roughly the virial temperature, the mean inflow
becomes highly subsonic, and most of the energy released by accretion is
transported to large radii by thermal conduction and the accretion rate became
much smaller than Bondi accretion rate. Moreover, they found that for a larger
values of conductive parameter, energy transport through thermal conduction
becomes the dominant energy transport mechanism at small radii. Ghanbari et
al. (2009) studied a two-dimensional advective accretion disc bathed in the
poloidal magnetic field of a central accretor in the presence of thermal
conduction. They did not consider toroidal component of magnetic field and for
simplicity assumed the resistivity as a constant. They studied induction
equation of magnetic field in a steady state that is not according to anti-
dynamo theorem (e.g. Cowling 1981) and is useful only in particular systems
where the magnetic dissipation time is much longer than the age of the system.
On the other hand, this assumption implies that the flow is in balance between
escape and creation of the magnetic field.
In this paper, we adopt the presented solutions by Tanaka & Menou (2006) and
Faghei (2011a). Thus, we will investigate the influences of thermal conduction
on a viscous and resistive ADAF in the presence of a toroidal magnetic field.
Moreover, it is assumed that magnetic diffusivity in the present model is not
constant, and escaping and creating of magnetic field are unbalanced. From
some aspects will be shown that the present model is in according with the
observations and the resistive MHD simulations. The paper is organized as
follows. In section 2, the basic equations of constructing a model for quasi-
spherical magnetized advection dominated accretion flow with thermal
conduction will be defined. In section 3, self- similar method for solving
equations which govern the behaviour of the accreting gas was utilized. The
summary of the model will appear in section 4.
## 2 Basic Equations
We suppose a rotating and accreting gas around a Schwarzschild black hole of
mass $M$. The flow is assumed to be in advection dominated stage, where
viscous and resistive heating are balanced by the advection cooling and
thermal conduction. We use a spherical coordinate ($r$, $\theta$, $\phi$)
centred on the accreting object. Furthermore, the flow is assumed to be steady
and axisymmetric ($\partial_{t}=\partial_{\phi}=0$), and the equations will be
considered in the equatorial plane, $\theta=\pi/2$. Thus, all flow variables
are a function of only $r$. For the sake of simplicity, the general
relativistic effects are ignored and Newtonian gravity is used. The magnetic
field in the present model has only a toroidal component. Under assumptions,
the model is described by the following equations:
The continuity equation with mass loss is
$\frac{1}{r^{2}}\frac{d}{dr}(r^{2}\rho v_{r})=\dot{\rho},$ (1)
where $\rho$ is density, $v_{r}$ is the radial infall velocity, and
$\dot{\rho}$ the mass-loss rate per unit volume.
The radial equation of momentum is
$v_{r}\frac{dv_{r}}{dr}=r(\Omega^{2}-\Omega^{2}_{K})-\frac{1}{\rho}\frac{d}{dr}(\rho
c^{2}_{s})-\frac{c^{2}_{A}}{r}-\frac{1}{2\rho}\frac{d}{dr}(\rho c^{2}_{A}),$
(2)
where $\Omega$ is the angular velocity of the flow,
$\Omega_{K}=\sqrt{GM/r^{3}}$ is the Keplerian angular velocity, $c_{s}$ the
isothermal sound speed, which is defined as $c_{s}^{2}=p_{gas}/\rho$,
$p_{gas}$ being the gas pressure, and $c_{A}$ is Alfven speed, which is
defined as $c_{A}^{2}=B_{\varphi}^{2}/4\pi\rho=2p_{mag}/\rho$, $p_{mag}$ being
the magnetic pressure. The angular momentum transfer equation is
$\rho v_{r}\frac{d}{dr}(r^{2}\Omega)=\frac{1}{r^{2}}\frac{d}{dr}\left[\nu\rho
r^{4}\frac{d\Omega}{\partial r}\right],$ (3)
where the right-hand side of above equation describes the effects of viscous
torques due to shear ($\nu$, here, is kinematic coefficient of viscosity). As
noted in the introduction, we assume both of the kinematic coefficient of
viscosity and the magnetic diffusivity due to turbulence in the accretion
flow. Thus, it is reasonable to use these parameters in analogy to the
$\alpha$-prescription of Shakura & Sunyaev (1973) for the turbulent,
$\nu=P_{m}\eta=\alpha\frac{c_{s}^{2}}{\Omega_{K}},$ (4)
where $P_{m}$ is the magnetic Prandtl number, which is assumed a constant of
order of unity, $\eta$ is the magnetic diffusivity, and $\alpha$ is a free
parameter less than unity.
The energy equation becomes
$\displaystyle\frac{v_{r}}{\gamma-1}\frac{d}{dr}(\rho
c^{2}_{s})+\frac{\gamma}{\gamma-1}\frac{\rho
c^{2}_{s}}{r^{2}}\frac{d}{dr}\left(r^{2}v_{r}\right)=$ $\displaystyle
Q_{diss}-Q_{rad}+Q_{cond},$ (5)
where $Q_{diss}=Q_{vis}+Q_{resis}$ is the dissipation rate by viscosity
$Q_{vis}$ and resistivity $Q_{resis}$, $Q_{rad}$ represents the energy loss
through radiative cooling, and $Q_{cond}$ is the energy transported by thermal
conduction. For the right-hand side of the energy equation, we can write
$Q_{adv}=Q_{diss}-Q_{rad}+Q_{cond}$ (6)
where $Q_{adv}$ is the advective transport of energy. We exploit the advection
factor, $f=1-Q_{rad}/Q_{diss}$, that describes the fraction of the dissipation
energy which is stored in the accretion flow and advected into the central
object rather than being radiated away. In general, the advection factor
depends on the details of the heating and radiative cooling mechanisms and
will vary by position (e.g. Watari 2006, 2007; Sinha et al. 2009). Here, we
assume a constant $f$ for simplicity. Clearly, the flow in the case of $f=1$
is in the extreme limit of no radiative cooling and in the limit of efficient
radiative cooling, we have $f=0$.
As mentioned, the inner regions of hot accretion flows are collisionless and
the electron mean free path due to Coulomb collision is larger than the
radius. This property is described as saturation (Cowie & McKee 1977). The
traditional equation of heat flux due to thermal conduction,
$F_{cond}=-\kappa\nabla T$ which $\kappa$ being the thermal conduction
coefficient, is not valid for such systems. Because this equation is suitable
for the collisional plasma, in which mean free path of electron energy
exchange is smaller than temperature scale height. Cowie & McKee (1977)
derived the heat flux for the collisionless plasma as
$F_{sat}=5\phi_{s}\rho c_{s}^{3}=5\phi_{s}p\sqrt{\frac{p}{\rho}},$ (7)
where $\phi_{s}$ is a factor is less than unity and is called as saturation
constant. Now, the viscous and resistive heating rates and the energy
transport by thermal conduction are expressed as
$Q_{vis}=\nu\rho r^{2}\left(\frac{\partial\Omega}{\partial r}\right)^{2}$ (8)
$Q_{resis}=\frac{\eta}{4\pi}{\bf J}^{2}$ (9)
$Q_{cond}=-\left|\frac{1}{r^{2}}\frac{\partial}{\partial
r}\left(r^{2}F_{sat}\right)\right|$ (10)
where ${\mathbf{J}}=\nabla\times{\mathbf{B}}$ is the current density,
${\mathbf{B}}$ being the magnetic field. We used a minus sign in equation
(10). Because we want thermal conduction as energy transport mechanism
outward. On the other hand, heat conduction flux will be behaved like a
cooling mechanism in accretion flow. Finally, since we consider only the
toroidal field, the induction equation with field escape can be written as
$\frac{1}{r}\frac{d}{dr}\left[rv_{r}B_{\varphi}-\eta\frac{d}{dr}(rB_{\varphi})\right]=\dot{B}_{\varphi}.$
(11)
where $B_{\varphi}$ is the toroidal component of magnetic field and
$\dot{B}_{\varphi}$ is the field escaping/creating rate due to a magnetic
instability or dynamo effect. This induction equation is rewritten as
$\displaystyle\frac{1}{r}\frac{d}{dr}\left[\sqrt{4\pi\rho
c^{2}_{A}}\left(rv_{r}-\frac{\alpha}{4\beta
P_{m}}\frac{1}{r\rho\Omega_{K}}\frac{d}{dr}(r^{2}\rho
c^{2}_{A})\right)\right]=$ $\displaystyle\dot{B}_{\varphi},$ (12)
where $\beta$ is the degree of magnetic pressure to the gas pressure and can
be defined by
$\beta=\frac{p_{mag}}{p_{gas}}=\frac{1}{2}\left(\frac{c_{A}}{c_{s}}\right)^{2}.$
(13)
In this paper, we apply the steady self-similar methods to solve the system
equations. Thus, this parameter will be constant throughout the disc. While,
Khesali & Faghei (2008, 2009) showed that it varies by position. In hot
accretion flows, typical value of $\beta$ is in the range $0.01$-$1$ (e.g. De
Villiers et al. 2003; Beckwith et al. 2008). Here, we will also consider the
magnetically dominated case ($\beta>1$). Because, when thermal instability
happens in an ADAF, the MHD numerical simulations imply that the thermal
pressure rapidly decreases while the magnetic pressure increases due to the
conservation of magnetic flux (Machida et al. 2006). This will result in large
$\beta$ and forms a magnetically dominated accretion flow (Bu et al. 2009).
## 3 Self-Similar Solutions
### 3.1 Analysis
The self-similar method is useful to understand of the physics of accretion
flows. This method is familiar due to its wide range of applications in many
research fields of astrophysics. Self-similar solution, although constituting
only a limited part of problem, is often useful to understand the basic
behaviour of the system. Thus, in order to seek similarity solutions for the
above equations, we seek solutions in the following form:
$v_{r}(r)=-c_{1}\alpha\sqrt{\frac{GM_{*}}{r}}$ (14)
$\Omega(r)=c_{2}\sqrt{\frac{GM_{*}}{r^{3}}}$ (15)
$c^{2}_{s}(r)=c_{3}\frac{GM_{*}}{r}$ (16)
$c^{2}_{A}(r)=\frac{B^{2}_{\varphi}}{4\pi\rho}=2\beta c_{3}\frac{GM_{*}}{r}$
(17)
where coefficients $c_{1}$, $c_{2}$, and $c_{3}$ are determined later. We
assume a power-law relation for density
$\rho(r)=\rho_{0}r^{s}$ (18)
where $\rho_{0}$ and $s$ are constant. By using above self-similar quantities,
the mass-loss rate and the field escaping/creating rate must have the
following form:
$\dot{\rho}(r)=\dot{\rho}_{0}r^{s-3/2}$ (19)
$\dot{B}_{\varphi}(r)=\dot{B}_{0}r^{\frac{s-4}{2}}$ (20)
where $\dot{\rho}_{0}$ and $\dot{B}_{0}$ are constant.
Substituting the above solutions in the continuity, momentum, angular
momentum, energy, and induction equations [(1)-(3), (5), and (11)], we can
obtain the following relations:
$\dot{\rho}_{0}=-\left(s+\frac{3}{2}\right)\alpha\rho_{0}c_{1}\sqrt{GM_{*}},$
(21)
$-\frac{1}{2}c^{2}_{1}\alpha^{2}=c^{2}_{2}-1-c_{3}\left[s-1+\beta(1+s)\right],$
(22)
$c_{1}=3(s+2)c_{3},$ (23) $\displaystyle-\alpha
c_{1}\left[\frac{2(s-1)+3\gamma}{\gamma-1}\right]=\alpha
f\left[\frac{9}{2}c^{2}_{2}+\frac{\beta}{P_{m}}c_{3}(1+s)^{2}\right]$
$\displaystyle-10\phi_{s}c_{3}^{1/2}\left|s+\frac{1}{2}\right|,$ (24)
$\dot{B}_{0}=-\frac{\alpha s}{2}GM_{*}\sqrt{2\pi\rho_{0}\beta
c_{3}}\left[2c_{1}+\frac{c_{3}}{P_{m}}(1+s)\right].$ (25)
Above equations express for $s=-3/2$, there is no mass loss, while for
$s>-3/2$ mass loss (wind) exists. After algebraic manipulations, we obtain an
algebraic equation for $c_{3}$:
$\displaystyle{\frac{81}{8}}\,{\alpha}^{3}(s+2)^{2}f{c_{{3}}}^{2}+c_{3}\Bigg{[}\frac{3\,\alpha(s+2)}{2}\times\frac{2(s-1)+3\gamma}{\gamma-1}$
$\displaystyle-\frac{9\alpha
f}{4}\,\Bigg{(}\frac{2\beta}{9P_{m}}(s+1)^{2}+(s+1)\beta+s-1\Bigg{)}\Bigg{]}$
$\displaystyle-5\phi_{s}\left|s+\frac{1}{2}\right|\sqrt{c_{{3}}}-\frac{9}{4}\,f\alpha=0,$
(26)
and the rest of the physical variables are
$\dot{\rho}_{0}=-3\alpha\rho_{0}\sqrt{GM_{*}}(s+2)\left(s+\frac{3}{2}\right)c_{3},$
(27) $c_{1}=3c_{3}(s+2),$ (28)
$c_{2}^{2}=1-\frac{9\alpha^{2}}{2}(s+2)^{2}c_{3}^{2}+c_{3}\left[(s+1)\beta+s-1\right],$
(29) $\displaystyle\dot{B}_{0}=-3\alpha
sGM_{*}(s+2)^{3/2}c_{3}^{3/2}\sqrt{6\pi\rho_{0}\beta}\times$
$\displaystyle\left[1+\frac{s+1}{6P_{m}(s+2)}\right].$ (30)
We can solve algebraic equation (26) numerically and clearly only real roots
which correspond to positive $c_{1}$ are physically acceptable. Without
thermal conduction, i.e. $\phi_{s}=0$, (26) and similarity solutions reduce to
Faghei (2011a). But our main algebraic equation includes thermal conduction.
Fig. 1 : Physical quantities of the flow as a function of saturation constant
for $\gamma=1.3$, $\alpha=0.5$, $f=1$, and $P_{m}=100$. Solid, dotted, dashed,
and dot-dashed lines represent $\beta=0.1,0.3,0.7$, and $1.5$.
### 3.2 Numerical Results
Now, we consider the behaviour of solutions in the presence of thermal
conduction. But in this paper, only case of no wind ($s=-3/2$) is considered
that $\dot{\rho}_{0}=0$ and $\dot{B}_{\varphi}\propto r^{-11/4}$. In addition
to introduced coefficients, we also define a new parameter of $c_{4}$ that is
the right-hand side of equation (24)
$c_{4}=\alpha
f\left[\frac{9}{2}c^{2}_{2}+\frac{\beta}{P_{m}}c_{3}(1+s)^{2}\right]-10\phi_{s}c_{3}^{1/2}\left|s+\frac{1}{2}\right|.$
(31)
Equations (5), (24), and (31) imply that the parameter $c_{4}$ is the
advection transport of energy. The behaviour of coefficients $c_{i}$ as a
function of $\phi_{s}$ are shown in Figures 1-3. Moreover, Figure 1 represents
the profiles of physical quantities for several values of the magnetic
pressure fraction, i. e. $\beta=0.1,0.3,0.7$, and $1.5$. The value of $\beta$
measures the strength of magnetic field, and a larger $\beta$ denotes a
stronger magnetic field. Figure 2 represents the profiles of physical
quantities for several values of magnetic Prandtl number, i. e.
$P_{m}=\infty,1,2/3$, and $1/2$. The smaller values of $P_{m}$ denotes a
stronger magnetic diffusivity, $\eta$. Figure 3 shows the profiles of physical
quantities for several values of adiabatic index, i. e. $\gamma=1.2,1.25,1.3$,
and $1.35$.
Fig. 2 : Same as Fig. 1, but $\beta=1.0$, and solid, dotted, dashed, and dot-
dashed lines represent $P_{m}=\infty,1,2/3$, and $1/2$.
Fig. 3 : Same as Fig. 1, but $\beta=1.0$, and solid, dotted, dashed, and dot-
dashed lines represent $\gamma=1.2,1.25,1.3$, and $1.35$.
The solutions in Figures 1-3 imply that the radial infall velocity, $c_{1}$,
and the sound speed, $c_{3}$, both decrease with the magnitude of conduction,
while the squared angular velocity, $c_{2}^{2}$, increases. These properties
are qualitatively consistent with dynamical analysis of Faghei (2011b). One
and two dimensional simulations of hot accretion flows have also shown that
thermal conduction reduces the flow temperature (Sharmal et al. 2008; Wu et
al. 2010). The profiles of advection transport of energy, $c_{4}$, in Figures
1-3 imply that thermal conduction behaves as a cooling mechanism, resulting in
a local decrease of the gas temperature relative to the original ADAF
solution. At the same time, the gas adjust its angular velocity (which
increases the level of viscous dissipation) and reduces its inflow speed.
Figure 1 shows the radial velocity, sound speed, and advection transport of
energy decrease by adding the magnetic pressure fraction, $\beta$. These
properties are qualitatively consistent with results of Bu et al. (2009) and
Faghei (2011a). Moreover, the angular velocity decreases with the magnitude of
magnetic field that is in according with results of Khesali & Faghei (2009)
and Faghei (2011a).
Figure 2 shows the magnetic diffusivity has the opposite effects of thermal
conduction on the physical variables. As, by adding the magnetic diffusivity,
$P_{m}^{-1}$, the radial velocity, sound speed, and advection transport of
energy increase, while the rotational velocity decrease. These results are
similar to resistive ADAF models without thermal conduction (e. g. Faghei
2011a).
Figure 3 represents the gas adiabatic index similar to magnetic diffusivity
has the opposite effects of thermal conduction on the physical quantities. As,
with the magnitude of $\gamma$, the inflow and sound speed increase, while the
rotational velocity decreases. This is in accord with dynamical study of hot
accretion flow (Faghei 2011b). Moreover, Figure 3 shows that the gas adiabatic
index contributes with thermal conduction to reduce advection transport of
energy. It can be due to decrease of rotational velocity by adding $\gamma$,
which reduces the level of viscous dissipation.
The studies of hot accretion flows (e. g. Shadmehri 2008) imply that the
solution for a given set of the input parameters reaches to a non-rotating
limit at a specific of $\phi_{s}$ which we denote it by $\phi_{s}^{c}$. With
zero insertion of $c_{2}$ in equations (22)-(24), $\phi_{s}^{c}$ can be
written as
$\displaystyle\phi_{s}^{c}=\frac{1}{60}\left[\frac{\beta
f}{P_{m}}-\frac{5/3-\gamma}{\gamma-1}\right]\times$
$\displaystyle\sqrt{2(\beta+5)\left[-1+\sqrt{1+\frac{18\alpha^{2}}{(\beta+5)^{2}}}\,\right]}.$
(32)
We can not extend the studies beyond $\phi_{s}^{c}$, because the right-hand
side of equation (29) becomes negative and a negative $c_{2}^{2}$ is clearly
unphysical. As, the rotational velocity profiles in Figures 1-3 show, we have
selected the input parameters that $c_{2}^{2}$ is positive. The behaviour of
critical saturation constant, $\phi_{s}^{c}$, as a function of $\beta$ for
several values of magnetic Prandtl number is shown in Figures 4. In left panel
of Figure 4, the viscosity parameter value is $\alpha=0.1$, and in right panel
is $\alpha=0.2$. The profiles of Figure 4 show the critical saturation
constant highly depends on magnetic pressure fraction, $\beta$, magnetic
Prandtl number, $P_{m}$, and viscosity parameter, $\alpha$. As, higher values
of $\beta$ corresponds to larger $\phi_{s}^{c}$. Critical saturation constant
also increases with higher value of $\alpha$. Since, magnetic diffusivity is
proportional to $P_{m}^{-1}$, Figure 4 shows that the magnetic diffusivity
similar to magnetic field increases $\phi_{s}^{c}$ value.
As mentioned in the introduction, Sharma et al. (2008) by resistive MHD
simulation studied a spherical accretion with thermal conduction. They found
for even modest thermal conductivities, conduction is the significant
mechanism of energy. Here, to compare conduction mechanism to others, we study
the ratio of energy transport by thermal conduction, $Q_{cond}$, to the gas
heating rate by viscosity $Q_{vis}$ and resistivity $Q_{resis}$. Such
solutions are shown in Figure 5 for two cases of non-resistive (left-panel)
and resistive (right-panel) flows. The solutions imply that thermal conduction
is the significant energy mechanism in the flow. This result confirms
simulation of Sharma et al. (2008). Moreover, the ratio of $Q_{cond}$ to
$Q_{diss}$ increases slightly by adding the magnetic field and does not change
for different values of magnetic Prandtl number. Because, a large fraction of
$Q_{diss}$ is generated by viscous dissipation.
## 4 Summary and Discussion
The collision timescale between ions and electrons in hot accretion flows is
longer than the inflow timescale. Thus, the inflow plasma is collisionless,
and transfer of energy by thermal conduction can be dynamically important. The
low collisional rate of the gas is confirmed by direct observation,
particularly in the case of the Galactic centre (Quataert 2004; Tanaka & Menou
2006) and in the intracluster medium of galaxy clusters (Sarazin 1986).
In this paper, the structure of a magnetized ADAF in the presence of
resistivity and thermal conduction is investigated. We assumed the magnetic
field has a purely toroidal component. We adopted the presented solutions by
Tanaka & Menou (2006) and Faghei (2011a). Thus, we assumed that angular
momentum transport is due to viscous turbulence and the $\alpha$\-
prescription is used for the kinematic coefficient of viscosity. We also
assumed the flow does not have a good cooling efficiency and so a fraction of
energy accretes along with matter on to the central object. In order to solve
the equations that govern the structure behaviour of magnetized ADAF with
thermal conduction, we have used steady self-similar solution.
Fig. 4 : The critical saturation constant as a function of the ratio of
magnetic pressure to gas pressure. Solid, dotted, dashed, and dot-dashed lines
represent $P_{m}=10,5,1$, and $0.5$. The input parameters are set to
$\gamma=5/3$, $f=1$, $s=-3/2$, and the viscous parameter $\alpha$ in left-
panel is $0.1$ and in right-panel is $0.2$.
The solutions showed the radial infall velocity and sound speed in the
presence of thermal conduction both decrease, while angular velocity increase.
These properties are consistent with dynamical study of hot accretion flow (e.
g. Faghei 2011b) and from some aspects also are in accord with simulations of
Sharma et al. (2008) and Wu et al. (2010). Moreover, the solutions represent
the magnetic diffusivity and thermal conduction have the opposite effects on
physical quantities. For a moderate thermal conduction, the solutions imply
that thermal conduction can play an important role in energy mechanism of the
system. This property is qualitatively consistent with non-ideal simulations
of Sharma et al. (2008).
In the present model, accretion flow is studied in one-dimensional approach
and ignored from latitudinal dependence of physical quantities. There are some
researches in two-dimensional approach that express the importance of such
studies (Tanaka & Menou 2006; Ghanbari et al. 2009; Wu et al. 2010 ). Thus,
latitudinal study of present model can be investigated in other research.
Here, we used a saturated heat conduction flux. While, there are some studies
with unsaturated heat flux (e. g. Shcherbakkov & baganoff 2010) that show a
good agreement with observations. Thus, the study of the present model in a
unsaturated case will be interesting.
Fig. 5 : The ratio of energy transport by thermal conduction, $Q_{cond}$, to
the gas heating rate by viscosity $Q_{vis}$ and resistivity $Q_{resis}$ as a
function of saturation constant. Solid, dotted, dashed, and dot-dashed lines
represent $\beta=0.0,1.0,3.0$, and $5.0$. The input parameters are set to
$\gamma=4/3$, $f=1$, $s=-3/2$, $\alpha=0.5$, and the magnetic Prandtl number
in left-panel is $\infty$ and in right-panel is $0.5$.
## Acknowledgements
I wish to thank the anonymous referee for very useful comments that helped me
to improve the initial version of the paper. I would also like to thank Roman
V. Shcherbakov for his helpful comments.
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|
arxiv-papers
| 2011-12-12T20:11:22 |
2024-09-04T02:49:25.205520
|
{
"license": "Creative Commons - Attribution - https://creativecommons.org/licenses/by/3.0/",
"authors": "Kazem Faghei",
"submitter": "Kazem Faghei",
"url": "https://arxiv.org/abs/1112.2678"
}
|
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